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warning/0002/math-ph0002009.html | ar5iv | text | # Szego kernels and a theorem of Tian
## 1. Introduction
A variety of results in complex geometry and mathematical physics depend upon the analysis of holomorphic sections of high powers $`L^N`$ of holomorphic line bundles $`LM`$ over compact Kähler manifolds (\[A\]\[Bis\]\[Bis.V\] \[Bou.1\]\[Bou.2\]\[B.G\]\[D\]\[Don\] \[G\]\[G.S\]\[K\]\[Ji\] \[T\] \[W\]). The principal tools have been Hörmander’s $`L^2`$-estimate on the $`\overline{}`$-operator over $`M`$ \[T\], the asymptotics of heat kernels $`k_N(t,x,y)`$ for associated Laplacians \[Bis\]\[Bis.V\]\[D\]\[Bou.1\]\[Bou.2\]\[G\], the method of stationary phase for formal functional integrals \[A\]\[W\] and the microlocal analysis of Szegö and Bergman kernels \[B.F.G\]\[B.S\]\[B.G\]. In this note we wish to apply the latter methods, specifically the Boutet de Monvel-Sjöstrand parametrix for the Szegö kernel, to a problem in complex geometry. Our purpose is to prove the following theorem:
###### Theorem 1.
Let $`M`$ be a compact complex manifold of dimension $`n`$ (over $`\mathrm{C}`$) and let $`(L,h)M`$ be a positive hermitian holomorphic line bundle. Let $`g`$ be the Kähler metric on $`M`$ corresponding to the Kähler form $`\omega _g:=Ric(h)`$. For each $`N𝐍`$, $`h`$ induces a hermitian metric $`h_N`$ on $`L^N`$. Let $`\{S_0^N,\mathrm{},S_{d_N}^N\}`$ be any orthonormal basis of $`H^0(M,L^N)`$, $`d_N=\mathrm{dimH}^0(\mathrm{M},\mathrm{L}^\mathrm{N})`$, with respect to the inner product $`s_1,s_2_{h_N}=_Mh_N(s_1(z),s_2(z))𝑑V_g`$. Here, $`dV_g=\frac{1}{n!}\omega _g^n`$ is the volume form of $`g`$. Then there exists a complete asymptotic expansion:
$$\underset{i=0}{\overset{d_N}{}}S_i^N(z)_{h_N}^2=a_0N^n+a_1(z)N^{n1}+a_2(z)N^{n2}+$$
for certain smooth coefficients $`a_j(z)`$ with $`a_0=1`$. More precisely, for any $`k`$
$$\underset{i=0}{\overset{d_N}{}}S_i^N(z)_{h_N}^2\underset{j<R}{}a_j(x)N^{nj}_{C^k}C_{R,k}N^{nR}.$$
Above, $`Ric(h)`$ is the Ricci curvature of $`h`$, given locally by $`\frac{\sqrt{1}}{2\pi }\overline{}a`$ where $`a=e_L_h`$ is the positive function locally representing $`h`$ in a local holomorphic frame $`e_L`$.
This theorem has a number of corollaries. First it implies that for sufficiently large $`N`$, there are no common zeroes of the sections $`\{S_0^N,\mathrm{},S_{d_N}^N\}`$. Hence one can define the holomorphic map
(1)
$$\varphi _N:M\text{ }\mathrm{C}𝐏^{d_N},z[S_0^N(z),\mathrm{},S_{d_N}^N(z)]$$
where $`[S_0^N(z),\mathrm{},S_{d_N}^N(z)]`$ denotes the line thru $`(S_0^N(z),\mathrm{},S_{d_N}^N(z))`$ as defined in a local holomorphic frame. Since all the components transform by the same scalar under a change of frame, the line is well-defined. As is well-known \[G.H\], $`\varphi _N`$ is equivalent to the invariantly defined map
(2)
$$\stackrel{~}{\varphi }_N:M𝐏H^0(M,L^N)^{},zH_z:=\{sH^0(M,L^N):s(z)=0\}.$$
Secondly it gives an asymptotic formula for the distortion function between the metrics $`h_N`$ and $`h_{FS,N}`$, where $`h_{FS,N}`$ is the Fubini-Study metric on $`L^N`$ induced by $`\stackrel{~}{\varphi }_N`$. The following result was simeltaneously proved in the case of abelian varieties by G.Kempf \[K\] and S.Ji \[Ji\] (with $`C^0`$ convergence) and for general projective varieties (with $`C^4`$ convergence) by G.Tian (\[T\], Lemma 3.2(i)). Heat kernel proofs were later found by T. Bouche \[Bou.1\]\[Bou.2\] and J.P.Demailly \[D\].
###### Corollary 2.
Let $`G`$ be any Riemannian metric on $`M`$, endow $`H^0(M,L^N)`$ with the Hermitian inner product induced by $`(G,h_N)`$ and define the map $`\stackrel{~}{\varphi }_N`$ as above. Identify $`L^N`$ with the pull-back $`\stackrel{~}{\varphi }_N^{}O(1)`$ of the hyperplane bundle $`O(1)𝐏H^0(M,L^N)`$, and let $`h_{FS,N}`$ be the pullback of the standard Hermitian metric on $`O(1).`$ Then
$$\frac{h_{N,z}}{h_{FS,N,z}}=(\frac{N}{2\pi })^n|\alpha _1(z)\mathrm{}\alpha _n(z)|+O(N^{n1}).$$
where $`\alpha _1(z),\mathrm{},\alpha _n(z)`$ are the eigenvalues $`Ric(h)`$ with respect to $`G`$.
If we take the background metric $`G`$ to be the Kähler metric associated to the Kähler form $`\omega _g=Ric(h)`$, then the curvature eigenvalues are all equal to one.
Third, it implies:
###### Corollary 3.
Let $`\omega _{FS}`$ denote the Fubini-Study form on $`\text{ }\mathrm{C}𝐏^{d_N}`$. Then:
$$\frac{1}{N}\varphi _N^{}(\omega _{FS})\omega _g_{C^k}=O(\frac{1}{N})$$
for any $`k`$.
This statement for $`k2`$ was the principal result of Tian (\[T\], Theorem A). The map $`\varphi _N`$ depends on the choice of $`\{S_0^N,\mathrm{},S_{d_N}^N\}`$ but it is easily seen that $`\varphi _N^{}(\omega _{FS})`$ does not. This result follows formally from the preceding one by taking the curvature of both sides of the asymptotic formula.
Our result strengthens the previous ones in two ways: First, it shows that the convergence of $`\frac{1}{N}\varphi _N^{}(\omega _{FS})\omega _g`$ takes place in the $`C^{\mathrm{}}`$-topology and not just in $`C^2`$. This was conjectured in \[T\]. Second, it shows that this convergence is just the first term of a complete asymptotic expansion. If only the principal terms are desired, then the proof could be simplified further: as in \[Z\], which treats an analogous problem in the context of Zoll manifolds, one could obtain the principal terms by the symbol calculus of Toeplitz operators. But we believe that the lower order terms should also be of interest.
The proof begins by expressing the maps $`\varphi _N`$ and $`\stackrel{~}{\varphi }_N`$ in terms of an associated equivariant map $`\mathrm{\Phi }_N`$ on the unit circle bundle $`X`$ of the dual line bundle $`L^{}`$ with respect to the induced metric $`h`$. Roughly, this converts the holomorphic geometry of $`L`$ to the CR geometry of $`X`$. Since $`X`$ is the boundary of the strictly pseudo-convex domain $`D=\{vL^{}:|v|_h<1\}L^{}`$ it has a Szegö kernel $`\mathrm{\Pi }(x,y)`$ which projects $`L^2(X)`$ to the Hardy space $`H^2(X)`$ of boundary values of holomorphic functions in $`D`$. Under the natural $`S^1`$ action of $`XM`$, $`H^2(X)`$ splits up into weight spaces $`H_N^2(X)`$ and one has a canonical isomorphism $`s\widehat{s}:H^0(M,L^N)H_N^2(X)`$, $`\widehat{s}(z,u)=u,s(z)`$ where we write a point of $`X`$ as $`(z,u),uL_z^{},|u|_h=1`$ and where $`,`$ is the pairing between $`L`$ and $`L^{}.`$ So the basis $`\{S_0^N,\mathrm{},S_{d_N}^N\}`$ of $`H^0(M,L^N)`$ corresponds to an orthonormal basis $`\{\widehat{S}_0^N,\mathrm{},\widehat{S}_{d_N}^N\}`$ of $`H_N^2(X)`$. One then gets an associated CR map of $`X`$ into $`H_N^2(X)^{}`$. Expressed invariantly in terms of the orthogonal projection $`\mathrm{\Pi }_N`$ onto $`H_N^2(X)`$ it is defined by:
(3)
$$\mathrm{\Phi }_N(x)=\mathrm{\Pi }_N(x,):XH_N^2(X)^{}.$$
Using the canonical isomorphism above we get an essentially equivalent map $`\stackrel{~}{\mathrm{\Phi }}_N:XH^0(M,L^N)^{}`$. (We note that these maps are well-defined even when the set $`Z_N`$ of common zeroes of the sections is non-empty.) Thus for $`N0`$ we get the diagram:
(4)
$$\begin{array}{ccc}X\hfill & \stackrel{\stackrel{~}{\mathrm{\Phi }}_N}{}\hfill & H^0(M,L^N)^{}0\hfill \\ \pi \hfill & & \rho \hfill \\ M\hfill & \stackrel{\stackrel{~}{\varphi }_N}{}\hfill & 𝐏H^0(M,L^N)^{}\hfill \end{array}$$
where $`\pi `$ and $`\rho `$ are the canonical projections. The top arrow is well-defined for all $`N`$.
After unravelling the identifications, we find (§1) that $`_{i=0}^{d_N}S_i^N(z)_{h_N}^2=\mathrm{\Pi }_N(x,x)`$ and that $`\frac{1}{N}\varphi _N^{}\omega _{FS}=\omega _g+\frac{i}{2\pi }\overline{}_b_b\mathrm{log}\mathrm{\Pi }_N(x,x)`$ for any $`x`$ with $`\pi (x)=z`$. The second term on the right side is an $`S^1`$ invariant form so we have identified it with a form on $`M`$. The theorem is therefore equivalent to the statement that $`\mathrm{\Pi }_N(x,x)`$ has a complete asymptotic expansion as $`N\mathrm{}`$ which can be differentiated any number of times. This will follow by applying the method of stationary phase to Boutet de Monvel-.Sjöstrand’s parametrix for the Szegö projector $`\mathrm{\Pi }(x,y)`$\[B.S\] (see §3).
The analysis of Szegö kernels should have other applications in complex geometry. By studying the off-diagonal of $`\mathrm{\Pi }_N(x,y)`$ one can show that the maps $`\mathrm{\Phi }_N`$ are embeddings, thus obtaining an analytic proof of the Kodaira embedding theorem. In a forthcoming paper \[S.Z\], B.Shiffman and the author also use the Szegö kernels to show (among other things) that the zeroes of a ‘random section’ of $`H^0(M,L^N)`$ become uniformly distributed as $`N\mathrm{}.`$ There are also some potential analogues in the almost complex setting. In a recent paper \[B.U\], Borthwick-Uribe conjecture some results on Szegö kernels in the almost complex setting which seem very close to what is proved here in the complex setting. In part their motivation ( as well as ours) was to reinterpret some constructions of Donaldson \[Don\] from the viewpoint of semiclassical analysis.
We thank B.Shiffman for help with the relevant complex geometry and for his and M. Zworski’s encouragement to publish this note.
## 2. From line bundle to circle bundle
The purpose of this section is to convert the statements of Theorem 1 and Corollaries 2 \- 3 into statements about $`\mathrm{\Pi }_N.`$ Before doing so let us recall why one exists in this context and establish some notation.
### 2.1. The CR setting
Let $`O(1)\text{ }\mathrm{C}𝐏^n`$ denote the hyperplane section line bundle and let $`,`$ denote its natural Hermitian metric. Let $`M\text{ }\mathrm{C}P^n`$ be a non-singular projective variety, let $`L`$ denote the restriction of $`O(1)`$ to $`M`$ and let $`h`$ denote the restriction of $`,`$ to $`L`$. The following proposition is well-known (it was originally observed by Grauert in the 50’s):
###### Proposition 4.
Let $`D=\{(m,v)L^{}:h(v,v)1\}`$. Then $`D`$ is a strictly pseudoconvex domain in $`L`$.
Here $`L^{}`$ is the dual line bundle to $`L`$. The boundary of $`D`$ is a principal $`S^1`$ bundle $`XM`$ whose defining function is given by
(5)
$$\rho :L^{}\mathrm{I}\mathrm{R},\rho (z,\nu )=1|\nu |_z^2$$
where $`\nu L_z^{}`$ and where $`|\nu |_z`$ is its norm in the metric induced by $`h`$. That is, $`D=\{\rho >0\}`$. In a local coframe $`e_L^{}`$ over $`UM`$ we may write $`\nu =\lambda e_L^{}`$ and then $`|\nu |_z^2=a(z)|\lambda |^2`$ where $`a(z)=|e_L^{}|_z^2`$ is a positive smooth function on $`U`$. Thus in local holomorphic coordinates $`(z,\lambda )`$ on $`L^{}`$ the defining function is given by $`\rho =1a(z)|\lambda |^2.`$ We will denote the $`S^1`$ action by $`r_\theta x`$ and its infinitesimal generator by $`\frac{}{\theta }`$. We note that $`\rho `$ is $`S^1`$-invariant.
Let us denote by $`T^{}D,T^{\prime \prime }DTD\text{ }\mathrm{C}`$ the holomorphic, resp. anti-holomorphic subspaces and define $`d^{}f=df|_T^{},d^{\prime \prime }f=df|_{T^{\prime \prime }}`$ for $`fC^{\mathrm{}}(D).`$ Then $`X`$ inherits a CR structure $`TX\text{ }\mathrm{C}=T^{}T^{\prime \prime }\text{ }\mathrm{C}\frac{}{\theta }`$. Here $`T^{}X`$ (resp. $`T^{\prime \prime }X`$) denotes the holomorphic (resp. anti-holomorphic vectors) of $`D`$ which are tangent to $`X`$. They are given in local coordinates by vector fields $`a_j\frac{}{\overline{z}_j}`$ such that $`a_j\frac{}{\overline{z}_j}\rho =0.`$ A local basis is given by the vector fields $`Z_j^k=\frac{}{\overline{z}_j}(\frac{\rho }{\overline{z}_k})^1(\frac{\rho }{\overline{z}_j})\frac{}{\overline{z}_k}`$ ($`jk.`$)
The Cauchy-Riemann operator on $`X`$ is defined by
(6)
$$\overline{}_b:C^{\mathrm{}}(X)C^{\mathrm{}}(X,(T^{\prime \prime })^{}),\overline{}_bf=df|_{T^{\prime \prime }}.$$
In terms of the local basis above, it is given by
(7)
$$\overline{}_bf=\underset{jk}{}Z_j^kfd\overline{z}_j|_{T^{\prime \prime }}.$$
Also associated to $`X`$ are
(8)
$$\begin{array}{cc}\hfill & \mathrm{the}\mathrm{contact}\mathrm{form}\alpha =\frac{1}{\mathrm{i}}\mathrm{d}^{}\rho |_\mathrm{X}=\frac{1}{\mathrm{i}}\mathrm{d}^{\prime \prime }\rho |_\mathrm{X}\hfill \\ & \\ \hfill & \mathrm{the}\mathrm{volume}\mathrm{form}\mathrm{d}\mu =\alpha (\mathrm{d}\alpha )^\mathrm{n}\hfill \\ & \\ \hfill & \mathrm{the}\mathrm{Levi}\mathrm{form}\mathrm{L}_\rho (\mathrm{z})=\frac{^2\rho }{\mathrm{z}_\mathrm{j}\overline{\mathrm{z}}_\mathrm{k}}\mathrm{z}_\mathrm{j}\overline{\mathrm{z}}_\mathrm{k}.\hfill \\ & \\ \hfill & \mathrm{the}\mathrm{Levi}\mathrm{form}\mathrm{on}\mathrm{X}\mathrm{L}_\mathrm{X}=\mathrm{L}_\rho |_{\mathrm{T}^{}\mathrm{T}^{\prime \prime }\mathrm{TX}}\hfill \end{array}$$
which are independent of the choice of $`\rho .`$ The Levi form on $`X`$ is related to $`d\alpha =\pi ^{}\omega _g`$ by: $`L_X(V,W)=d\alpha (V,\overline{W}).`$ Since $`\omega _g`$ is Kähler, $`D`$ is a strictly pseudoconvex domain.
The Hardy space $`H^2(X)`$ is the space of boundary values of holomorphic functions on $`D`$ which are in $`L^2(X)`$, or equivalently $`H^2=(\mathrm{ker}\overline{}_\mathrm{b})\mathrm{L}^2(\mathrm{X}).`$ The $`S^1`$ action commutes with $`\overline{}_b`$, hence $`H^2(X)=_{N=1}^{\mathrm{}}H_N^2(X)`$ where $`H_N^2(X)=\{fH^2(X):f(r_\theta x)=e^{iN\theta }f(x)\}.`$
A section $`s`$ of $`L`$ determines an equivariant function $`\widehat{s}`$ on $`L^{}0`$ by the rule: $`\widehat{s}(z,\lambda )=\lambda ,s(z)`$ ( $`zM,\lambda L_z^{}.`$) It is clear that if $`\tau \text{ }\mathrm{C}^{}`$ then $`\widehat{s}(z,\tau \lambda )=\tau \widehat{s}.`$ We will usually restrict $`\widehat{s}`$ to $`X`$ and then the equivariance property takes the form: $`\widehat{s}(r_\theta x)=e^{i\theta }\widehat{s}(x).`$ Similarly, a section $`s_N`$ of $`L^N`$ determines an equivariant function $`\widehat{s}_N`$ on $`L^{}0`$: put $`\widehat{s}_N(z,\lambda )=\lambda ^N,s_N(z)`$ where $`\lambda ^N=\lambda \lambda \mathrm{}\lambda .`$ The following proposition is well-known:
###### Proposition 5.
The map $`s\widehat{s}`$ is a unitary equivalence between $`H^0(M,L^N)`$ and $`H_N^2(X).`$
As above, we let $`\mathrm{\Pi }_N:L^2(X)H_N^2(X)`$ denote the orthogonal projection. Its kernel is defined by
(9)
$$\mathrm{\Pi }_Nf(x)=_X\mathrm{\Pi }_N(x,y)f(y)𝑑\mu (y).$$
This definition differs from that of \[B.S\] in using $`d\mu `$ as the reference density.
### 2.2. Line bundles and maps to projective space
Since the definitions of the various maps $`\varphi _N,\stackrel{~}{\varphi }_N,\mathrm{\Phi }_N,\stackrel{~}{\mathrm{\Phi }}_N`$ involve some identifications, we pause to recall some basic facts about maps to projective space \[\[G.H\], I.4\].
Let $`EM`$ denote a holomorphic line bundle. Since we are interested in $`E=L^N`$ for large $`N`$ we may assume that not all sections $`sH^0(M,E)`$ vanish at any point $`zM.`$ Then the space of sections vanishing at $`z`$ forms a hyperplane $`H_z`$ in $`H^0(M,E)`$ and one can define a map $`\iota _E:M𝐏(H^0(M,E))^{}`$ by $`zH_z.`$ Here $`𝐏(H^0(M,E))^{}`$ denotes the dual projective space of linear functionals on $`H^0(M,E)`$ modulo scalar multiplication.
Now equip $`E`$ with a Hermitian metric $`h`$ and $`M`$ with a volume form, and let $`,`$ denote the induced inner product on $`H^0(M,E)`$. Then choose an orthonormal basis $`\{s_0,\mathrm{},s_m\}`$ of $`H^0(M,E)`$ with respect to $`,.`$ Also, choose a local holomorphic frame $`e_E`$ and write $`s_j=f_je_E.`$ Then the point of $`𝐏^m`$ with homogeneous coordinates $`[f_0(z),\mathrm{},f_m(z)]`$ is independent of $`e_E`$ and defines a map $`\varphi _E:M\text{ }\mathrm{C}𝐏^m`$. The same basis also gives an identification $`𝐏^m𝐏(H^0(M,E))^{}`$ by writing a linear functional in the dual basis $`\{s_0^{},\mathrm{},s_m^{}\}`$ of $`H^0(M,E)^{}.`$ We observe that under this identification, $`\varphi _E\iota _E:`$ for $`H_z=\{_ja_js_j:_ja_jf_j(z)=0\}=\mathrm{ker}(_\mathrm{j}\mathrm{f}_\mathrm{j}(\mathrm{z})\mathrm{s}_\mathrm{j}^{})[\mathrm{f}_0(\mathrm{z}),\mathrm{},\mathrm{f}_\mathrm{m}(\mathrm{z})].`$
Next, recall that the Kähler form $`\omega _{FS}`$ of the Fubini-Study metric $`g_{FS}`$ on $`\text{ }\mathrm{C}P^m`$ is given in homogeneous coordinates $`[w_0,\mathrm{},w_m]`$ by $`\omega _{FS}=\frac{\sqrt{1}}{2\pi }\overline{}\mathrm{log}(_{j=0}^m|w_i|^2).`$ Hence
(10)
$$\varphi _E^{}\omega _{FS}=\frac{\sqrt{1}}{2\pi }\overline{}\mathrm{log}(\underset{j=0}{\overset{m}{}}|f_j|^2).$$
It is easy to see that this form is independent of the choice of orthonormal basis.
### 2.3. The maps $`\varphi _N,\stackrel{~}{\varphi }_N,\mathrm{\Phi }_N,\stackrel{~}{\mathrm{\Phi }}_N`$
Now let us return to our setting. We fix a local holomorphic section $`e_L`$ of $`L`$ over $`UM`$. It induces sections $`e_L^N`$ of $`L^N|_U`$ and we write $`S_i^N(z)=f_i^N(z)e_L^N(z)`$ for a holomorphic function $`f_i^N`$ on $`U`$. By the above we have:
(11)
$$\stackrel{~}{\varphi }_N^{}(\omega _{FS})=\varphi _N^{}(\omega _{FS})=\frac{\sqrt{1}}{2\pi }\overline{}\mathrm{log}(\underset{j=0}{\overset{d_N}{}}|f_j^N|^2).$$
The definition is independent of the choice of $`e_L^N.`$
Since $`S_i^N`$ is a holomorphic section of $`L^N`$, $`\widehat{S}_i^N`$ is an equivariant CR function of level $`N`$, i.e. $`\widehat{S}_i^NH_N^2(X).`$ We will need a series of formulae relating expressions in $`S_i^N`$ to expressions in $`\widehat{S}_i^N.`$ Let us introduce local coordinates $`(z,\theta )`$ on $`X`$ on the domain $`U`$ of the unitary frame $`\frac{e_L}{|e_L|}`$ by $`(z,\theta )(z,r_\theta \frac{e_L}{|e_L|}).`$
###### Proposition 6.
$`S_j^N(z)_{h_N}^2=|\widehat{S}_i^N(x)|^2`$ for any $`x`$ with $`\pi (x)=z.`$
Proof:
By definition,
(12)
$$\widehat{S}_i^N(z,u)=u^N,S_i^N(z)=f_i^N(z)u^N,e_L^N(z)=f_i^N(z)a^{\frac{N}{2}}(z)u,\frac{e_L}{|e_L|}(z)^N.$$
In the above local coordinates, we get $`\widehat{S}_i^N(z,\theta )=f_i^N(z)a(z)^{\frac{N}{2}}e^{iN\theta }.`$ Hence $`|\widehat{S}_i^N(z,\theta )|^2=a(z)^N|f_i^N(z)|^2.`$ This obviously equals $`S_j^N(z)_{h_N}^2`$. ∎
###### Proposition 7.
$`\{\widehat{S}_N^i\}`$ is an orthonormal basis of $`H_N^2(X).`$
Proof:
Let $`dV_g=\omega _g^n`$ be the volume form of $`(M,g)`$. Then we have:
$$\begin{array}{c}S_\alpha ^N,S_\beta ^N:=_Mh_N(S_\alpha ^N,S_\beta ^N)𝑑V_g=_Ma^N(z)f_i^N(z)\overline{f}_j^N(z)𝑑V_g\hfill \\ \\ =_X\widehat{S}_\alpha ^N\widehat{S}_\beta ^N𝑑\mu \hfill \end{array}$$
where $`d\mu =\alpha d\alpha ^n`$. In the latter step we use that $`\alpha d\alpha ^n=d\theta \pi ^{}\omega _g^n.`$ This follows from the fact $`d\alpha =\pi ^{}\omega `$ and that $`\alpha =d\theta +\eta `$ where $`\eta `$ only involves $`dz,d\overline{z}.`$
We further have:
###### Proposition 8.
$`\frac{1}{N}\varphi _N^{}\omega _{FS}=\omega _g+\frac{\sqrt{1}}{2\pi N}\overline{}\mathrm{log}(_{j=0}^{d_N}S_j^N_{h_N}^2)=\frac{\sqrt{1}}{2\pi N}_b\overline{}_b\mathrm{log}(_{j=0}^{d_N}|\widehat{S}_j^N|^2)+\omega _g.`$
Proof:
The first statement follows by writing $`S_j^N(z)_{h_N}^2=a^N(z)|f_j^N(z)|`$ and using that $`\frac{\sqrt{1}}{2\pi N}\overline{}\mathrm{log}a^N=\omega _g.`$ To prove the second statement, we note that $`_{j=0}^{d_N}|\widehat{S}_j^N|^2`$ and $`_b\overline{}_b\mathrm{log}(_{j=0}^{d_N}|\widehat{S}_j^N|^2)`$ are $`S^1`$-invariant and hence may be identified with functions on $`M`$. In the latter case, this uses the fact that the $`S^1`$ action commutes with $`_b,`$ i.e. acts by CR automorphisms. The statement then follows from the general fact that $`\pi _{}_b\pi ^{}f=f`$ for any $`fC^{\mathrm{}}(M),`$ where $`\pi _{}F`$ denotes the function on $`M`$ corresponding to an $`S^1`$-invariant function $`F`$ on $`X`$. ∎
Now let us rewrite these relations in terms of the Szegö projectors $`\mathrm{\Pi }_N`$.
###### Proposition 9.
$`\frac{1}{N}\varphi _N^{}\omega _{FS}=\omega _g+\frac{\sqrt{1}}{2\pi N}_b\overline{}_b\mathrm{log}\mathrm{\Pi }_N(x,x).`$
Proof:
We first observe that
(13)
$$\mathrm{\Pi }_N(x,y)=\underset{i=0}{\overset{d_N}{}}\widehat{S}_i^N(x)\widehat{S}^N(y)$$
or, in local coordinates,
(14)
$$\mathrm{\Pi }_N(z,\theta ,w,\theta ^{})=a(z)^{\frac{N}{2}}a(w)^{\frac{N}{2}}e^{iN(\theta \theta ^{})}\underset{i=0}{\overset{d_N}{}}f_i^N(z)\overline{f}_i^N(w).$$
Hence we have
(15)
$$\begin{array}{cc}(a)\hfill & _{j=0}^{d_N}S_j^N(z)_{h_N}^2=\mathrm{\Pi }_N(z,0,z,0)\hfill \\ & \\ (b)\hfill & \frac{\sqrt{1}}{2\pi N}\overline{}\mathrm{log}(_{j=0}^{d_N}S_j^N_{h_N}^2)=\frac{\sqrt{1}}{2\pi N}_b\overline{}_b\mathrm{log}\mathrm{\Pi }_N(z,0,w,0).\hfill \end{array}$$
Together with Proposition 8 this completes the proof. ∎
###### Corollary 10.
The statement of Corollary 2 is equivalent to: $`\overline{}_b_b\mathrm{log}\mathrm{\Pi }_N(x,x)_{C^k}=O(1).`$
## 3. Parametrix for the Cauchy-Szegö kernel
Now we recall the necessary background on the Szegö kernel $`\mathrm{\Pi }(x,y)`$ for a strictly pseudoconvex domain. The following theorem states that it is a Fourier integral operator of complex type, or more precisely a Toeplitz operator in the sense of Boutet de Monvel-Guillemin \[B.G\]. The notation below differs from \[B.S\] in that $`n+1=\mathrm{dim}_{\text{ }\mathrm{C}}\mathrm{D}.`$
###### Theorem 11.
(\[B.S\], Theorem 1.5 and §2.c. Let $`\mathrm{\Pi }(x,y)`$ be the Szegö kernel of the boundary $`X`$ of a bounded strictly pseudo-convex domain $`\mathrm{\Omega }`$ in a complex manifold $`L`$. Then there exists a symbol $`sS^n(X\times X\times \mathrm{I}\mathrm{R}^+)`$ of the type
$$s(x,y,t)\underset{k=0}{\overset{\mathrm{}}{}}t^{nk}s_k(x,y)$$
so that
$$\mathrm{\Pi }(x,y)=_0^{\mathrm{}}e^{it\psi (x,y)}s(x,y,t)𝑑t$$
where the phase $`\psi C^{\mathrm{}}(D\times D)`$ is determined by the following properties:
$``$ $`\psi (x,x)=\frac{1}{i}\rho (x)`$ where $`\rho `$ is the defining function of $`X.`$
$``$ $`d_x^{\prime \prime }\psi `$ and $`d_y^{}\psi `$ vanish to infinite order along the diagonal.
$``$ $`\psi (x,y)=\overline{\psi (y,x)}.`$
More precisely, the phase is determined up to a function which vanishes to infinite order at $`x=y.`$ The integrals are regularized by taking the principal value (see \[B.S\]). The second condition states that $`\psi (x,y)`$ is almost analytic. Roughly speaking, $`\psi `$ is obtained by Taylor expanding $`\rho (z,\overline{z})`$ and replacing all the $`\overline{z}`$’s by $`\overline{w}`$’s. More precisely, the Taylor expansion of $`\psi `$ near the diagonal is given by
$$\psi (x+h,x+k)=\frac{1}{i}\frac{^{\alpha +\beta }\rho }{z^\alpha \overline{z}^\beta }(x)\frac{h^\alpha }{\alpha !}\frac{\overline{k}^\beta }{\beta !}.$$
We note that Theorem 1.5 of \[B.S\] is stated only in the case where $`L=\text{ }\mathrm{C}^n`$. But in \[B.S\]( §2.c, especially (2.17)) it is extended to general complex manifolds.
The simplest example is that of the unit ball in $`\text{ }\mathrm{C}^{n+1}`$ in which case the above formula has the form
$$K_B(z,w)=\frac{1}{(1z,w)^{n+1}}=_0^{\mathrm{}}e^{it\psi _B(z,w)}t^n𝑑t$$
with $`\psi _B(z,w)=1z,w.`$
The above result states that $`\mathrm{\Pi }`$ is a Fourier integral operator with complex phase, in the class $`I_c^0(X\times X,C^+)`$, where $`C^+`$ is the canonical relation $`C^+T^{}X\times T^{}X`$ generated by the phase $`t\psi (x,y)`$ on $`X\times X\times \mathrm{I}\mathrm{R}^+.`$ Its critical points are the solutions of $`\frac{d}{dt}(t\psi )=0`$, i.e. $`\psi =0`$ and on the diagonal of $`X\times X`$ one has
(16)
$$d_x\psi =d_y\psi =\frac{1}{i}d^{}\rho |_X.$$
In particular, the real points of $`C^+`$ consist in the diagonal of $`\mathrm{\Sigma }^+\times \mathrm{\Sigma }^+`$ where $`\mathrm{\Sigma }^+=\{(x,r\alpha ):r>0\}`$ is the cone generated by the contact form $`\alpha =\frac{1}{i}d^{}\rho .`$ In the terminology of \[B.G\], $`\mathrm{\Pi }`$ is a Toeplitz structure on the symplectic cone $`\mathrm{\Sigma }^+.`$
The principal term $`s_0(x,y)`$ was also determined in \[B.S (4.10)\], using that $`\mathrm{\Pi }`$ is a projection. On the diagonal one has
(17)
$$s_0(x,x)d\mu (x)=\frac{1}{4\pi ^n}(\mathrm{det}\mathrm{L}_\mathrm{X})\mathrm{d}\rho \mathrm{dx}$$
where $`L_X=L_\rho |_{T^{}T^{\prime \prime }TX}`$ is the restriction of the Levi form to the maximal complex subspace of $`TX.`$
## 4. Proof of the Theorem 1
The weight space projections $`\mathrm{\Pi }_N`$ are Fourier coefficients of $`\mathrm{\Pi }`$ and hence may be expressed as:
(18)
$$\mathrm{\Pi }_N(x,y)=_0^{\mathrm{}}_{S^1}e^{iN\theta }e^{it\psi (r_\theta x,y)}s(r_\theta x,y,t)𝑑t𝑑\theta $$
where $`r_\theta `$ denotes the $`S^1`$ action on $`X`$. Changing variables $`tNt`$ gives
(19)
$$\mathrm{\Pi }_N(x,y)=N_0^{\mathrm{}}_{S^1}e^{iN(\theta +t\psi (r_\theta x,y)}s(r_\theta x,y,tN)𝑑t𝑑\theta .$$
From the fact that Im$`\psi (x,y)C(d(x,X)+d(y,X)+|xy|^2+O(|xy|^3)`$ (see \[B.S\], Corollary (1.3)) it follows that the phase
(20)
$$\mathrm{\Psi }(t,\theta ;x,y)=t\psi (r_\theta x,y)\theta .$$
has positive imaginary part. Here, $`d(x,X)`$ is the distance from $`x`$ to $`X`$ and $`|xy|`$ is a local Euclidean metric. It follows that the integral is a complex oscillatory integral. Before analysing its asymptotics we simplify the phase. As above, we choose a local holomorphic co-frame $`e_L^{}`$, put $`a(z)=|e_L^{}|_z^2,`$ and write $`\nu L_z^{}`$ as $`\nu =\lambda e_L^{}`$. In the associated coordinates $`(x,y)=(z,\lambda ,w,\mu )`$ on $`X\times X`$ we have:
(21)
$$\rho (z,\lambda )=a(z)|\lambda |^2,\psi (z,\lambda ,w,\mu )=\frac{1}{i}a(z,w)\lambda \overline{\mu }$$
where $`a(z,w)`$ is an almost analytic function on $`M\times M`$ satisfying $`a(z,z)=a(z)`$. On $`X`$ we have $`a(z)|\lambda |^2=1`$ so we may write $`\lambda =a(z)^{\frac{1}{2}}e^{i\varphi }`$. Similarly for $`\mu `$. So for $`(x,y)=(z,\varphi ,w,\varphi ^{})X\times X`$ we have
(22)
$$\psi (z,\varphi ,w,\varphi ^{})=\frac{1}{i}(1\frac{a(z,w)}{i\sqrt{a(z)}\sqrt{a(w)}})e^{i(\varphi \varphi ^{})}.$$
### 4.1. Proof of Theorem 1
On the diagonal $`x=y`$ we have $`\psi (r_\theta x,x)=\frac{1}{i}(1\frac{a(z,z)}{a(z)}e^{i\theta })=\frac{1}{i}(1e^{i\theta }).`$ So
(23)
$$\mathrm{\Psi }(t,\theta ;x,x)=\frac{t}{i}(1e^{i\theta })\theta .$$
We have
(24)
$$\begin{array}{c}d_t\mathrm{\Psi }=\frac{1}{i}(1e^{i\theta })\hfill \\ d_\theta \mathrm{\Psi }=te^{i\theta }1\hfill \end{array}$$
so the critical set is $`𝒞=\{(x,t,\theta ):\theta =0,t=1\}.`$ The Hessian $`\mathrm{\Psi }^{\prime \prime }`$ on the critical set equals
$$\left(\begin{array}{cc}0\hfill & 1\hfill \\ 1\hfill & i\hfill \end{array}\right)$$
so the phase is non-degenerate and the Hessian operator is given by $`L_\mathrm{\Psi }=(\mathrm{\Psi }^{\prime \prime }(1,0)^1D,D=2\frac{^2}{t\theta }i\frac{^2}{t^2}.`$ It follows by the stationary phase method for complex oscillatory intgrals that
(25)
$$\mathrm{\Pi }_N(x,x)N\frac{1}{\sqrt{\mathrm{det}(\mathrm{N}\mathrm{\Psi }^{\prime \prime }(1,0)/2\pi \mathrm{i})}}\underset{j,k=0}{\overset{\mathrm{}}{}}N^{nkj}L_js_k(x,x)$$
where $`L_j`$ is a differential operator of order $`2j`$ defined by
(26)
$$L_js_k(x,x)=\underset{\nu \mu =j}{}\underset{2\nu 3\mu }{}\frac{1}{2^\nu i^j\mu !\nu !}L_\mathrm{\Psi }^\nu [ts_k(r_\theta x,x)g^\mu (t,\theta )]|_{t=1,\theta =0}$$
with $`g(t,\theta )`$ the third order remainder in the Taylor expansion of $`\mathrm{\Psi }`$ at $`(t,\theta )=(1,0).`$ More precisely, for any $`m0`$, one has by \[H I, Theorem 7.7.5\] that
(27)
$$\begin{array}{c}|\mathrm{\Pi }_N(x,x)N\frac{1}{\sqrt{\mathrm{det}(\mathrm{N}\mathrm{\Psi }^{\prime \prime }(1,0)/2\pi \mathrm{i})}}_{j+k<R}^{\mathrm{}}N^{nkj}L_js_k(x,x)|\hfill \\ \\ CN^{nR}_{k<R,|\alpha |2R2k}D^\alpha s_k_{\mathrm{}}.\hfill \end{array}$$
Note that the hypotheses of \[H I, loc.cit.\] are satisfied since the phase has a non-negative imaginary part and since its critical points are real and independent of $`x`$. Note also that the expansion can be differentiated any number of times. After some rearrangement, the series has the form
(28)
$$\mathrm{\Pi }_N(x,x)=N^nC_ns_0(x,x)+N^{n1}a_1(x,x)+\mathrm{}$$
where $`C_n`$ is a universal constant depending only on $`n`$ and where the coefficients $`s_0(x,x),a_1(x,x)\mathrm{}`$ depend only on the jets of the terms $`s_k`$ along the diagonal. From the description above of the leading coefficient $`s_0(x,x)`$ we have (for some other universal constant $`C_n`$):
(29)
$$\mathrm{\Pi }_N(x,x)d\mu (x)=N^nC_n\alpha \omega ^n+O(N^{n1}).$$
Relative to the Riemannian volume measure $`dV_g`$, the coefficient is a (non-zero) constant $`a_0`$ times $`N^n`$. Comparing to the leading term of the Riemann-Roch polynomial gives that $`a_0=1`$, concluding the proof of (a). ∎
### 4.2. Proof Corollary 2
The new element here is that an arbitrary metric $`G`$ on $`M`$, or more precisely its volume form $`dV_G`$, is used to define orthogonality of sections. We may express $`dV_G=J_G\omega ^n`$ for some positive $`J_GC^{\mathrm{}}(M)`$ and then express $`d\mu _G:=d\theta \pi ^{}dV_G=J_G\alpha d\alpha ^n.`$ Let $`\mathrm{\Pi }^G:L^2(X,d\mu _G)H^2(X,d\mu _G)`$ denote the corresponding orthogonal projection and let $`\mathrm{\Pi }_N^G`$ denote the Fourier components. The Boutet de Monvel-Sjöstrand parametrix construction applies to $`\mathrm{\Pi }^G`$ just as well as to $`\mathrm{\Pi }`$, the only difference lying in their symbols. The principal symbol for $`\mathrm{\Pi }^G`$ equals $`d\rho detL_X^G`$ where $`detL_X^G`$ is the determinant relative to $`d\mu _G`$, that is, $`detL_X^G=\frac{\omega ^n}{d\mu _G}=J_G^n.`$ Clearly this is equal to the determinant of $`\omega =Ric(h)`$ relative to $`dV_G.`$ Hence we get
(30)
$$\mathrm{\Pi }_N^G(x,x)=C_nN^ndet_GRic(h)(x)[1+O(\frac{1}{N})].$$
Corollary 2 follows immediately from this and from
(31)
$$|\xi |_{FS,N}^2=\frac{|\xi |^2}{|S_0^N(x)|^2+\mathrm{}+|S_N^N(x)|^2}\mathrm{for}\xi \mathrm{E}_\mathrm{x}^\mathrm{N}.$$
It is equivalent to the statement (cf.\[Bou.1, Theoreme Principal\] \[D, §4\])
$$|S_j^N(x)|^2N^n(2\pi )^n|\alpha _1(x)\mathrm{}\alpha _n(x)|$$
where $`\alpha _j(z)`$ are the eigenvalues of $`ic(L)=Ric(h)`$ relative to $`G`$
### 4.3. Proof of Corollary 3
Because $`\mathrm{\Phi }_N`$ is a CR map, the asymptotics of the derivatives follow immediately from the asymptotics of $`\mathrm{\Pi }_N(x,x)`$. Indeed, $`_b\overline{}_b\mathrm{log}\mathrm{\Pi }_N(x,x)=_b\overline{}_b\mathrm{log}\mathrm{\Pi }_N(x,y)|_{y=x}.`$
By (a) we have
(32)
$$\begin{array}{c}\mathrm{log}\mathrm{\Pi }_N(x,x)=\mathrm{log}(N^ns_0(x,x)[1+N^1\frac{s_1}{s_0}+\mathrm{}])\hfill \\ =n\mathrm{log}N+\mathrm{log}s_0(x,x)+\mathrm{log}[1+N^1\frac{a_1}{s_0}+\mathrm{}])=n\mathrm{log}N+\mathrm{log}s_0(x,x)+O(\frac{1}{N}).\hfill \end{array}$$
By differentiating the expansion we get
(33)
$$_b\overline{}_b\mathrm{log}\mathrm{\Pi }_N(x,x)=_b\overline{}_b\mathrm{log}s_0(x,x)+O(\frac{1}{N})=O(1).$$ |
warning/0002/hep-ph0002058.html | ar5iv | text | # DESY 99-189 ISSN 0418-9833 MPI/PhT/99-59 hep-ph/0002058 December 1999 Virtual Sfermion Effects on Vector-Boson Pair Production at 𝑒⁺𝑒⁻ Colliders
## Abstract
We study the quantum effects on vector-boson pair production in $`e^+e^{}`$ annihilation induced by the sleptons and squarks of the minimal supersymmetric extension of the standard model (MSSM) in the one-loop approximation. We list full analytic results, and quantitatively analyze the resulting deviation from the standard-model prediction of $`e^+e^{}W^+W^{}`$ for the supergravity-inspired MSSM. The latter can be rendered small throughout the whole parameter space by an appropriate choice of renormalization scheme.
PACS numbers: 12.60.Jv, 13.10.+q, 14.80.Ly
The production of $`W`$-boson pairs in $`e^+e^{}`$ annihilation offers a unique opportunity to probe the nonabelian gauge structure of the standard model (SM) at the tree level, which manifests itself in a distinctive cancellation between the $`s`$-, $`t`$-, and $`u`$-channel scattering amplitudes. This process is being studied experimentally with high precision at the CERN Large Electron-Positron Collider (LEP2) in the centre-of-mass (CM) energy range $`2M_W\mathrm{\Gamma }<\sqrt{s}\mathrm{\Gamma }<\mathrm{\hspace{0.17em}205}`$ GeV. At the same time, the related reactions $`e^+e^{}\gamma \gamma ,\gamma Z,ZZ`$ are being measured there, too. With a future $`e^+e^{}`$ linear supercollider, such as JLC, NLC, or TESLA, these measurements can be extended to higher energies, way up to the TeV range, and rendered more precise. On the theoretical side, enormous effort has been invested into the computation of the one-loop radiative corrections to the cross sections of these processes in the SM, both for on- and off-shell vector bosons, and useful low- and high-energy approximations have been elaborated; for a comprehensive review, see Ref. .
Significant deviations of the measured cross sections from their SM predictions could signal physics beyond the SM. Since the gauge couplings of the electron are so tightly constrained by low-energy and LEP1 data, such deviations should mainly originate from the triple gauge-boson couplings (TGC’s). Generic parameterizations of the so-called anomalous TGC’s were introduced and applied to the processes $`e^+e^{}W^+W^{},ZZ`$ in Ref. . In order to explain the physical origin of anomalous TGC’s, it is necessary to consider specific new-physics scenarios. From the theoretical point of view, renormalizable extensions of the SM are most satisfactory. As a rule, the deviations are then induced through loop effects of new particles, which affect not only the TGC’s, but also the vector-boson propagators and the renormalizations of the parameters and wave functions of the tree-level amplitudes. Such deviations were investigated at the one-loop level in Refs. for models with a modified lepton sector including Majorana neutrinos and were found to be generally small. In Ref. , the one-loop radiative corrections to the anomalous parameters $`\mathrm{\Delta }\kappa _V`$ and $`\lambda _V`$ of the $`VW^+W^{}`$ TGC’s, with $`V=\gamma ,Z`$, were studied in the minimal supersymmetric extension of the SM (MSSM), using the pinch technique to render them gauge independent. The sfermion contributions were found to generally dominate the Higgs and gaugino contributions. This may be understood by observing that mass splittings between the up and down components of the sfermion doublets give rise to significant contributions and that the sfermions come in large numbers, due to their multiplicities in flavour and colour.
In this letter, we calculate the sfermion contributions to the cross sections of $`e^+e^{}V_1V_2`$, with $`V_1V_2=\gamma \gamma ,\gamma Z,ZZ,W^+W^{}`$, at one loop in the MSSM. Preliminary results of this study were published in Ref. . Our calculation proceeds along the lines of Ref. , which gives full analytic results. We use the conventions of Ref. and list only those formulas which need to be substituted therein. In a way, our analysis extends Ref. , where the sfermion-induced radiative corrections to the processes $`e^+e^{}Zh^0`$ and $`Z\gamma h^0`$, with $`h^0`$ being the lightest CP-even Higgs boson, were calculated at one loop in the MSSM. The authors of Ref. did not list analytic results that could be compared with ours. In Refs. , the one-loop radiative correction to $`e^+e^{}W^+W^{}`$ in the MSSM was considered for the full supersymmetric particle spectrum, under the simplistic assumption that the mass matrix of each sfermion flavour is proportional to the unit matrix, so that the two weak eigenstates are mass eigenstates with a common mass. We shall compare the sfermion loop correction of Refs. with our result below. As we shall see later, the size of the correction is significantly affected by the sfermion mass splittings.
The Higgs sector of the MSSM is made up by two complex Higgs isodoublet of opposite hypercharge and accommodates five physical Higgs bosons: the neutral CP-even $`h^0`$ and $`H^0`$ bosons, the neutral CP-odd $`A^0`$ boson, and the charged $`H^\pm `$-boson pair. At the tree level, it has two free parameters, which are usually taken to be the mass $`m_A`$ of the $`A^0`$ boson and the ratio $`\mathrm{tan}\beta =v_2/v_1`$ of the vacuum expectation values of the two Higgs doublets. For each of these Higgs bosons and each SM fermion and gauge boson there is a supersymmetric partner. Thus, the spectrum of states is more than doubled if one passes from the SM to the MSSM, which gives rise to a proliferation of parameters and weakens the predictive power of the theory. A canonical method to reduce the number of parameters is to embed the MSSM into a grand unified theory (GUT), e.g., a suitable supergravity (SUGRA) model, in such a way that it is recovered in the low-energy limit. The MSSM thus constrained is described by the following parameters at the GUT scale, which come in addition to $`\mathrm{tan}\beta `$ and $`m_A`$: the universal scalar mass $`m_0`$, the universal gaugino mass $`m_{1/2}`$, the trilinear Higgs-sfermion coupling $`A`$, the bilinear Higgs coupling $`B`$, and the Higgs-higgsino mass parameter $`\mu `$. Notice that $`m_A`$ is then not an independent parameter anymore, but it is fixed through the renormalization group equation. The number of parameters can be further reduced by making additional assumptions. Unification of the tau and bottom Yukawa couplings at the GUT scale leads to a correlation between $`m_t`$ and $`\mathrm{tan}\beta `$. Furthermore, if the electroweak symmetry is broken radiatively, then $`B`$ and $`\mu `$ are determined up to the sign of $`\mu `$. Finally, it turns out that the MSSM parameters are nearly independent of the value of $`A`$, as long as $`|A|\mathrm{\Gamma }<\mathrm{\hspace{0.17em}500}`$ GeV at the GUT scale.
We now present our analytic results. We denote the four-momenta of $`e^+`$, $`e^{}`$, and the two produced vector bosons, $`V_1`$ and $`V_2`$, by $`p_+`$, $`p_{}`$, $`k_1`$, and $`k_2`$, and define the Mandelstam variables as $`s=(p_++p_{})^2`$, $`t=(p_+k_1)^2`$, and $`u=(p_+k_2)^2`$. Neglecting the electron mass, we have $`s+t+u=M_1^2+M_2^2`$, where $`M_1`$ and $`M_2`$ are the masses of $`V_1`$ and $`V_2`$, respectively. In this limit, also the $`s`$-channel contributions due to Higgs-boson exchanges vanish. Because each of the four processes $`e^+e^{}V_1V_2`$ has more than one tree-level diagram, it is convenient to introduce helicity amplitudes $`^\kappa (\lambda _1,\lambda _2,s,t)`$, where $`\kappa `$, $`\lambda _1`$, and $`\lambda _2`$ denote the helicities of $`e^{}`$, $`V_1`$, and $`V_2`$ in the CM frame, respectively. The $`e^+`$ helicity is then $`\kappa `$. The helicity amplitudes $`^\kappa `$ can be decomposed into the standard matrix elements $`_i^\kappa `$ ($`i=0,\mathrm{},9`$) , which are written down in Appendix B of Ref. . In addition to those for $`i=0,\mathrm{},3`$, which already appear at the tree level, we only need $`_9^\kappa `$ for the present analysis.
The tree-level cross sections of $`e^+e^{}V_1V_2`$ are well known and may be found in Eq. (3.1) of Ref. . The sfermion-induced one-loop corrections receive contributions from diagrams containing self-energy corrections, vertex corrections, and counterterm insertions. We work in the Fermi-constant ($`G_F`$) formulation of the electroweak on-shell renormalization scheme, which is explained in the context of Eq. (4.1) in Ref. . Specifically, starting from the results in the pure on-shell renormalization scheme, which uses Sommerfeld’s fine-structure constant $`\alpha `$ and the physical particle masses as basic parameters, we fix $`\alpha =\sqrt{2}G_F\mathrm{sin}^2\theta _wM_W^2/\pi `$, where $`\theta _w`$ is the weak mixing angle, and supplement the radiative corrections with the term $`2\mathrm{\Delta }r`$, where $`\mathrm{\Delta }r`$ contains those radiative corrections to the muon lifetime which the SM or its extensions introduce on top of the purely photonic corrections from within the Fermi model. The sfermion contribution to $`\mathrm{\Delta }r`$ in the MSSM was examined in Ref. . All the formulas listed in Section III and Appendix D of Ref. carry over to the sfermion case, except for Eqs. (3.5) and (3.8), which give the transverse parts of the vector-boson vacuum polarizations $`\mathrm{\Pi }_T^{V_1V_1}`$ and the proper vertex corrections $`\delta _V^\kappa `$, respectively. The relevant Feynman rules for the MSSM sfermion sector are summarized in Appendix A of Ref. . For each fermion flavour $`Q=U,D`$, where $`U=\nu _e,\nu _\mu ,\nu _\tau ,u,c,t`$ and $`D=e,\mu ,\tau ,d,s,b`$, there is a corresponding sfermion flavour, denoted by a tilde. Except for the sneutrinos, which we assume to be left handed, $`\stackrel{~}{Q}`$ comes in two mass eigenstates $`a=1,2`$. The masses $`M_{\stackrel{~}{Q}a}`$ of the sfermions and their trilinear and quartic couplings, $`\stackrel{~}{V}_{Q_aQ_b^{}}^{V_i}`$ and $`\stackrel{~}{U}_{Q_aQ_b^{}}^{V_iV_j}`$, respectively, to the vector bosons $`V_i=\gamma ,Z,W`$ are defined in Appendix A of Ref. . In the absence of flavour-changing neutral currents, we have $`Q=Q^{}`$ in $`\stackrel{~}{V}_{Q_aQ_b^{}}^{V_i}`$ and $`\stackrel{~}{U}_{Q_aQ_b^{}}^{V_iV_j}`$ if $`V_i,V_j=\gamma ,Z`$, which explains the notation $`\stackrel{~}{V}_{Qab}^{V_i}`$ and $`\stackrel{~}{U}_{Qab}^{V_iV_j}`$ used in Ref. . As in Ref. , we neglect the Cabibbo-Kobayashi-Maskawa mixing, so that we may write $`\stackrel{~}{V}_{UaDb}^W`$ and $`\stackrel{~}{U}_{Qab}^{WW}`$. The sfermion contributions to the $`\mathrm{\Pi }_T^{V_1V_2}`$ functions read
$`\mathrm{\Pi }_T^{V_1V_2}(p^2)`$ $`=`$ $`{\displaystyle \frac{1}{48\pi ^2}}{\displaystyle \underset{Q,a,b}{}}N_{\mathrm{col}}^Q\stackrel{~}{V}_{Q_aQ_b}^{V_1}\stackrel{~}{V}_{Q_bQ_a}^{V_2}\{[s2(M_{\stackrel{~}{Q}_a}^2+M_{\stackrel{~}{Q}_b}^2)+{\displaystyle \frac{\left(M_{\stackrel{~}{Q}_a}^2M_{\stackrel{~}{Q}_b}^2\right)^2}{s}}]`$ (1)
$`\times B_0(p,M_{\stackrel{~}{Q}_a},M_{\stackrel{~}{Q}_b})+M_{\stackrel{~}{Q}_a}^2\left(2{\displaystyle \frac{M_{\stackrel{~}{Q}_a}^2M_{\stackrel{~}{Q}_b}^2}{s}}\right)B_0(0,M_{\stackrel{~}{Q}_a},M_{\stackrel{~}{Q}_a})`$
$`+M_{\stackrel{~}{Q}_b}^2(2{\displaystyle \frac{M_{\stackrel{~}{Q}_b}^2M_{\stackrel{~}{Q}_a}^2}{s}})B_0(0,M_{\stackrel{~}{Q}_b},M_{\stackrel{~}{Q}_b})+{\displaystyle \frac{2}{3}}s{\displaystyle \frac{\left(M_{\stackrel{~}{Q}_a}^2M_{\stackrel{~}{Q}_b}^2\right)^2}{s}}\},`$
where $`p`$ is the external four-momentum, $`N_{\mathrm{col}}^Q=1`$ (3) for sleptons (squarks), and the standard two-point scalar function $`B_0`$ is defined in Eq. (C.2) of Ref. . If $`V_1=W^{}`$ and $`V_2=W^+`$, then $`Q_a=U_a`$, $`Q_b=D_b`$, and it is summed over $`(U,D)`$ instead of $`Q`$. The sfermion contribution to $`\delta _V^\kappa `$ is found to be
$`\delta _V^\kappa `$ $`=`$ $`{\displaystyle \frac{1}{2\pi ^2}}{\displaystyle \underset{B=\gamma ,Z}{}}{\displaystyle \frac{g_{eeB}^\kappa }{sM_B^2}}{\displaystyle \underset{Q,a,b,c}{}}N_{\mathrm{col}}^Q\left(\stackrel{~}{V}_{Q_cQ_b}^B\stackrel{~}{V}_{Q_bQ_a}^{V_1}\stackrel{~}{V}_{Q_aQ_c}^{V_2}\stackrel{~}{V}_{Q_cQ_a}^{V_2}\stackrel{~}{V}_{Q_aQ_b}^{V_1}\stackrel{~}{V}_{Q_bQ_c}^B\right)`$ (2)
$`\times \left[_1^\kappa \left(C_2^0+C_3^{01}+C_3^{02}\right)+_2^\kappa C_3^{01}+_3^\kappa C_3^{02}_9^\kappa \left(C_2^{12}+C_3^{12}+C_3^{21}\right)\right],`$
where the $`C`$ functions are the Lorentz coefficients of the standard three-point tensor integrals defined in Eq. (C4) of Ref. . In Eq. (2), we have suppressed their common argument $`(k_1,k_2,M_{\stackrel{~}{Q}_a},M_{\stackrel{~}{Q}_b},M_{\stackrel{~}{Q}_c})`$. Deviating from Ref. , the electron gauge couplings appearing in Eq. (2) are defined as
$$g_{ee\gamma }^\pm =e,g_{eeZ}^+=\frac{e\mathrm{sin}\theta _w}{\mathrm{cos}\theta _w},g_{eeZ}^{}=\frac{e}{\mathrm{cos}\theta _w\mathrm{sin}\theta _w}\left(\frac{1}{2}\mathrm{sin}^2\theta _w\right),$$
(3)
where $`e=\sqrt{4\pi \alpha }`$. We caution the reader that the Feynman rules used in Refs. differ in the sign of $`\mathrm{sin}\theta _w`$. Consequently, we need to multiply the expression for $`\mathrm{\Pi }_T^{\gamma Z}`$ in Eq. (1) with an extra minus sign when we insert it into the relevant formulas, Eqs. (3.6) and (D2), of Ref. . Notice that Eq. (2) includes the contributions from both the direct and crossed triangle diagrams, which are proportional to $`\stackrel{~}{V}_{Q_cQ_b}^B\stackrel{~}{V}_{Q_bQ_a}^{V_1}\stackrel{~}{V}_{Q_aQ_c}^{V_2}`$ and $`\stackrel{~}{V}_{Q_cQ_a}^{V_2}\stackrel{~}{V}_{Q_aQ_b}^{V_1}\stackrel{~}{V}_{Q_bQ_c}^B`$, respectively. If $`V_1=W^{}`$ and $`V_2=W^+`$, then the first term contributes for $`Q_a=U_a`$, $`Q_b=D_b`$, and $`Q_c=D_c`$ and the second one for $`Q_a=D_a`$, $`Q_b=U_b`$, and $`Q_c=U_c`$.
At this point, we should compare our results with those published in Refs. . To that end, we put $`M_{\stackrel{~}{Q}_1}=M_{\stackrel{~}{Q}_2}`$ and nullify the mixing angle relating the weak and mass eigenstates for each sfermion flavour $`\stackrel{~}{Q}`$. Then, our Eq. (1) agrees with Eqs. (C7), (D5), (E9), and (F8) in Ref. , up to an overall minus sign, if we eliminate the factor 1/2 multiplying $`T_{3f}^i`$ in Eq. (D5) and the sum over $`i`$ in Eq. (F8). As for the $`\gamma W^+W^{}`$ vertex correction, our Eq. (2) is in accordance with Eqs. (65) and (B13)–(B18) in Ref. if we replace the first two appearances of $`C_{36}`$ in Eq. (B14) by $`C_{35}`$, substitute $`C_{24}`$ in Eq. (B18) by $`C_{22}`$, and include an overall minus sign in Eqs. (B17) and (B18). As for the $`ZW^+W^{}`$ vertex correction, we find agreement with Eqs. (97) and (C48)–(C53) in Ref. if we alter the overall signs of Eqs. (C52) and (C53).
Now, we explore the phenomenological implications of our results. We concentrate on the case of $`e^+e^{}W^+W^{}`$ because, for $`\sqrt{s}\mathrm{\Gamma }>\mathrm{\hspace{0.17em}180}`$ GeV, it has the largest cross section of the four processes under consideration and it is the only one involving TGC’s at the tree level in the SM. The SM input parameters for our numerical analysis are taken to be $`G_F=1.1663910^5`$ GeV<sup>-2</sup> , $`m_W=80.385`$ GeV, $`m_Z=91.1871`$ GeV, $`m_t=174.3`$ GeV , and $`m_b=4.7`$ GeV. We vary $`\mathrm{tan}\beta `$ and $`m_A`$ in the ranges $`1<\mathrm{tan}\beta <35m_t/m_b`$ and 100 GeV$`<m_A<600`$ TeV, respectively. As for the GUT parameters, we choose $`m_{1/2}=150`$ GeV, $`A=0`$, and $`\mu <0`$, and tune $`m_0`$ so as to be consistent with the desired value of $`m_A`$. All other MSSM parameters are then determined according to the SUGRA-inspired scenario as implemented in the program package SUSPECT . We checked that the results obtained from the program package ISAJET 7.49 , where the electroweak-symmetry-breaking scale is fixed to be $`Q=\sqrt{M_{\stackrel{~}{t}_L}M_{\stackrel{~}{t}_R}}`$, agree with those from SUSPECT within typically 5% or less if the same scale convention is implemented in the latter. In our analysis, we adopt the SUSPECT default value $`Q=M_Z`$. We do not impose the unification of the tau and bottom Yukawa couplings at the GUT scale, which would just constrain the allowed $`\mathrm{tan}\beta `$ range without any visible effect on the results for these values of $`\mathrm{tan}\beta `$. We exclude solutions which do not comply with the present experimental lower mass bounds of the sfermions, charginos, neutralinos, and Higgs bosons .
In Fig. 1, the sfermion-induced correction $`\delta (\theta )`$ in the relationship $`d\sigma /d\mathrm{cos}\theta =`$$`(d\sigma /d\mathrm{cos}\theta )_{\mathrm{Born}}[1+\delta (\theta )]`$ between the one-loop-corrected and tree-level cross sections of $`e^+e^{}W^+W^{}`$ is shown as a function of the scattering angle $`\theta `$, enclosed between the $`e^+`$ and $`W^+`$ three-momenta in the CM frame, for $`\sqrt{s}=200`$, 500 and 1000 GeV assuming $`\mathrm{tan}\beta =10`$ and $`m_A=250`$ GeV. We observe that $`\delta (\theta )`$ has a typical size of order 0.1% or less and can be of either sign. In the backward direction, it strongly depends on the CM energy, while the energy dependence is rather feeble in the forward direction. We emphasize that the smallness of $`\delta (\theta )`$ is a special feature of the $`G_F`$ scheme. In the $`\alpha `$ scheme, the correction is given by $`\delta (\theta )+2\mathrm{\Delta }r`$ and thus shifted to negative values because we have $`\mathrm{\Delta }r0.11\%`$, as indicated in Fig. 1. Next, we study the correction $`\mathrm{\Delta }`$ to the integrated cross section, defined by $`\sigma =\sigma _{\mathrm{Born}}(1+\mathrm{\Delta })`$, for $`\sqrt{s}=200`$ GeV. In Fig. 2, the $`\mathrm{tan}\beta `$ dependences of $`\mathrm{\Delta }`$ and $`\mathrm{\Delta }r`$ are shown for $`m_A=100`$, 250, and 600 GeV, while, in Fig. 3, the $`m_A`$ dependences are shown for $`\mathrm{tan}\beta =3`$, 10, and 30. These dependences are implicit in the sense that our formulas for $`\mathrm{\Delta }`$ and $`\mathrm{\Delta }r`$ do not contain $`\mathrm{tan}\beta `$ or $`m_A`$. In fact, $`\mathrm{\Delta }`$ and $`\mathrm{\Delta }r`$ only depend on $`\mathrm{tan}\beta `$ or $`m_A`$ via the sfermion masses and gauge couplings, the latter being affected through the mixing angles which rotate the weak eigenstates of the sfermion into their mass eigenstates. We note that the SUGRA-inspired MSSM with our choice of input parameters does not permit $`\mathrm{tan}\beta `$ and $`m_A`$ to be simultaneously small, due to the experimental selectron mass lower bound . This explains why the curves for $`m_A=100`$ GeV in Fig. 2 only start at $`\mathrm{tan}\beta 11`$ and those for $`\mathrm{tan}\beta =3`$ in Fig. 3 at $`m_A240`$ GeV. For large $`m_A`$, the experimental $`m_h`$ lower bound enforces $`\mathrm{tan}\beta \mathrm{\Gamma }>\mathrm{\hspace{0.17em}3}`$. On the other hand, the experimental lower bounds on the chargino and neutralino masses induce an upper limit on $`\mathrm{tan}\beta `$, which depends on $`m_A`$. From Fig. 2 we observe that the $`\mathrm{tan}\beta `$ dependence of $`\mathrm{\Delta }r`$ for fixed $`m_A`$ is modest for intermediate values of $`\mathrm{tan}\beta `$, while $`\mathrm{\Delta }r`$ increases in magnitude towards the edges of the allowed $`\mathrm{tan}\beta `$ range. The stau and tau-sneutrino contributions dominate for large $`\mathrm{tan}\beta `$ and small $`m_A`$, while the sbottom and stop contributions dominate for small $`\mathrm{tan}\beta `$ and large $`m_A`$. The contributions due to the sfermions of the first and second generations are insignificant for all values of $`\mathrm{tan}\beta `$ and $`m_A`$. It is interesting to investigate the mixings between the left- and right-handed components of the charged sfermions in the third generation. The mixing is strongest for stop, especially for small $`\mathrm{tan}\beta `$ and small $`m_A`$. For stau and sbottom, the mixings are generally feeble for large $`\mathrm{tan}\beta `$, independently of $`m_A`$. The magnitude of $`\mathrm{\Delta }r`$ may reach several tenths of percent if $`m_A`$ is small to medium and $`\mathrm{tan}\beta `$ is close to its lower or upper limits. For $`\mathrm{tan}\beta 30`$, $`\mathrm{\Delta }r`$ is almost independent of $`m_A`$, while for smaller (larger) values of $`\mathrm{tan}\beta `$, the size of $`\mathrm{\Delta }r`$ monotonically decreases (increases) as $`m_A`$ increases. These features are also nicely illustrated in Fig. 3. We learn from Figs. 2 and 3 that $`\mathrm{\Delta }`$ is insignificant, below 0.02% in size, for all considered values of $`\mathrm{tan}\beta `$ and $`m_A`$. We stress that this happens by virtue of the $`G_F`$ scheme.
In summary, we derived analytic results for the sfermion-induced radiative corrections to the cross sections of $`e^+e^{}\gamma \gamma ,\gamma Z,ZZ,W^+W^{}`$ at one loop in the MSSM and presented a phenomenological discussion for the most interesting case, $`e^+e^{}W^+W^{}`$, adopting a SUGRA-inspired scenario. In the latter case, the correction can essentially be quenched by adopting the $`G_F`$ scheme, which could not be anticipated without explicit calculation. On the other hand, the sfermions are likely to generate the bulk of the MSSM correction to $`e^+e^{}W^+W^{}`$ because of their multiplicities in flavour and colour. This expectation is substantiated by a study of the MSSM corrections to the $`\gamma W^+W^{}`$ and $`ZW^+W^{}`$ TGC’s . We conclude that significant deviations of the measured cross section of $`e^+e^{}W^+W^{}`$ from its SM predictions will not point towards the SUGRA-inspired MSSM.
Note added
After the completion of this work, we received a preprint which reports on the MSSM sfermion corrections to the cross section of $`e^+e^{}W^+W^{}`$ in the modified minimal-subtraction scheme. The analytic results for the vector-boson vacuum polarizations and the proper vertex corrections given in Eqs. (B.1), (B.3)–(B.5), (B.8) and (B.9) of Ref. agree with our Eqs. (1) and (2), respectively. The agreement was also established numerically to very high precision.
Acknowledgements
We are grateful to Ralf Hempfling for his collaboration in earlier stages of this work. We thank Jean-Loic Kneur, Gilbert Moultaka, Joannis Papavassiliou, and Jay Watson for useful comments on Ref. . The work of A.A.B.B. was supported by the Friedrich-Ebert-Stiftung through Grant No. 219747. The II. Institut für Theoretische Physik is supported by the Bundesministerium für Bildung und Forschung under Contract No. 05 HT9GUA 3, and by the European Commission through the Research Training Network Quantum Chromodynamics and the Deep Structure of Elementary Particles under Contract No. ERBFMRXCT980194. |
warning/0002/nucl-th0002011.html | ar5iv | text | # The 𝐹/𝐷 ratio and meson-baryon couplings from QCD sum rules-II
## I INTRODUCTION
The QCD sum rule is often used to determine hadronic parameters from QCD. In this framework, an interpolating field appropriate for the hadron of concern is introduced using quark and gluon fields and used to construct an appropriate correlation function. The correlation function is calculated on the one hand by the operator product expansion (OPE) at the deep Euclidean region of the correlator momentum $`q^2\mathrm{}`$ using QCD degrees of freedom. On the other hand, its phenomenological form is constructed using hadronic degrees of freedom. In most cases of practical QCD sum rule calculations, the phenomenological form is analytic in the complex $`q^2`$ plane except along the positive real axis. Through a dispersion relation, the nonanalytic structure along the positive real axis is matched to the QCD representation of the correlator at $`q^2\mathrm{}`$ and the hadron parameter of concern is extracted in terms of QCD parameters.
The QCD sum rule framework has been widely used to calculate various hadronic properties . Among various applications, determining meson-baryon couplings is of particular interest because meson-baryon couplings are important ingredients for analyzing baryon-baryon interactions. Their values determined from QCD may provide important constraints in constructing baryon-baryon potentials . One main feature of meson-baryon couplings is SU(3) symmetry as it provides a systematic classification of the couplings in terms of the two parameters, the $`\pi NN`$ coupling and the $`F/D`$ ratio. This systematic classification of the couplings is a basis for making realistic potential models for hyperon-baryon interactions . In this approach, however, implementing the SU(3) breaking in the couplings is somewhat limited because the models used rely solely on hadronic degrees of freedom and thus the way of introducing the SU(3) breaking terms in the model may not be unique. Moreover, the baryon-baryon scattering data used in the fitting processes are not precise enough to pick out a specific mesonic channel. Therefore, it would be useful to constrain each model of meson-baryon coupling directly from other non-perturbative methods of QCD, such as QCD sum rules.
Recently, the two-point correlation function of the nucleon interpolating fields with an external pion field
$`\mathrm{\Pi }(q,p)=i{\displaystyle d^4xe^{iqx}0|T[J_N(x)\overline{J}_N(0)]|\pi (p)}.`$ (1)
has been extensively used to calculate the pion-nucleon coupling within the conventional QCD sum rule method . Another approach relying on the three-point function gives results that may contain non-negligible contributions from the higher resonances $`\pi (1300)`$ and $`\pi (1800)`$ . Moreover, using the two-point correlation function, the sum rule can be easily extended to other meson-baryon couplings and the SU(3) limit can be easily taken to identify the $`F/D`$ ratio. Among various developments using Eq. (1), one interesting attempt is to calculate the coupling beyond the chiral limit by considering the Dirac structure $`i\gamma _5`$ at the order $`p^2=m_\pi ^2`$. As a consistent chiral counting, linear terms in the quark mass $`m_q`$ have been included in the OPE side. The pion-nucleon coupling obtained by combining this sum rule with the nucleon chiral-odd sum rule seems to be quite satisfactory.
This sum rule beyond the chiral limit has been recently applied to other meson-baryon couplings such as $`\eta NN`$, $`\pi \mathrm{\Xi }\mathrm{\Xi }`$, $`\eta \mathrm{\Xi }\mathrm{\Xi }`$, $`\pi \mathrm{\Sigma }\mathrm{\Sigma }`$ and $`\eta \mathrm{\Sigma }\mathrm{\Sigma }`$ . The OPE of each sum rule is found to satisfy the SU(3) relations for the couplings proposed in Ref. , which enables us to identify the terms responsible for the $`F/D`$ ratio. This nontrivial observation for the $`F/D`$ ratio, which is a natural consequence of using the SU(3) symmetric interpolating fields, was possible because the sum rules are constructed beyond the chiral limit. The sum rules, if constructed in the soft-meson limit for example, provide the OPE trivially satisfying the SU(3) relations with $`F/D=0`$. Therefore, going beyond the chiral limit especially in the $`i\gamma _5`$ sum rules is important for obtaining nontrivial value of the $`F/D`$ ratio. From this sum rule, the ratio obtained is $`F/D0.2`$ substantially smaller than what it has been known from SU(6) consideration. Furthermore, meson-baryon couplings after taking into account the SU(3) breaking in the OPE as well as in the phenomenological part undergo huge changes from their SU(3) symmetric values. This finding is not consistent with Nijmegen potentials or with common assumptions in studying hypernuclei.
On the other hand, similar sum rule calculations can be performed for the couplings using Dirac structures other than the $`i\gamma _5`$ structure. As discussed in Ref. , the correlation function Eq. (1) contains the other Dirac structures, $`i\gamma _5\widehat{p}`$ ($`\widehat{p}\gamma _\mu p^\mu `$ ) and $`\gamma _5\sigma _{\mu \nu }q^\mu p^\nu `$, which can also be used to construct sum rules for the $`\pi NN`$ coupling beyond the soft-pion limit. Of course, it is straightforward to extend the $`i\gamma _5\widehat{p}`$ and $`\gamma _5\sigma _{\mu \nu }q^\mu p^\nu `$ sum rules and calculate the couplings $`\eta NN`$, $`\pi \mathrm{\Xi }\mathrm{\Xi }`$, $`\eta \mathrm{\Xi }\mathrm{\Xi }`$, $`\pi \mathrm{\Sigma }\mathrm{\Sigma }`$and $`\eta \mathrm{\Sigma }\mathrm{\Sigma }`$. In the SU(3) limit, the OPE should satisfy the SU(3) relations for the couplings, which will allow us to identify the OPE terms responsible for the $`F/D`$ ratio. This is our primary purpose of this work. In particular, we will see if the identifications made from these Dirac structures are consistent with the previous identification made in Ref. , or if the Dirac structure dependence of the sum rule results still persists in these identifications. Once the sum rules are constructed, it will be straightforward to introduce the SU(3) breaking within this framework and see how the SU(3) breaking affects the couplings.
As presented in Ref. , the $`\gamma _5\sigma _{\mu \nu }q^\mu p^\nu `$ Dirac structure has nice features in calculating the $`\pi NN`$ coupling. To be specific, this sum rule provides a coupling independent of the models employed to construct its phenomenological side and the result is rather stable against the variation of the continuum parameter. Indeed, this structure has been further applied to the study of the couplings $`g_{NK\mathrm{\Lambda }}`$ and $`g_{NK\mathrm{\Sigma }}`$ and other pion-baryon couplings . Even though the calculated $`\pi NN`$ coupling is a bit smaller than the empirical value, it may be useful to investigate this sum rule further and improve the previous result. In this work, we will revisit the previous $`\pi NN`$ calculations and improve the OPE calculation in the highest dimensions. We then apply the framework to other meson-baryon couplings in the SU(3) sector to identify the OPE responsible for the $`F/D`$ ratio.
The other Dirac structure $`i\gamma _5\widehat{p}`$, as discussed in Ref. , is found not to be reliable for calculating the $`\pi NN`$ coupling as it contains large contributions from the continuum. The extracted value is highly sensitive to the continuum threshold. Nevertheless, we will study this sum rule again for the completeness first of all. We will improve the OPE calculation by including higher dimensional operators. Then, by extending the sum rule to the other couplings, we will see how the identification for the $`F/D`$ ratio from this Dirac structure is different from the ones obtained from other Dirac structures.
The paper is organized as follows. In Section II, we will revisit the $`i\gamma _5\widehat{p}`$ and $`\gamma _5\sigma _{\mu \nu }q^\mu p^\nu `$ sum rules for $`\pi NN`$ and refine the OPE calculation. We then extend this framework to construct $`\eta NN`$ sum rules in Section III and identify the OPE responsible for the $`F/D`$ ratio. In Section IV, we construct sum rules for the couplings, $`\pi \mathrm{\Xi }\mathrm{\Xi }`$, $`\eta \mathrm{\Xi }\mathrm{\Xi }`$, $`\pi \mathrm{\Sigma }\mathrm{\Sigma }`$ and $`\eta \mathrm{\Sigma }\mathrm{\Sigma }`$. We confirm that the OPE of each sum rule satisfies the SU(3) relations with the identification of the $`F/D`$ ratio made in Section III. In Section V, we present our analysis in the SU(3) symmetric limit. The analysis beyond the SU(3) limit is given in Section VI but only for the $`\gamma _5\sigma _{\mu \nu }q^\mu p^\nu `$ sum rules.
## II Pion-nucleon coupling determined beyond the soft-pion limit
In this section, we construct QCD sum rules for the $`\pi ^0pp`$ coupling at the first order of the pion momentum $`𝒪(p_\mu )`$, to predict the coupling beyond the soft-pion limit. To do this, we use the two-point correlation function with a pion,
$`\mathrm{\Pi }(q,p)=i{\displaystyle d^4xe^{iqx}0|T[J_p(x)\overline{J}_p(0)]|\pi ^0(p)}{\displaystyle d^4xe^{iqx}\mathrm{\Pi }(x,p)}.`$ (2)
Here $`J_p`$ is the proton interpolating field suggested by Ioffe ,
$`J_p=ϵ_{abc}[u_a^TC\gamma _\mu u_b]\gamma _5\gamma ^\mu d_c,`$ (3)
where $`a,b,c`$ are color indices, $`T`$ denotes the transpose with respect to the Dirac indices, $`C`$ the charge conjugation. Using this correlator, we construct the sum rules beyond the soft-pion limit by considering the Dirac structures, $`i\gamma _5\widehat{p}`$ ($`\widehat{p}\gamma _\mu p^\mu `$) and $`\gamma _5\sigma _{\mu \nu }q^\mu p^\nu `$. After taking out the Dirac structures containing one power of the pion momentum, we take the limit $`p_\mu 0`$ in the rest of the correlator. These sum rules have been considered in Ref. but we revisit them here to improve the OPE calculations. By considering the sum rules for $`\pi ^0pp`$ instead of $`\pi ^+pn`$, we can easily extend the formalism to the other diagonal meson-baryon couplings, $`\eta NN`$, $`\pi \mathrm{\Xi }\mathrm{\Xi }`$, $`\eta \mathrm{\Sigma }\mathrm{\Sigma }`$ and so on.
The QCD side of the correlator Eq.(2) is calculated via the operator product expansion (OPE) at the deep spacelike region $`q^2\mathrm{}`$. In our calculations, we use the vacuum saturation hypothesis to factor out quark-antiquark component with a pion (denoted by $`D_{ab}^q`$ below) from the correlator. The rest of the correlator is basically time-ordered products of quark fields which are normally evaluated by background-field techniques . Accordingly, it is straightforward to write the correlator in the coordinate space,
$`\mathrm{\Pi }(x,p)=iϵ_{abc}ϵ_{a^{}b^{}c^{}}\{`$ $`\gamma _5\gamma ^\mu D_{cc^{}}^d\gamma ^\nu \gamma _5\mathrm{Tr}\left[iS_{aa^{}}(x)(\gamma _\nu C)^TiS_{bb^{}}^T(x)(C\gamma _\mu )^T\right]`$ (4)
$``$ $`\gamma _5\gamma ^\mu D_{cc^{}}^d\gamma ^\nu \gamma _5\mathrm{Tr}\left[iS_{ab^{}}(x)\gamma _\nu CiS_{ba^{}}^T(x)(C\gamma _\mu )^T\right]`$ (5)
$``$ $`\gamma _5\gamma ^\mu iS_{cc^{}}(x)\gamma ^\nu \gamma _5\mathrm{Tr}\left[iS_{ab^{}}(x)\gamma _\nu C(D_{ba^{}}^u)^T(C\gamma _\mu )^T\right]`$ (6)
$`+`$ $`\gamma _5\gamma ^\mu iS_{cc^{}}(x)\gamma ^\nu \gamma _5\mathrm{Tr}\left[iS_{aa^{}}(x)(\gamma _\nu C)^T(D_{bb^{}}^u)^T(C\gamma _\mu )^T\right]`$ (7)
$`+`$ $`\gamma _5\gamma ^\mu iS_{cc^{}}(x)\gamma ^\nu \gamma _5\mathrm{Tr}\left[D_{aa^{}}^u(\gamma _\nu C)^TiS_{bb^{}}^T(x)(C\gamma _\mu )^T\right]`$ (8)
$``$ $`\gamma _5\gamma ^\mu iS_{cc^{}}(x)\gamma ^\nu \gamma _5\mathrm{Tr}\left[D_{ab^{}}^u\gamma _\nu CiS_{ba^{}}^T(x)(C\gamma _\mu )^T\right]\}.`$ (9)
The quark propagators $`iS(x)`$ See Ref. for a detailed expression for the quark propagator. inside the traces are the u-quark propagators and the ones outside of the traces are the d-quark propagators. For the time being, we postpone discussions on the gluonic contributions which are obtained by moving a gluon tensor from a quark propagator into the quark-antiquark component with a pion (constituting thus the three-particle pion wave functions according to the nomenclature commonly used in the light-cone QCD sum rules ). Then, it is possible to write the quark-antiquark component with a pion in terms of three matrix elements involving a pion. Namely, for $`q=u\mathrm{or}d`$-quark, we have
$`(D_{aa^{}}^q)^{\alpha \beta }`$ $``$ $`0|q_a^\alpha (x)\overline{q}_a^{}^\beta (0)|\pi ^0(p)`$ (10)
$`=`$ $`{\displaystyle \frac{\delta _{aa^{}}}{12}}(\gamma ^\mu \gamma _5)^{\alpha \beta }0|\overline{q}(0)\gamma _\mu \gamma _5q(x)|\pi ^0(p)+{\displaystyle \frac{\delta _{aa^{}}}{12}}(i\gamma _5)^{\alpha \beta }0|\overline{q}(0)i\gamma _5q(x)|\pi ^0(p)`$ (12)
$`{\displaystyle \frac{\delta _{aa^{}}}{24}}(\gamma _5\sigma ^{\mu \nu })^{\alpha \beta }0|\overline{q}(0)\gamma _5\sigma _{\mu \nu }q(x)|\pi ^0(p).`$
The pseudoscalar pion matrix element $`0|\overline{q}(0)i\gamma _5q(x)|\pi ^0(p)`$ contributes only to the $`i\gamma _5`$ structure of the correlator, not participating in the sum rules of the Dirac structures $`i\gamma _5\widehat{p}`$ and $`\gamma _5\sigma _{\mu \nu }q^\mu p^\nu `$. The rest two components differ by their chirality and both contribute to the $`i\gamma _5\widehat{p}`$ and $`\gamma _5\sigma _{\mu \nu }q^\mu p^\nu `$ sum rules.
At the first order in $`p_\mu `$, the pseudovector matrix element for the $`u`$-quark up to twist-4 can be calculated as provided in Appendix A, namely,
$`A_\mu ^u(\pi ^0)0|\overline{u}(0)\gamma _\mu \gamma _5u(x)|\pi ^0(p)if_\pi p_\mu i{\displaystyle \frac{5}{18}}f_\pi \delta ^2\left({\displaystyle \frac{1}{2}}x^2p_\mu {\displaystyle \frac{1}{5}}x_\mu xp\right),`$ (13)
where $`f_\pi =93`$ MeV and the twist-4 parameter $`\delta ^2=0.2`$ GeV<sup>2</sup> according to Ref. . The $`d`$-quark matrix element has the opposite sign from the $`u`$-quark element,
$`A_\mu ^d(\pi ^0)0|\overline{d}(0)\gamma _\mu \gamma _5d(x)|\pi ^0(p)if_\pi p_\mu +i{\displaystyle \frac{5}{18}}f_\pi \delta ^2\left({\displaystyle \frac{1}{2}}x^2p_\mu {\displaystyle \frac{1}{5}}x_\mu xp\right),`$ (14)
because the pion is an isovector particle. In the local limit (that is, $`x_\mu =0`$), this equation becomes just the PCAC relation. By taking the divergence of the local operator and applying the soft-pion theorem to the resulting matrix element, one can easily derive the well-known Gell-Mann$``$Oakes$``$Renner relation,
$`2m_q\overline{q}q=m_\pi ^2f_\pi ^2.`$ (15)
This implies that the $`m_q`$ terms in the OPE are not the same order as the first order in $`p_\mu `$. In other words, the $`m_q`$ terms should not be included in the OPE when the sum rules are constructed at the first order in $`p_\mu `$. This tricky point does not matter in the $`\pi pp`$ or $`\eta pp`$ sum rules because the $`u`$ or $`d`$ quark masses are small anyway. However, in cases when strange quarks are involved, the quark-mass corrections could be compatible with other OPE terms and therefore this point should be kept in mind: the sum rules at the first order in $`p_\mu `$ should not include quark-mass terms in the OPE.
The other pion matrix element contributing to our sum rules is the pseudotensor type $`B_{\mu \nu }^q(\pi ^0)0|\overline{q}(0)\gamma _5\sigma _{\mu \nu }q(x)|\pi ^0(p)`$ with $`q=u,d`$, which up to leading order in $`p_\mu `$ can be written,
$`B_{\mu \nu }^u(\pi ^0)0|\overline{u}(0)\gamma _5\sigma _{\mu \nu }u(x)|\pi ^0(p)`$ $``$ $`i(p_\mu x_\nu p_\nu x_\mu ){\displaystyle \frac{\overline{u}u}{6f_\pi }},`$ (16)
$`B_{\mu \nu }^d(\pi ^0)0|\overline{d}(0)\gamma _5\sigma _{\mu \nu }d(x)|\pi ^0(p)`$ $``$ $`+i(p_\mu x_\nu p_\nu x_\mu ){\displaystyle \frac{\overline{d}d}{6f_\pi }}.`$ (17)
This will be derived in Appendix B. In the sum rules of the Dirac structures $`i\gamma _5\widehat{p}`$ and $`\gamma _5\sigma _{\mu \nu }q^\mu p^\nu `$, the matrix elements $`A_\mu ^q`$ and $`B_{\mu \nu }^q`$ each multiplied by the corresponding Dirac matrix according to Eq. (12) contribute.
The OPE diagrams that we are considering up to dimension 7 are given in figure 1. Each blob in figures 1 \[except the figures (e) and (f)\] denotes either $`A_\mu ^q`$ or $`B_{\mu \nu }^q`$. For example, if we take the term containing $`A_\mu ^q`$ in Eq. (12) for the pion matrix element, figure 1 (a) contribute to the $`i\gamma _5\widehat{p}`$ sum rule. But in the case of figure 1 (b), even if we take $`A_\mu ^q`$ for the blob, because the chirality is flipped by the disconnected quark line \[namely by the $`\overline{q}q`$ condensate\], this diagram contributes to the $`\gamma _5\sigma _{\mu \nu }q^\mu p^\nu `$ sum rule.
Figures 1 (e) (f), where a gluon from a quark propagator interacts with the quark-antiquark element with a pion, are new in this work, not properly considered in our previous calculations . In these cases, the blob denotes the three-particle pion wave functions commonly used in the light-cone QCD sum rules . At order $`p_\mu `$, the only three-particle wave function contributing to our sum rules is
$`0|q(x)_a^\alpha ig_s[\stackrel{~}{G}^A(0)]^{\sigma \rho }\overline{q}(0)_b^\beta |\pi ^0(p)`$ $``$ $`{\displaystyle \frac{1}{16}}t_{ab}^A(\gamma _\theta )^{\alpha \beta }(p^\rho g^{\theta \sigma }p^\sigma g^{\theta \rho })A_G^q`$ (18)
$``$ $`\pm t_{ab}^A(\gamma _\theta )^{\alpha \beta }{\displaystyle \frac{f_\pi \delta ^2}{48}}(p^\rho g^{\theta \sigma }p^\sigma g^{\theta \rho }),`$ (19)
where $`\stackrel{~}{G}_{\alpha \beta }^A`$ is the dual of the gluon strength tensor, $`\stackrel{~}{G}_{\alpha \beta }^A=\frac{1}{2}ϵ_{\alpha \beta \sigma \rho }(G^A)^{\sigma \rho }`$. The plus sign is for the $`d`$-quark and the minus sign is for the $`u`$-quark. For its derivation, see Appendix C.
Other diagrams contributing to our sum rules but not explicitly shown in figures 1 are when a gluonic tensor coming from the disconnected quark line combines with the quark condensate to form the quark-gluon mixed condensate. These are basically obtained from figure 1 (b) by expanding the disconnected quark line in $`x_\mu `$. We have not drawn those diagrams since they are basically the same kind as figure 1 (c). Similarly, the OPE coming from the coordinate expansion of the quark-antiquark component with a pion, namely the diagram obtained by expanding the blob in figure 1 (a) or (b), is the same kind as figure 1 (e) or (f) and therefore is not explicitly shown.
We now collect the OPE contributing to the $`i\gamma _5\widehat{p}`$ sum rule first in the coordinate space and take the Fourier transformation afterward. Figure 1 (a) contributes to the $`i\gamma _5\widehat{p}`$ sum rule when the pion matrix element $`A_\mu ^q`$ is taken. It is
$`{\displaystyle \frac{4i}{\pi ^4}}\gamma _5\left[{\displaystyle \frac{xA^d}{x^8}}\widehat{x}+{\displaystyle \frac{\widehat{A}^u}{x^6}}{\displaystyle \frac{xA^u}{x^8}}\widehat{x}\right]`$ $`=`$ $`{\displaystyle \frac{4i}{\pi ^4}}\gamma _5{\displaystyle \frac{\widehat{A}^u}{x^6}}`$ (20)
$`\stackrel{\mathrm{F}.\mathrm{T}.}{}`$ $`i\gamma _5\widehat{p}{\displaystyle \frac{f_\pi }{2\pi ^2}}\left[q^2ln(q^2)+\delta ^2ln(q^2)\right],`$ (21)
where we have used the isospin relation $`A_\mu ^u(\pi ^0)=A_\mu ^d(\pi ^0)`$ in the first step and then Eq. (13) afterward before the Fourier transformation is taken. In the first step, the $`d`$-quark contribution was canceled by the corresponding $`u`$-quark contribution. It should be noted however that, in the case of the $`\eta NN`$ coupling, such a cancellation does not occur as $`\eta `$ is an isoscalar. The contribution from figure 1 (d) is similarly obtained,
$`{\displaystyle \frac{\alpha _s}{\pi }}𝒢^2`$ $`{\displaystyle \frac{i\gamma _5}{144\pi ^2}}\left[A_\alpha ^d{\displaystyle \frac{1}{x^4}}(x^2\gamma ^\alpha x^\alpha \widehat{x})A_\alpha ^u{\displaystyle \frac{1}{x^4}}(4x^2\gamma ^\alpha +x^\alpha \widehat{x})\right]`$ (22)
$`\stackrel{\mathrm{F}.\mathrm{T}.}{}`$ $`i\gamma _5\widehat{p}{\displaystyle \frac{f_\pi }{12}}{\displaystyle \frac{\alpha _s}{\pi }}𝒢^2\left[{\displaystyle \frac{1}{q^2}}{\displaystyle \frac{\delta ^2}{9q^4}}\right].`$ (23)
As in Eq. (21), the isospin relation $`A_\mu ^u(\pi ^0)=A_\mu ^d(\pi ^0)`$ has been used. Other diagrams contributing to the $`i\gamma _5\widehat{p}`$ sum rule are calculated straightforwardly,
$`\mathrm{Fig}.\text{1}(\mathrm{f})`$ $`i\gamma _5\widehat{p}{\displaystyle \frac{3}{4\pi ^2}}ln(q^2)[A_G^d+A_G^u]=0,`$ (24)
$`\mathrm{Fig}.\text{1}(\mathrm{b})`$ $`i\gamma _5\widehat{p}{\displaystyle \frac{2}{9f_\pi }}\overline{u}u^2{\displaystyle \frac{1}{q^2}},`$ (25)
$`\mathrm{Fig}.\text{1}(\mathrm{c})`$ $`i\gamma _5\widehat{p}{\displaystyle \frac{m_0^2\overline{u}u^2}{36f_\pi }}{\displaystyle \frac{1}{q^4}}.`$ (26)
In obtaining Eqs. (25) and (26), the pseudotensor pion matrix element $`B_{\mu \nu }^q`$ has been used for the blobs in the figures. The chirality change due to taking $`B_{\mu \nu }^q`$ is compensated by the disconnected quark line so that total chirality is the same as Eq. (23). Note that the quark-gluon mixed condensate in the last equation is parametrized as $`\overline{u}g_s\sigma 𝒢um_0^2\overline{u}u`$ with $`m_0^2=0.8`$ GeV<sup>2</sup> .
The phenomenological side for the $`i\gamma _5\widehat{p}`$ sum rule takes the form,
$`{\displaystyle \frac{g_{\pi N}\lambda _N^2m_N}{(q^2m_N^2)^2}}+\mathrm{}.`$ (27)
The ellipsis includes the continuum whose spectral density is parametrized by a step function, and the nucleon single-pole term. The single-pole term comes from $`NN^{}`$ transitions as well as the PS and PV scheme-dependent $`NN`$ . As the pion momentum is taken out along with its Dirac structure $`i\gamma _5\widehat{p}`$, the rest correlator contains only the correlator momentum $`q^2`$. We thus employ a single-variable dispersion relation in matching the two sides of the sum rule.
In QCD sum rules with external fields, a double dispersion relation is sometime used. (See for example Ref. .) But as suggested in Ref. , sum rules invoking the double dispersion relation may contain spurious terms. As discussed in Ref , the spurious terms, at least in a few examples considered in Ref , give rise to unphysical poles in the spectral density located at the continuum threshold. Indeed, it was demonstrated in Ref. that if sum rules are constructed for each power of the external momentum as we did in this work, the sum rules coming from the double dispersion relation are identical to that coming from the single dispersion relation, provided the spurious terms are eliminated. Thus, even though there are on-going discussions , we believe that the single-variable dispersion relation leads to the correct sum rules constructed at the order $`p_\mu `$.
Taking the Borel transformation with respect to the correlator momentum $`q^2`$, we obtain the $`i\gamma _5\widehat{p}`$ sum rule,
$`g_{\pi N}m_N\lambda _N^2(1+C_{\pi N}M^2)e^{m_N^2/M^2}={\displaystyle \frac{f_\pi }{2\pi ^2}}\left[M^6E_1(x)+\delta ^2M^4E_0(x)\right]`$ (28)
$`+M^2\left[{\displaystyle \frac{f_\pi }{12}}{\displaystyle \frac{\alpha _s}{\pi }}𝒢^2+{\displaystyle \frac{2\overline{q}q^2}{9f_\pi }}\right]{\displaystyle \frac{\delta ^2f_\pi }{108}}{\displaystyle \frac{\alpha _s}{\pi }}𝒢^2+{\displaystyle \frac{m_0^2\overline{q}q^2}{36f_\pi }}.`$ (29)
As the isospin symmetry is always imposed, $`\overline{q}q`$ denotes either $`\overline{u}u`$ or $`\overline{d}d`$. The contribution from the nucleon single pole term whose residue is not known has been denoted by $`C_{\pi N}`$. The continuum contribution is denoted by the factor, $`E_n(xS_0/M^2)=1(1+x+\mathrm{}+x^n/n!)e^x`$ where $`S_0`$ is the continuum threshold. Comparing the corresponding sum rule in the previous calculation , we note that the last two terms are new in this work. They however belong to the highest dimension and their magnitudes are suppressed by the large numerical factors involved.
We now construct the $`\gamma _5\sigma ^{\mu \nu }q_\mu p_\nu `$ sum rule. This differs from the $`i\gamma _5\widehat{p}`$ sum rule by its chirality. In order to sort out the OPE diagrams contributing to this sum rule, the total chirality should be counted carefully. For example, figure 1 (a) contributes to the sum rule when the pseudotensor matrix element $`B_{\mu \nu }^q`$ is taken for the pion matrix element. Since the other quark propagators do not change the chirality, the total chirality is given by the term containing the matrix element $`B_{\mu \nu }^q`$. On the other hand, figure 1 (b) contributes to the sum rule when the pseudovector element $`A_\mu ^q`$ is taken for the pion matrix element because the chirality change due to this choice is compensated by the disconnected quark line. Collecting terms contributing to the $`\gamma _5\sigma ^{\mu \nu }q_\mu p_\nu `$ structure, we have
$`\mathrm{Fig}.1(\mathrm{a})`$ $`\stackrel{\mathrm{F}.\mathrm{T}.}{}`$ $`{\displaystyle \frac{\overline{d}d}{12\pi ^2f_\pi }}\gamma _5\sigma ^{\mu \nu }q_\mu p_\nu ln(q^2),`$ (30)
$`\mathrm{Fig}.1(\mathrm{b})`$ $``$ $`{\displaystyle \frac{2i}{3\pi ^2}}\overline{d}d{\displaystyle \frac{A_\alpha ^u}{x^4}}x_\beta \gamma _5\sigma ^{\alpha \beta }`$ (31)
$`\stackrel{\mathrm{F}.\mathrm{T}.}{}`$ $`{\displaystyle \frac{4f_\pi }{3}}\overline{d}d\gamma _5\sigma ^{\mu \nu }q_\mu p_\nu \left({\displaystyle \frac{1}{q^2}}{\displaystyle \frac{5\delta ^2}{9q^4}}\right),`$ (32)
$`\mathrm{Fig}.1(\mathrm{c})`$ $``$ $`{\displaystyle \frac{im_0^2}{48\pi ^2}}\overline{d}dA_\alpha ^u\gamma _5\sigma ^{\alpha \beta }{\displaystyle \frac{x_\beta }{x^2}}`$ (33)
$`\stackrel{\mathrm{F}.\mathrm{T}.}{}`$ $`{\displaystyle \frac{f_\pi m_0^2}{6}}\overline{d}d\gamma _5\sigma ^{\mu \nu }q_\mu p_\nu {\displaystyle \frac{1}{q^4}},`$ (34)
$`\mathrm{Fig}.1(\mathrm{d})`$ $`\stackrel{\mathrm{F}.\mathrm{T}.}{}`$ $`{\displaystyle \frac{1}{216f_\pi }}\overline{d}d{\displaystyle \frac{\alpha _s}{\pi }}𝒢^2\gamma _5\sigma ^{\mu \nu }q_\mu p_\nu {\displaystyle \frac{1}{q^4}},`$ (35)
$`\mathrm{Fig}.1(\mathrm{e})`$ $`\stackrel{\mathrm{F}.\mathrm{T}.}{}`$ $`{\displaystyle \frac{2\overline{d}d}{3}}A_G^u\gamma _5\sigma ^{\mu \nu }q_\mu p_\nu {\displaystyle \frac{1}{q^4}}`$ (36)
$``$ $`{\displaystyle \frac{2f_\pi }{9}}\overline{d}d\delta ^2\gamma _5\sigma ^{\mu \nu }q_\mu p_\nu {\displaystyle \frac{1}{q^4}}.`$ (37)
Note in some cases, we have specified the steps before final expressions as they are useful for the extension to other couplings. It should be remarked at this stage that, since we have taken into account the contributions from operators with two gluon lines, we also have to take into account the $`\alpha _s`$ correction in the leading term of the OPE. However, that is a formidable task that will introduce only a small correction to our estimate and we will not consider it here. Rather, here we will concentrate on the effects from the higher twist component of the pion wave function.
The phenomenological side of this sum rule takes the form
$`{\displaystyle \frac{g_{\pi N}\lambda _N^2}{(q^2m_N^2)^2}}\gamma _5\sigma ^{\mu \nu }q_\mu p_\nu +\mathrm{}.`$ (38)
This expression is independent of the pseudoscalar-pseudovector coupling schemes and the ellipsis includes only the transitions $`NN^{}`$ . Matching the two expressions and subsequent Borel transformation lead to
$`g_{\pi N}\lambda _N^2(1+D_{\pi N}M^2)e^{m_N^2/M^2}`$ (39)
$`={\displaystyle \frac{\overline{q}q}{f_\pi }}\left[{\displaystyle \frac{M^4E_0(x)}{12\pi ^2}}+{\displaystyle \frac{4}{3}}f_\pi ^2M^2+{\displaystyle \frac{\alpha _s}{\pi }}𝒢^2{\displaystyle \frac{1}{216}}{\displaystyle \frac{m_0^2f_\pi ^2}{6}}+{\displaystyle \frac{26}{27}}f_\pi ^2\delta ^2\right].`$ (40)
The contribution from $`NN^{}`$ has been denoted by $`D_{\pi N}`$. The last term containing $`\delta ^2`$ is new in this work, coming from figures 1 (b) and (e). This new term cancels the other term containing $`m_0^2`$ and makes the OPE larger. Without this term, we previously reported $`g_{\pi N}10`$. Thus, the new term should make better agreement with the empirical $`\pi NN`$ coupling.
## III $`\eta NN`$ coupling and the $`F/D`$ ratio
Extensions of the $`\pi ^0pp`$ sum rules to the $`\eta NN`$ coupling is straightforward. In this case, we consider the correlation function with an $`\eta `$,
$`i{\displaystyle d^4xe^{iqx}0|T[J_p(x)\overline{J}_p(0)]|\eta (p)}.`$ (41)
In the SU(3) symmetric limit, there is no $`\eta \eta ^{}`$ mixing and, within the OPE dimension that we are considering, the strange quark component of $`\eta `$ does not participate in the sum rule. The difference from the $`\pi ^0pp`$ case is that, since $`\eta `$ is isoscalar, the $`d`$-quark matrix element couples to an $`\eta `$ with the same sign as the $`u`$-quark element does. Thus, the quark-antiquark elements with an $`\eta `$ contributing to our sum rules have the following relations,
$`0|\overline{u}(0)\gamma _\mu \gamma _5u(x)|\eta (p)`$ $`=`$ $`0|\overline{d}(0)\gamma _\mu \gamma _5d(x)|\eta (p)`$ (42)
$``$ $`{\displaystyle \frac{if_\eta }{\sqrt{3}}}\left[p_\mu {\displaystyle \frac{5}{18}}\delta ^2\left({\displaystyle \frac{1}{2}}x^2p_\mu {\displaystyle \frac{1}{5}}x_\mu xp\right)\right],`$ (43)
$`0|\overline{u}(0)\gamma _5\sigma _{\mu \nu }u(x)|\eta (p)`$ $`=`$ $`0|\overline{d}(0)\gamma _5\sigma _{\mu \nu }d(x)|\eta (p)`$ (44)
$``$ $`i(p_\mu x_\nu p_\nu x_\mu ){\displaystyle \frac{\overline{u}u}{6\sqrt{3}f_\eta }},`$ (45)
$`0|u(x)_a^\alpha ig_s(\stackrel{~}{G}^A(0))^{\sigma \rho }\overline{u}(0)_b^\beta |\eta (p)`$ $`=`$ $`0|d(x)_a^\alpha ig_s(\stackrel{~}{G}^A(0))^{\sigma \rho }\overline{d}(0)_b^\beta |\eta (p)`$ (46)
$``$ $`t_{ab}^A(\gamma _\theta )^{\alpha \beta }{\displaystyle \frac{f_\eta \delta ^2}{48\sqrt{3}}}(p^\rho g^{\theta \sigma }p^\sigma g^{\theta \rho }).`$ (47)
Similarly for the $`\pi ^0pp`$ case, these matrix elements are the basic ingredients in calculating the OPE. In this construction, the SU(3) breaking effects are driven only by $`f_\eta f_\pi `$.
Keeping in mind the sign change of the matrix elements involving the $`d`$-quark, it is straightforward to obtain the OPE for the $`\eta pp`$ coupling directly from Eq.(20) - Eq.(26) for the $`i\gamma _5\widehat{p}`$ sum rule, and from Eq.(30) - Eq.(37) for the $`\gamma _5\sigma ^{\mu \nu }q_\mu p_\nu `$ sum rule. For example in the case of figure 1 (a), we have the similar expression as Eq. (20) but now the $`d`$-quark element adds up with the $`u`$-quark element to yield the OPE
$`{\displaystyle \frac{if_\eta }{3\sqrt{3}\pi ^2}}\gamma _5\widehat{p}\left[q^2ln(q^2)+\delta ^2ln(q^2)\right].`$ (48)
By similarly calculating for the other OPE, we obtain the $`i\gamma _5\widehat{p}`$ sum rule,
$`g_{\eta N}m_N\lambda _N^2(1+`$ $`C_{\eta N}`$ $`M^2)e^{m_N^2/M^2}={\displaystyle \frac{f_\eta }{3\sqrt{3}\pi ^2}}[M^6E_1+{\displaystyle \frac{5}{2}}\delta ^2M^4E_0]`$ (49)
$`+`$ $`{\displaystyle \frac{M^2}{\sqrt{3}}}\left[{\displaystyle \frac{f_\eta }{9}}{\displaystyle \frac{\alpha _s}{\pi }}𝒢^2+{\displaystyle \frac{2}{9f_\eta }}\overline{q}q^2\right]{\displaystyle \frac{2f_\eta \delta ^2}{81\sqrt{3}}}{\displaystyle \frac{\alpha _s}{\pi }}𝒢^2+{\displaystyle \frac{m_0^2\overline{q}q^2}{36\sqrt{3}f_\eta }}.`$ (50)
The RHS in the SU(3) limit ($`f_\pi =f_\eta `$) should satisfy the SU(3) relation for the $`\eta NN`$ coupling ,
$`{\displaystyle \frac{g_{\eta N}}{g_{\pi N}}}{\displaystyle \frac{1}{\sqrt{3}}}(4\alpha 1),`$ (51)
where $`\alpha =\frac{F}{F+D}`$. To identify the OPE corresponding to $`\alpha `$, we neglect the nucleon single pole terms $`C_{\pi N}`$ and $`C_{\eta N}`$ for the time being. A full analysis including them will be given in later sections. By taking the ratio of Eqs.(29) and (50) and comparing with Eq.(51), we find a rather complicated expression for $`\alpha `$,
$`4\alpha `$ $``$ $`\{{\displaystyle \frac{5}{6\pi ^2}}M^6E_1(x)+{\displaystyle \frac{4\delta ^2}{3\pi ^2}}M^4E_0(x)+{\displaystyle \frac{\alpha _s}{\pi }}𝒢^2({\displaystyle \frac{7}{36}}M^2{\displaystyle \frac{11}{324}}\delta ^2)+{\displaystyle \frac{4}{9f_\pi ^2}}\overline{q}q^2M^2`$ (53)
$`+{\displaystyle \frac{m_0^2\overline{q}q^2}{18f_\pi ^2}}\}`$
$`\times `$ $`\{{\displaystyle \frac{1}{2\pi ^2}}M^6E_1(x)+{\displaystyle \frac{\delta ^2}{2\pi ^2}}M^4E_0(x)+{\displaystyle \frac{1}{12}}{\displaystyle \frac{\alpha _s}{\pi }}𝒢^2(M^2{\displaystyle \frac{\delta ^2}{9}})+{\displaystyle \frac{2\overline{q}q^2}{9f_\pi ^2}}M^2`$ (55)
$`+{\displaystyle \frac{m_0^2\overline{q}q^2}{36f_\pi ^2}}\}^1.`$
Of course, this equation as is written here should not be used to determine the $`F/D`$ ratio because the important contribution from the unknown nucleon single-pole terms has been neglected in the construction. This equation only provides a consistency check for the OPE as the QCD side of the sum rule should satisfy the SU(3) relations for the couplings.
By performing a similar calculation, we obtain for the $`\gamma _5\sigma ^{\mu \nu }q_\mu p_\nu `$ sum rule,
$`g_{\eta N}\lambda _N^2(1+D_{\eta N}M^2)e^{m_N^2/M^2}`$ (56)
$`={\displaystyle \frac{\overline{q}q}{\sqrt{3}f_\eta }}\left[{\displaystyle \frac{M^4E_0(x)}{12\pi ^2}}{\displaystyle \frac{4}{3}}f_\eta ^2M^2+{\displaystyle \frac{\alpha _s}{\pi }}𝒢^2{\displaystyle \frac{1}{216}}+{\displaystyle \frac{m_0^2f_\eta ^2}{6}}{\displaystyle \frac{26}{27}}f_\eta ^2\delta ^2\right].`$ (57)
Note the first and third terms, as they come from the OPE containing the $`d`$-quark elements with $`\eta `$, have the opposite sign from the corresponding terms in Eq. (40). Because of this, there seems a delicate cancellation between the OPE terms and the overall OPE strength becomes small.
Combining Eqs. (40) and (57) and neglecting again the nucleon single-pole terms, $`D_{\pi N}`$ and $`D_{\eta N}`$, we identify $`\alpha `$ in the SU(3) limit ($`f_\pi =f_\eta `$),
$`4\alpha `$ $``$ $`\left\{{\displaystyle \frac{8f_\pi }{3}}M^2{\displaystyle \frac{f_\pi m_0^2}{3}}+{\displaystyle \frac{52f_\pi \delta ^2}{27}}\right\}`$ (58)
$`\times `$ $`\left\{{\displaystyle \frac{M^4E_0(x)}{12\pi ^2f_\pi }}+{\displaystyle \frac{4f_\pi }{3}}M^2+{\displaystyle \frac{26}{27}}f_\pi \delta ^2{\displaystyle \frac{f_\pi m_0^2}{6}}+{\displaystyle \frac{1}{216f_\pi }}{\displaystyle \frac{\alpha _s}{\pi }}𝒢^2\right\}^1.`$ (59)
The quark condensate is canceled in the ratio. The expression is clearly different from Eq. (55). In our previous work, we obtain from the $`i\gamma _5`$ sum rules beyond the chiral limit ,
$`4\alpha `$ $``$ $`{\displaystyle \frac{4m_qm_0^2\overline{q}q^2}{3f_\pi }}`$ (60)
$`\times `$ $`\{m_\pi ^2M^4E_0(x)[{\displaystyle \frac{\overline{q}q}{12\pi ^2f_\pi }}+{\displaystyle \frac{3f_{3\pi }}{4\sqrt{2}\pi ^2}}]+{\displaystyle \frac{2m_q}{f_\pi }}\overline{q}q^2M^2+{\displaystyle \frac{m_\pi ^2}{72f_\pi }}\overline{q}q{\displaystyle \frac{\alpha _s}{\pi }}𝒢^2`$ (61)
$`+`$ $`{\displaystyle \frac{2m_q}{3f_\pi }}m_0^2\overline{q}q^2\}^1.`$ (62)
Through the Gell-Mann$``$Oakes$``$Renner relation, the $`m_q`$ as well as the $`m_\pi ^2`$ dependence will be canceled in the ratio. This expression is also not consistent with Eq.(55) or Eq.(59). Therefore, depending on Dirac structures, we clearly have different expressions for $`\alpha `$. Later sections, we will discuss the numerical value of $`\alpha `$ when we analyze the sum rules including the unknown single-pole terms.
## IV $`\pi \mathrm{\Xi }\mathrm{\Xi }`$, $`\eta \mathrm{\Xi }\mathrm{\Xi }`$, $`\pi \mathrm{\Sigma }\mathrm{\Sigma }`$ and $`\eta \mathrm{\Sigma }\mathrm{\Sigma }`$
Having identified the OPE corresponding to $`\alpha `$ in the SU(3) limit, we apply, as consistency checks, the sum rule framework introduced above to the couplings, $`\pi \mathrm{\Xi }\mathrm{\Xi }`$, $`\eta \mathrm{\Xi }\mathrm{\Xi }`$, $`\pi \mathrm{\Sigma }\mathrm{\Sigma }`$ and $`\eta \mathrm{\Sigma }\mathrm{\Sigma }`$. In this extension, the strange quark enters to the sum rules as an active degree of freedom. One important constraint to be satisfied always is the SU(3) relations for the couplings . This means that the identification of $`\alpha `$ made in Eq. (55) for the $`i\gamma _5\widehat{p}`$ sum rule and in Eq.(59) for the $`\gamma _5\sigma ^{\mu \nu }q_\mu p_\nu `$ should be separately satisfied whenever SU(3) symmetry is imposed on the OPE. This statement is obvious because the interpolating fields used for $`\mathrm{\Xi }`$ and $`\mathrm{\Sigma }`$ are constructed from the nucleon interpolating field via the SU(3) rotation. However, it is often the case that this SU(3) constraint is not properly imposed in QCD sum rule constructions of meson-baryon couplings in the SU(3) sector.
One important limitation in this extension is related to quark mass corrections. Our sum rules are constructed at the order of $`𝒪(p_\mu )`$ in the expansion of the meson momentum. The quark-mass terms are of higher orders and thus should not be included in this sum rule. However, when the massive $`s`$-quark is involved, it is questionable whether the sum rules truncated at the first order in $`p_\mu `$ is reliable: the sum rule at the order $`p_\mu `$ may not properly represent the total strength of the correlators. In this sense, our sum rules in this extension are somewhat limited and a more systematic procedure which does not rely on the expansion of the meson momentum may be required for realistic prediction for the couplings. Indeed, the light-cone sum rule may be useful for that purpose but in this case, QCD inputs contain the meson wave functions at a specific point instead of their integrated strength as in our case. Therefore, predictions may depend on a specific ‘ansatz’ for the wave functions . Moreover, there is an issue in applying QCD duality in the construction of the light-cone sum rule and the usual application of QCD duality in QCD sum rules with external fields may not be well satisfied . In future a more systematic method to overcome these difficulties is anticipated. Nevertheless, our sum rules at the linear order in $`p_\mu `$ are reliable as long as the SU(3) symmetric limit is imposed, because, in that limit, the $`s`$-quark mass is small. Therefore, the discussion regarding the $`F/D`$ ratio is reasonable. When we discuss the couplings beyond the SU(3) limit, however, this limitation should be noted.
We calculate the meson-baryon couplings for $`\pi \mathrm{\Xi }\mathrm{\Xi }`$, $`\eta \mathrm{\Xi }\mathrm{\Xi }`$, $`\pi \mathrm{\Sigma }\mathrm{\Sigma }`$ and $`\eta \mathrm{\Sigma }\mathrm{\Sigma }`$ from a correlation function of the type,
$`i{\displaystyle d^4xe^{iqx}0|T[J_B(x)\overline{J}_B(0)]|M(p)},`$ (63)
where $`J_B`$ is the corresponding baryon interpolating field and $`|M(p)`$ is the meson state of concern. For $`\mathrm{\Xi }`$ and $`\mathrm{\Sigma }`$, we use the interpolating fields
$`J_\mathrm{\Xi }`$ $`=`$ $`ϵ_{abc}[s_a^TC\gamma _\mu s_b]\gamma _5\gamma ^\mu u_c,`$ (64)
$`J_\mathrm{\Sigma }`$ $`=`$ $`ϵ_{abc}[u_a^TC\gamma _\mu u_b]\gamma _5\gamma ^\mu s_c,`$ (65)
respectively obtained from the nucleon interpolating field via the SU(3) rotations. Calculation can be performed similarly as before. But, since the baryon interpolating fields have the similar structure as the nucleon interpolating field, we can easily obtain the OPE for each sum rule by making simple replacements from the $`\pi NN`$ or $`\eta NN`$ sum rules. To be more specific, the OPE for the $`\pi \mathrm{\Xi }\mathrm{\Xi }`$ coupling is obtained from that of the $`\pi pp`$ coupling by replacing the quark fields, $`us`$ and $`du`$. The same replacements are required to obtain the OPE for the $`\eta \mathrm{\Xi }\mathrm{\Xi }`$ from that of $`\eta pp`$. For the $`\pi \mathrm{\Sigma }\mathrm{\Sigma }`$ and $`\eta \mathrm{\Sigma }\mathrm{\Sigma }`$ couplings, we need to replace $`ds`$ from the corresponding nucleon sum rules. A new ingredient in this extension is the strange quark-antiquark component with the specific meson of concern. The strange quark-antiquark component does not couple to a pion within the OPE dimension that we are considering. However, in the case of the $`\eta `$-baryon couplings, there is nonzero strength between the strange quark-antiquark operators and an $`\eta `$. Its strength relative to the $`u`$ or $`d`$-quark operators can be read off from the SU(3) Gell-Mann matrix. Namely, we have
$`A_\mu ^s(\eta )0|\overline{s}(0)\gamma _\mu \gamma _5s(x)|\eta (p)`$ $``$ $`{\displaystyle \frac{2}{\sqrt{3}}}if_\eta p_\mu +i{\displaystyle \frac{5}{9\sqrt{3}}}f_\eta \delta ^2\left({\displaystyle \frac{1}{2}}x^2p_\mu {\displaystyle \frac{1}{5}}x_\mu xp\right),`$ (66)
$`B_{\mu \nu }^s(\eta )0|\overline{s}(0)\gamma _5\sigma _{\mu \nu }s(x)|\eta (p)`$ $``$ $`+i(p_\mu x_\nu p_\nu x_\mu ){\displaystyle \frac{\overline{s}s}{3\sqrt{3}f_\eta }},`$ (67)
$`0|s(x)_a^\alpha ig_s(\stackrel{~}{G}^A(0))^{\sigma \rho }\overline{s}(0)_b^\beta |\eta (p)`$ $``$ $`+t_{ab}^A(\gamma _\theta )^{\alpha \beta }{\displaystyle \frac{f_\eta \delta ^2}{24\sqrt{3}}}(p^\rho g^{\theta \sigma }p^\sigma g^{\theta \rho }).`$ (68)
Compared with Eqs.(14), (17),(19), the strange quark-antiquark elements with an $`\eta `$ have the overall sign consistent with the d-quark components with a pion but the magnitude has been multiplied by the factor $`2/\sqrt{3}`$ as it should be. In Eqs. (66) and (68), the SU(3) breaking is reflected only in $`f_\eta `$ but in Eq. (67), there is another SU(3) breaking source, the strange quark condensate. As we know how to go to the SU(3) symmetric limit from these two breaking sources, the SU(3) relations for the couplings can be easily investigated in this approach.
With these differences in mind, we can straightforwardly calculate the OPE for each coupling. In the case of the $`i\gamma _5\widehat{p}`$ Dirac structure, we obtain the $`\pi \mathrm{\Xi }\mathrm{\Xi }`$ sum rule,
$`g_{\pi \mathrm{\Xi }}m_\mathrm{\Xi }\lambda _\mathrm{\Xi }^2(1+C_{\pi \mathrm{\Xi }}M^2)e^{m_\mathrm{\Xi }^2/M^2}`$ $`=`$ $`{\displaystyle \frac{f_\pi }{12\pi ^2}}\left[M^6E_12\delta ^2M^4E_0\right]`$ (69)
$`+`$ $`{\displaystyle \frac{f_\pi }{72}}{\displaystyle \frac{\alpha _s}{\pi }}𝒢^2\left(M^2{\displaystyle \frac{5\delta ^2}{9}}\right).`$ (70)
Again, neglecting the unknown single pole term $`C_{\pi \mathrm{\Xi }}`$, it is easy to see that in the SU(3) limit the RHS satisfies the SU(3) relation,
$`{\displaystyle \frac{g_{\pi \mathrm{\Xi }}}{g_{\pi N}}}2\alpha 1,`$ (71)
if we identify $`\alpha `$ as Eq. (55). This suggests that our approach makes sense at least in retrieving consistently the SU(3) relation for the coupling. For the other couplings, we similarly proceed the calculations. The LHS side of each sum rule has the similar structure as that of Eq.(70) but now we have different baryon mass, the coupling, and the strength of the interpolating field to the baryon of concern. The RHS of each sum rule for the $`i\gamma _5\widehat{p}`$ structure is obtained as follows
$`\eta \mathrm{\Xi }\mathrm{\Xi }`$ $`:{\displaystyle \frac{1}{\sqrt{3}}}\{{\displaystyle \frac{11f_\eta }{12\pi ^2}}M^6E_1+{\displaystyle \frac{7f_\eta \delta ^2}{6\pi ^2}}M^4E_0+{\displaystyle \frac{f_\eta }{72}}{\displaystyle \frac{\alpha _s}{\pi }}𝒢^2(13M^2{\displaystyle \frac{17\delta ^2}{9}})`$ (72)
$`+`$ $`{\displaystyle \frac{4\overline{s}s^2}{9f_\eta }}M^2+{\displaystyle \frac{m_0^2\overline{s}s^2}{18f_\eta }}\},`$ (73)
$`\pi \mathrm{\Sigma }\mathrm{\Sigma }`$ $`:{\displaystyle \frac{f_\pi }{12\pi ^2}}\left[5M^6E_1+8\delta ^2M^4E_0\right]+{\displaystyle \frac{f_\pi }{72}}{\displaystyle \frac{\alpha _s}{\pi }}𝒢^2\left(7M^2{\displaystyle \frac{11\delta ^2}{9}}\right)`$ (74)
$`+`$ $`{\displaystyle \frac{2\overline{q}q^2}{9f_\pi }}M^2+{\displaystyle \frac{m_0^2\overline{q}q^2}{36f_\pi }},`$ (75)
$`\eta \mathrm{\Sigma }\mathrm{\Sigma }`$ $`:{\displaystyle \frac{1}{\sqrt{3}}}\{{\displaystyle \frac{7f_\eta }{12\pi ^2}}M^6E_1+{\displaystyle \frac{f_\eta \delta ^2}{3\pi ^2}}M^4E_0+{\displaystyle \frac{f_\eta }{72}}{\displaystyle \frac{\alpha _s}{\pi }}𝒢^2(5M^2{\displaystyle \frac{\delta ^2}{9}})`$ (76)
$`+`$ $`{\displaystyle \frac{2\overline{q}q^2}{9f_\eta }}M^2+{\displaystyle \frac{m_0^2\overline{q}q^2}{36f_\eta }}\}.`$ (77)
Again it is straightforward to show that in the SU(3) limit ($`f_\pi =f_\eta `$, $`\overline{s}s=\overline{q}q`$ and $`\lambda _N=\lambda _\mathrm{\Xi }=\lambda _\mathrm{\Sigma }`$) the RHS of each sum rule satisfies the SU(3) relations for the couplings,
$`{\displaystyle \frac{g_{\eta \mathrm{\Xi }}}{g_{\pi N}}}{\displaystyle \frac{1}{\sqrt{3}}}(1+2\alpha );{\displaystyle \frac{g_{\pi \mathrm{\Sigma }}}{g_{\pi N}}}2\alpha ;{\displaystyle \frac{g_{\eta \mathrm{\Sigma }}}{g_{\pi N}}}{\displaystyle \frac{2}{\sqrt{3}}}(1\alpha ),`$ (78)
with $`\alpha `$ given in Eq. (55). Therefore the consistency check has been made for these sum rules.
In the case of the $`\gamma _5\sigma ^{\mu \nu }q_\mu p_\nu `$ Dirac structure, the LHS side of the sum rule for meson($`m_i=\pi \mathrm{or}\eta `$)-baryon ($`B_j=\mathrm{\Xi }\mathrm{or}\mathrm{\Sigma }`$) coupling takes the form
$`g_{m_iB_j}\lambda _{B_j}^2(1+D_{m_iB_j}M^2)e^{m_{B_j}^2/M^2}.`$ (79)
In the RHS, we have the following set of the sum rules,
$`\pi \mathrm{\Xi }\mathrm{\Xi }:`$ $`{\displaystyle \frac{\overline{q}q}{12\pi ^2f_\pi }}M^4E_0+{\displaystyle \frac{\overline{q}q}{216f_\pi }}{\displaystyle \frac{\alpha _s}{\pi }}𝒢^2,`$ (80)
$`\eta \mathrm{\Xi }\mathrm{\Xi }:`$ $`{\displaystyle \frac{\overline{q}q}{\sqrt{3}}}\left\{{\displaystyle \frac{M^4E_0}{12\pi ^2f_\eta }}+{\displaystyle \frac{8f_\eta M^2}{3}}+{\displaystyle \frac{52f_\eta \delta ^2}{27}}{\displaystyle \frac{f_\eta m_0^2}{3}}+{\displaystyle \frac{1}{216f_\eta }}{\displaystyle \frac{\alpha _s}{\pi }}𝒢^2\right\},`$ (81)
$`\pi \mathrm{\Sigma }\mathrm{\Sigma }:`$ $`\overline{s}s\left\{{\displaystyle \frac{4f_\pi }{3}}M^2{\displaystyle \frac{26f_\pi \delta ^2}{27}}+{\displaystyle \frac{f_\pi m_0^2}{6}}\right\},`$ (82)
$`\eta \mathrm{\Sigma }\mathrm{\Sigma }:`$ $`{\displaystyle \frac{\overline{s}s}{\sqrt{3}}}\left\{{\displaystyle \frac{M^4E_0}{6\pi ^2f_\eta }}{\displaystyle \frac{4f_\eta }{3}}M^2{\displaystyle \frac{26f_\eta \delta ^2}{27}}{\displaystyle \frac{1}{108f_\eta }}{\displaystyle \frac{\alpha _s}{\pi }}𝒢^2+{\displaystyle \frac{f_\eta m_0^2}{6}}\right\}.`$ (83)
It can be easily checked that these RHS of the sum rules satisfy the SU(3) relations for the couplings Eqs. (71) and (78) if we identify $`\alpha `$ as given in Eq.(59), again making sure the consistency with SU(3) symmetry.
## V Analysis in the SU(3) symmetric limit
We now analyze the sum rules provided in the previous sections within the SU(3) symmetric limit. What we have demonstrated so far is that we have different identifications for the $`F/D`$ ratio depending on the Dirac structure considered. Each set of sum rules satisfies the SU(3) relations for the couplings with different identifications for the $`F/D`$ ratio. In this section, we calculate numerical values of the $`F/D`$ ratio from the $`i\gamma _5\widehat{p}`$ and $`\gamma _5\sigma _{\mu \nu }q^\mu p^\nu `$ sum rules separately, and see if they are consistent with the previous calculation beyond the chiral limit . In the SU(3) limit, all the baryon masses are the same $`m_N=m_\mathrm{\Xi }=m_\mathrm{\Sigma }`$, and we also have the relations, $`\overline{s}s=\overline{q}q`$ and $`f_\eta =f_\pi `$. Moreover, from the baryon mass sum rules , the strength of each interpolating field to the low-lying baryon should be the same in this limit, $`\lambda _N=\lambda _\mathrm{\Xi }=\lambda _\mathrm{\Sigma }`$. We take the conventional values for the QCD parameters in this analysis,
$`\overline{q}q`$ $`=`$ $`(0.23\mathrm{GeV})^3;{\displaystyle \frac{\alpha _s}{\pi }}𝒢^2=(0.33\mathrm{GeV})^4`$ (84)
$`\delta ^2`$ $`=`$ $`0.2\mathrm{GeV}^2;m_0^2=0.8\mathrm{GeV}^2.`$ (85)
We start with the $`i\gamma _5\widehat{p}`$ sum rules in the SU(3) limit. This structure for the $`\pi NN`$ coupling has been studied in Ref. . By revisiting them here, we want to see whether the higher dimensional operators included in this work change the previous results. To proceed, we rearrange the sum rule equations for the $`i\gamma _5\widehat{p}`$ structure, Eqs. (29),(50), (70),(73), (75),(77) into the form,
$`a+bM^2=f(M^2),`$ (86)
by transferring baryon masses and the exponential factors to the RHS of the sum rules. Thus, in the case of the $`\pi NN`$ sum rule, Eq. (29), the parameters represent that
$`a=g_{\pi N}\lambda _N^2;b=g_{\pi N}\lambda _N^2C_{\pi N},`$ (87)
and similarly for the other couplings.
In figures 2 (a) (b), we plot the RHS of the sum rules $`f(M^2)`$ for the couplings. The continuum threshold is set to $`S_0=2.07`$ GeV<sup>2</sup> corresponding to the Roper resonance and is used in obtaining the solid lines. To see the sensitivity to the continuum threshold, the Borel curves with $`S_0=2.57`$ GeV<sup>2</sup> are also shown with the dashed lines. The $`\pi NN`$ Borel curves are almost the same as the ones presented in Ref. indicating that the higher dimensional operators included in this work do not change the previous results. By linearly fitting each Borel curve within an appropriate Borel window, we extract the two phenomenological parameters $`a`$ and $`b`$. The parameter $`a`$ is given by the intersection of the vertical axis ($`M^2=0`$) with the best fitting straight line, and the parameter $`b`$ is given by the slope of the line. As one can see from the figures, in most sum rules, there is huge sensitivity to the continuum threshold, which prevents us to extract reliably the parameters of the concern. At $`M^21`$ GeV<sup>2</sup>, the $`\pi NN`$ Borel curve undergoes 14 % change due to the continuum parameter, which however yields rather different value of $`a`$ as shown in table I. For the other sum rules, we can also see from the table that the extracted parameters are highly sensitive to the continuum. One of the reasons may be, as suggested in Ref. , because higher resonances with different parities add up in forming the continuum or the unknown single pole terms. The SU(3) parameter $`\alpha =F/(F+D)`$ extracted from table I is $`\alpha 1.24`$ when the continuum parameter $`S_0=2.07`$ GeV<sup>2</sup> is used. This gives $`F/D=5.17`$. But with $`S_0=2.57`$ GeV<sup>2</sup>, we have totally different value, $`\alpha =0.288`$, which yields $`F/D=0.4`$. Therefore, the $`i\gamma _5\widehat{p}`$ sum rules may not be useful in determining the $`F/D`$ ratio.
Figure 3 shows the Borel curves for the Dirac structure $`\gamma _5\sigma _{\mu \nu }q^\mu p^\nu `$. The sum rules, Eqs. (40), (57), (80), (81),(82) and (83) are arranged into the form of Eq. (86) and the RHS of that is plotted for each coupling in the figure. Hence, the best fitting straight line within a Borel window will provide us with the parameters $`a`$ and $`b`$, which represent the same quantities as before. In contrast to the $`i\gamma _5\widehat{p}`$ sum rules, the Borel curves in this case are rather insensitive to the continuum parameter $`S_0`$. The sensitivity of the Borel curves to the continuum at $`M^21`$ GeV<sup>2</sup> is about 2 % level. Therefore, this $`\gamma _5\sigma _{\mu \nu }q^\mu p^\nu `$ structure may provide a useful constraint for the $`F/D`$ ratio.
The $`\pi NN`$ sum rule in this work has been improved from the previous calculations of Ref. by including gluonic contributions combined with the external pion state. They are from the three-particle pion wave functions, which produces the term involving $`\delta ^2`$ in Eq. (40). The new term appears in the highest dimension and cancels the other OPE at the same dimension containing the quark-gluon mixed parameter $`m_0^2`$. Thus, the total OPE is well saturated by the first two OPE terms. (Note that the term involving the gluon condensate in the same dimension contributes negligibly to the total OPE.) Combining Eq. (40) with the chiral-odd nucleon sum rule and taking the standard sum rule analysis, we obtain,
$`g_{\pi N}1314.`$ (88)
The errors are coming from how we choose the Borel window. This is certainly consistent with its empirical value as well as the one obtained from the sum rule beyond the chiral limit . This also means that the gluonic contributions which were not included in our previous study are important in stabilizing the sum rules and in obtaining the coupling agreeing with the phenomenology.
Other Borel curves for the sum rules, Eqs. (57), (80), (81), (82) and (83) are plotted in figure 3. As we have demonstrated, each OPE satisfies the SU(3) relation if we identify $`\alpha `$ as given in Eq. (59). Thus, as far as the OPE is concerned, all the sum rules in the SU(3) limit are related by the SU(3) rotations. This means that the same Borel window should to be used for the other couplings. Table II shows our results from the $`\gamma _5\sigma _{\mu \nu }q^\mu p^\nu `$ sum rules. In the $`\pi \mathrm{\Sigma }\mathrm{\Sigma }`$ case, there is no dependence on the continuum threshold. The ratios given in the fourth column are directly related to the SU(3) relations for the couplings. From them, we consistently obtain $`\alpha =0.44`$, which yields
$`F/D0.78.`$ (89)
This is a factor of 4 larger than our previous value $`F/D0.2`$ determined beyond the chiral limit , clearly indicating the Dirac structure dependence of a sum rule result. It is even larger than the value from the SU(6) symmetry $`F/D2/3`$ or the recent value $`F/D0.57`$ . Using the empirical value for the $`\pi NN`$ coupling, $`g_{\pi N}=13.4`$, we obtain the following couplings in the SU(3) limit (indicated by the superscript below),
$`g_{\eta N}^{(S)}`$ $`=`$ $`5.76;g_{\pi \mathrm{\Xi }}^{(S)}=1.61;g_{\eta \mathrm{\Xi }}^{(S)}=14.47,`$ (90)
$`g_{\pi \mathrm{\Sigma }}^{(S)}`$ $`=`$ $`11.79;g_{\eta \mathrm{\Sigma }}^{(S)}=8.58.`$ (91)
These values are in contrast with the ones determined beyond the chiral limit,
$`g_{\eta N}^{(S)}`$ $`=`$ $`2.3;g_{\pi \mathrm{\Xi }}^{(S)}=8.7;g_{\eta \mathrm{\Xi }}^{(S)}=10.5,`$ (92)
$`g_{\pi \mathrm{\Sigma }}^{(S)}`$ $`=`$ $`4.7;g_{\eta \mathrm{\Sigma }}^{(S)}=12.8.`$ (93)
Once again, the Dirac structure dependence of a sum rule result is clearly exhibited.
What could be the reasons for this Dirac structure dependence ? Since the $`i\gamma _5`$ structure has the same chirality as the $`\gamma _5\sigma _{\mu \nu }q^\mu p^\nu `$ structure, the dependence can not be attributed only to the higher resonance contributions . It seems rather due to the use of Ioffe current for baryon interpolating fields. In fact, Ioffe current is a specific choice for the nucleon interpolating field from a more general current of the type,
$`J_N=2ϵ_{abc}[(u_a^\mathrm{T}Cd_b)\gamma _5u_c+t(u_a^\mathrm{T}C\gamma _5d_b)u_c].`$ (94)
That is, when $`t=1`$, this current reduces to Ioffe current. One speculation that we can think of is that the choice with $`t=1`$ is not optimal for the nucleon. Other speculation is the following. Since we obtain the right strength for the $`\pi NN`$ coupling at least from the $`i\gamma _5`$ and $`\gamma _5\sigma _{\mu \nu }q^\mu p^\nu `$ sum rules, it could be that the Ioffe current is fine for the nucleon but its SU(3) rotated versions may not be optimal for the hyperons. Further studies are necessary to understand this problem in future.
## VI Meson-baryon couplings from the pseudotensor structure
In the previous section, we studied the SU(3) relations for the couplings in the SU(3) limit from the $`i\gamma _5\widehat{p}`$ and $`\gamma _5\sigma _{\mu \nu }q^\mu p^\nu `$ sum rules respectively. The $`i\gamma _5\widehat{p}`$ sum rules contain strong dependence on the continuum parameter and may not be relevant for our purpose of obtaining the $`F/D`$ ratio. On the other hand, the $`\gamma _5\sigma _{\mu \nu }q^\mu p^\nu `$ sum rules have been found to provide a reasonable $`\pi NN`$ coupling with less sensitivity to the continuum threshold. The obtained valued for the $`F/D`$ ratio does not however agree with the previous result beyond the chiral limit. Thus, the $`\gamma _5\sigma _{\mu \nu }q^\mu p^\nu `$ sum rules provide another set of couplings when the calculation is performed beyond the SU(3) limit. As we have emphasized, however, our results for strangeness baryons should be taken with some caution because our sum rules constructed at the order $`𝒪(p_\mu )`$ may draw a doubt when the strange quark is involved. The strange quark mass, as it is higher than $`𝒪(p_\mu )`$, should not be included in our approach. A question therefore is whether or not the sum rules truncated at order of $`𝒪(p_\mu )`$ make sense when the massive strange quark is involved. A more systematic method may be needed in future to verify (or refute) our approach here. But in our standpoint, there is no such a method at present.
With this limitation in mind, we present our results for the $`\gamma _5\sigma _{\mu \nu }q^\mu p^\nu `$ sum rules beyond the SU(3) limit. In our analysis, we will use the SU(3) breaking parameters
$`f_\eta =1.2f_\pi ;\overline{s}s=0.8\overline{q}q.`$ (95)
The $`f_\eta `$ value is from Ref. . In addition, phenomenological parameters such as baryon masses and the strengths $`\lambda _{B_j}`$ will change as we move away from the SU(3) limit. We take the empirical values for the masses and, for the strengths, we will discuss them as we move along. We ignore the SU(3) breaking driven by the singlet-octet mixing.
Figure 4 shows the Borel curves for the $`\pi \mathrm{\Sigma }\mathrm{\Sigma }`$, $`\eta \mathrm{\Sigma }\mathrm{\Sigma }`$, $`\pi \mathrm{\Xi }\mathrm{\Xi }`$ and $`\eta \mathrm{\Xi }\mathrm{\Xi }`$. Around the resonance masses ($`M^21.41`$ GeV<sup>2</sup> for $`\mathrm{\Sigma }`$ and $`M^21.73`$ GeV<sup>2</sup> for $`\mathrm{\Xi }`$), they are well-fitted by a straight line, suggesting that the dependence on the chosen Borel window is marginal. As we have discussed, the dependence on the continuum threshold is also small. The $`\eta NN`$ curve (not shown) basically has the similar features as the case in the SU(3) symmetric limit but is shifted up slightly. The best fitting parameters are given in table III. Due to the unknown strengths $`\lambda _{B_j}`$ ($`B_j=N,\mathrm{\Sigma },\mathrm{\Xi }`$), we here present ratios of the couplings obtained from table III,
$`{\displaystyle \frac{g_{\eta N}^{(B)}}{g_{\pi N}^{(B)}}}`$ $`=`$ $`0.55;{\displaystyle \frac{g_{\eta \mathrm{\Xi }}^{(B)}}{g_{\pi \mathrm{\Xi }}^{(B)}}}=4.2;{\displaystyle \frac{g_{\eta \mathrm{\Sigma }}^{(B)}}{g_{\pi \mathrm{\Sigma }}^{(B)}}}=0.94.`$ (96)
In comparison with the ratios in the SU(3) limit,
$`{\displaystyle \frac{g_{\eta N}^{(S)}}{g_{\pi N}^{(S)}}}`$ $`=`$ $`0.43;{\displaystyle \frac{g_{\eta \mathrm{\Xi }}^{(S)}}{g_{\pi \mathrm{\Xi }}^{(S)}}}=8.9;{\displaystyle \frac{g_{\eta \mathrm{\Sigma }}^{(S)}}{g_{\pi \mathrm{\Sigma }}^{(S)}}}=0.74.`$ (97)
the SU(3) breaking effects are huge for the meson-$`\mathrm{\Xi }`$ couplings but not so large for the other couplings.
Let’s us now compare our results to that of other works. In table IV, our results in the SU(3) limit and that from Refs. are shown. The couplings in Ref. are based on the assumption of the hyperon-nucleon potentials obeying SU(3) symmetry. Except for the results on $`g_{\pi \mathrm{\Xi }}`$ and $`g_{\eta \mathrm{\Sigma }}`$, our results qualitatively agree with that of Ref. . Another approach based on the QCD parametrization method gives results (see the 5th column of table IV.) not so different from Ref. . Comparing our results in SU(3) to that in Ref. , we find qualitative agreement for $`\pi \mathrm{\Sigma }\mathrm{\Sigma }`$ but large discrepancy in $`\pi \mathrm{\Xi }\mathrm{\Xi }`$. To see the SU(3) breaking directly reflected in the couplings, we simply use the scaling proposed by Ref. , $`\lambda _{B_j}^2m_{B_j}^6`$, and calculate each coupling. As shown in the 3rd column of table IV, most couplings change noticeably as we turn on the SU(3) breaking effects. As a result, they do not agree with the ones from other works. However, our results with the broken SU(3) should be taken with cautions. Since our sum rules are constructed at the order $`p_\mu `$, the couplings obtained are the ones at the kinematical point $`p^2=0`$. But the physical couplings are defined at the kinematical point $`p^2=m_{\pi ,\eta }^2`$. Therefore, in the $`\eta BB`$ cases, one can expect some changes in this extrapolation. Furthermore our formalism should be improved by including higher effects of the broken SU(3) more systematically.
## VII Summary
In this work, we have developed QCD sum rules beyond the soft-meson limit for the diagonal meson-baryon couplings, $`\pi NN`$, $`\eta NN`$, $`\pi \mathrm{\Xi }\mathrm{\Xi }`$, $`\eta \mathrm{\Xi }\mathrm{\Xi }`$, $`\pi \mathrm{\Sigma }\mathrm{\Sigma }`$ and $`\eta \mathrm{\Sigma }\mathrm{\Sigma }`$. The Dirac structures $`i\gamma _5\widehat{p}`$ and $`\gamma _5\sigma _{\mu \nu }q^\mu p^\nu `$ are separately considered in constructing the sum rules. In the first stage, we have improved the previous calculations of the $`\pi NN`$ coupling by including three-particle pion wave functions mediated by the gluonic tensor. The $`\gamma _5\sigma _{\mu \nu }q^\mu p^\nu `$ sum rule in this revision provides the $`\pi NN`$ coupling closed to its empirical value, while no critical change has been observed for the $`i\gamma _5\widehat{p}`$ sum rules. By extend the $`\pi NN`$ sum rules to the other couplings, we have studied the SU(3) relations for the couplings. Depending on the Dirac structure considered, we have reported different identifications of the $`F/D`$ ratio. Therefore, our findings support the previous claim of the Dirac structure dependence of a sum rule result . In the sum rule analysis, the $`i\gamma _5\widehat{p}`$ sum rules were found to give the results very sensitive to the continuum threshold. On the other hand, stable results are obtained from the $`\gamma _5\sigma _{\mu \nu }q^\mu p^\nu `$ sum rule, which however is not consistent with the previous results obtained from the sum rule beyond the chiral limit . We have therefore provided a different set of the couplings in the SU(3) limit and beyond the SU(3) limit using the Dirac structure $`\gamma _5\sigma _{\mu \nu }q^\mu p^\nu `$. The obtained $`F/D`$ ratio from the $`\gamma _5\sigma _{\mu \nu }q^\mu p^\nu `$ sum rules is $`0.78`$, slightly larger than the SU(6) value of $`2/3`$. We have also discussed the SU(3) breaking effects in the couplings.
###### Acknowledgements.
This work is supported in part by the Grant-in-Aid for JSPS fellow, and the Grant-in-Aid for scientific research (C) (2) 11640261 of the Ministry of Education, Science, Sports and Culture of Japan. The work of H. Kim is also supported by Research Fellowships of the Japan Society for the Promotion of Science. The work of S. H. Lee is supported by the KOSEF grant number 1999-2-111-005-5 and by the BK 21 project of the Korean Ministry of Education.
## A Derivation of the pseudovector matrix elements
Here we evaluate the matrix elements given in Eqs. (13) and (14) by considering
$`0|\overline{d}(0)\gamma _\mu \gamma _5u(x)|\pi ^+(p)`$ (A1)
to leading order in the pion momentum $`p_\mu `$. We consider the case with a charged pion to use some informations from Ref. for twist-4 element. The matrix elements for the neutral pion, which are of our concern, can be obtained simply by isospin rotations afterward. By expanding $`u(x)`$ in $`x_\mu `$, we have
$`u(x)=u(0)+x_\mu ^\mu u(0)+{\displaystyle \frac{1}{2}}x_\mu x_\nu ^\mu ^\nu u(0)+\mathrm{}.`$ (A2)
In the fixed-point gauge ($`x_\mu A^\mu =0`$), the partial derivative can be replaced by the covariant derivative, $`_\mu D_\mu _\mu ig_sA_\mu `$. In this expansion of Eq.(A1), the first term is given by the PCAC,
$`0|\overline{d}(0)\gamma _\mu \gamma _5u(0)|\pi ^+(p)=i\sqrt{2}f_\pi p_\mu ,`$ (A3)
where $`f_\pi =93`$ MeV. The second term in the expansion, as it contains one covariant derivative, is linear in the quark mass, $`𝒪(m_q)`$, which is higher order than $`p_\mu `$. The third term in the expansion containing two covariant derivatives can be written
$`{\displaystyle \frac{1}{2}}x^\alpha x^\beta 0|\overline{d}(0)\gamma _\mu \gamma _5D_\alpha D_\beta u(0)|\pi ^+(p)`$ (A4)
$`={\displaystyle \frac{1}{4}}x^\alpha x^\beta 0|\overline{d}(0)\gamma _\mu \gamma _5(D_\alpha D_\beta +D_\beta D_\alpha )u(0)|\pi ^+(p).`$ (A5)
The matrix element in the RHS sandwiched by the vacuum and the pion state is symmetric under $`\alpha \beta `$. Hence, it should be built by symmetrically combining the available four vector $`p_\mu `$ and the metric $`g_{\mu \nu }`$. At the first order in $`p_\mu `$, it is easy to see that
$`0|\overline{d}(0)\gamma _\mu \gamma _5{\displaystyle \frac{1}{2}}(D_\alpha D_\beta +D_\beta D_\alpha )u(0)|\pi ^+(p)=g_{\alpha \beta }p_\mu B+(g_{\alpha \mu }p_\beta +g_{\beta \mu }p_\alpha )C.`$ (A6)
Note, other possible combinations are higher order than $`p_\mu `$. To determine the invariant functions $`B`$ and $`C`$, first let us multiply $`g_{\alpha \beta }`$ on both sides of Eq. (A6) to obtain,
$`0|\overline{d}(0)\gamma _\mu \gamma _5D^2u(0)|\pi ^+(p)=(4B+2C)p_\mu .`$ (A7)
Other constraint equation can be obtained by multiplying $`g_{\alpha \mu }`$ on Eq. (A6),
$`0|\overline{d}(0)\gamma _5{\displaystyle \frac{1}{2}}(\widehat{D}D_\beta D_\beta \widehat{D})u(0)|\pi ^+(p)=(B+5C)p_\beta ,`$ (A8)
where $`\widehat{D}\gamma _\mu D^\mu `$. The LHS is linear in the quark mass $`m_q`$ and higher order in chiral counting than the order $`p_\mu `$. Thus, to leading order in $`p_\mu `$, the LHS of Eq (A8) should be zero, which yields the relation,
$`B=5C.`$ (A9)
Combining this with Eq. (A7), we have
$`Bp_\mu `$ $`=`$ $`{\displaystyle \frac{5}{18}}0|\overline{d}(0)\gamma _\mu \gamma _5D^2u(0)|\pi ^+(p),`$ (A10)
$`Cp_\mu `$ $`=`$ $`{\displaystyle \frac{1}{18}}0|\overline{d}(0)\gamma _\mu \gamma _5D^2u(0)|\pi ^+(p).`$ (A11)
According to Ref. , the matrix element in the RHS is given by
$$0|\overline{d}(0)\gamma _\mu \gamma _5D^2u(0)|\pi ^+(p)=i\sqrt{2}f_\pi \delta ^2p_\mu ,$$
(A12)
with $`\delta ^2=0.2`$ GeV<sup>2</sup>. Therefore, we have
$`B=i{\displaystyle \frac{5\sqrt{2}}{18}}f_\pi \delta ^2;C=i{\displaystyle \frac{\sqrt{2}}{18}}f_\pi \delta ^2.`$ (A13)
Using this result in Eq. (A6) and putting into Eq. (A5), we obtain
$`{\displaystyle \frac{1}{2}}x^\alpha x^\beta 0|\overline{d}(0)\gamma _\mu \gamma _5D_\alpha D_\beta u(0)|\pi ^+(p)=i{\displaystyle \frac{5}{18}}\sqrt{2}f_\pi \delta ^2\left({\displaystyle \frac{1}{2}}x^2p_\mu {\displaystyle \frac{1}{5}}x_\mu xp\right).`$ (A14)
Thus, up to twist-4 but to leading order in $`p_\mu `$, we have the expansion
$`0|\overline{d}(0)\gamma _\mu \gamma _5u(x)|\pi ^+(p)=i\sqrt{2}f_\pi p_\mu i{\displaystyle \frac{5}{18}}\sqrt{2}f_\pi \delta ^2\left({\displaystyle \frac{1}{2}}x^2p_\mu {\displaystyle \frac{1}{5}}x_\mu xp\right).`$ (A15)
Note that higher twists have been neglected as usual in QCD sum rules. Invoking the isospin symmetry, we directly obtain
$`0|\overline{u}(0)\gamma _\mu \gamma _5u(x)|\pi ^0(p)`$ $`=`$ $`if_\pi p_\mu i{\displaystyle \frac{5}{18}}f_\pi \delta ^2\left({\displaystyle \frac{1}{2}}x^2p_\mu {\displaystyle \frac{1}{5}}x_\mu xp\right),`$ (A16)
$`0|\overline{d}(0)\gamma _\mu \gamma _5d(x)|\pi ^0(p)`$ $`=`$ $`if_\pi p_\mu +i{\displaystyle \frac{5}{18}}f_\pi \delta ^2\left({\displaystyle \frac{1}{2}}x^2p_\mu {\displaystyle \frac{1}{5}}x_\mu xp\right).`$ (A17)
## B Derivation of the pseudotensor matrix elements
Here we calculate the matrix element
$`0|\overline{u}(0)\gamma _5\sigma _{\alpha \beta }u(x)|\pi ^0(p)`$ (B1)
to leading order in $`p_\mu `$. By expanding $`u(x)`$ in $`x_\mu `$, we have
$`0|\overline{u}(0)\gamma _5\sigma _{\alpha \beta }u(0)|\pi ^0(p)+x^\mu 0|\overline{u}(0)\gamma _5\sigma _{\alpha \beta }D_\mu u(0)|\pi ^0(p)+\mathrm{}.`$ (B2)
The dots are higher orders than $`𝒪(p_\mu )`$ . The first term in the expansion is zero simply because it is not possible to make antisymmetric combinations with respect to $`\alpha `$ and $`\beta `$ using the available vector $`p_\mu `$ and the metric $`g_{\mu \nu }`$. The matrix element in the second term can be written
$`0|\overline{u}(0)\gamma _5\sigma _{\alpha \beta }D_\mu u(0)|\pi ^0(p)=i(p_\alpha g_{\beta \mu }p_\beta g_{\alpha \mu })A.`$ (B3)
No other combinations are allowed at order $`𝒪(p_\mu )`$. To get the scalar function $`A`$, we multiply both sides with $`g_{\mu \beta }`$ and get,
$`i0|\overline{u}(0)\gamma _5D_\alpha u(0)|\pi ^0(p)+𝒪(m_q)=3ip_\alpha A.`$ (B4)
The $`m_q`$ term is higher chiral order than $`p_\mu `$ and can be neglected at the order that we are interested in. The other term in the LHS is proportional to the first moment of the twist-3 pion wave function,
$`0|\overline{u}(0)\gamma _5D_\alpha u(0)|\pi ^0(p)={\displaystyle \frac{\overline{u}u}{f_\pi }}p_\alpha {\displaystyle _0^1}𝑑uu\phi _p(u)={\displaystyle \frac{\overline{u}u}{2f_\pi }}p_\alpha .`$ (B5)
Note that the zeroth moment of the wave function, that is $`_0^1𝑑u\phi _p(u)=1`$, is fixed solely by the soft-pion theorem. The first moment $`_0^1𝑑uu\phi _p(u)=1/2`$ has been used according to Ref. . Hence,
$$A=\frac{\overline{u}u}{6f_\pi },$$
(B6)
which yields the tensor matrix element at the order of $`𝒪(p_\mu )`$,
$`0|\overline{u}(0)\gamma _5\sigma _{\alpha \beta }u(x)|\pi ^0(p)=i(p_\alpha x_\beta p_\beta x_\alpha ){\displaystyle \frac{\overline{u}u}{6f_\pi }}.`$ (B7)
Due to the isopin symmetry, the $`d`$-quark element is given by
$`0|\overline{d}(0)\gamma _5\sigma _{\alpha \beta }d(x)|\pi ^0(p)=i(p_\alpha x_\beta p_\beta x_\alpha ){\displaystyle \frac{\overline{d}d}{6f_\pi }}.`$ (B8)
The sign different from the u-quark element can be directly seen by using the soft-pion theorem.
## C Derivation of the three-particle pion matrix elements
In this appendix, we derive the three-particle pion matrix element, Eq. (19), contributing to our sum rules to leading order in $`p_\mu `$. This matrix element can be obtained by inserting a gluonic tensor from a quark propagator into the quark-antiquark component with a pion. Among various possibilities, the one that survives to leading order in $`p_\mu `$ can be written,
$`0|q(x)_a^\alpha ig_s[\stackrel{~}{G}^A(0)]^{\sigma \rho }\overline{q}(0)_b^\beta |\pi ^0(p)=t_{ab}^A(\gamma _\theta )^{\alpha \beta }i(A_G^q)^{\theta \sigma \rho }.`$ (C1)
Here $`\stackrel{~}{G}_{\alpha \beta }^A`$ is the dual of the gluon strength tensor, $`\stackrel{~}{G}_{\alpha \beta }^A=\frac{1}{2}ϵ_{\alpha \beta \sigma \rho }(G^A)^{\sigma \rho }`$, and the color matrix $`t^A`$ is related to the Gell-Mann matrices via $`t^A=\lambda ^A/2`$. Other possibilities in combining a gluonic tensor with the quark-antiquark component are at least the second order in $`p_\mu `$ . On multiplying $`t_{ba}^A(\gamma ^\delta )^{\beta \alpha }`$ on both sides, Eq. (C1) becomes,
$`0|\overline{q}(0)\gamma ^\delta ig_s\stackrel{~}{𝒢}^{\sigma \rho }(0)q(x)|\pi ^0(p)=16i(A_G^q)^{\delta \sigma \rho },`$ (C2)
where $`\stackrel{~}{𝒢}_{\sigma \rho }=t^A\stackrel{~}{G}_{\sigma \rho }^A`$. At order $`p_\mu `$, the matrix element of the LHS contributes at the local limit ($`x_\mu =0`$). Among all possible antisymmetric combinations with respect to the indices $`\sigma `$ and $`\rho `$, the only possibility is
$`0|\overline{q}(0)\gamma ^\delta ig_s\stackrel{~}{𝒢}^{\sigma \rho }(0)q(0)|\pi ^0(p)=(p^\rho g^{\sigma \delta }p^\sigma g^{\rho \delta })A_G^q.`$ (C3)
To determine the scalar function $`A_G^q`$, we multiply both sides with $`g^{\delta \sigma }`$. Then after some manipulations, the LHS becomes
$`0|\overline{q}(0)\gamma _\sigma ig_s\stackrel{~}{𝒢}^{\sigma \rho }(0)q(0)|\pi ^0(p)`$ $`=`$ $`0|\overline{q}(0)\gamma ^\rho \gamma _5iD^2q(0)|\pi ^0(p)`$ (C4)
$`=`$ $`\pm \delta ^2f_\pi p^\rho `$ (C5)
where the plus sign is for the $`u`$-quark and the minus sign is for the $`d`$-quark. From Eqs. (C3) and (C5), we have $`A_G^q=\pm \delta ^2f_\pi /3`$, which from Eq.(C2) yields
$`i(A_G^q)^{\delta \sigma \rho }={\displaystyle \frac{1}{48}}(p^\rho g^{\sigma \delta }p^\sigma g^{\rho \delta })\delta ^2f_\pi .`$ (C6)
Therefore, putting into Eq. (C1), we have Eq. (19). |
warning/0002/hep-th0002193.html | ar5iv | text | # Noncommutative instantons and twistor transform
## 1. Physical motivation
In this section we explain the physical motivation for studying instantons on a noncommutative $`^4`$. Readers uninterested in the motivation may skip most of this section and proceed directly to subsection 1.5. Likewise, readers familiar with the way noncommutative instantons arise in string theory may start with subsection 1.5.
### 1.1. Instanton equations
Let $`E`$ be a vector bundle with structure group $`G`$ on an oriented Riemannian 4-manifold $`X,`$ and let $`A`$ be a connection on $`E`$. Instanton equation is the equation
(1)
$$F_A^+=0,$$
where $`F_A`$ is the curvature of $`A,`$ and $`F_A^+`$ denotes the self-dual (SD) part of $`F_A`$. Solutions of this equation are called instantons, or anti-self-dual (ASD) connections. The second Chern class of $`E`$ is known in the physics literature as the instanton number. Instantons automatically satisfy the Yang-Mills equation $`d_A(F)=0,`$ where $`d_A:\mathrm{\Omega }^p\mathrm{End}(E)\mathrm{\Omega }^{p+1}\mathrm{End}(E)`$ is the covariant differential, and $`:\mathrm{\Omega }^p\mathrm{\Omega }^{4p}`$ is the Hodge star operator.
There are several physical reasons to be interested in instantons. If one is studying quantum gauge theory on a Riemannian 3-manifold $`M`$ (space), then instantons on $`X=M\times `$ describe quantum-mechanical tunneling between different classical vacua. The possibility of such tunneling has drastic physical effects, some of which can be experimentally observed. If one is studying classical gauge theory on a 5-dimensional space-time $`X\times ,`$ then instantons on $`X`$ can be interpreted as solitons, i.e. as static solutions of the Yang-Mills equations of motion. In fact, instantons are the absolute minima of the Yang-Mills energy function of the 5-dimensional theory (with fixed second Chern class).
Both interpretations arise in string theory, but to explain this we need to make a digression and discuss D-branes.
### 1.2. D-branes
It has been discovered in the last few years that string theory describes, besides strings, extended objects (branes) of various dimensions. These extended objects should be regarded as static solutions of (as yet poorly understood) stringy equations of motion. D-branes are a particularly manageable class of branes. Recall that ordinary closed oriented superstrings, known as Type II strings, are described by maps from a Riemann surface $`\mathrm{\Sigma }`$ (“worldsheet”) to a 10-dimensional manifold $`Z`$ (“target”). The physical definition of a D-brane is “a submanifold of $`Z`$ on which strings can end.” This means that if a D-brane is present, then one needs to consider maps from a Riemann surface with boundaries to $`Z`$ such that the boundaries are mapped to a certain submanifold $`XZ`$. In this case one says that there is a D-brane wrapped on $`X`$. If $`X`$ is connected and has dimension $`p+1,`$ then one says that one is dealing with a Dp-brane. In general, $`X`$ can have several components with different dimensions, and then each component corresponds to a D-brane.
In perturbative string theory, the role of equations of motion is played by the condition that a certain auxiliary quantum field theory on the Riemann surface $`\mathrm{\Sigma }`$ is conformally invariant. When D-branes are present, $`\mathrm{\Sigma }`$ has boundaries, and the auxiliary theory must be supplemented with boundary conditions. The requirement that the boundary conditions preserve conformal invariance imposes constraints on the submanifold $`X`$. This constraints should be regarded as equations of motion for D-branes. For example, if we consider a D0-brane wrapped on a 1-dimensional submanifold $`X,`$ then conformal invariance requires that $`X`$ be a geodesic in $`Z`$. This is the usual equation of motion for a relativistic particle moving in $`Z`$.
An important subtlety is that to specify fully the boundary conditions for the auxiliary theory on $`\mathrm{\Sigma }`$ it is not sufficient to specify $`X;`$ one should also specify a unitary vector bundle $`E`$ on $`X`$ and a connection on it. In the simplest case this bundle has rank $`1,`$ but one can also have “multiple” D-branes, described by bundles of rank $`r>1`$. Such bundles describe $`r`$ coincident D-branes wrapped on the same submanifold $`X`$. Using the requirement of conformal invariance of the auxiliary two-dimensional quantum field theory, one can derive equations of motion for the Yang-Mills connection on $`E`$. In the low-energy approximation, the equations of motion are the usual Yang-Mills equations $`d_A(F_A)=0.`$ In particular, instantons are solutions of these equations.
### 1.3. Instantons and D-branes
Let $`Z`$ be $`^{10}`$ with a flat metric, and let $`XZ`$ be $`^5=^4\times `$ linearly embedded in $`Z`$. We regard $`^4`$ as space and $``$ as time. Consider $`r`$ D4-branes wrapped on $`X`$. This physical system is described by the Yang-Mills action on $`^5=^4\times `$. If one is looking for static solutions of the equations of motion, one needs to consider the minima of the Yang-Mills energy function
$$W[A]=_^4F_A^2,$$
where $`F_A`$ is the curvature of a $`U(r)`$ connection $`A,`$ and $`||F_A||^2=\mathrm{Tr}(F_AF_A).`$ The instanton number of $`A`$ is defined by
(2)
$$c_2=\frac{1}{8\pi ^2}_^4\mathrm{Tr}\left(F_AF_A\right).$$
If the Yang-Mills energy evaluated on $`A`$ is finite, then the bundle $`E`$ and the connection $`A`$ extend to $`𝕊^4`$, the one-point compactification of $`^4`$ (see for details). In this case $`c_2`$ is the second Chern class of $`E`$ and is therefore an integer.
Solutions of instanton equations on $`^4`$ are precisely the absolute minima of the Yang-Mills energy function. These solutions should be regarded as composed of identical particle-like objects (instantons) on $`X,`$ their number being $`c_2`$. Since the energy of the instanton is proportional to $`c_2,`$ all “particles” have the same mass. Since the solution is static, the particles neither repel nor attract. This is actually a consequence of supersymmetry: Type II string theory is supersymmetric, and D4-branes with instantons on them leave part of supersymmetry unbroken.
In string theory one may also consider $`k`$ D0-branes present simultaneously with $`r`$ D4-branes. More specifically, we will consider D0-branes which are at rest, i.e. the corresponding one-dimensional manifolds are straight lines parallel to the time axis. Such a configuration of branes is also supersymmetric, and consequently there are no forces between any of the branes. The positions of D0-branes are not constrained by anything, so their moduli space is $`(^9)^k`$. More precisely, since D0-branes are indistinguishable, the moduli space is $`Sym^k(^9)`$.
It turns out that an instanton with instanton number $`k`$ and $`k`$ D0-branes are related: they can be deformed into each other without any cost in energy. A convenient point of view is the following. In the presence of D4-branes wrapped on $`X`$ the moduli space of D0-branes has two branches: a branch where their positions are unconstrained and D0-branes are point-like (this branch is isomorphic to $`Sym^k(^9)`$), and the branch where they are constrained to lie on $`X`$. The latter branch is isomorphic to the moduli space $`_{r,k}`$ of $`U(r)`$ instantons on $`X=^4`$ with $`c_2=k`$.
The dimension of $`_{r,k}`$ is known to be $`4rk`$ for $`r>1`$ (see for example ). For $`r=1`$ instantons do not exist. The translation group of $`^4`$ acts freely on $`_{r,k},`$ and the quotient space describes the relative positions and sizes of instantons. Thus D0-branes are point-like objects when they are away from D4-branes, but when they bind to D4-branes they can acquire finite size.
The “instanton” branch touches the “point-like” branch at submanifolds where some or all of the instantons shrink to zero size. These are the submanifolds where the instanton moduli space is singular. At these submanifolds the point-like instantons can detach from D4-branes and start a new life as D0-branes. This lowers the second Chern class of the bundle on D4-branes. Thus from the string theory perspective it is natural to glue together the moduli spaces of instantons with different Chern classes along singular submanifolds.
### 1.4. Noncommutative geometry and D-branes
Instanton equations (and, more generally, Yang-Mills equations) arise in the low-energy limit of string theory, or equivalently for large string tension. Recently, another kind of low-energy limit of string theory was discussed in the literature . Consider a trivial $`U(r)`$-bundle on $`X=^4`$ with a connection $`A`$ whose curvature $`F_A`$ is of the form $`1f,`$, where $`1`$ is the unit section of $`\mathrm{End}(E)`$, and $`f`$ is a constant nondegenerate 2-form. For small $`f`$ the D4-branes are described by the ordinary Yang-Mills action, but for large $`F_A`$ the stringy equations of motion get complicated. It turns out that the equations of motion simplify again in the limit when both $`F_A`$ and the string tension are taken to infinity, with a certain combination of the two kept fixed (one also has to scale the metric appropriately, see ). We will call this limit the Seiberg-Witten limit. In this limit the D4-branes are described by Yang-Mills equations on a certain noncommutative deformation of $`^4`$ (see and references therein).
There is another description of the Seiberg-Witten limit, which is gauge-equivalent to the previous one. Type II string theory reduces at low energies to Type II supergravity in 10 dimensions. The bosonic fields of this low-energy theory include a symmetric rank-two tensor (metric) and a 2-form $`B`$. $`^{10}`$ with a flat Lorenzian metric and a constant $`B`$ is a solution of supergravity equations of motion, as well as full stringy equations of motion. A constant $`B`$ can be gauged away, so this is not a very interesting solution. Life gets more interesting if there are D-branes present. For example, consider $`r`$ coincident flat D4-branes embedded in $`^{10}`$ with a constant $`B`$-field. It turns out that one can gauge away a constant $`B`$-field only at the expense of introducing a constant $`F_A`$ of the form $`1f`$ where $`f`$ is equal to the pull-back of $`B`$ to the worldvolume of the D4-branes. Thus the solution with zero $`F_A`$ and nonzero $`B`$ is equivalent to the solution with nonzero $`F_A`$ and zero $`B`$. Therefore the Seiberg-Witten limit can be described as the limit in which both the $`B`$-field and the string tension become infinite.
The idea that D-branes in a nonzero B-field are described Yang-Mills theory on a noncommutative space was first put forward in for the case of D-branes wrapped on tori.
### 1.5. Instanton equations on a noncommutative $`^4`$
The deformed $`^4`$ that one obtains in the Seiberg-Witten limit is completely characterized by its algebra of functions $`𝒜`$. It is a noncommutative algebra whose underlying space is a certain subspace of $`C^{\mathrm{}}`$ functions on $`^4`$. The product is the so-called Wigner-Moyal product formally given by
(3)
$$(fg)(x)=\underset{yx}{lim}\mathrm{exp}\left(\frac{1}{2}\mathrm{}\theta _{ij}\frac{^2}{x_iy_j}\right)f(x)g(y).$$
Here $`\theta `$ is a purely imaginary matrix, and $`\mathrm{}`$ is a real parameter (“Planck constant”) which is introduced to emphasize that the Wigner-Moyal product is a deformation of the usual product. In the string theory context $`\theta `$ is proportional to $`f^1.`$
Of course, to make sense of this definition we must specify a subspace in the space of $`C^{\mathrm{}}`$ functions which is closed under the Wigner-Moyal product. Leaving this question aside for a moment,<sup>1</sup><sup>1</sup>1String theory considerations do not shed light on this problem. one can define the exterior differential calculus over $`𝒜`$. Differential geometry of noncommutative spaces has been developed by A. Connes . In our situation Connes’ general theory is greatly simplified. For example, the sheaf of 1-forms $`\mathrm{\Omega }^1(𝒜)`$ is simply a bimodule $`𝒜^4`$ (the relation of this definition with the general theory is explained in subsection 8.11). The elements of $`\mathrm{\Omega }^1(𝒜)`$ will be denoted $`_if^i(x)dx_i,`$ or simply $`f^i(x)dx_i`$. The exterior differential $`d`$ is a vector space morphism
$$d:𝒜\mathrm{\Omega }^1(𝒜),f\frac{f}{x_i}dx_i.$$
The exterior differential $`d`$ satisfies the Leibniz rule
$$d(f_1f_2)=df_1f_2+f_1df_2.$$
This makes sense because $`\mathrm{\Omega }^1(𝒜)`$ is a bimodule.
The sheaf of 2-forms over $`𝒜`$ is a bimodule $`\mathrm{\Omega }^2(𝒜)=𝒜^6`$ (see subsection 8.11). The definition of the exterior differential can be extended to $`\mathrm{\Omega }^1(A)`$ in an obvious manner.
Complex conjugation acts as an anti-linear anti-homomorphism of $`𝒜`$, i.e. $`\overline{(fg)}=\overline{g}\overline{f}.`$ Thus $`𝒜`$ has a natural structure of a $``$-algebra. We will denote the $``$-conjugate of $`f𝒜`$ by $`f^{}`$.
A trivial bundle over the noncommutative $`^4`$ is defined as a free $`𝒜`$-module $`E`$. A trivial unitary bundle over the noncommutative $`^4`$ is defined as a free module $`V_{}𝒜,`$ where $`V`$ is a Hermitean vector space. A connection on a trivial bundle $`E`$ is defined as a map
$$:EE_𝒜\mathrm{\Omega }^1(𝒜),$$
which is a vector space morphism satisfying the Leibniz rule
$$(mf)=(m)f+mdf.$$
This formula makes use of the bimodule structure on $`\mathrm{\Omega }^1(𝒜)`$.
The curvature $`F_{}=[,]`$ is a morphism of $`𝒜`$-modules
$$F_{}:EE_𝒜\mathrm{\Omega }^2(𝒜).$$
As in the commutative case, a connection on a trivial bundle $`E`$ can be written in terms of a connection 1-form $`A\mathrm{End}_𝒜(E)_𝒜\mathrm{\Omega }^1(𝒜)`$:
$$(m)=dm+Am.$$
This formula uses the bimodule structure on $`m`$. If $`E`$ is a unitary bundle, and we have $`A^{}=A`$, then we say that $`A`$ is a unitary connection.
The curvature is given in terms of $`A`$ by the usual formula
$$F_{}:=F_A=dA+AA.$$
Here it is understood that
$$f^idx_ig^jdx_j=f^ig^jdx_idx_j.$$
The instanton equation on $`𝒜`$ is again given by (1), and the instanton number is defined by (2).
The most obvious choice of the space of functions closed under the Wigner-Moyal product is the space of polynomial functions. However, this choice is not suitable for our purposes because it precludes the decrease of $`F_A`$ at infinity which is necessary for the instanton action to converge. In the commutative case, components of an instanton connection are rational functions , so we would like our class of functions to include rational functions on $`^4`$. A possible choice for the underlying set of $`𝒜`$ is the set of $`C^{\mathrm{}}`$ functions on $`^4`$ all of whose derivatives are polynomially bounded. Then we face the question of the convergence of the series (3). To avoid dealing with this issue, we modify our definition of the Wigner-Moyal product (see Appendix for details). The modified product makes the space of $`C^{\mathrm{}}`$ functions all of whose derivatives are polynomially bounded into an algebra over $`,`$ and agrees with (3) on polynomial functions.
Polynomial functions form a subalgebra of $`𝒜`$. This subalgebra is isomorphic to the algebra generated by four variables $`x_i,i=1,2,3,4`$ with relations
$$[x_i,x_j]=\mathrm{}\theta _{ij}.$$
This algebra is usually called the Weyl algebra.
To summarize, there is a limit of string theory in which D4 branes are described by Yang-Mills equations on the noncommutative $`^4`$ ($`=𝒜`$). D0-branes bound to D4-branes are described in this limit by the instanton equations on the noncommutative $`^4`$. One can show that, unlike in the commutative case, instantons cannot be deformed to point-like D0-branes without a cost in energy. Thus it is natural to suspect that the moduli space of instantons on the noncommutative $`^4`$ is metrically complete.
## 2. Review of the ADHM construction and summary
All instantons on the commutative $`^4`$ arise from the so-called ADHM construction. Recently N. Nekrasov and A. Schwarz introduced a modification of this construction which produces instantons on the noncommutative $`^4`$.<sup>2</sup><sup>2</sup>2As in the commutative case, one may consider different classes of functions on the noncommutative $`^4`$: polynomial, $`C^{\mathrm{}}`$ functions rapidly decreasing at infinity, $`C^{\mathrm{}}`$ functions all of whose derivatives are polynomially bounded, etc. Our class of functions differs somewhat from that adopted in . In the commutative case the completeness of the ADHM construction can proved using the twistor transform of R. Penrose, so one could hope that that the same approach could work in the noncommutative case. In this paper we show that the deformed ADHM data of describe holomorphic bundles on certain noncommutative algebraic varieties and interpret the deformed ADHM construction in terms of noncommutative twistor transform. In this subsection we review both ordinary and deformed ADHM constructions and make a summary of our results.
### 2.1. Review of the ADHM construction of instantons
First let us outline the ADHM construction of $`U(r)`$ instantons on the commutative $`^4`$ following . We assume that the constant metric $`G`$ on $`^4`$ has been brought to the standard form $`G=\mathrm{diag}(1,1,1,1)`$ by a linear change of basis. To construct a $`U(r)`$ instanton with $`c_2=k`$ one starts with two Hermitean vector spaces $`V^k`$ and $`W^r`$. The ADHM data consist of four linear maps $`B_1,B_2\mathrm{Hom}(V,V),I\mathrm{Hom}(W,V),J\mathrm{Hom}(V,W)`$ which satisfy the following two conditions:
(i) $`\mu _c=[B_1,B_2]+IJ=0,\mu _r=[B_1,B_1^{}]+[B_2,B_2^{}]+II^{}J^{}J=0.`$
(ii) For any $`\xi =(\xi _1,\xi _2)^2^4`$ the linear map $`𝒟_\xi \mathrm{Hom}(VVW,VV)`$ defined by
(4)
$$𝒟_\xi =\left(\begin{array}{ccc}B_1\xi _1\hfill & \hfill B_2+\xi _2& I\hfill \\ B_2^{}\overline{\xi }_2\hfill & \hfill B_1^{}\overline{\xi }_1& J^{}\hfill \end{array}\right)$$
is surjective.
The equations $`\mu _c=\mu _r=0`$ are called the ADHM equations. They are invariant with respect to the action of the group of unitary transformations of $`V`$. Solutions of these equations are called ADHM data. The space of ADHM data modulo $`U(V)`$ transformations has dimension $`4rk`$ and carries a natural hyperkähler metric. ADHM construction identifies this moduli space with the moduli space of $`U(r)`$ instantons with $`c_2=k`$ and fixed trivialization at infinity. The role of the condition (ii) above is to remove submanifolds in this moduli space where the hyperkähler metric becomes singular (these are point-like instanton singularities mentioned in subsection 1.3). As a result the moduli space of the ADHM data is metrically incomplete.
The instanton connection can be reconstructed from the ADHM data as follows. The condition (ii) implies that the family $`𝐊𝐞𝐫𝒟_\xi `$ forms a trivial subbundle of $`VVW`$ of rank $`r`$. Let $`v(\xi )`$ be its trivialization, i.e. a linear map $`v(\xi ):^rVVW`$ smoothly depending on $`\xi `$ such that $`𝒟_\xi v(\xi )=0`$ for all $`\xi `$, and $`\rho (\xi )=v(\xi )^{}v(\xi )`$ is an isomorphism for all $`\xi `$. We set
$$A(\xi )=\rho (\xi )^1v(\xi )^{}dv(\xi ).$$
The matrix-valued one-form $`A`$ is a connection on a trivial unitary bundle of rank $`r`$. One can show that its curvature $`F_A`$ is ASD (see ). However, it does not satisfy $`A^{}=A`$, because we are not using a unitary gauge. Instead $`A`$ satisfies
$$A^{}(\xi )=(\rho (\xi )A(\xi )\rho (\xi )^1+\rho (\xi )d\rho (\xi )^1).$$
To go to a unitary gauge, we must make a gauge transformation
$$A^{}(\xi )=g(\xi )A(\xi )g(\xi )^1+g(\xi )dg(\xi )^1,$$
where $`g(\xi )`$ is a function taking values in Hermitean $`r\times r`$ matrices and satisfying $`g(\xi )^2=\rho (\xi ).`$
We now explain, following , how to modify the ADHM construction so that it produces rank $`r`$ instantons on the noncommutative $`^4`$ defined in the previous section. It proves convenient to apply an orthogonal transformation which brings the matrix $`\theta `$ in (3) to the standard form
$$\theta =\sqrt{1}\left(\begin{array}{cccc}0& a& 0& 0\\ a& 0& 0& 0\\ 0& 0& 0& b\\ 0& 0& b& 0\end{array}\right).$$
We will assume that $`a+b0.`$ Since $`\theta `$ enters only in the combination $`\mathrm{}\theta ,`$ we can set $`a+b=1`$ without loss of generality. The relation between the affine coordinates $`\xi _1,\xi _2`$ on $`^2`$ and affine coordinates $`x_1,x_2,x_3,x_4`$ on $`^4`$ is chosen as follows:
$$\xi _1=x_4\sqrt{1}x_3,\xi _2=x_2+\sqrt{1}x_1.$$
Then $`\xi _1,\xi _2,\overline{\xi }_1,\overline{\xi }_2`$ obey the Weyl algebra relations
$$[\xi _1,\overline{\xi }_1]=2\mathrm{}b,[\xi _2,\overline{\xi }_2]=2\mathrm{}a,[\xi _1,\xi _2]=[\xi _1,\overline{\xi }_2]=[\overline{\xi }_1,\xi _2]=[\overline{\xi }_1,\overline{\xi }_2]=0.$$
The modified ADHM data consist of the same four maps which now satisfy
$$\mu _c=0,\mu _r=2\mathrm{}(a+b)1_{k\times k}.$$
The instanton connection is given by essentially the same formulas as in the commutative case. The operator $`𝒟`$ is given by the same formula as $`𝒟_\xi ,`$ but is now regarded as an element of
$$\mathrm{Hom}_𝒜((VVW)_{}𝒜,(VV)_{}𝒜).$$
The module $`𝐊𝐞𝐫𝒟`$ is a projective module over $`𝒜`$. Following , we assume that it is isomorphic to a free module of rank $`r`$, and $`v`$ is the corresponding isomorphism $`v:𝒜^r𝐊𝐞𝐫𝒟.`$ We further assume that the morphism
$$\mathrm{\Delta }=𝒟𝒟^{}\mathrm{End}_𝒜((VV)𝒜)$$
is an isomorphism.<sup>3</sup><sup>3</sup>3One can show that the latter assumption is always valid provided $`\mathrm{}0.`$ As for the former one, it is not known what constraints should be imposed on the deformed ADHM data to ensure that $`𝐊𝐞𝐫𝒟`$ is a free $`𝒜`$–module of rank $`r.`$ For $`r=1`$ $`𝐊𝐞𝐫𝒟`$ is never free . Then it is easy to see that $`\rho =v^{}v\mathrm{End}_𝒜(^r𝒜)`$ is an isomorphism too. We set
(5)
$$A=\rho ^1v^{}dv.$$
(The multiplication here and below is understood to be the Wigner-Moyal multiplication.) This formula defines a connection 1-form on a trivial unitary bundle on $`𝒜`$ of rank $`r`$. The curvature of this connection is given by
$$F_A=\rho ^1dv^{}(1v\rho ^1v^{})dv.$$
A short computation (essentially the same as in the commutative case) shows that the curvature can be written in the form
$$F_A=\rho ^1v^{}d𝒟^{}\mathrm{\Delta }^1d𝒟v.$$
Furthermore, since $`𝒟`$ and $`𝒟^{}`$ are linear in $`\xi _i,\overline{\xi }_i,`$ their exterior derivatives have a very simple form:
$$d𝒟=\left(\begin{array}{ccc}d\xi _1& d\xi _2& 0\\ d\overline{\xi }_2& d\overline{\xi }_1& 0\end{array}\right),d𝒟^{}=\left(\begin{array}{cc}d\overline{\xi }_1& d\xi _2\\ d\overline{\xi }_2& d\xi _1\\ 0& 0\end{array}\right).$$
Note also that by virtue of the deformed ADHM equations $`\mathrm{\Delta }`$ has a block-diagonal form:
$$\mathrm{\Delta }=\left(\begin{array}{cc}\delta & 0\\ 0& \delta \end{array}\right),$$
where $`\delta \mathrm{End}_𝒜(V𝒜)`$ is an isomorphism. Using this fact, one can easily see that $`F_A`$ is proportional to the 2-forms
$$d\xi _1d\overline{\xi }_1+d\xi _2d\overline{\xi }_2,d\xi _1d\overline{\xi }_2,d\xi _2d\overline{\xi }_1,$$
which are anti-self-dual.
As in the commutative case, the connection $`A`$ does not satisfy $`A^{}=A`$. To go to a unitary gauge one has to perform a gauge transformation
$$A^{}=gAg^1+gdg^1.$$
Here $`g\mathrm{Aut}_𝒜(^r𝒜)`$ should be found from the conditions $`g^{}=g,`$ $`gg=\rho .`$ The existence of such $`g`$ is an additional assumption.
### 2.2. Summary of results
In the commutative case there is a one-to-one correspondence between the following three classes of objects:
$`A.`$ Rank $`r`$ holomorphic bundles on $`^2`$ with $`c_2=k`$ and a fixed trivialization on the line at infinity.
$`B.`$ The set of ADHM data modulo the natural action of $`U(k)`$.
$`C.`$ Rank $`r`$ holomorphic bundles on $`^3`$ with $`c_2=k,`$ a trivialization on a fixed line, vanishing $`H^1(E(2)),`$ and satisfying a certain reality condition.
$`D.`$ $`U(r)`$ instantons on $`^4`$ with $`c_2=k`$.
The correspondence between $`C`$ and $`D`$ is a particular instance of twistor transform . The correspondence between $`B`$ and $`C`$ has been proved by Atiyah, Hitchin, Drinfeld, and Manin . Together these two results imply that all instantons on $`^4`$ arise from the ADHM construction. The correspondence between $`A`$ and $`B`$ has been proved by Donaldson . One can also prove the correspondence between $`A`$ and $`D`$ directly .
The goal of this paper is to extend some of these results to the noncommutative case. We show that there is a natural one-to-one correspondence between the isomorphism classes of the following objects:
$`A^{}.`$ Algebraic bundles on a noncommutative deformation of $`^2`$ with $`c_2=k`$ and a fixed trivialization on the line at infinity.
$`B^{}.`$ Deformed ADHM data of Nekrasov and Schwarz modulo the natural $`U(k)`$ action.
$`C^{}.`$ Certain complexes of sheaves on a noncommutative deformation of $`^3`$ satisfying reality conditions.
The moduli space of the deformed ADHM data has a natural hyperkähler metric, and the other two moduli spaces inherit this metric.
Furthermore, we reinterpret the deformed ADHM construction of Nekrasov and Schwarz in terms of a noncommutative deformation of the twistor transform.
It is interesting to note that H. Nakajima studied the same linear algebra data as Nekrasov and Schwarz and showed that their moduli space coincides with the moduli space of torsion free sheaves on a commutative $`^2`$ with a trivialization on a fixed line. On the other hand, we show that the same data describe algebraic bundles on a noncommutative $`^2`$. As shown below, the interpretation in terms of complexes of sheaves on a noncommutative $`^3`$ provides a geometric reason for this “coincidence.” We prove that the two moduli spaces are isomorphic as hyperkähler manifolds, but the natural complex structures on them differ by an $`\mathrm{SO}(3)`$ rotation.
The rest of the paper is organized as follows. In Section 3 we define noncommutative deformations of certain commutative projective varieties ($`^2,`$ $`^3,`$ and a quadric in $`^5`$). Section 4 is an algebraic preparation for the study of bundles on noncommutative projective spaces. In Section 5 we study the cohomological properties of sheaves on noncommutative $`^2`$ and $`^3`$ and define locally free sheaves (i.e. bundles). In Section 6 we show that any bundle on a noncommutative $`^2`$ trivial on the commutative line at infinity arises as a cohomology of a monad. In Section 7 we exhibit bijections between $`A^{},`$ $`B^{},`$ and $`C^{}`$ and explain the relation with Nakajima’s results. In Section 8 we construct a noncommutative deformation of Grassmannians and flag manifolds and describe a noncommutative version of the twistor transform. We also describe a nice class of noncommutative projective varieties associated with a Yang-Baxter operator and define differential forms on these varieties. In section 9 we consider a more general deformation of $`^4`$ (a $`q`$– deformed $`^4`$) whose physical significance is obscure at present. We propose an ADHM–like construction of instantons on this space and outline its relation to noncommutative algebraic geometry. In the Appendix we define the Wigner-Moyal product on the space of $`C^{\mathrm{}}`$ functions on $`^n`$ all of whose derivatives are polynomially bounded, and prove that the Wigner-Moyal product provides this space with a structure of an algebra over $`.`$
## 3. Geometry of noncommutative varieties
### 3.1. Algebraic preliminaries
Let $`𝚔`$ be a base field (we will be dealing only with $`𝚔=`$ or $`𝚔=`$ in this paper). Let $`A`$ be an algebra over $`𝚔`$. It is called right (left) noetherian if every right (left) ideal is finitely generated, and it is called noetherian if it is both right and left noetherian.
Let $`A=\underset{i0}{}A_i`$ be a graded noetherian algebra. We denote by $`mod(A)`$ the category of finitely generated right $`A`$–modules, by $`gr(A)`$ the category of finitely generated graded right $`A`$–modules, and by $`tors(A)`$ the full subcategory of $`gr(A)`$ which consists of finite dimensional graded $`A`$–modules.
An important role will be played by the quotient category $`qgr(A)=gr(A)/tors(A)`$. It has the following explicit description. The objects of $`qgr(A)`$ are the objects of $`gr(A)`$ (we denote by $`\stackrel{~}{M}`$ the object in $`qgr(A)`$ which corresponds to a module $`M`$). The morphisms in $`qgr(A)`$ are given by
$$\mathrm{Hom}_{qgr}(\stackrel{~}{M},\stackrel{~}{N})=\underset{\stackrel{}{M^{}}}{lim}\mathrm{Hom}_{gr}(M^{},N)$$
where $`M^{}`$ runs over submodules of $`M`$ such that $`M/M^{}`$ is finite dimensional.
On the category $`gr(A)`$ there is a shift functor: for a given graded module $`M=_{i0}M_i`$ the shifted module $`M(r)`$ is defined by $`M(r)_i=M_{r+i}`$. The induced shift functor on the quotient category $`qgr(A)`$ sends $`\stackrel{~}{M}`$ to $`\stackrel{~}{M}(r)=\stackrel{~}{M(r)}`$.
Similarly, we can consider the category $`Gr(A)`$ of all graded right $`A`$–modules. It contains the subcategory $`Tors(A)`$ of torsion modules. Recall that a module $`M`$ is called torsion if for any element $`xM`$ one has $`xA_s=0`$ for some $`s,`$ where $`A_s=\underset{is}{}A_i`$. We denote by $`QGr(A)`$ the quotient category $`Gr(A)/Tors(A)`$. The category $`QGr(A)`$ contains $`qgr(A)`$ as a full subcategory. Sometimes it is convenient to work in $`QGr(A)`$ instead of $`qgr(A)`$.
Henceforth, all graded algebras will be noetherian algebras generated by the first component $`A_1`$ with $`A_0=𝚔`$.
Sometimes we use subscripts $`R`$ or $`L`$ for categories $`gr(A),qgr(A),`$ etc., to specify whether right or left modules are considered. If the subscript is omitted, the modules are taken to be right modules. For the same reason for an $`A`$–bimodule $`M`$ we sometimes write $`M_A`$ or $`{}_{A}{}^{}M`$ to specify whether the right or left module structure is considered.
### 3.2. Noncommutative varieties
A variety in commutative geometry is a topological space with a sheaf of functions (continuous, smooth, analytic, algebraic, etc.) which is, obviously, a sheaf of algebras. One of the main objects in geometry (algebraic or differential) is a bundle or, more generally, a sheaf. To any variety $`X`$ we can associate an abelian category of sheaves of modules (maybe with some additional properties) over the sheaf of algebras of functions. Given a sheaf of modules on $`X,`$ the space of its global sections is a module over the algebra of global functions on $`X.`$ Thus the functor of global sections associates to every $`X`$ an algebra and a certain category of modules over it. Under favorable circumstances, much of the information about the geometry of $`X`$ is contained in this purely algebraic datum. Let us give a few examples.
If $`X`$ is a compact Hausdorff topological space, then the category of vector bundles over $`X`$ is equivalent to the category of finitely generated projective modules over the algebra of continuous functions on $`X`$ . The equivalence is given by the functor which maps a vector bundle to the module of its global sections.
It is well known that if $`A`$ is a commutative noetherian algebra, the category of coherent sheaves on the noetherian affine scheme $`Spec(A)`$ is equivalent to the category of finitely generated modules over $`A`$. The equivalence is again given by the functor which attaches to a coherent sheaf the module of its global sections.
In the case of projective varieties the only global functions are constants, so one has to act somewhat differently. Since a projective variety $`X`$ is by definition a subvariety of a projective space, it inherits from it the line bundle $`𝒪_X(1)`$ and its tensor powers $`𝒪_X(i)`$. We can consider a graded algebra
$$\mathrm{\Gamma }(X)=\underset{i0}{}H^0(X,𝒪_X(i)).$$
This algebra is called the homogeneous coordinate algebra of $`X`$. Furthermore, for any sheaf $``$ we can define a graded $`A`$–module
$$\mathrm{\Gamma }()=\underset{i0}{}H^0(X,(i)).$$
It can be checked that $`\mathrm{\Gamma }`$ is a functor from the category of coherent sheaves on $`X`$ $`coh(X)`$ to $`gr(\mathrm{\Gamma }(X))`$. In a brilliant paper , J-P. Serre described the category of coherent sheaves on a projective scheme $`X`$ in terms of graded modules over the graded algebra $`\mathrm{\Gamma }(X)`$. He proved that the category $`coh(X)`$ is equivalent to the quotient category $`qgr(\mathrm{\Gamma }(X))=gr(\mathrm{\Gamma }(X))/tors(\mathrm{\Gamma }(X))`$. The equivalence is given by the composition of the functor $`\mathrm{\Gamma }`$ with the projection $`\pi :gr(A)qgr(A)`$. On other hand, let $`A=\underset{i0}{}A_i`$ be a graded commutative algebra generated over $`𝚔`$ by the first component (which is assumed to be finite dimensional). We can associate to $`A`$ a projective scheme $`X=Proj(A)`$. Serre proved that the category $`coh(X)`$ is equivalent to the category $`qgr(A)`$.
The equivalence also holds for the category of quasicoherent sheaves on $`X`$ and the category $`QGr(A)=Gr(A)/Tors(A)`$.
In all of the above examples it turned out that the natural category of sheaves or bundles on a variety is equivalent to a certain category defined in terms of (graded) modules over some (graded) algebra. On the other hand, “as A. Grothendieck taught us, to do geometry you really don’t need a space, all you need is a category of sheaves on this would-be space” (, p.83).
For this reason, in the realm of algebraic geometry it is natural to regard a noncommutative noetherian algebra as a coordinate algebra of a noncommutative affine variety; then the category of finitely generated right modules over this algebra is identified with the category of coherent sheaves on the corresponding variety. Similarly, a noncommutative graded noetherian algebra is regarded as a homogeneous coordinate algebra of a noncommutative projective variety. The category of finitely generated graded right modules over this algebra modulo torsion modules is identified with the category of coherent sheaves on this variety (see , , ).
A different approach to noncommutative geometry has been pursued by A. Connes .
### 3.3. Noncommutative deformations of commutative varieties
Many important noncommutative varieties arise as deformations of commutative ones.
Let $`X`$ be a commutative variety (affine or projective). Let $`A`$ be the corresponding commutative (graded) algebra. A noncommutative deformation of $`X`$ is a deformation of the algebra structure on $`A,`$ that is, a deformation of the multiplication law. Usually it is not easy to write down an explicit formula for the deformed product.
There is a more algebraic way to describe noncommutative deformations of commutative varieties. Assume that the algebra $`A`$ is given in terms of generators and relations. This means that $`A`$ is given as a quotient $`A=T(V)/R,`$ where $`V`$ is the vector space spanned by the generators, $`T(V)`$ is the tensor algebra of $`V,`$ and $`R`$ is a two-sided ideal in $`T(V)`$ generated by a subspace of relations $`RT(V)`$. Assume that $`R_{\mathrm{}}T(V)`$ is a one-parameter deformation of the subspace $`R`$. Then $`A_{\mathrm{}}=T(V)/R_{\mathrm{}}`$ is a one-parameter deformation of $`A`$. (If $`A`$ is graded, then we assume that $`R`$ is a graded subspace, and the deformation preserves the grading).
We denote by $`X_{\mathrm{}}`$ the noncommutative variety corresponding to the algebra $`A_{\mathrm{}}`$. Thus $`X_{\mathrm{}}`$ is a noncommutative one-parameter deformation of $`X`$.
If $`X`$ is projective and $`A`$ is a graded algebra, then we denote by $`coh(X_{\mathrm{}})`$ the category $`qgr(A_{\mathrm{}})`$. Furthermore, as in the commutative case, we will write $`𝒪(r)`$ for the object $`\stackrel{~}{A_{\mathrm{}}}(r)`$.
Now we define noncommutative varieties which are going to be used in this paper.
### 3.4. Noncommutative $`^4`$
Denote by $`A(^4)`$ the algebra of polynomial functions on $`^4`$. Let $`\theta `$ be a skew-symmetric $`4\times 4`$ matrix.
Let us define the algebra $`A(_{\mathrm{}}^4)`$ as an algebra over $``$ generated by $`x_i`$ ($`i=1,2,3,4`$) with relations $`[x_i,x_j]=\mathrm{}\theta _{ij}`$:
(6)
$$A(_{\mathrm{}}^4)=\mathrm{T}(x_1,x_2,x_3,x_4)/[x_i,x_j]=\mathrm{}\theta _{ij}_{1i,j4}.$$
We will regard $`A(_{\mathrm{}}^4)`$ as the algebra of polynomial functions on a noncommutative affine variety $`_{\mathrm{}}^4`$.
### 3.5. Noncommutative 4-dimensional quadric
Let $`G`$ be a $`4\times 4`$ symmetric nondegenerate matrix. Consider a graded algebra $`Q_{\mathrm{}}=\underset{i0}{}Q_i`$ over $``$ generated by the elements $`X_1,X_2,X_3,X_4,D,T`$ of degree $`1`$ with the following quadratic relations:
(7)
$$\begin{array}{cc}[T,D]=[T,X_i]=0,\hfill & \\ [X_i,X_j]=\mathrm{}\theta _{ij}T^2,\hfill & \\ [D,X_i]=2\mathrm{}\underset{lk}{}\theta _{il}G^{lk}X_kT,\hfill & \\ \underset{ij}{}G^{ij}X_iX_j=DT.\hfill & \end{array}$$
We denote by $`_{\mathrm{}}^4`$ the noncommutative projective variety corresponding to the algebra $`Q_{\mathrm{}}`$. It is evident that $`_{\mathrm{}}^4`$ is a deformation of a 4-dimensional commutative quadric $`^4=\{_{ij}G^{ij}X_iX_j=DT\}^5`$.
### 3.6. Embedding $`_{\mathrm{}}^4_{\mathrm{}}^4`$
Let $`Q_{\mathrm{}}[T^1]`$ be the localization of the algebra $`Q_{\mathrm{}}`$ with respect to $`T`$. Elements of degree $`0`$ in $`Q_{\mathrm{}}[T^1]`$ form a subalgebra which will be denoted by $`Q_{\mathrm{}}[T^1]_0`$.
###### Lemma 3.1.
The map $`x_iT^1X_i`$ $`(i=1,2,3,4)`$ induces an isomorphism of the algebra $`A(_{\mathrm{}}^4)`$ with the algebra $`Q_{\mathrm{}}[T^1]_0`$.
###### Proof.
Obvious. ∎
This means that $`_{\mathrm{}}^4`$ can be identified with the open subset $`\{T0\}`$ in $`_{\mathrm{}}^4`$. For this reason, $`_{\mathrm{}}^4`$ may be regarded as a compactification of $`_{\mathrm{}}^4`$ which is compatible with the bilinear form $`G`$. Note also that the complement of $`_{\mathrm{}}^4`$ in $`_{\mathrm{}}^4`$ corresponds to the algebra
$$Q_{\mathrm{}}/T=\mathrm{T}(X_1,X_2,X_3,X_4,D)/[X_i,X_j]=[D,X_i]=0,\underset{ij}{}G^{ij}X_iX_j=0.$$
Since this algebra is commutative, the complement is the usual 3-dimensional commutative quadratic cone. Thus one may say that $`_{\mathrm{}}^4`$ is obtained from $`_{\mathrm{}}^4`$ by adding a cone “at infinity”. This is in complete analogy with the commutative case.
### 3.7. Noncommutative $`_{\mathrm{}}^2`$ and $`_{\mathrm{}}^3`$
Noncommutative deformations of the projective plane have been classified in , , . We will need one of them, namely the one whose homogeneous coordinate algebra is a graded algebra $`PP_{\mathrm{}}=\underset{i0}{}PP_{\mathrm{}}^{}{}_{i}{}^{}`$ over $``$ generated by the elements $`w_1,w_2,w_3`$ of degree 1 with the relations:
(8)
$$\begin{array}{c}[w_3,w_i]=0\text{ for any }i=1,2,3,\hfill \\ [w_1,w_2]=2\mathrm{}w_3^2.\hfill \end{array}$$
We will also need a noncommutative deformation of the 3–dimensional projective space, whose homogeneous coordinate algebra will be denoted $`PS_{\mathrm{}}=\underset{i0}{}PS_{\mathrm{}}^{}{}_{i}{}^{}`$. It is a graded algebra over $``$ generated by $`PS_{\mathrm{}}^{}{}_{1}{}^{}=U,`$ where the vector space $`U`$ is spanned by elements $`z_1,z_2,z_3,z_4`$ obeying the relations
(9)
$$\begin{array}{c}[z_3,z_i]=[z_4,z_i]=0\text{ for any }i=1,2,3,4,\hfill \\ [z_1,z_2]=2\mathrm{}z_3z_4.\hfill \end{array}$$
The noncommutative projective varieties corresponding to $`PP_{\mathrm{}}`$ and $`PS_{\mathrm{}}`$ will be denoted $`_{\mathrm{}}^2`$ and $`_{\mathrm{}}^3,`$ respectively.
Note that for $`\mathrm{}0`$ all algebras $`PS_{\mathrm{}}`$ are isomorphic, and therefore the varieties $`_{\mathrm{}}^3`$ are the same for all $`\mathrm{}0`$. The same is true for $`_{\mathrm{}}^2.`$
### 3.8. Subvarieties in $`_{\mathrm{}}^3`$ and $`_{\mathrm{}}^2`$
If $`IPS_{\mathrm{}}`$ is a graded two-sided ideal, then the quotient algebra $`PS_{\mathrm{}}/I`$ corresponds to a closed subvariety $`X(I)_{\mathrm{}}^3`$. Let us describe some of them.
Let $`J`$ be the graded two-sided ideal generated by $`z_3`$ and $`z_4`$. Then $`PS_{\mathrm{}}/J=\mathrm{T}(z_1,z_2)/[z_1,z_2]=0,`$ hence $`X(J)`$ is the commutative projective line.
For each point $`p=(\lambda :\mu )^1`$ let $`J_p`$ denote the graded two-sided ideal generated by $`\lambda z_3+\mu z_4`$. If $`p=(0:1)`$ or $`p=(1:0),`$ then it is easy to see that $`X(J_p)`$ is the commutative projective plane. For all other $`p^1`$ we have
$$PS_{\mathrm{}}/J_p=\mathrm{T}(z_1,z_2,z_3)/[z_1,z_3]=[z_2,z_3]=0,[z_1,z_2]=2\mathrm{}\frac{\lambda }{\mu }z_3^2,$$
hence $`X(J_p)`$ is a noncommutative projective plane isomorphic to $`_{\mathrm{}}^2`$.
We have $`J_pJ`$ for all $`p^1,`$ hence all planes $`X(J_p)`$ pass through the line $`X(J)`$. Thus we see that $`_{\mathrm{}}^3`$ is a pencil of noncommutative projective planes passing through a fixed commutative projective line.
Similarly, the two-sided ideal generated by $`w_3`$ in $`PP_{\mathrm{}}`$ corresponds to a commutative projective line $`l=\{w_3=0\}_{\mathrm{}}^2`$.
## 4. Properties of algebras $`PS_{\mathrm{}}`$ and $`PP_{\mathrm{}}`$ and the resolution of the diagonal
This section is a preparation for the study of sheaves on $`_{\mathrm{}}^3`$ and $`_{\mathrm{}}^2`$. We show that the algebras $`PS_{\mathrm{}}`$ and $`PP_{\mathrm{}}`$ are regular and Koszul and construct the resolution of the diagonal, which will enable us to associate monads to certain bundles on $`_{\mathrm{}}^2`$.
### 4.1. Quadratic algebras
A graded algebra $`A=\underset{i0}{}A_i`$ over a field $`𝚔`$ is called quadratic if it is connected (i.e. $`A_0=𝚔`$), is generated by the first component $`A_1,`$ and the ideal of relations is generated by the subspace of quadratic relations $`R(A)A_1A_1`$.
Therefore the algebra $`A`$ can be represented as $`T(A_1)/R(A),`$ where $`T(A_1)`$ is a free tensor algebra generated by the space $`A_1`$.
The algebras $`PS_{\mathrm{}}`$ and $`PP_{\mathrm{}}`$ are quadratic algebras. For example, $`PS_{\mathrm{}}`$ can be represented as $`\mathrm{T}(U)/W,`$ where $`U=PS_{\mathrm{}}^{}{}_{1}{}^{}`$ is a 4-dimensional vector space and $`W`$ is the 6–dimensional subspace of $`UU`$ spanned by the relations (9).
### 4.2. The dual algebra
For any quadratic algebra $`A=T(A_1)/R(A)`$ we can define its dual algebra which is also quadratic.
Let us identify $`A_1^{}A_1^{}`$ with $`(A_1A_1)^{}`$ by $`(lm)(ab)=m(a)l(b)`$. Denote by $`R(A)^{}`$ the annulator of $`R(A)`$ in $`A_1^{}A_1^{},`$ i.e. the subspace which consists of such $`q(A_1^{})^2`$ that $`q(r)=0`$ for any $`rR(A)`$.
###### Definition 4.1.
() The algebra $`A^!=T(A_1^{})/R(A)^{}`$ is called the dual algebra of $`A`$.
###### Example 4.2.
Let $`\{\stackrel{ˇ}{z}_i\},i=1,2,3,4,`$ be the basis of $`PS_{\mathrm{}}^{}{}_{1}{}^{!}=U^{}`$ which is dual to $`\{z_i\}`$. By definition, $`PS_{\mathrm{}}^{}{}_{}{}^{!}`$ is generated by $`\{\stackrel{ˇ}{z}_i\}`$ with defining relations
$$\begin{array}{c}\stackrel{ˇ}{z}_i^2=0\text{ for all}i=1,\mathrm{},4;\hfill \\ \stackrel{ˇ}{z}_i\stackrel{ˇ}{z}_j+\stackrel{ˇ}{z}_j\stackrel{ˇ}{z}_i=0\text{ for all}i<j,(i,j)(3,4);\hfill \\ \stackrel{ˇ}{z}_3\stackrel{ˇ}{z}_4+\stackrel{ˇ}{z}_4\stackrel{ˇ}{z}_3=\mathrm{}[\stackrel{ˇ}{z}_1,\stackrel{ˇ}{z}_2]=2\mathrm{}\stackrel{ˇ}{z}_1\stackrel{ˇ}{z}_2.\hfill \end{array}$$
In the commutative case the dual algebra of the symmetric algebra $`S^{}(U)`$ is isomorphic to the exterior algebra $`\mathrm{\Lambda }^{}(U^{})`$. Obviously, the algebras $`PS_{\mathrm{}}^{}{}_{}{}^{!}`$ and $`PP_{\mathrm{}}^{}{}_{}{}^{!}`$ are deformations of exterior algebras. For example, the vector space $`PS_{\mathrm{}}^{}{}_{k}{}^{!}`$ is spanned by the elements $`\stackrel{ˇ}{z}_{i_1}\mathrm{}\stackrel{ˇ}{z}_{i_k}`$ with $`i_1<\mathrm{}<i_k`$. In particular, the dimension of the vector space $`PS_{\mathrm{}}^{}{}_{k}{}^{!}`$ is equal to $`\left(\genfrac{}{}{0pt}{}{4}{k}\right)`$. Similarly, the dimension of $`PP_{\mathrm{}}^{}{}_{k}{}^{!}`$ is equal to $`\left(\genfrac{}{}{0pt}{}{3}{k}\right)`$.
###### Proposition 4.3.
Let $`A`$ be $`PS_{\mathrm{}}`$ or $`PP_{\mathrm{}},`$ and let $`n`$ be $`4`$ or $`3,`$ respectively. The multiplication map $`A_k^!A_{nk}^!A_n^!`$ is a non-degenerate pairing. Hence the dual algebra $`A^!`$ is a Frobenius algebra, i.e. $`(A^!)_{A^!}({}_{A^!}{}^{}A_{}^{!})^{}`$ as right $`A^!`$–modules.
###### Proof.
The proposition holds for the exterior algebra, and therefore also for the algebra $`A^!,`$ since the latter is a “small” deformation of the exterior algebra. ∎
### 4.3. The Koszul complex
Consider right $`A`$–modules $`(A_k^!)^{}A`$. The following complex $`K_{}(A)`$ is called the (right) Koszul complex of a quadratic algebra:
$$\mathrm{}\stackrel{d}{}(A_3^!)^{}A(3)\stackrel{d}{}(A_2^!)^{}A(2)\stackrel{d}{}(A_1^!)^{}A(1)\stackrel{d}{}(A_0^!)^{}A0,$$
where the map $`d:(A_k^!)^{}A(A_{k1}^!)^{}A`$ is a composition of two natural maps:
$$(A_k^!)^{}A(A_k^!)^{}A_1^!A_1A(A_k^!)^{}A.$$
Here the first arrow sends $`\alpha a`$ to $`\alpha ea`$ with $`e`$ defined as
$$e=\underset{i}{}y_ix_iA_1^!A_1,$$
and $`\{x_i\}`$ and $`\{y_i\}`$ being the dual bases of $`A_1`$ and $`A_1^!,`$ respectively. The second map is determined by the algebra structures on $`A^!`$ and $`A`$.
It is a well–known fact that $`d^2=0`$ (see, for example, ).
Let $`𝚔_A`$ be the trivial right $`A`$-module. The Koszul complex $`K_{}(A)`$ possesses a natural augmentation $`K_{}\stackrel{\epsilon }{}𝚔_A0`$.
###### Definition 4.4.
(see ) A quadratic algebra $`A=\underset{i0}{}A_i`$ is called a Koszul algebra if the augmented Koszul complex $`K_{}(A)\stackrel{\epsilon }{}𝚔_A0`$ is exact.
In the same manner one can define the left Koszul complex of a quadratic algebra. It is well known that the exactness of the right Koszul complex is equivalent to the exactness of the left Koszul complex (see, for example, ).
###### Proposition 4.5.
The algebras $`PS_{\mathrm{}}`$ and $`PP_{\mathrm{}}`$ are Koszul algebras.
###### Proof.
For $`\mathrm{}=0`$ this is a well-known fact about the symmetric algebra $`S^{}(U)`$. Since the augmented Koszul complex is exact for $`\mathrm{}=0,`$ it is also exact for small $`\mathrm{},`$ and consequently for all $`\mathrm{}`$. ∎
Since the dual algebras $`PS_{\mathrm{}}^{}{}_{}{}^{!}`$ and $`PP_{\mathrm{}}^{}{}_{}{}^{!}`$ are finite, the Koszul resolutions for the algebras $`PS_{\mathrm{}}`$ and $`PP_{\mathrm{}}`$ are finite too and have the same form as the resolutions for ordinary symmetric algebras. For example, the Koszul resolution for $`A=PP_{\mathrm{}}`$ is:
$$\{0(A_3^!)^{}A(3)(A_2^!)^{}A(2)(A_1^!)^{}A(1)(A_0^!)^{}A\}.$$
### 4.4. Resolution of the diagonal
Consider a bigraded vector space
$$K_{}^2(A)=\underset{k,l0}{}K_{k,l}^2(A)\text{with}K_{k,l}^2(A)=A(k)(A_{lk}^!)^{}A(l).$$
Consider morphisms $`d_R:K_{k,l}^2K_{k,l1}^2`$ and $`d_L:K_{k,l}^2K_{k+1,l}^2`$ given by the following compositions
$$\begin{array}{c}d_R:A(A_k^!)^{}AA(A_k^!)^{}A_1^!A_1AA(A_{k1}^!)^{}A,\hfill \\ d_L:A(A_k^!)^{}AAA_1A_1^!(A_k^!)^{}AA(A_{k1}^!)^{}A.\hfill \end{array}$$
Here the leftmost maps are given by
$$e_R=\underset{i}{}y_ix_iA_1^!A_1\text{and}e_L=\underset{i}{}x_iy_iA_1A_1^!,$$
where $`\{x_i\}`$ and $`\{y_i\}`$ are the dual bases of $`A_1`$ and $`A_1^!,`$ respectively, while the rightmost maps are induced by the algebra structures of $`A^!`$ and $`A`$. It is easy to show that
$$d_R^2=d_L^2=0\text{and}d_Rd_L=d_Ld_R,$$
hence $`K_{}^2(A)`$ is a bicomplex. It is called the double Koszul bicomplex of the quadratic algebra $`A`$.
The topmost part of the bicomplex looks as follows:
$$\begin{array}{ccccccc}\mathrm{}& \stackrel{d_R}{}& A(A_{l+1}^!)^{}A(1l)& \stackrel{d_R}{}& A(A_l^!)^{}A(l)& \stackrel{d_R}{}& \mathrm{}\\ & & d_L& & d_L& & \\ \mathrm{}& \stackrel{d_R}{}& A(1)(A_l^!)^{}A(1l)& \stackrel{d_R}{}& A(1)(A_{l1}^!)^{}A(l)& \stackrel{d_R}{}& \mathrm{}\end{array}$$
Each term of the bicomplex $`K_{}^2(A)`$ has an obvious structure of a bigraded $`A`$-bimodule, and it is clear that the differentials are morphisms of bigraded $`A`$-bimodules.
Let
$$𝒦_l(A)=\mathrm{Ker}d_L:K_{0,l}^2(A)K_{1,l}^2(A).$$
Then $`𝒦_{}(A)`$ is a complex of bigraded $`A`$-bimodules (with respect to the differential $`d_R`$).
Consider a bigraded algebra $`\mathrm{\Delta }=_{i,j}\mathrm{\Delta }_{ij}`$ with $`\mathrm{\Delta }_{ij}=A_{i+j}`$ and with the multiplication induced from $`A`$. The algebra $`\mathrm{\Delta }`$ is called the diagonal bigraded algebra of $`A`$. Note that the multiplication map induces a surjective morphism of $`A`$-bimodules $`\delta :AA\mathrm{\Delta }`$.
###### Lemma 4.6.
The map
$$\delta :𝒦_0(A)=AA\mathrm{\Delta }$$
gives an augmentation of the complex $`𝒦_{}(A)`$.
###### Proof.
We have to check that $`\delta d_R:𝒦_1(A)A`$ vanishes. Note that $`K_{0,1}^2(A)=AA_1A(1),`$ and that the differentials $`d_R`$ and $`d_L`$ restricted to $`K_{0,1}^2(A)`$ coincide with the multiplication maps $`m_{1,2}`$ and $`m_{2,3},`$ respectively. Thus we have the following commutative diagram:
$$\begin{array}{ccccc}𝒦_1(A)& \stackrel{d_R}{}& 𝒦_0(A)& \stackrel{\delta }{}& \mathrm{\Delta }\\ & & & & & & \\ AA_1A(1)& \stackrel{m_{1,2}}{}& AA& \stackrel{\delta }{}& \mathrm{\Delta }\\ m_{2,3}& & \\ A(1)A(1)\end{array}$$
Now the Lemma follows because $`\delta m_{1,2}=\delta m_{2,3}`$ (associativity) obviously annihilates $`𝐊𝐞𝐫m_{2,3}=𝒦_1(A)`$. ∎
###### Proposition 4.7.
If $`A`$ is Koszul, then $`𝒦_{}(A)\stackrel{\delta }{}\mathrm{\Delta }`$ is exact.
###### Proof.
The $`(p,q)`$–bigraded component of $`K_{k,l}^2(A)`$ is equal to $`A_{p+k}(A_{lk}^!)^{}A_{ql},`$ hence the $`(p,q)`$– bigraded component of the bicomplex $`K_{}^2(A)`$ vanishes for $`l<k`$ or $`l>q`$. Thus the $`(p,q)`$–bigraded component of the bicomplex $`K_{}^2(A)`$ is bounded. Therefore both spectral sequences of the bicomplex $`K_{}^2(A)`$ converge to the cohomology of the total complex $`\mathrm{Tot}(K_{}^2(A))`$. The first term of the first spectral sequence reads
$$E_{k,l}^1=\{\begin{array}{cc}A(l)𝚔(l),\hfill & \text{if }k=l\hfill \\ 0,\hfill & \text{otherwise}\hfill \end{array}$$
Hence the spectral sequence degenerates in the first term, and we have
$$H^0(\mathrm{Tot}(K_{}^2(A)))=\underset{l=0}{\overset{\mathrm{}}{}}A(l)𝚔(l),H^0(\mathrm{Tot}(K_{}^2(A)))=0.$$
On the other hand, the first term of the second spectral sequence reads
$$E_{k,l}^1=\{\begin{array}{cc}𝚔(l)A(l),\hfill & \text{if }k=l>0\hfill \\ 𝒦_l(A),\hfill & \text{if }k=0\hfill \\ 0,\hfill & \text{otherwise}\hfill \end{array}$$
Hence the spectral sequence degenerates in the second term, and we have
$$H^0(\mathrm{Tot}(K_{}^2(A)))=H^0(𝒦_{}(A))\left(\underset{l=1}{\overset{\mathrm{}}{}}𝚔(l)A(l)\right),H^l(\mathrm{Tot}(K_{}^2(A)))=H^l(𝒦_{}(A)).$$
Therefore $`H^0(𝒦_{}(A))=0,`$ and we have an exact sequence
$$0H^0(𝒦_{}(A))\underset{l=0}{\overset{\mathrm{}}{}}A(l)𝚔(l)\underset{l=1}{\overset{\mathrm{}}{}}𝚔(l)A(l)0.$$
Looking at $`(p,q)`$–bigraded component of this sequence we see that
$$(H^0(𝒦_{}(A)))_{p,q}=\{\begin{array}{cc}A_{p+q},\hfill & \text{if }p,q0\hfill \\ 0,\hfill & \text{otherwise}\hfill \end{array}$$
Thus $`H^0(𝒦_{}(A))=\mathrm{\Delta }`$. ∎
###### Definition 4.8.
Define the left $`A`$–module $`\mathrm{\Omega }^k`$ as the cohomology of the left Koszul complex, truncated in the term $`K_k`$. In particular, $`\mathrm{\Omega }^1`$ is defined by the so-called Euler sequence
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$$0\mathrm{\Omega }^1A(1)A_1\stackrel{m}{}A\stackrel{\epsilon }{}𝚔0.$$
In section 8.11 we will show that for noncommutative projective spaces the sheaves corresponding to the modules $`\mathrm{\Omega }^k`$ can be regarded as sheaves of differential forms.
###### Proposition 4.9.
We have $`𝒦_k(A)=\mathrm{\Omega }^k(k)A(k)`$.
###### Proof.
This follows immediately from the definition of $`\mathrm{\Omega }^k`$ and $`𝒦_k(A)`$. ∎
Combining Propositions 4.7 and 4.9, we obtain the following resolution of the diagonal:
(11)
$$\mathrm{}\mathrm{\Omega }^2(2)A(2)\mathrm{\Omega }^1(1)A(1)AA\mathrm{\Delta }0.$$
### 4.5. Cohomological properties of the algebras $`PS_{\mathrm{}}`$ and $`PP_{\mathrm{}}`$
First we note that both algebras $`PS_{\mathrm{}}`$ and $`PP_{\mathrm{}}`$ are noetherian. This follows from the fact that they are Ore extensions of commutative polynomial algebras (see for example, ). For the same reason the algebras $`PS_{\mathrm{}}`$ and $`PP_{\mathrm{}}`$ have finite right (and left) global dimension, which is equal to 4 and 3, respectively (see , p. 273).
We remind that the global dimension of a ring $`A`$ is the minimal number $`n`$ (if it exists) such that for any two modules $`M`$ an $`N`$ we have $`\mathrm{Ext}_A^{n+1}(M,N)=0`$.
In the paper the notion of a regular algebra has been introduced. Regular algebras have many good properties (see , , , etc.).
###### Definition 4.10.
A graded algebra $`A`$ is called regular of dimension $`d`$ if it satisfies the following conditions:
| (1) | $`A`$ has global dimension $`d,`$ |
| --- | --- |
| (2) | $`A`$ has polynomial growth, i.e. $`dimA_ncn^\delta `$ for some $`c,\delta ,`$ |
| (3) | $`A`$ is Gorenstein, meaning that $`\mathrm{Ext}_A^i(𝚔,A)=0`$ if $`id,`$ |
| | and $`\mathrm{Ext}_A^d(𝚔,A)=𝚔(l)`$ for some $`l`$. |
Here $`\mathrm{Ext}_A`$ stands for the Ext functor in the category $`mod(A)`$.
It is easy to see that these properties are verified for $`PS_{\mathrm{}}`$ and $`PP_{\mathrm{}}`$. Property (2) holds because our algebras grow as ordinary polynomial algebras. Property (3) follows from the fact that $`PS_{\mathrm{}}`$ and $`PP_{\mathrm{}}`$ are Koszul algebras and the dual algebras are Frobenius resolutions. In this case the Gorenstein parameter $`l`$ in (3) is equal to the global dimension $`d`$. Thus we have
###### Proposition 4.11.
The algebras $`PS_{\mathrm{}}`$ and $`PP_{\mathrm{}}`$ are noetherian regular algebras of global dimension $`4`$ and $`3,`$ respectively. For these algebras the Gorenstein parameter $`l`$ coincides with the global dimension $`d`$.
## 5. Cohomological properties of sheaves on $`_{\mathrm{}}^2`$ and $`_{\mathrm{}}^3`$
### 5.1. Ampleness and cohomology of $`𝒪(i)`$
Let $`A`$ be a graded algebra and $`X`$ be the corresponding noncommutative projective variety. Consider the sequence of sheaves $`\{𝒪(i)\}_i`$ in the category $`coh(X)qgr(A),`$ where $`𝒪(i)=\stackrel{~}{A(i)}`$.
This sequence is called ample if the following conditions hold:
| (a) | For every coherent sheaf $``$ there are integers $`k_1,\mathrm{},k_s`$ and an epimorphism |
| --- | --- |
| | $`\underset{i=1}{\overset{s}{}}𝒪(k_i)`$. |
| (b) | For every epimorphism $`𝒢`$ the induced map |
| | $`\mathrm{Hom}(𝒪(n),)\mathrm{Hom}(𝒪(n),𝒢)`$ is surjective for $`n0`$. |
It is proved in that the sequence $`\{𝒪(i)\}`$ is ample in $`qgr(A)`$ for a graded right noetherian $`𝚔`$–algebra $`A`$ if it satisfies the extra condition:
$$(\chi _1):dim_𝚔\mathrm{Ext}_A^1(𝚔,M)<\mathrm{}$$
for any finitely generated graded $`A`$–module $`M`$.
This condition can be verified for all noetherian regular algebras (see , Theorem 8.1). In particular, the categories $`coh(_{\mathrm{}}^3),`$ $`coh(_{\mathrm{}}^2)`$ have ample sequences.
For any sheaf $`qgr(A)`$ we can define a graded module $`\mathrm{\Gamma }()`$ by the rule:
$$\mathrm{\Gamma }():=\underset{i0}{}\mathrm{Hom}(𝒪(i),)$$
It is proved in that for any noetherian algebra $`A`$ that satisfies the condition $`\chi _1`$ the correspondence $`\mathrm{\Gamma }`$ is a functor from $`qgr(A)`$ to $`gr(A)`$ and the composition of $`\mathrm{\Gamma }`$ with the natural projection $`\pi :gr(A)gqr(A)`$ is isomorphic to the identity functor (see , ch. 3,4).
Now we formulate a result about the cohomology of sheaves on noncommutative projective spaces. This result is proved in for a general regular algebra and parallels the commutative case.
###### Proposition 5.1.
(Theorem 8.1. ) Let $`A`$ be $`PS_{\mathrm{}}`$ or $`PP_{\mathrm{}},`$ and $`X`$ be $`_{\mathrm{}}^3`$ or $`_{\mathrm{}}^2,`$ respectively. Denote by $`n`$ the dimension of $`X`$ (in our case $`n=3`$ or $`n=2,`$ respectively). Then
1) The cohomological dimension of $`coh(X)`$ is equal to $`dim(X),`$ i.e. for any two coherent sheaves $``$ and $`𝒢`$ $`\mathrm{Ext}^i(,𝒢)`$ vanishes if $`i>n`$.
2) There are isomorphisms
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$$H^p(X,𝒪(i))=\{\begin{array}{cc}A_k\hfill & \text{for }p=0,i0\hfill \\ A_{i1n}^{}\hfill & \text{ for }p=n,in1\hfill \\ 0\hfill & \text{otherwise}\hfill \end{array}$$
This proposition and the ampleness of the sequence $`\{𝒪(i)\}`$ implies the following corollary:
###### Corollary 5.2.
Let $`X`$ be either $`_{\mathrm{}}^3`$ or $`_{\mathrm{}}^2`$. Then for any sheaf $`coh(X)`$ and for all sufficiently large $`i0`$ we have
$$\mathrm{Hom}(,𝒪(i))=0.$$
###### Proof.
By ampleness a sheaf $``$ can be covered by a finite sum of sheaves $`𝒪(k_j)`$. Now the statement follows from the Proposition, because $`\mathrm{Hom}(𝒪(k_j),𝒪(i))=0`$ for all $`i<k_j`$. ∎
###### Corollary 5.3.
Let $`X`$ be either $`_{\mathrm{}}^3`$ or $`_{\mathrm{}}^2`$. Then for any sheaf $`coh(X)`$ and for all sufficiently large $`i0`$ we have
$$H^k(X,(i))=0$$
for all $`k1`$.
###### Proof.
The group $`H^k(X,(i))`$ coincides with $`\mathrm{Ext}^k(𝒪(i),)`$. Let $`k`$ be the maximal integer (it exists because the global dimension is finite) such that for some $``$ there exists arbitrarily large $`i`$ such that $`\mathrm{Ext}^k(𝒪(i),)0`$. Assume that $`k1`$. Choose an epimorphism $`\underset{j=1}{\overset{s}{}}𝒪(k_j)`$. Let $`_1`$ denote its kernel. Then for $`i>\mathrm{max}\{k_j\}`$ we have $`\mathrm{Ext}^{>0}(𝒪(i),\underset{j=1}{\overset{s}{}}𝒪(k_j))=0,`$ hence $`\mathrm{Ext}^k(𝒪(i),)0`$ implies $`\mathrm{Ext}^{k+1}(𝒪(i),)0`$. This contradicts the assumption of the maximality of $`k`$. ∎
### 5.2. Serre duality and the dualizing sheaf
A very useful property of commutative smooth projective varieties is the existence of the dualizing sheaf. Recall that a sheaf $`\omega `$ is called dualizing if for any $`coh(X)`$ there are natural isomorphisms of $`𝚔`$–vector spaces
$$H^i(X,)\mathrm{Ext}^{ni}(,\omega )^{},$$
where $``$ is denotes the $`𝚔`$–dual. The Serre duality theorem asserts the existence of the dualizing sheaf for smooth projective varieties. In this case the dualizing sheaf is a line bundle and coincides with the sheaf of differential forms $`\mathrm{\Omega }_X^n`$ of top degree.
Since the definition of $`\omega `$ is given in abstract categorical terms, it can be extend to the noncommutative case. More precisely, we will say that $`qgr(A)`$ satisfies classical Serre duality if there is an object $`\omega qgr(A)`$ together with natural isomorphisms
$$\mathrm{Ext}^i(𝒪,)\mathrm{Ext}^{ni}(,\omega )^{}.$$
Our noncommutative varieties $`_{\mathrm{}}^3`$ and $`_{\mathrm{}}^2`$ satisfy classical Serre duality, with dualizing sheaves being $`𝒪__{\mathrm{}}^3(4)`$ and $`𝒪__{\mathrm{}}^2(3),`$ respectively. This follows from the paper , where the existence of a dualizing sheaf in $`qgr(A)`$ has been proved for a general class of algebras which includes all noetherian regular algebras. In addition, the authors of showed that the dualizing sheaf coincides with $`\stackrel{~}{A}(l),`$ where $`l`$ is the Gorenstein paramenter for $`A`$ (see condition (3) of Definition 4.10).
### 5.3. Bundles on noncommutative projective spaces
To any graded right $`A`$-module $`M`$ one can attach a left $`A`$-module $`M^{}=\mathrm{Hom}_A(M,A)`$ which is also graded. Note that under this correspondence the right module $`A_A(r)`$ goes to the left module $`{}_{A}{}^{}A(r)`$.
It is known that if $`A`$ is a noetherian regular algebra, then $`\mathrm{Hom}_A(,A)`$ is a functor from the category $`gr(A)_R`$ to the category $`gr(A)_L`$. Moreover, its derived functor $`𝐑\mathrm{Hom}_A^{}(,A)`$ gives an anti-equivalence between the derived categories of $`gr(A)_R`$ and $`gr(A)_L`$ (see , , ).
If we assume that the composition of the functor $`\mathrm{Hom}_A(,A)`$ with the projection $`gr(A)_Lqgr(A)_L`$ factors through the projection $`gr(A)_Rqgr(A)_R,`$ then we obtain a functor from $`qgr(A)_R`$ to $`qgr(A)_L`$ which is denoted by $`\underset{¯}{om}(,𝒪)`$. This functor is not right exact and has right derived functors $`\underset{¯}{xt}^i(,𝒪),i>0,`$ from $`qgr(A)_R`$ to $`qgr(A)_L`$.
For a noetherian regular algebra the functor $`\underset{¯}{om}(,𝒪)`$ and its right derived functors exist. This follows from the fact that the functors $`\mathrm{Ext}_A^i(,A)`$ send a finite dimensional module to a finite dimensional module (see condition (3) of Definition 4.10).
Moreover, in this case the functor $`\underset{¯}{om}(,𝒪)`$ can be represented as the composition of the functor $`\mathrm{\Gamma }:qgr(A)_Rgr(A)_R,`$ the functor $`\mathrm{Hom}_A(,A):gr(A)_Rgr(A)_L,`$ and the projection $`\pi :gr(A)_Lqgr(A)_L`$. This can be illustrated by the following commutative diagram:
(13)
$$\begin{array}{ccc}gr(A)_R& \stackrel{\mathrm{Hom}_A(,A)}{–-\to }& gr(A)_L\\ \pi \mathrm{\Gamma }& & \pi \\ qgr(A)_R& \stackrel{\underset{¯}{om}(,𝒪)}{–-\to }& qgr(A)_L\end{array}$$
For a noetherian regular algebra the functor $`𝐑\mathrm{Hom}_A^{}(,A)`$ is an anti-equivalence between the derived categories of $`gr(A)_R`$ and $`gr(A)_L`$ and takes complexes of finite dimensional modules over $`gr(A)_R`$ to complexes of finite dimensional modules over $`gr(A)_L`$. This implies that the functor $`𝐑\underset{¯}{om}^{}(,𝒪)`$ gives an anti-equivalence between the derived categories of $`qgr(A)_R`$ and $`qgr(A)_L`$. (Note that for derived functors $`𝐑\mathrm{Hom}_A(,A)`$ and $`𝐑\underset{¯}{om}(,𝒪)`$ there is also a commutative diagram like (13)).
The functors $`\underset{¯}{xt}^j(,𝒪)`$ can be described more explicitly. Let $`M`$ be an $`A`$–bimodule. Regarding it as a right module, we see that for any $`QGr(A)_R`$ the groups $`\mathrm{Ext}^i(,\stackrel{~}{M})`$ have the structure of left $`A`$–modules. We can project them to $`QGr(A)_L`$. Thus each bimodule $`M`$ defines functors from $`QGr(A)_R`$ to $`QGr(A)_L,`$ which will be denoted by $`\pi \mathrm{Ext}^i(,\stackrel{~}{M})`$.
Now, using $`\pi \mathrm{\Gamma }=id`$ and the commutativity of the diagram (13) for the derived functors $`\mathrm{Ext}_A^j(,A)`$ and $`\underset{¯}{xt}^j(,𝒪),`$ we obtain isomorphisms
(14)
$$\underset{¯}{xt}^j(,𝒪)\pi \mathrm{Ext}_A^j(\mathrm{\Gamma }(),A)\pi \mathrm{Ext}_{gr(A)}^j(\mathrm{\Gamma }(),\underset{i0}{}A(i))\pi \mathrm{Ext}^j(,\underset{i0}{}𝒪(i))$$
for any sheaf $`qgr(A)_R`$.
###### Definition 5.4.
We call a coherent sheaf $`qgr(A)_R`$ locally free (or a bundle) if $`\underset{¯}{xt}^j(,𝒪)=0`$ for any $`j0`$.
Remark. In the commutative case this definition is equivalent to the usual definition of a locally free sheaf.
###### Definition 5.5.
The dual sheaf $`\underset{¯}{om}(,𝒪)qgr(A)_L`$ will be denoted by $`^{}`$ .
If $`qgr(A)_L`$ is a bundle, then the dual sheaf $`^{}`$ is a bundle in $`qgr(A)_L,`$ because $`𝐑\underset{¯}{om}^{}(^{},𝒪)=`$ in the derived category, and $`\underset{¯}{xt}^j(^{},𝒪)=0`$ for $`j0`$.
Thus we have a good definition of locally free sheaves on $`_{\mathrm{}}^3`$ and $`_{\mathrm{}}^2`$. Since the derived functor $`𝐑\underset{¯}{om}(,𝒪)`$ gives an anti-equivalence between the derived categories of $`qgr(A)_R`$ and $`qgr(A)_L,`$ there is an isomorphism:
(15)
$$\mathrm{Hom}(,𝒢)\mathrm{Hom}(𝒢^{},^{})$$
for any two bundles $``$ and $`𝒢`$ on $`_{\mathrm{}}^3`$ or $`_{\mathrm{}}^2`$.
## 6. Bundles on $`_{\mathrm{}}^2`$
### 6.1. Bundles on $`_{\mathrm{}}^2`$ with a trivialization on the commutative line
In this section we study bundles on $`_{\mathrm{}}^2`$. By definition, a bundle is an object $`coh(_{\mathrm{}}^2)`$ satisfying the additional condition $`\underset{¯}{xt}^i(,𝒪)=0`$ for all $`i>0`$ (see (5.4)).
The noncommutative plane $`_{\mathrm{}}^2`$ contains the commutative projective line $`l^1`$ given by the equation $`w_3=0`$. If $`M`$ is a $`PP_{\mathrm{}}`$-module, then the quotient module $`M/Mw_3`$ is a $`PP_{\mathrm{}}/w_3`$-module. This gives a functor $`coh(_{\mathrm{}}^2)coh(^1),`$ $`|_l`$. The sheaf $`|_l`$ is referred to as the restriction of $``$ to the line $`l`$.
###### Lemma 6.1.
If $``$ is a bundle, there is an exact sequence:
(16)
$$0(1)\stackrel{w_3}{}|_l0.$$
###### Proof.
To prove this we only need to show that multiplication by $`w_3`$ is a monomorphism. If $``$ is a bundle, it can be embedded into a direct sum $`\underset{i=1}{\overset{s}{}}𝒪(k_i),`$ because by ampleness the dual bundle $`^{}`$ is covered by a direct sum of line bundles. Now, since the morphism $`𝒪(k_i1)\stackrel{w_3}{}𝒪(k_i)`$ is mono for any $`i,`$ the same is true for the morphism $`(1)\stackrel{w_3}{}`$. ∎
###### Lemma 6.2.
Let $``$ be a bundle on $`_{\mathrm{}}^2`$ such that its restriction $`|_l`$ to the commutative line $`l`$ is isomorphic to a trivial bundle $`𝒪_l^r`$. Then
$$H^0(_{\mathrm{}}^2,(1))=H^0(_{\mathrm{}}^2,(2))=H^2(_{\mathrm{}}^2,(1))=H^2(_{\mathrm{}}^2,(2))=0.$$
###### Proof.
We have the following exact sequence in the category $`coh(_{\mathrm{}}^2)`$:
(17)
$$0(2)(1)(1)|_l0.$$
Since $`(1)|_l𝒪_l(1)^r,`$ we have $`H^0((1)|_l)=0`$.
Assume that $`(1)`$ has a nontrivial section. Then $`(2)`$ has a nontrivial section too. For the same reason $`(3)`$ has a nontrivial section, and so on. Thus for any $`n<0`$ the bundle $`(n)`$ has a nontrivial section. By (15) we have isomorphisms:
$$H^0((n))\mathrm{Hom}(𝒪(n),)\mathrm{Hom}(^{},𝒪(n)).$$
On the other hand, by Corollary 5.2 the last group is trivial for $`n0`$. Hence $`H^0((n))=0`$ for all $`n0,`$ and consequently $`H^0((2))=H^0((1))=0`$.
Further, assume that $`H^2((2))`$ is nontrivial. Since $`H^1((i)|_l)=0`$ for all $`i1`$ we have from the exact sequence (16) with $`=(i)`$ that $`H^2((i))`$ is nontrivial too for all $`i1`$. But this contradicts Corollary 5.3. Therefore $`H^2((2))=H^2((1))=0`$. This completes the proof. ∎
### 6.2. Monads on $`_{\mathrm{}}^2`$ and $`_{\mathrm{}}^3`$
As in the commutative case, a non-degenerate monad on $`_{\mathrm{}}^2`$ or $`_{\mathrm{}}^3`$ is a complex over $`coh(_{\mathrm{}}^2)`$
$$0H𝒪(1)\stackrel{m}{}K𝒪\stackrel{n}{}L𝒪(1)0$$
for which the map $`n`$ is an epimorphism and $`m`$ is a monomorphism. (Note that there is another more restrictive definition of a monad, according to which the dual map $`(m)^{}`$ has to be an epimorphism, see ). The coherent sheaf
$$E=𝐊𝐞𝐫(n)/𝐈𝐦(m)$$
is called the cohomology of a monad. A morphism between two monads is a morphism of complexes. The following lemma is proved in (Lemma 4.1.3) in the commutative case, but the proof is categorical and applies to the noncommutative case as well.
###### Lemma 6.3.
Let $`X`$ be either $`_{\mathrm{}}^2`$ or on $`_{\mathrm{}}^3,`$ and let $`E`$ and $`E^{}`$ be the cohomology bundles of two monads
$$\begin{array}{cc}M:\hfill & 0H𝒪(1)\stackrel{m}{}K𝒪\stackrel{n}{}L𝒪(1)0,\hfill \\ M^{}:\hfill & 0H^{}𝒪(1)\stackrel{m^{}}{}K^{}𝒪\stackrel{n^{}}{}L^{}𝒪(1)0\hfill \end{array}$$
on $`X`$. Then the natural mapping
$$\mathrm{Hom}(M,M^{})\mathrm{Hom}(E,E^{})$$
is bijective.
The proof is based on the fact that
$$\mathrm{Ext}^j(𝒪,𝒪(1))=\mathrm{Ext}^j(𝒪(1),𝒪(1))=\mathrm{Ext}^j(𝒪(1),𝒪)=0$$
for all $`j`$ (see , Lemma 4.1.3).
### 6.3. Non-degeneracy conditions
In the definition of a monad we require that the map $`n`$ be an epimorphism. In the commutative case this condition must be verified pointwise. In the noncommutative case the situation is simpler in some sense, because the complement of the commutative line $`l`$ does not have points.
###### Lemma 6.4.
If the restriction of a sheaf $`coh(_{\mathrm{}}^2)`$ to the projective line $`l`$ is the zero object, then $``$ is also the zero object.
###### Proof.
Let $`M`$ be a finitely generated graded $`PP_{\mathrm{}}`$-module such that $`\stackrel{~}{M}`$. Consider an exact sequence:
$$M\stackrel{w_3}{}M(1)N0.$$
Since $`\stackrel{~}{N}=(1)|_l=0,`$ the module $`N`$ is finite dimensional. This implies that for $`i0`$ the map $`M_i\stackrel{w_3}{}M_{i+1}`$ is surjective. Moreover, these maps are isomorphisms for $`i0,`$ because all $`M_i`$ are finite dimensional vector spaces. Let us identify all $`M_i`$ for $`i0`$ with respect to these isomorphisms. Using the $`A`$-module structure on $`M`$, we obtain a representation of the Weyl algebra $`\mathrm{T}(X,Y)/[X,Y]=2\mathrm{}`$ on the vector space $`M_i`$. But it is well known that the Weyl algebra does not have finite dimensional representations. Thus $`M_i=0`$ for all $`i0,`$ and $`M`$ is finite dimensional. Therefore $`=0.`$
The following corollary is an immediate consequence of the Lemma.
###### Corollary 6.5.
Let $`f:𝒢`$ be a morphism in $`coh(_{\mathrm{}}^2)`$. Suppose its restriction $`\overline{f}:|_l𝒢|_l`$ is an epimorphism. Then $`f`$ is an epimorphism too.
### 6.4. From the resolution of the diagonal to a monad
Let $`M`$ be an $`A`$–bimodule. Regarding it as a left module, we see that for any $`QGr(A)_L`$ the groups $`\mathrm{Ext}^i(,\stackrel{~}{M})`$ have the structure of right $`A`$–modules. We can project them to $`QGr(A)_R`$. Thus each bimodule $`M`$ defines functors $`\pi \mathrm{Ext}^i(,\stackrel{~}{M})`$ from $`QGr(A)_L`$ to $`QGr(A)_R`$.
Let $``$ be a bundle on $`_{\mathrm{}}^2`$ such that its restriction to the line $`l`$ is a trivial bundle. Let us consider the bundle $`^{}(1)qgr(PP_{\mathrm{}})_L`$ and the resolution of the diagonal $`𝒦_{}(PP_{\mathrm{}}),`$ which has only three terms:
$$\{0PP_{\mathrm{}}(1)PP_{\mathrm{}}(2)\mathrm{\Omega }^1(1)PP_{\mathrm{}}(1)PP_{\mathrm{}}PP_{\mathrm{}}\}\mathrm{\Delta }.$$
The resolution of the diagonal is a complex of bimodules. It induces a complex $`\stackrel{~}{𝒦}_{}`$ over $`QGr(PP_{\mathrm{}})_L`$:
(18)
$$\{0𝒪(1)PP_{\mathrm{}}(2)\mathrm{\Omega }^1(1)PP_{\mathrm{}}(1)𝒪PP_{\mathrm{}}\}\stackrel{~}{\mathrm{\Delta }}$$
where $`\mathrm{\Omega }^1`$ is a sheaf on $`_{\mathrm{}}^2`$ corresponding to the $`PP_{\mathrm{}}`$–module $`\mathrm{\Omega }^1`$.
As described above, each $`A`$–bimodule $`M`$ gives the functors $`\pi \mathrm{Ext}^i(,\stackrel{~}{M})`$ from $`QGr(A)_L`$ to $`QGr(A)_R`$. In particular, each object of the resolution of the diagonal induces such functors.
First we calculate these functors for the object $`\stackrel{~}{\mathrm{\Delta }}`$. Note that the object $`\stackrel{~}{\mathrm{\Delta }}`$ coincides with $`\underset{i0}{}𝒪(i)`$. Hence by (14) we have
$$\pi \mathrm{Ext}^j(^{}(1),\stackrel{~}{\mathrm{\Delta }})=0$$
if $`j>0,`$ while $`\pi \mathrm{Ext}^0(^{}(1),\stackrel{~}{\mathrm{\Delta }})(1).`$
The resolution of the diagonal (18) gives us a spectral sequence with the $`E_1`$ term
$$E_1^{pq}=\pi \mathrm{Ext}^q(^{}(1),\stackrel{~}{𝒦}_p)\pi \mathrm{Ext}^{p+q}(^{}(1),\stackrel{~}{\mathrm{\Delta }}),$$
which converges to
$$E_{\mathrm{}}^i=\{\begin{array}{cc}(1)\hfill & \text{if }i=0\hfill \\ 0\hfill & \text{otherwise}\hfill \end{array}$$
Now we describe all terms $`E_1^{pq}`$ of this spectral sequence. First we have
$$\pi \mathrm{Ext}^j(^{}(1),𝒪PP_{\mathrm{}})\mathrm{Ext}^j(^{}(1),𝒪)\stackrel{~}{PP_{\mathrm{}}}\mathrm{Ext}^j(^{}(1),𝒪)𝒪H^j(_{\mathrm{}}^2,(1))𝒪.$$
By Lemma 6.2, these groups are trivial for $`j1`$. For the same reason we have
$$\pi \mathrm{Ext}^j(^{}(1),𝒪(1)PP_{\mathrm{}}(2))=H^j(_{\mathrm{}}^2,(2))𝒪(2)=0$$
for $`j1`$ and
$$\pi \mathrm{Ext}^1(^{}(1),𝒪(1)PP_{\mathrm{}}(2))H^1(_{\mathrm{}}^2,(2))𝒪(2).$$
Now let us consider the functors which are associated with the object $`\mathrm{\Omega }^1(1)PP_{\mathrm{}}(1)`$. We have
$$\pi \mathrm{Ext}^j(^{}(1),\mathrm{\Omega }^1(1)PP_{\mathrm{}}(1))\mathrm{Ext}^j(^{},\mathrm{\Omega }^1)𝒪(1).$$
It follows from the Koszul complex that the sheaf $`\mathrm{\Omega }^1`$ can be included in two exact sequences:
$$\begin{array}{c}0\mathrm{\Omega }^1𝒪(1)PP_{\mathrm{}}^{}{}_{1}{}^{}𝒪0,\\ 0𝒪(3)𝒪(2)(PP_{\mathrm{}}^{}{}_{1}{}^{})^{}\mathrm{\Omega }^10.\end{array}$$
Applying the functor $`\mathrm{Hom}(^{},)`$ to the first sequence and taking into account that $`\mathrm{Hom}(^{},𝒪(1))=0,`$ we obtain $`\mathrm{Hom}(^{},\mathrm{\Omega }^1)=0`$. Similarly, we deduce from the second sequence that $`\mathrm{Ext}^2(^{},\mathrm{\Omega }^1)=0,`$ because $`\mathrm{Ext}^2(^{},𝒪(2))=0`$. This implies that the object $`\pi \mathrm{Ext}^j(^{}(1),\mathrm{\Omega }^1(1)PP_{\mathrm{}}(1))`$ is trivial for all $`j1`$.
Thus our spectral sequence is nothing more than the complex
$$\pi \mathrm{Ext}^1(^{}(1),\stackrel{~}{𝒦}_2)\pi \mathrm{Ext}^1(^{}(1),\stackrel{~}{𝒦}_1)\pi \mathrm{Ext}^1(^{}(1),\stackrel{~}{𝒦}_0),$$
which is isomorphic to the complex
$$H^1(_{\mathrm{}}^2,(2))𝒪(2)\mathrm{Ext}^1(^{},\mathrm{\Omega }^1)𝒪(1)H^1(_{\mathrm{}}^2,(1))𝒪.$$
It has only one cohomology which coincides with $`(1)`$.
###### Theorem 6.6.
Let $``$ be a bundle on $`_{\mathrm{}}^2`$ such that its restriction to the commutative line $`l`$ is isomorphic to the trivial bundle $`𝒪_l^r`$. Then $``$ is the cohomology of a monad
$$0H𝒪(1)\stackrel{m}{}K𝒪\stackrel{n}{}L𝒪(1)0$$
with $`H=H^1(_{\mathrm{}}^2,(2)),L=H^1(_{\mathrm{}}^2,(1)),`$ and such monad is unique up to an isomorphism. Moreover, in this case the vector spaces $`H`$ and $`L`$ have the same dimension.
###### Proof.
The existence of such a monad was proved above. The uniqueness follows from Lemma 6.3. The equality of dimensions of $`H`$ and $`L`$ follows immediately from the exact sequence (17). ∎
### 6.5. Barth description of monads
Now following Barth , we give a description of the moduli space of vector bundles on $`_{\mathrm{}}^2`$ trivial on the line $`l`$ in terms of linear algebra (see also ).
Denote by $`𝔐_{\mathrm{}}(r,0,k)`$ the moduli space of bundles on the noncommutative $`_{\mathrm{}}^2`$ trivial on the line $`l`$ and with a fixed trivialization there (i.e. with a fixed isomorphism $`|_l𝒪_l^r`$). Let $`dimH^1(_{\mathrm{}}^2,(1))=k`$. As in the commutative case, the numbers $`r,0,k`$ can be regarded as the rank, first Chern class, and second Chern class of $`,`$ respectively.
The following theorem gives a description of this moduli space which is similar to the description given by Barth in the commutative case.
###### Theorem 6.7.
Let $`\{(b_1,b_2;\mathrm{j},\mathrm{i})\}`$ be the set of quadruples of matrices $`b_1,b_2M_{k\times k}(),\mathrm{j}M_{r\times k}(),\mathrm{i}M_{k\times r}(),`$ which satisfy the condition
$$[b_1,b_2]+\mathrm{ij}+2\mathrm{}1_{k\times k}=0.$$
Then the space $`𝔐_{\mathrm{}}(r,0,k)`$ is the quotient of this set with respect to the following free action of $`\mathrm{GL}(k,)`$:
$$b_igb_ig^1,\mathrm{j}\mathrm{j}g^1,\mathrm{i}g\mathrm{i},\text{where }g\mathrm{GL}(k,).$$
###### Proof.
Let $``$ be a bundle on $`_{\mathrm{}}^2`$ trivial on the line $`l`$. We showed above that any such bundle comes from a monad unique up to an isomorphism. Conversely, suppose we have a monad
(19)
$$0H𝒪(1)\stackrel{m}{}K𝒪\stackrel{n}{}L𝒪(1)0$$
with $`dimH=dimL=k`$ such that its restriction to the line $`l`$ is a monad with the cohomology $`𝒪_l^r`$. Then the cohomology of this monad is a bundle on $`_{\mathrm{}}^2`$ which belongs to $`𝔐_{\mathrm{}}(r,0,k)`$. Indeed, the cohomologies of the dual complex
$$0𝒪(1)L^{}\stackrel{n^{}}{}𝒪K^{}\stackrel{m^{}}{}𝒪(1)H^{}0$$
coincide with $`\underset{¯}{om}(,𝒪)`$ and $`\underset{¯}{xt}^1(,𝒪)`$. Hence, to prove that $``$ is a bundle, it is sufficient to show that the dual complex is a monad too, i.e. that the map $`m^{}`$ is an epimorphism. The restriction of the dual complex to $`l`$ is a monad which is dual to the restriction of the monad (19) to $`l`$. Hence the restriction of $`m^{}`$ on $`l`$ is an epimorphism. Then, by Lemma 6.5, $`m^{}`$ is an epimorphism as well. Thus to describe the moduli space $`𝔐_{\mathrm{}}(r,0,k)`$ we have to decsribe the space of all monads (19) modulo isomorphisms preserving trivialization on $`l`$.
Consider a monad
$$0H𝒪(1)\stackrel{m}{}K𝒪\stackrel{n}{}L𝒪(1)0$$
with $`dimH=dimL=k`$ and $`dimK=2k+r`$. Denote by $``$ its cohomology bundle.
The maps $`m`$ and $`n`$ can be regarded as elements of $`H^{}KW`$ and $`K^{}LW,`$ respectively, where $`W=H^0(_{\mathrm{}}^2,𝒪(1))`$ is the vector space spanned by $`w_1,w_2,w_3`$. The maps $`m`$ and $`n`$ can be written as
$$m_1w_1+m_2w_2+m_3w_3,n_1w_1+n_2w_2+n_3w_3,$$
where $`m_i:HK`$ and $`n_i:KL`$ are constant linear maps.
Let us restrict the monad to the line $`l`$. The monadic condition $`nm=0`$ implies now:
$$n_1m_2+n_2m_1=0,n_1m_1=0,n_2m_2=0.$$
Moreover, since the restriction of $``$ to $`l`$ is trivial, the composition $`n_1m_2`$ is an isomorphism (see , Lemma 4.2.3). We can choose bases for $`H,K,L`$ so that $`n_1m_2=1_{k\times k}`$ (the identity matrix) and
$$m_1=\left(\begin{array}{c}1_{k\times k}\\ 0_{k\times k}\\ 0_{r\times k}\end{array}\right),m_2=\left(\begin{array}{c}0_{k\times k}\\ 1_{k\times k}\\ 0_{r\times k}\end{array}\right),n_1=\left(\begin{array}{ccc}0_{k\times k}& 1_{k\times k}& 0_{k\times r}\end{array}\right),n_2=\left(\begin{array}{ccc}1_{k\times k}& 0_{k\times k}& 0_{k\times r}\end{array}\right).$$
Using the equations $`n_3m_1+n_1m_3=0`$ and $`n_3m_2+n_2m_3=0`$ we can write:
$$m_3=\left(\begin{array}{c}b_1\\ b_2\\ \mathrm{j}\end{array}\right),n_3=\left(\begin{array}{ccc}b_2& b_1& \mathrm{i}\end{array}\right).$$
Now the monadic condition $`nm=0`$ can be written as:
$$(n_3m_3)w_3^2+1_{k\times k}[w_1,w_2]=0.$$
Therefore we obtain the following matrix equation:
$$[b_1,b_2]+\mathrm{ij}+2\mathrm{}1_{k\times k}=0.$$
Note that the last $`r`$ basis vectors of $`K`$ give us a trivialization of the restriction of $``$ to the line $`l`$. It is easy to check that any isomorphism of a monad which preserves trivialization on $`l`$ and the choice of the bases of $`H,K,L`$ made above has the form
$$b_igb_ig^1,\mathrm{j}\mathrm{j}g^1,\mathrm{i}g\mathrm{i},\text{where }g\mathrm{GL}(k,).$$
This proves the theorem. ∎
## 7. The noncommutative variety $`_{\mathrm{}}^3`$ as a twistor space
### 7.1. Real structures
A $``$–algebra is, by definition, an algebra over $``$ with an anti-linear anti-homomorphism $``$ satisfying $`^2=id`$. A $``$–structure on a (graded) algebra is regarded as a real structure on the corresponding (projective) noncommutative variety.
Let us introduce real structures on the complex varieties $`_{\mathrm{}}^4`$ and $`_{\mathrm{}}^4`$ defined in section 3. Assume that in (6), (7) the skew-symmetric matrix $`\theta `$ is purely imaginary and $`\mathrm{}`$ is real. Then there is a unique $``$–structure on the algebra $`A(_{\mathrm{}}^4)`$ such that $`x_i^{}=x_i`$. We denote the corresponding noncommutative variety by $`_{\mathrm{}}^4`$.
Assume in addition that the symmetric matrix $`G`$ in (7) is real and positive definite. There is a unique $``$–structure on the algebra $`Q_{\mathrm{}}`$ such that $`X_i^{}=X_i,`$$`D^{}=D,`$ and $`T^{}=T`$. The corresponding noncommutative real variety will be called the noncommutative sphere and denoted by $`𝕊_{\mathrm{}}^4`$. The embedding of $`_{\mathrm{}}^4`$ into $`_{\mathrm{}}^4`$ induces an embedding $`_{\mathrm{}}^4𝕊_{\mathrm{}}^4`$. Recall that the complement of $`_{\mathrm{}}^4`$ in $`_{\mathrm{}}^4`$ is a commutative quadratic cone $`\underset{kl}{}G^{kl}X_kX_l=0`$ which has only one real point. Thus $`𝕊_{\mathrm{}}^4`$ can be regarded as a one-point compactification of $`_{\mathrm{}}^4`$.
By a linear change of basis one can bring the pair $`(G,\theta )`$ to the standard form
(20)
$$G=\left(\begin{array}{cccc}1& 0& 0& 0\\ 0& 1& 0& 0\\ 0& 0& 1& 0\\ 0& 0& 0& 1\end{array}\right),\theta =\sqrt{1}\left(\begin{array}{cccc}0& a& 0& 0\\ a& 0& 0& 0\\ 0& 0& 0& b\\ 0& 0& b& 0\end{array}\right).$$
Furthermore, since $`\mathrm{}`$ and $`\theta `$ enter only in the combination $`\mathrm{}\theta ,`$ and we asssume that $`a+b0,`$ we can set $`a+b=1`$ without loss of generality.
### 7.2. Realification of $`_{\mathrm{}}^3`$
Recall that the noncommutative projective space $`_{\mathrm{}}^3`$ corresponds to the algebra $`PS_{\mathrm{}}`$ with generators $`z_i,i=1,2,3,4,`$ and relations (9). Consider an algebra $`\widehat{PS_{\mathrm{}}}`$ with generators $`z_i,\overline{z}_i,i=1,2,3,4,`$ and relations
(21)
$$\begin{array}{c}\begin{array}{ccc}[z_1,z_2]=2\mathrm{}(a+b)z_3z_4,\hfill & [z_1,\overline{z}_1]=2\mathrm{}bz_3\overline{z}_32\mathrm{}az_4\overline{z}_4,\hfill & [z_1,\overline{z}_2]=0,\hfill \\ [\overline{z}_1,\overline{z}_2]=2\mathrm{}(a+b)\overline{z}_3\overline{z}_4,\hfill & [z_2,\overline{z}_2]=2\mathrm{}az_3\overline{z}_32\mathrm{}bz_4\overline{z}_4,\hfill & [z_2,\overline{z}_1]=0,\hfill \end{array}\hfill \\ [z_i,z_j]=[z_i,\overline{z}_j]=[\overline{z}_i,z_j]=[\overline{z}_i,\overline{z}_j]=0\text{for all}i=3,4;j=1,2,3,4\hfill \end{array}$$
There is a unique $``$–structure on this algebra such that $`z_i^{}=\overline{z}_i,`$$`\overline{z}_i^{}=z_i`$. We denote the corresponding real variety $`_{\mathrm{}}^3()`$. This variety can be considered a realification of $`_{\mathrm{}}^3`$.
Remark. In contrast to the commutative situation, a noncommutative complex variety in general has many different realifications. We have an ambiguity in the choice of relations involving both $`z_i`$ and $`\overline{z}_j`$. The realification (21) is distinguished by the fact that it is the twistor space of the noncommutative sphere $`𝕊_{\mathrm{}}^4,`$ as explained below.
In the commutative case there is a map from $`^3()`$ to the sphere $`𝕊^4`$ which is a $`^1`$ fibration. The corresponding $`^1`$ bundle is the projectivization of a spinor bundle on $`𝕊^4`$. This map is known as the Penrose map. In the noncommutative case we have a similar picture. The analogue of the Penrose map is a map $`\mathrm{\Pi }:_{\mathrm{}}^3()𝕊_{\mathrm{}}^4`$ which is associated with the homomorphism of $``$-algebras $`Q_{\mathrm{}}\widehat{PS_{\mathrm{}}}`$:
$$\begin{array}{ccccccc}X_1\hfill & \hfill & \frac{\sqrt{1}}{2}& (z_1\overline{z}_4\overline{z}_1z_4\overline{z}_2z_3+z_2\overline{z}_3),D\hfill & \hfill & \frac{1}{2}(z_1\overline{z}_1+\overline{z}_1z_1+z_2\overline{z}_2+\overline{z}_2z_2),\hfill & \\ X_2\hfill & \hfill & \frac{1}{2}& (z_1\overline{z}_4+\overline{z}_1z_4\overline{z}_2z_3z_2\overline{z}_3),T\hfill & \hfill & (z_3\overline{z}_3+z_4\overline{z}_4),\hfill & \\ X_3\hfill & \hfill & \frac{\sqrt{1}}{2}& (\overline{z}_1z_3z_1\overline{z}_3+z_2\overline{z}_4\overline{z}_2z_4),\hfill & & & \\ X_4\hfill & \hfill & \frac{1}{2}& (z_1\overline{z}_3+\overline{z}_1z_3+\overline{z}_2z_4+z_2\overline{z}_4).\hfill & & & \end{array}$$
Note that for $`\mathrm{}=0`$ we obtain the homomorphism of commutative algebras which corresponds to the usual Penrose map. This means that $`_{\mathrm{}}^3()`$ is the twistor space of $`𝕊_{\mathrm{}}^4`$.
The variety $`_{\mathrm{}}^3()`$ is a twistor space in yet another sense. For the commutative $`^4`$ the complex structures compatible with the symmetric bilinear form $`G`$ and orientation are parametrized by points of a $`^1`$. This remains true in the noncommutative case. A complex structure (resp. orientation) on $`_{\mathrm{}}^4`$ is defined as a complex structure (resp. orientation) on the real vector space $`U`$ spanned by $`x_1,\mathrm{},x_4`$. We will choose an orientation on $`U`$ and require that the complex structure be compatible with it. All such complex structures are parametrized by points of a $`^1`$.
Recall now that $`_{\mathrm{}}^3`$ is a pencil of noncommutative projective planes passing through the commutative line. Let us pick any one of them. The realification of $`_{\mathrm{}}^3`$ defined above induces a realification of the noncommutative projective plane. It is easy to see that the complement of the commutative line $`w_3=\overline{w}_3=0`$ in the realified projective plane is isomorphic to $`_{\mathrm{}}^4`$. Furthermore, the complement carries a natural complex structure defined by
$$w_3^1w_i\sqrt{1}w_3^1w_i,\overline{w}_3^1\overline{w}_i\sqrt{1}\overline{w}_3^1\overline{w}_i,i=1,2.$$
The Penrose map induces an identification between the complement and $`_{\mathrm{}}^4𝕊_{\mathrm{}}^4,`$ and therefore induces a complex structure on the latter. Varying the noncommutative projective plane, one obtains all possible complex structures on $`_{\mathrm{}}^4`$ compatible with a particular orientation. This is completely analogous to the commutative case.
### 7.3. Connection between sheaves on commutative and noncommutative planes
In this subsection we are going to connect the moduli space $`𝔐_{\mathrm{}}(r,0,k)`$ of bundles on $`_{\mathrm{}}^2`$ with a trivialization on the line $`l`$ with the moduli space $`𝔐(r,0,k)`$ of torsion free sheaves on the commutative $`^2`$ with a trivialization on a fixed line. The bridge between bundles on $`_{\mathrm{}}^2`$ and torsion free sheaves on $`^2`$ is provided by the twistor variety $`_{\mathrm{}}^3`$. This gives a geometrical interpretation of Nakajima’s results (the description of the moduli space $`𝔐(r,0,k)`$ by the deformed ADHM data ). We will construct a hyperkähler manifold $``$ parametrizing certain complexes on $`_{\mathrm{}}^3`$ which is isomorphic to $`𝔐(r,0,k)`$ (which is also a hyperkähler manifold ). The isomorphism is given by the restriction of complexes to one of the commutative $`^2`$’s. On the other hand, the restriction of complexes to a noncommutative plane $`_{\mathrm{}}^2`$ yields an isomorphism between $``$ with a particular choice of complex structure and the moduli space $`𝔐_{\mathrm{}}(r,0,k)`$. Thus $`𝔐_{\mathrm{}}(r,0,k)`$ can be obtained from $`𝔐(r,0,k)`$ by a rotation of complex structure.
Consider complexes $`𝒞^{}`$ on $`_{\mathrm{}}^3`$ of the form
(22)
$$0H𝒪(1)\stackrel{M}{}K𝒪\stackrel{N}{}L𝒪(1)0$$
with $`dimH=dimL=k,`$ $`dimK=2k+r,`$ which satisfies the condition that its restriction to the line $`l`$ has only one cohomology which is a trivial bundle (with a fixed trivialization). This condition implies that $`M`$ is a monomorphism. Note that $`N`$ is not an epimorphism in general, so (22) is not a monad. But the restriction of the complex (22) to any noncommutative plane is a monad by Corollary 6.5. Thus $`N`$ can fail to be surjective only on the commutative planes $`z_3=0`$ and $`z_4=0`$.
Now we introduce a real structure on $`_{\mathrm{}}^3`$ (this is different from the real structure on the realification of $`_{\mathrm{}}^3`$ defined above). Assume that $`\mathrm{}`$ is a real number. Consider an anti-linear anti-homomorphism $`\overline{𝒥}`$ of $`PS_{\mathrm{}}`$ defined by
$$\overline{𝒥}(z_1)=z_2,\overline{𝒥}(z_2)=z_1,\overline{𝒥}(z_3)=z_4,\overline{𝒥}(z_4)=z_3,\overline{𝒥}(\lambda )=\overline{\lambda },\lambda .$$
Thus $`\overline{𝒥}`$ is a homomorphism of $``$–algebras from $`PS_{\mathrm{}}`$ to the opposite algebra $`PS_{\mathrm{}}^{}{}_{}{}^{op}`$. (The notation $`\overline{𝒥}`$ is used by analogy with the commutative case, where this anti-homomorphism is a composition of a complex structure $`J`$ with complex conjugation .)
The anti-homomorphism $`\overline{𝒥}`$ induces a functor $`\overline{𝒥}^{}`$ from $`qgr(PS_{\mathrm{}})_R`$ to $`qgr(PS_{\mathrm{}}^{}{}_{}{}^{op})_R`$. The latter category is naturally identified with the category $`qgr(PS_{\mathrm{}})_L`$. Using this identification we can consider the composition of $`\overline{𝒥}^{}`$ with the dualization functor $`\underset{¯}{om}(,𝒪)`$ as a functor from $`qgr(PS_{\mathrm{}})_R`$ to itself. For any bundle $``$ we denote by $`\overline{𝒥}^{}()^{}`$ its image under this functor. The functor can be extended to complexes of bundles. It takes the complex $`𝒞^{}`$ (22) to the complex $`\overline{𝒥}^{}(𝒞^{})^{}`$
$$0\overline{L}^{}𝒪(1)\stackrel{\overline{𝒥}^{}(N)^{}}{}\overline{K}^{}𝒪\stackrel{\overline{𝒥}^{}(M)^{}}{}\overline{H}^{}𝒪(1)0.$$
Let us consider complexes $`𝒞^{}`$ on $`_{\mathrm{}}^3`$ with an isomorphism
(23)
$$\overline{𝒥}^{}(𝒞^{})^{}𝒞^{}$$
and trivialization on the line $`l`$. Then the space $`K`$ acquires a hermitian metric and $`L`$ becomes isomorphic to $`\overline{H}^{}`$. The reasoning of section 6 shows that we can represent the maps $`M`$ and $`N`$ as
$$M_1z_1+M_2z_2+M_3z_3+M_4z_4,N_1z_1+N_2z_2+N_3z_3+N_4z_4,$$
where $`M_i`$ and $`N_i`$ are constant maps. By a suitable choice of bases we can put these maps into the form
(24)
$$M_1=\left(\begin{array}{c}1\\ 0\\ 0\end{array}\right),M_2=\left(\begin{array}{c}0\\ 1\\ 0\end{array}\right),M_3=\left(\begin{array}{c}B_1\\ B_2\\ J\end{array}\right),M_4=\left(\begin{array}{c}B_{1}^{}{}_{}{}^{^{}}\\ B_{2}^{}{}_{}{}^{^{}}\\ J^{}\end{array}\right),$$
$$N_1=\left(\begin{array}{ccc}0& 1& 0\end{array}\right),N_2=\left(\begin{array}{ccc}1& 0& 0\end{array}\right),N_3=\left(\begin{array}{ccc}B_2& B_1& I\end{array}\right),N_4=\left(\begin{array}{ccc}B_{2}^{}{}_{}{}^{^{}}& B_{1}^{}{}_{}{}^{^{}}& I^{^{}}\end{array}\right).$$
Using the reality conditions $`\overline{𝒥}^{}(N)^{}=M`$ and $`\overline{𝒥}^{}(M)^{}=N`$ we find that
(25)
$$B_{1}^{}{}_{}{}^{^{}}=B_{2}^{}{}_{}{}^{},B_{2}^{}{}_{}{}^{^{}}=B_{1}^{}{}_{}{}^{},J^{^{}}=I^{},I^{^{}}=J^{}.$$
Finally the condition $`NM=0`$ gives
$$\begin{array}{cc}a)\hfill & \mu _c=[B_1,B_2]+\mathrm{𝐼𝐽}=0,\hfill \\ b)\hfill & \mu _r=[B_1,B_{1}^{}{}_{}{}^{}]+[B_2,B_{2}^{}{}_{}{}^{}]+\mathrm{𝐼𝐼}^{}J^{}J=2\mathrm{}1_{k\times k}.\hfill \end{array}$$
These matrix equations are invariant under the following action of $`U(k)`$:
(26)
$$B_igB_ig^1,IgI,JJg^1,\text{where }gU(k).$$
Denote by $`𝐌`$ the vector space of complex matrices $`(B_1,B_2,I,J)`$. It has a structure of a quaternionic vector space defined by
$$(B_1,B_2,I,J)(B_{2}^{}{}_{}{}^{},B_{1}^{}{}_{}{}^{},J^{},I^{}),$$
and, moreover, it is a flat hyperkähler manifold (see ). The map $`\mu =(\mu _r,\mu _c)`$ is a hyperkähler moment map for the action of $`U(k)`$ defined in (26) (see ). Since the action of $`U(k)`$ on $`\mu _c^1(0)\mu _r^1(2\mathrm{}1)`$ is free, the quotient $`=\mu _c^1(0)\mu _r^1(2\mathrm{}1)/U(k)`$ is a smooth hyperkähler manifold. This manifold parametrizes complexes (22) with a real structure (23) and a trivialization on the line $`l`$.
On the other hand, it was proved in that the moduli space $`𝔐(r,0,k)`$ of torsion free sheaves on the commutative $`^2`$ with a trivialization on a fixed line can be identified with $``$.
This identification can be described geometrically as follows. Let us assume that $`\mathrm{}`$ is positive. It can be checked that in this case the map $`N`$ can fail to be surjective only on the plane $`z_4=0`$. We can restrict the complex (22) to the commutative plane $`z_3=0`$. The restriction is a monad and its cohomology sheaf is a torsion free sheaf. It is easy to see that this yields a complex isomorphism from $``$ to $`𝔐(r,0,p)`$.
The restriction of the complex (22) to a noncommutative plane is a monad as well. This yields a map from $``$ to the moduli space $`𝔐_{\mathrm{}}(r,0,k)`$ of bundles on the noncommutative plane. Let us show that this map is an isomorphism. To this end we note that on the level of the linear algebra data this map sends a quadruple $`(B_1,B_2,I,J)`$ to the quadruple $`(b_1,b_2,\mathrm{i},\mathrm{j})`$ with
$$\begin{array}{c}b_1=B_1B_{2}^{}{}_{}{}^{},b_2=B_2+B_{1}^{}{}_{}{}^{},\mathrm{i}=IJ^{},\mathrm{j}=J+I^{}.\hfill \end{array}$$
Further, note that the equations $`\mu _c=0,`$ $`\mu _r=2\mathrm{}1`$ are equivalent to the equation $`[b_1,b_2]+\mathrm{i}\mathrm{j}+2\mathrm{}1=0`$ and the vanishing of the moment map for the action of the group $`U(k)`$ on the space of quadruples $`(b_1,b_2,\mathrm{i},\mathrm{j})`$. Now it follows from the theorem of Kempf and Ness (, ) that the map $`𝔐_{\mathrm{}}(r,0,k)`$ is a diffeomorphism. It becomes a complex isomorphism if we replace the natural complex structure of the space $``$ with another one within the $`^1`$ of complex structures on $``$.
Thus we have
###### Theorem 7.1.
The moduli space $`𝔐_{\mathrm{}}(r,0,k)`$ is a smooth hyperkähler manifold of real dimension $`4rk,`$ and as a hyperkähler manifold it is isomorphic to the moduli space $`𝔐(r,0,k)`$ of torsion free sheaves on the commutative $`^2`$ with a trivialization on a fixed line. As a complex manifold $`𝔐_{\mathrm{}}(r,0,k)`$ is obtained from $`𝔐(r,0,k)`$ by a rotation of the complex structure.
The above discussion shows that there are natural bijections between
$`A^{}.`$ Bundles on $`_{\mathrm{}}^2`$ with a trivialization on the commutative line $`l`$ and $`c_2=k.`$
$`B^{}.`$ Solutions of the equations $`\mu _c=0,\mu _r=2\mathrm{}1`$ modulo the action of $`U(k).`$
$`C^{}.`$ Complexes of sheaves on $`_{\mathrm{}}^3`$ of the form (22) with a trivialization on the commutative line $`l`$ satisfying the reality condition (23).
One can show that for $`r>1`$ a generic complex (22) is a monad and its cohomology is a bundle $``$ on $`_{\mathrm{}}^3`$ such that
(27)
$$H^1(_{\mathrm{}}^3,(2))=0,\overline{𝒥}^{}()^{}.$$
Moreover, it can be shown that any bundle $``$ satisfying the conditions (27) can be represented as a cohomology of a monad of the form (22).
## 8. Noncommutative twistor transform
### 8.1. Review of the twistor transform
In the commutative case the ADHM construction of instantons has the following geometric interpretation. Consider the double fibration
(28)
$$\begin{array}{ccccc}𝐆(2;4)& \stackrel{p}{}& \mathrm{𝐅𝐥}(1,2;4)& \stackrel{q}{}& ^3,\end{array}$$
where $`𝐆(2;4)`$ is the Grassmannian and $`\mathrm{𝐅𝐥}(1,2;4)`$ is the partial flag variety. The Grassmannian $`𝐆(2;4)`$ has a real structure with $`𝕊^4`$ as the set of real points. For any bundle $``$ on $`^3`$ its twistor transform is defined as a sheaf $`p_{}q^{}`$ on $`𝐆(2;4)`$. Given ADHM data we have a monad on $`^3`$ whose cohomology is a bundle. It can be shown that the restriction of its twistor transform to the sphere $`𝕊^4`$ coincides with the instanton bundle corresponding to these ADHM data. The instanton connection can also be reconstructed from the bundle on $`^3`$ (see for details).
In this section we show that one can consider the noncommutative quadric introduced in section 3 as a noncommutative Grassmannian $`𝐆(2;4)`$. We also construct a noncommutative flag variety $`\mathrm{𝐅𝐥}(1,2;4)`$ and projections $`p,`$ $`q`$ giving a noncommutative analogue of the twistor diagram (28). The twistor transform can be defined in the same way as above. It produces a bundle on the noncommutative sphere from the deformed ADHM data. We show that this bundle is precisely the kernel of the map $`𝒟`$ defined in section 2.
It should also be possible to construct the instanton connection on the noncommutative $`^4`$ from the complex of sheaves on $`_{\mathrm{}}^3`$. To do this, one needs to develop differential geometry of noncommutative affine and projective varieties. We go some way in this direction by defining differential forms and spinors.
Since the goal of this section is mainly illustrative, we limit ourselves to stating the results. An interested reader should be able to fill in the proofs.
### 8.2. Tensor categories
A good way to construct noncommutative varieties with properties similar to those of commutative varieties is to start with a tensor category (see ). Let $`𝒯`$ be an abelian tensor category. Consider a tensor functor $`\mathrm{\Phi }:𝒯𝒱ect`$ to the abelian tensor category of vector spaces compatible with the associativity constraint but not compatible with the commutativity constraint. If $`A`$ is a commutative algebra in the tensor category $`𝒯,`$ then $`\mathrm{\Phi }(A)`$ is a noncommutative algebra in the tensor category $`𝒱ect`$. If $`M𝒯`$ is a right $`A`$-module, then $`\mathrm{\Phi }(M)`$ is a right $`\mathrm{\Phi }(A)`$-module. Any right $`A`$–module (in the category $`𝒯`$) has a natural structure of a left $`A`$–module (and hence an $`A`$–bimodule). Thus any right $`\mathrm{\Phi }(A)`$–module of the form $`\mathrm{\Phi }(M)`$ has a natural structure of a $`\mathrm{\Phi }(A)`$–bimodule.
Consider the category $`𝒞omm_𝒯`$ of all finitely generated (graded) commutative algebras in the tensor category $`𝒯`$. Then under $`\mathrm{\Phi }`$ the category $`𝒞omm_𝒯`$ is mapped to a subcategory of the category of finitely generated (graded) algebras. This subcategory enjoys many properties of the category of commutative (graded) algebras. For example, for all $`A,B𝒞omm_𝒯`$ there is a natural algebra structure on $`\mathrm{\Phi }(A)\mathrm{\Phi }(B)`$ coming from the algebra structure on $`AB.`$ The corresponding subcategory in the category of noncommutative affine (resp. projective) varieties shares a lot of properties with the category of commutative varieties. For example, if $`X`$ and $`Y`$ are varieties in this category, then using the tensor product of the corresponding algebras one can define the “Carthesian” product $`X\times Y`$. More generally, given a pair of morphisms $`XZ`$ and $`YZ`$ one can define the fiber product $`X\times _ZY`$. Further, starting from the module of differential forms of $`A`$ one can construct the sheaf of differential forms on the corresponding noncommutative variety.
The category $`qgr(\mathrm{\Phi }(A))`$ has a nice subcategory which consists of modules of the form $`\mathrm{\Phi }(M),`$ where $`M𝒯`$ is an $`A`$–module. To any object $`\mathrm{\Phi }(M)`$ of this subcategory one can associate its symmetric and exterior powers. The symmetric powers of $`\mathrm{\Phi }(M)`$ form a noncommutative graded algebra. This enables one to define the projectivization of the sheaf corresponding to the module $`\mathrm{\Phi }(M).`$
### 8.3. Yang-Baxter operators
One way to construct an abelian tensor category $`𝒯`$ with a functor $`\mathrm{\Phi }:𝒯𝒱ect`$ is to consider a Yang-Baxter operator (see , ).
A Yang-Baxter operator on a vector space $`V`$ is an operator $`R:VVVV,`$ such that
(29)
$$R^2=\mathrm{𝐢𝐝}_{VV},(R\mathrm{𝐢𝐝}_V)(\mathrm{𝐢𝐝}_VR)(R\mathrm{𝐢𝐝}_V)=(\mathrm{𝐢𝐝}_VR)(R\mathrm{𝐢𝐝}_V)(\mathrm{𝐢𝐝}_VR).$$
A Yang-Baxter operator induces an action of the permutation group $`𝔖_n`$ on the tensor power $`V^n,`$ where the transposition $`(i,i+1)𝔖_n`$ acts as the operator
$$R_{i,i+1}=\mathrm{𝐢𝐝}_{V^{(i1)}}R\mathrm{𝐢𝐝}_{V^{(ni1)}}:V^nV^n.$$
Equations (29) ensure that operators $`R_{i,i+1}`$ satisfy the relations between the transpositions $`(i,i+1)`$ in the group $`𝔖_n`$.
If $`R`$ is a Yang-Baxter operator on a vector space $`V,`$ then the dual operator $`R^{}:V^{}V^{}V^{}V^{}`$ is also a Yang-Baxter operator.
Given a Yang-Baxter operator $`R:VVVV,`$ one can construct an abelian tensor category $`𝒯_R`$ and a functor $`\mathrm{\Phi }_R:𝒯_R𝒱ect`$ such that $`V`$ is a $`\mathrm{\Phi }_R`$–image of some object of $`𝒯_R`$, and the commutativity morphism in the category $`𝒯_R`$ is mapped by $`\mathrm{\Phi }_R`$ to $`R`$ . As mentioned above, given any two objects $`A,B`$ of the category $`𝒞omm_{𝒯_R},`$ one has a natural algebra structure on the vector space $`\mathrm{\Phi }(A)\mathrm{\Phi }(B)`$. This algebra will be denoted $`\mathrm{\Phi }(A)\underset{𝑅}{}\mathrm{\Phi }(B)`$ and called the $`R`$-tensor product of $`\mathrm{\Phi }(A)`$ and $`\mathrm{\Phi }(B)`$.
It is well known that there is a one-to-one correspondence between irreducible representations of the group $`𝔖_n`$ and partitions of $`n`$ (Young diagrams). Under this correspondence the trivial partition $`(n)`$ corresponds to the sign representation, while the maximal partition $`(\underset{n\text{ times}}{\underset{}{1,1,\mathrm{},1}})`$ corresponds to the identity representation. Given a partition $`(k_1,\mathrm{},k_r)`$ of $`n`$ $`(k_1k_2\mathrm{}k_r)`$ we denote by $`(k_1,\mathrm{},k_r)`$ the corresponding irreducible representation and by $`\mathrm{\Sigma }_R^{(k_1,\mathrm{},k_r)}V`$ (resp. $`\mathrm{\Sigma }_R^{(k_1,\mathrm{},k_r)}V^{}`$) the $`(k_1,\mathrm{},k_r)`$-isotypical component of $`V^n`$ (resp. $`(V^{})^n`$), i.e. the sum of all subrepresentations of $`V^n`$ (resp. $`(V^{})^n`$) isomorphic to $`(k_1,\mathrm{},k_r)`$. We also put $`\mathrm{\Lambda }_R^nV=\mathrm{\Sigma }_R^{(n)}V,`$ $`\mathrm{\Lambda }_R^nV^{}=\mathrm{\Sigma }_R^{(n)}V^{}`$ for brevity.
Remark. The subspaces $`\mathrm{\Sigma }_R^\lambda VV^n`$ are the $`\mathrm{\Phi }_R`$-images of some objects of the category $`𝒯_R`$.
Let $`\lambda ,`$ $`\mu `$ be partitions of $`n`$ and $`m`$ respectively. It is clear that the action of the permutation $`\sigma _{n,m}𝔖_{n+m}`$
$$\sigma _{n,m}(i)=\{\begin{array}{cc}i+m,\hfill & \text{if }1in\hfill \\ in,\hfill & \text{if }n+1in+m\hfill \end{array}$$
gives an isomorphism
$$R_{n,m}:\mathrm{\Sigma }_R^\lambda V\mathrm{\Sigma }_R^\mu V\mathrm{\Sigma }_R^\mu V\mathrm{\Sigma }_R^\lambda V.$$
Remark. This isomorphism is the image of an isomorphism in the category $`𝒯_R`$.
The trivial example of a Yang-Baxter operator is the usual transposition
$$R_0(v_1v_2)=v_2v_1.$$
We will say that $`R`$ is a deformation-trivial Yang-Baxter operator if $`R`$ is an algebraic deformation of $`R_0`$ in the class of Yang-Baxter operators. For a deformation-trivial Yang-Baxter operator $`R`$ we have
$$dim\mathrm{\Sigma }_R^\lambda V=dim\mathrm{\Sigma }_{R_0}^\lambda V$$
for any partition $`\lambda `$.
### 8.4. The noncommutative projective space
Let $`R`$ be a deformation-trivial Yang-Baxter operator on the vector space $`V^{}`$. Then the graded algebra
$$S_R^{}V^{}=T(V^{})/\mathrm{\Lambda }_R^2V^{}$$
is a noncommutative deformation of the coordinate algebra of the projective space $`(V)`$. We denote by $`_R(V)`$ the corresponding noncommutative variety. Thus $`_R(V)`$ is a noncommutative deformation of the projective space $`(V)`$.
###### Example 8.1.
The operator
(30)
$$\begin{array}{c}R(z_iz_j)=z_jz_i,\text{if }(i,j)(1,2),\text{ }(2,1)\hfill \\ R(z_1z_2)=z_2z_1+2\mathrm{}(az_3z_4+bz_4z_3)\hfill \\ R(z_2z_1)=z_1z_22\mathrm{}(bz_3z_4+az_4z_3).\hfill \end{array}$$
is a deformation trivial Yang-Baxter operator on the $`4`$-dimensional vector space $`Z^{}`$ with the basis $`\{z_1,z_2,z_3,z_4\}`$. By definition the homogeneous coordinate algebra of $`_R(Z)`$ is generated by $`z_1,z_2,z_3,z_4`$ with relations $`(\text{9})`$ (we set $`a+b=1`$ as before). Hence $`_R(Z)`$ is isomorphic to the noncommutative projective space $`_{\mathrm{}}^3`$ defined in section 3. The space $`Z^{}`$ was denoted $`U`$ in that section.
The above example shows that part of the data encoded in the Yang-Baxter operator $`R`$ is lost in the structure of the corresponding noncommutative projective space. We will see below that this data appears in the structure of other noncommutative varieties associated with $`R`$.
### 8.5. Noncommutative Grassmannians
It is well known that the homogeneous coordinate algebra of the Grassmann variety $`𝐆(k;V)`$ is a graded quadratic algebra with $`\mathrm{\Lambda }^kV^{}`$ as the space of generators and
$$𝐊𝐞𝐫\left(\mathrm{\Lambda }^kV^{}\mathrm{\Lambda }^kV^{}(V^{})^{2k}\mathrm{\Sigma }^{(k,k)}V^{}\right)$$
as the space of relations. This description justifies the following definition.
###### Definition 8.2.
Let $`R`$ be a Yang-Baxter operator on the space $`V^{}`$. The noncommutative Grassmann variety $`𝐆_R(k;V)`$ is the noncommutative projective variety corresponding to the quadratic algebra
$$𝒢_R(k;V)=T(\mathrm{\Lambda }_R^kV^{})/𝐊𝐞𝐫(\mathrm{\Lambda }_R^kV^{}\mathrm{\Lambda }_R^kV^{}\mathrm{\Sigma }_R^{(k,k)}V^{})$$
The algebra $`𝒢_R(k;V)`$ is the $`\mathrm{\Phi }_R`$-image of a commutative algebra in the category $`𝒯_R`$.
If $`R`$ is deformation-trivial, then $`𝐆_R(k;V)`$ is a noncommutative deformation of $`𝐆(k;V)`$. Note that $`𝐆_R(1;V)=_R(V)`$ by definition.
###### Example 8.3.
Consider the noncommutative Grassmannian $`𝐆_R(2;Z)`$ corresponding to the Yang-Baxter operator $`(\text{30})`$. Let
$$z_{ij}=\frac{1}{2}((z_iz_jz_jz_i)R(z_iz_jz_jz_i))\mathrm{\Lambda }_R^2Z^{}.$$
Then it is easy to check that $`𝒢_R(2;Z)`$ is generated by the elements
$$Y_1=z_{13},Y_2=z_{24},Y_3=z_{23},Y_4=z_{14},D=z_{12},T=z_{34},$$
with relations
(31)
$$\begin{array}{c}\begin{array}{cccccc}[Y_1,Y_2]\hfill & =& 2\mathrm{}aT^2,\hfill & [Y_3,Y_4]\hfill & =& 2\mathrm{}bT^2,\hfill \\ [D,Y_1]\hfill & =& 2\mathrm{}aY_1T,\hfill & [D,Y_2]\hfill & =& 2\mathrm{}aY_2T,\hfill \\ [D,Y_3]\hfill & =& 2\mathrm{}bY_3T,\hfill & [D,Y_4]\hfill & =& 2\mathrm{}bY_4T,\hfill \end{array}\hfill \\ DT=\frac{1}{2}\left(Y_1Y_2+Y_2Y_1+Y_3Y_4+Y_4Y_3\right),\hfill \end{array}$$
$`[Y_i,Y_j]=[T,Y_j]=[T,D]=0`$ for all $`i=3,4,j=1,2,3,4.`$ Comparing with $`(\text{7})`$ one can see that the algebra $`𝒢_R(2;Z)`$ is isomorphic to $`_{\mathrm{}}`$ with $`G`$ and $`\theta `$ given by
$$G=\frac{1}{2}\left(\begin{array}{cccc}0& 1& 0& 0\\ 1& 0& 0& 0\\ 0& 0& 0& 1\\ 0& 0& 1& 0\end{array}\right),\theta =2\mathrm{}\left(\begin{array}{cccc}0& a& 0& 0\\ a& 0& 0& 0\\ 0& 0& 0& b\\ 0& 0& b& 0\end{array}\right).$$
Note that the variables $`X_i,i=1,2,3,4,`$ used in section 7 to describe the quadric are related to $`Y_i,i=1,2,3,4,`$ by the following formulas:
(32)
$$\begin{array}{cc}Y_1=X_2+\sqrt{1}X_1,\hfill & Y_2=X_2+\sqrt{1}X_1,\hfill \\ Y_3=X_4+\sqrt{1}X_3,\hfill & Y_4=X_4+\sqrt{1}X_3.\hfill \end{array}$$
### 8.6. Products of Grassmannians and flag varieties
Let $`R`$ be a Yang-Baxter operator on the vector space $`V^{}`$. Consider a sequence $`k_1,\mathrm{},k_r`$ of integers. Let $`^r`$ be the free abelian group with $`r`$ generators $`e_1,\mathrm{},e_r`$. The $`R`$-tensor product
$$𝒢_R(k_1;V)\underset{𝑅}{}\mathrm{}\underset{𝑅}{}𝒢_R(k_r;V)$$
is a $`^r`$-graded algebra generated by the vector spaces $`\mathrm{\Lambda }_R^{k_i}V^{}`$ in degree $`e_i`$, with relations
$$𝐊𝐞𝐫\left(\mathrm{\Lambda }_R^{k_i}V^{}\mathrm{\Lambda }_R^{k_i}V^{}\mathrm{\Sigma }_R^{(k_i,k_i)}V^{}\right)$$
in degree $`2e_i`$ for all $`i`$ and
$$\begin{array}{ccc}𝐊𝐞𝐫((\mathrm{\Lambda }_R^{k_i}V^{}\mathrm{\Lambda }_R^{k_j}V^{})(\mathrm{\Lambda }_R^{k_j}V^{}\mathrm{\Lambda }_R^{k_i}V^{})& \stackrel{(\mathrm{𝐢𝐝},R_{k_j,k_i})}{}& \mathrm{\Lambda }_R^{k_i}V^{}\mathrm{\Lambda }_R^{k_j}V^{})\end{array}$$
in degree $`e_i+e_j`$ for all $`i>j`$. For any increasing sequence $`k_1,\mathrm{},k_r`$ we define also a $`^r`$-graded algebra $`_R(k_1,\mathrm{},k_r;V).`$ It has the same generators as the algebra $`𝒢_R(k_1;V)\underset{𝑅}{}\mathrm{}\underset{𝑅}{}𝒢_R(k_r;V),`$, subject to the same relations in degrees $`2e_i`$ and to relations
$$\begin{array}{ccccc}𝐊𝐞𝐫((\mathrm{\Lambda }_R^{k_i}V^{}\mathrm{\Lambda }_R^{k_j}V^{})(\mathrm{\Lambda }_R^{k_j}V^{}\mathrm{\Lambda }_R^{k_i}V^{})& \stackrel{(\mathrm{𝐢𝐝},R_{k_j,k_i})}{}& \mathrm{\Lambda }_R^{k_i}V^{}\mathrm{\Lambda }_R^{k_j}V^{}& & \mathrm{\Sigma }_R^{(k_i,k_j)}V^{})\end{array}$$
in degree $`e_i+e_j`$ for all $`i>j`$. This definition is suggested by the Borel-Weil-Bott theorem (see ). In particular, for $`R=R_0`$ we get the algebra corresponding to the commutative flag variety.
We define the $`R`$-Carthesian product $`𝐆_R(k_1;V)\underset{𝑅}{\times }\mathrm{}\underset{𝑅}{\times }𝐆_R(k_r;V)`$ and the noncommutative flag variety $`\mathrm{𝐅𝐥}_R(k_1,\mathrm{},k_r;V)`$ as the noncommutative varieties corresponding to the algebras $`𝒢_R(k_1;V)\underset{𝑅}{}\mathrm{}\underset{𝑅}{}𝒢_R(k_r;V)`$ and $`_R(k_1,\mathrm{},k_r;V)`$ respectively.
To make this compatible with our definition of a noncommutative variety, we consider instead of a $`^r`$-graded algebra its diagonal subalgebra. The diagonal subalgebra is a graded algebra whose $`n`$-th graded component is the $`n(e_1+\mathrm{}+e_r)`$-graded component of the $`^r`$-graded algebra. Thus according to section 3 the category of coherent sheaves on the $`R`$-Cartesian product of Grassmannians (or the flag variety) is the category $`qgr`$ of the corresponding diagonal subalgebra.
The algebra $`_R(k_1,\mathrm{},k_r;V)`$ is the $`\mathrm{\Phi }_R`$-image of a commutative algebra in the category $`𝒯_R`$. Hence one can define the $`R`$-Carthesian product of several flag varieties.
If $`R`$ is deformation-trivial, then $`𝐆_R(k_1;V)\underset{𝑅}{\times }\mathrm{}\underset{𝑅}{\times }𝐆_R(k_r;V)`$ and $`\mathrm{𝐅𝐥}_R(k_1,\mathrm{},k_r;V)`$ are noncommutative deformations of the corresponding commutative varieties.
Note that we have a canonical embedding of the graded algebra $`𝒢_R(k_i;V)`$ into the graded algebra $`_R(k_1,\mathrm{},k_i,\mathrm{},k_r;V)`$ inducing the canonical projections
$$p_i:\mathrm{𝐅𝐥}_R(k_1,\mathrm{},k_i,\mathrm{},k_r;V)𝐆_R(k_i;V).$$
On the other hand, by definition $`_R(k_1,\mathrm{},k_r;V)`$ is a quotient algebra of the algebra $`𝒢_R(k_1;V)\underset{𝑅}{}\mathrm{}\underset{𝑅}{}𝒢_R(k_r;V)`$. Hence $`\mathrm{𝐅𝐥}_R(k_1,\mathrm{},k_r;V)`$ can be regarded as a closed subvariety in $`𝐆_R(k_1;V)\underset{𝑅}{\times }\mathrm{}\underset{𝑅}{\times }𝐆_R(k_r;V)`$.
###### Example 8.4.
The algebra $`𝒢_R(1;Z)\underset{𝑅}{}𝒢_R(2;Z)`$ corresponding to the Yang-Baxter operator $`(\text{30})`$ is generated by the elements $`z_1,z_2,z_3,z_4,Y_1,Y_2,Y_3,Y_4,D,T`$ with relations $`(\text{9})`$, $`(\text{31}),`$ and
$$\begin{array}{cccccc}[z_1,Y_2]\hfill & =& 2\mathrm{}az_4T,\hfill & [z_2,Y_1]\hfill & =& 2\mathrm{}az_3T,\hfill \\ [z_1,Y_3]\hfill & =& 2\mathrm{}bz_3T,\hfill & [z_2,Y_4]\hfill & =& 2\mathrm{}bz_4T,\hfill \\ [z_1,D]\hfill & =& 2\mathrm{}bz_3Y_42\mathrm{}az_4Y_1,\hfill & [z_2,D]\hfill & =& 2\mathrm{}az_3Y_22\mathrm{}bz_4Y_3,\hfill \end{array}$$
$`[z_1,Y_1]=[z_2,Y_2]=0,`$ $`[z_3,Y_i]=[z_3,D]=0,`$ $`[z_4,Y_i]=[z_4,D]=0,`$ $`[z_i,T]=0`$ for all $`i=1,2,3,4`$. The algebra $`_R(1,2;Z)`$ is given by the same generators subject to the same relations, as well as the additional relations
(33)
$$\left(\begin{array}{cccc}0& T& Y_2& Y_3\\ T& 0& Y_4& Y_1\\ Y_2& Y_4& 0& D\mathrm{}(a+b)T\\ Y_3& Y_1& D\mathrm{}(a+b)T& 0\end{array}\right)\left(\begin{array}{c}z_1\\ z_2\\ z_3\\ z_4\end{array}\right)=\left(\begin{array}{c}0\\ 0\\ 0\\ 0\end{array}\right).$$
As explained above, we have projections
$$\begin{array}{ccccccccc}_{\mathrm{}}& =& 𝐆_R(2;Z)& \stackrel{p}{}& \mathrm{𝐅𝐥}_R(1,2;Z)& \stackrel{q}{}& _R(Z)& =& _{\mathrm{}}^3\end{array}$$
and a closed embedding
$$\mathrm{𝐅𝐥}_R(1,2;Z)𝐆_R(2;Z)\underset{𝑅}{\times }_R(Z)=_{\mathrm{}}\underset{𝑅}{\times }_{\mathrm{}}^3.$$
### 8.7. Tautological bundles
Let $`𝒱`$ (resp. $`𝒱^{},`$ $`\mathrm{\Sigma }_R^\lambda 𝒱,`$ $`\mathrm{\Sigma }_R^\lambda 𝒱^{}`$) denote the coherent sheaf on $`𝐆_R(k;V)`$ corresponding to the free right $`𝒢_R(k;V)`$-module $`V𝒢_R(k;V)`$ (resp. $`V^{}𝒢_R(k;V),`$ $`\mathrm{\Sigma }_R^\lambda V𝒢_R(k;V),`$ $`\mathrm{\Sigma }_R^\lambda V^{}𝒢_R(k;V)`$). Since the space of global sections of the sheaf $`𝒪(1)`$ on the Grassmannian $`𝐆_R(k;V)`$ is $`\mathrm{\Lambda }_R^kV^{}`$, the maps $`\mathrm{\Lambda }_R^{k1}V^{}V\mathrm{\Lambda }_R^kV^{}`$ and $`\mathrm{\Lambda }_R^{k+1}V^{}V^{}\mathrm{\Lambda }_R^kV^{}`$ induce morphisms of sheaves
$$\begin{array}{ccccc}\mathrm{\Lambda }_R^{k1}𝒱^{}(1)& \stackrel{\varphi }{}& 𝒱\text{and}\mathrm{\Lambda }_R^{k+1}𝒱^{}(1)& \stackrel{\psi }{}& 𝒱^{}.\end{array}$$
We put $`S=𝐈𝐦\varphi ,`$ $`𝒱/S=𝐂𝐨𝐤𝐞𝐫\varphi `$, $`S^{}=𝐈𝐦\psi ,`$ $`𝒱^{}/S^{}=𝐂𝐨𝐤𝐞𝐫\psi `$.
Remark. For $`k=1`$ we have $`S=𝒪(1),`$ $`𝒱^{}/S^{}=𝒪(1)`$.
One can show that these sheaves are locally free. We refer to them as tautological bundles.
The free $`𝒢_R(k;V)`$-modules, corresponding to the sheaves $`\mathrm{\Sigma }_R^\lambda 𝒱`$, $`\mathrm{\Sigma }_R^\lambda 𝒱^{}`$ are the $`\mathrm{\Phi }_R`$-images of free modules over the corresponding algebra in the category $`𝒯_R`$. Furthermore, the morphisms $`\varphi `$ and $`\psi `$ are $`\mathrm{\Phi }_R`$-images. This implies that the $`𝒢_R(k;V)`$-modules corresponding to the tautological bundles are $`\mathrm{\Phi }_R`$-images as well. Therefore they all have a natural structure of $`𝒢_R(k;V)`$-bimodules. This allows to define $`R`$-symmetric powers $`S_R^k()`$ (resp. $`R`$-exterior powers $`\mathrm{\Lambda }_R^k()`$) of the tautological bundles as the corresponding $`\mathrm{\Phi }_R`$-images.
One can check that we have canonical isomorphisms of bimodules
$$𝒱^{}/S^{}S^{},S^{}(𝒱/S)^{}.$$
###### Example 8.5.
Let $`R`$ be the Yang-Baxter operator $`(\text{30})`$ and $`k=2`$. Let $`\stackrel{ˇ}{z}_1,\stackrel{ˇ}{z}_2,\stackrel{ˇ}{z}_3,\stackrel{ˇ}{z}_4`$ be the dual basis of $`Z`$. Then the twisted maps
$`\varphi (1):Z^{}𝒪_{𝐆_R}Z𝒪_{𝐆_R}(1),`$
$`\psi (1):Z𝒪_{𝐆_R}\mathrm{\Lambda }_R^3Z^{}𝒪_{𝐆_R}Z^{}𝒪_{𝐆_R}(1)`$
are given by
$$\begin{array}{cccc}\varphi (1):\left(\begin{array}{c}z_1\\ z_2\\ z_3\\ z_4\end{array}\right)& & \left(\begin{array}{cccc}0& D+\mathrm{}(ab)T& Y_1& Y_4\\ D\mathrm{}(ab)T& 0& Y_3& Y_2\\ Y_1& Y_3& 0& T\\ Y_4& Y_2& T& 0\end{array}\right)& \left(\begin{array}{c}\stackrel{ˇ}{z}_1\\ \stackrel{ˇ}{z}_2\\ \stackrel{ˇ}{z}_3\\ \stackrel{ˇ}{z}_4\end{array}\right)\\ & & & \\ \psi (1):\left(\begin{array}{c}\stackrel{ˇ}{z}_1\\ \stackrel{ˇ}{z}_2\\ \stackrel{ˇ}{z}_3\\ \stackrel{ˇ}{z}_4\end{array}\right)& & \left(\begin{array}{cccc}0& T& Y_2& Y_3\\ T& 0& Y_4& Y_1\\ Y_2& Y_4& 0& D\mathrm{}(a+b)T\\ Y_3& Y_1& D\mathrm{}(a+b)T& 0\end{array}\right)& \left(\begin{array}{c}z_1\\ z_2\\ z_3\\ z_4\end{array}\right)\end{array}$$
Note that $`\psi (1)\varphi =0`$ and $`\varphi (1)\psi =0`$. Hence we have isomorphisms
$$S^{}(1)𝒱/S,S(1)S^{}.$$
Note also that on the open subset $`T0`$ elements $`(z_3,z_4)`$ give a trivialization of the tautological bundle $`S^{}`$. More precisely, the restriction of the sections $`z_1,z_2`$ of $`S^{}`$ can be expressed as
(34)
$$z_1=y_4z_3y_1z_4,z_2=y_2z_3y_3z_4$$
where $`y_i=T^1Y_i`$. Similarly, the elements$`(\stackrel{ˇ}{z}_1,\stackrel{ˇ}{z}_2)`$ give a trivialization of $`𝒱/S`$ on $`T0`$. Thus the restrictions of all tautological bundles to the open subset $`T0`$ correspond to the free rank two bimodule over the Weyl algebra $`A(_{\mathrm{}}^4).`$
### 8.8. Pull-back and push-forward
Recall that we have canonical projections $`p_i:\mathrm{𝐅𝐥}_R(k_1,k_2;V)𝐆_R(k_i;V)`$ ($`i=1,2`$). Given a right graded $`𝒢_R(k_i;V)`$-module $`E`$ we consider the right bigraded $`_R(k_1,k_2;V)`$-module $`E_{𝒢_R(k_i;V)}_R(k_1,k_2;V)`$. The diagonal subspace of this module is a graded module over the diagonal subalgebra of $`_R(k_1,k_2;V)`$. This gives the pull-back functor
$$p_i^{}:coh(𝐆_R(k_i;V))coh(\mathrm{𝐅𝐥}_R(k_1,k_2;V)).$$
The pull-back functor is exact and takes a $`\mathrm{\Phi }_R`$-image to a $`\mathrm{\Phi }_R`$-image. In particular, the pull-backs of the tautological bundles have a canonical bimodule structure.
The pull-back functor $`p_i^{}`$ admits a right adjoint functor $`p_i:coh(\mathrm{𝐅𝐥}_R(k_1,k_2;V))coh(𝐆_R(k_i;V)),`$ called the push-forward functor. It also takes a $`\mathrm{\Phi }_R`$-image to a $`\mathrm{\Phi }_R`$-image.
The line bundles $`p_1^{}𝒪(i)`$ and $`p_2^{}𝒪(j)`$ on the flag variety $`\mathrm{𝐅𝐥}_R(k_1,k_2;V)`$ are $`\mathrm{\Phi }_R`$-images, hence they have a canonical bimodule structure. Therefore, we have a well-defined tensor product
$$𝒪(i,j)=p_1^{}𝒪(i)p_2^{}𝒪(j).$$
The line bundle $`𝒪(i,j)`$ is also a $`\mathrm{\Phi }_R`$-image and has a canonical bimodule structure. The $`n`$-th graded component of the corresponding module over the diagonal subalgebra of $`_R(k_1,k_2;V)`$ is the $`((n+i)e_1+(n+j)e_2)`$-graded component of the algebra $`_R(k_1,k_2;V)`$.
One can check that the push-forward of the line bundle $`𝒪(j_1,j_2)`$ with respect to $`p_2`$ is given by the formula
$$p_2𝒪(j_1,j_2)=S_R^{j_1}(S^{})(j_2).$$
### 8.9. $`\mathrm{𝐅𝐥}_R(1,2;Z)`$ as the projectivization of the tautological bundle
The $`R`$–symmetric powers of the tautological bundle form a sheaf of graded algebras on the Grassmannian $`𝐆_R(k;V)`$
$$S_R^{}(S^{})=T(S^{})/\mathrm{\Lambda }_R^2S^{}.$$
The corresponding $`𝒢_R(k;V)`$-module
$$\underset{i,j=0}{\overset{\mathrm{}}{}}\mathrm{\Gamma }(𝐆_R(k;V),S_R^j(S^{})(i))$$
is a bigraded module with a structure of a bigraded algebra. One can check that this bigraded algebra is isomorphic to the bigraded algebra $`_R(1,k;V)`$. Thus we can regard the flag variety $`\mathrm{𝐅𝐥}_R(1,k;V)`$ as the projectivization of the tautological bundle $`S`$ on the Grassmannian $`𝐆_R(k;V)`$. In particular, $`\mathrm{𝐅𝐥}_R(1,2;Z)`$ is the projectivization of the tautological bundle $`S`$ on the Grassmannian $`𝐆_R(2;Z)`$.
### 8.10. Noncommutative twistor transform
If $`E`$ is a coherent sheaf on the noncommutative projective space $`_R(Z)=_{\mathrm{}}^3,`$ we define its twistor transform as the sheaf $`p_{}q^{}E`$ on $`𝐆_R(2;Z)=_{\mathrm{}}`$, where $`q`$ is the projection $`\mathrm{𝐅𝐥}_R(1,2;Z)_R(Z)=_{\mathrm{}}^3`$ and $`p`$ is the projection $`\mathrm{𝐅𝐥}_R(1,2;Z)𝐆_R(2;Z)=_{\mathrm{}}`$. Similarly, we can define the twistor transform of a complex of sheaves on $`_{\mathrm{}}^3`$. Actually, it is more natural to consider the derived twistor transform, i.e. the derived functor of the ordinary twistor transform.
Consider a complex $`𝒞^{}`$ of the form
$$0H𝒪(1)\stackrel{M}{}K𝒪\stackrel{N}{}L𝒪(1)0$$
on the projective space $`_{\mathrm{}}^3`$. One can check that under the twistor transform one has
$$𝒪__{\mathrm{}}^3(1)0,𝒪__{\mathrm{}}^3𝒪_{𝐆_R},𝒪__{\mathrm{}}^3(1)S^{}.$$
In fact, for these sheaves the derived twistor transform coincides with the ordinary one. Thus the (derived) twistor transform takes the complex $`𝒞^{}`$ to the complex
$$0K𝒪\stackrel{𝒩}{}LS^{}0.$$
Let $``$ denote the middle cohomology of the complex $`𝒞^{}`$. It follows that the twistor transform takes $``$ to the kernel of the map $`𝒩:K𝒪LS^{}`$.
One can describe $`𝒩`$ without reference to the twistor transform. The morphism $`N`$ is the same thing as a vector space morphism
(35)
$$N_1z_1+N_2z_2+N_3z_3+N_4z_4:KZ^{}L.$$
Here the maps $`N_i`$ are given in terms of the deformed ADHM data according to (24) and (25). The map $`𝒩`$ is a composition of two maps
$$K𝒪_{𝐆_R}LZ^{}𝒪_{𝐆_R}LS^{},$$
where the first map is given by (35), while the second map comes from the canonical projection $`Z^{}𝒪_{𝐆_R}S^{}`$. (We remind that $`S^{}`$ is the cokernel of the map $`\psi :Z𝒪_{𝐆_R}(1)Z^{}𝒪_{𝐆_R}.`$)
Recall that on the open subset $`\{T0\}`$ the bundle $`S^{}`$ is trivial, and the elements $`(z_3,z_4)`$ give its trivialization (see (34)). Hence the restriction of the twistor transform of the complex (22) to this open subset is isomorphic to the complex
(36)
$$\begin{array}{ccccc}0& & K𝒪& \stackrel{\left(\begin{array}{c}N_3+y_4N_1y_2N_2\\ N_4y_1N_1y_3N_2\end{array}\right)}{}& (LL)𝒪0.\end{array}$$
Assume now that the complex (22) is given by the deformed ADHM data $`(B_1,B_2,I,J)`$ (see section (7)). Applying the formulas (24) and (25), we see that with respect to the chosen bases of $`L`$ and $`K`$ the map $`𝒩`$ is given by the matrix
$$\left(\begin{array}{ccc}B_2+y_2& B_1+y_4& I\\ B_{1}^{}{}_{}{}^{}+y_3& B_{2}^{}{}_{}{}^{}y_1& J^{}\end{array}\right).$$
It is evident that this operator is related to the operator $`𝒟`$ in (4) by a change of basis. In particular, the Nekrasov-Schwarz coordinates $`\xi _1,\xi _2,\overline{\xi }_1,\overline{\xi }_2`$ (see section 2) can be expressed through $`x_i=T^1X_i`$ as follows:
$$\begin{array}{cc}\xi _1=y_4=x_4\sqrt{1}x_3,\hfill & \xi _2=y_2=x_2+\sqrt{1}x_1,\hfill \\ \overline{\xi }_1=y_3=x_4+\sqrt{1}x_3,\hfill & \overline{\xi }_2=y_1=x_2\sqrt{1}x_1.\hfill \end{array}$$
Thus the twistor transform of the complex corresponding to the deformed ADHM data coincides with the instanton bundle corresponding to these data (see section 2). This gives a geometric interpretation of the deformed ADHM construction of the noncommutative instanton bundle.
### 8.11. Differential forms
Let an algebra $`A`$ be the $`\mathrm{\Phi }_R`$–image of a commutative algebra in the category $`𝒯_R`$. This means that there exists an operator $`R:A^2A^2`$ compatible with the multiplication law of $`A`$. Above we have defined the $`R`$–tensor product $`A\underset{𝑅}{}A`$ which is also an algebra with a Yang-Baxter operator. Explicitly, the multiplication law of $`A\underset{𝑅}{}A`$ is defined as follows. Let $`m`$ be the multiplication map from $`AA`$ to $`A`$. Then the multiplication map from $`(AA)(AA)`$ to $`AA`$ is given by $`m_{12}m_{34}R_{23}`$ in the obvious notation. It is easy to see that the multiplication map $`m`$ is a homomorphism of algebras.
Let $`I`$ denote the kernel of the map $`m:A\underset{𝑅}{}AA`$. Then $`I`$ is a two-sided ideal of the algebra $`A\underset{𝑅}{}A`$.
###### Definition 8.6.
We define the bimodule of $`R`$–differential forms of the algebra $`A`$ by
$$\mathrm{\Omega }_A^1=I/I^2.$$
For a motivation of this definition, see . Furthermore, suppose $`A`$ is a graded algebra. Consider the total grading of the bigraded algebra $`A\underset{𝑅}{}A`$. The two-sided ideal $`I`$ inherits the grading. Therefore the bimodule $`\mathrm{\Omega }_A^1`$ is graded too.
In the graded case, besides $`\mathrm{\Omega }_A^1,`$ we can define the module of projective differential forms of $`A`$ in the following way. Let $`\chi :A\underset{𝑅}{}AA\underset{𝑅}{}A`$ be the linear operator which acts on the $`(p,q)`$–th graded component of the algebra $`A\underset{𝑅}{}A`$ as a scalar multiplication by $`q`$. Since $`\chi `$ is a derivation, we have $`\chi (I^2)I`$. Therefore $`m(\chi (I^2))=0`$. Furthermore, the induced map $`\mathrm{\Omega }_A^1=I/I^2\stackrel{m\chi }{}A`$ is a morphism of graded $`A`$–bimodules.
###### Definition 8.7.
We define the $`A`$–bimodule of projective differential forms of the algebra $`A`$ by
$$\widehat{\mathrm{\Omega }}_A^1=𝐊𝐞𝐫(\mathrm{\Omega }_A^1\stackrel{m\chi }{}A).$$
First, let us apply this construction of differential forms to the noncommutative affine variety $`_{\mathrm{}}^4`$ (subsection 3.4). The algebra $`A(_{\mathrm{}}^4)`$ of polynomial functions on $`_{\mathrm{}}^4`$ is the Weyl algebra:
$$A(_{\mathrm{}}^4)=\mathrm{T}(x_1,x_2,x_3,x_4)/[x_i,x_j]=\mathrm{}\theta _{ij}_{1i,j4}.$$
Let us define the Yang-Baxter operator on the tensor square of the subspace of $`A(_{\mathrm{}}^4)`$ spanned by $`1,x_1,x_2,x_3,x_4`$ by the formula
$$1x_ix_i1,x_i11x_i,x_ix_jx_jx_i+\mathrm{}\theta _{ij}11\text{for all}1i,j4.$$
This Yang-Baxter operator has a unique extension to the whole $`A(_{\mathrm{}}^4)`$ compatible with the multiplication law.
There is another way to look at this Yang-Baxter operator. Recall that $`_{\mathrm{}}^4`$ is an open subset $`T0`$ in the noncommutative Grassmannian $`𝐆_R(2;Z)`$ where $`R`$ is defined by (30). The Yang-Baxter operator on the quadratic algebra $`𝒢_R(2;Z)`$ has the property that $`R(Ta)=aT`$ for any $`a𝒢_R(2;Z)`$. Hence it descends to a Yang-Baxter operator on $`A(_{\mathrm{}}^4)`$. It is easy to see that it acts on the tensor square of the subspace spanned by $`1,x_1,x_2,x_3,x_4`$ in the above manner.
We define the sheaf of differential forms $`\mathrm{\Omega }__{\mathrm{}}^4^1`$ as the bimodule of $`R`$–differential forms of the algebra $`A(_{\mathrm{}}^4)`$. It is easy to check that $`\mathrm{\Omega }__{\mathrm{}}^4^1`$ is isomorphic to the bimodule $`A(_{\mathrm{}}^4)^4`$. Futhermore, we can take any $`R`$–exterior power of $`\mathrm{\Omega }__{\mathrm{}}^4^1`$ and thereby define $`\mathrm{\Omega }__{\mathrm{}}^4^p`$. This enables us to define a connection and its curvature on any bundle on the noncommutative affine space. The relevant formulas were written above (see subsection 1.5).
Second, we define the sheaf of differential forms $`\mathrm{\Omega }_{𝐆_R}^1`$ on the noncommutative Grassmannian $`𝐆_R(k;V)`$ as the sheaf corresponding to the module of projective differential forms $`\widehat{\mathrm{\Omega }}_{𝒢_R}^1`$.
It can be shown that as in the commutative case we have an isomorphism of coherent sheaves on the noncommutative Grassmannian $`𝐆_R(k;V)`$:
$$\mathrm{\Omega }_{𝐆_R}^1SS^{}.$$
It follows that for $`k=1`$ that we have an exact sequence
$$0\mathrm{\Omega }_{_R(V)}^1𝒱^{}(1)𝒪0.$$
Thus this definition of the sheaf of differential forms $`\mathrm{\Omega }_{_R(V)}^1`$ is consistent with Definition 4.8.
Similarly, one can define the sheaf of differential forms $`\mathrm{\Omega }_{\mathrm{𝐅𝐥}_R}^1`$ on the noncommutative flag variety $`\mathrm{𝐅𝐥}_R(k_1,\mathrm{},k_r;V)`$. One can check that the projection
$$p_i:\mathrm{𝐅𝐥}_R(k_1,\mathrm{},k_i,\mathrm{},k_r;V)𝐆_R(k_i;V)$$
induces a morphism of bundles $`p_i^{}:\mathrm{\Omega }_{𝐆_R}^1\mathrm{\Omega }_{\mathrm{𝐅𝐥}_R}^1`$.
In the commutative case the ADHM construction of the instanton connection can be interpreted in terms of twistor transform (see for details). We believe that this can be done in the noncommutative case as well. It appears that the most convenient definition of connection on a bundle on a noncommutative projective variety is in terms of jet bundles (see, for example, ).
## 9. Instantons on a $`q`$– deformed $`^4`$
In this paper we have focused on a particular noncommutative deformation of $`^4`$ related to the Wigner-Moyal product (3). This is the only deformation of $`^4`$ which is known to arise in string theory. But most of our constructions work for more general deformations which do not have a clear physical interpretation. For example, let us replace $`_{\mathrm{}}^4`$ with a noncommutative affine variety whose coordinate ring is generated by $`z_1,z_2,z_3,z_4`$ subject to the following quadratic relations:
$$qz_1z_2q^1z_2z_1=\mathrm{},qz_3z_4q^1z_4z_3=\mathrm{},[z_1,z_3]=[z_1,z_4]=[z_2,z_3]=[z_2,z_4]=0.$$
We will denote this noncommutative affine variety by $`_{q,\mathrm{}}^4,`$ and its coordinate algebra by $`𝒜_{q,\mathrm{}}`$. If $`\mathrm{}`$ and $`q`$ are real, we can define a $``$-operation on $`𝒜_{q,\mathrm{}}`$ by $`z_1^{}=z_2,z_3^{}=z_4`$. The corresponding real noncommutative affine variety will be denoted by $`_{q,\mathrm{}}^4.`$
Consider now the following deformation of the ADHM equations:
(37)
$$[B_1,B_2]_{q^1}+IJ=0,[B_1,B_1^{}]_{q^1}+[B_2,B_2^{}]_q+II^{}J^{}J=2\mathrm{}1_{k\times k}.$$
Here $`B_1,B_2\mathrm{Hom}(V,V),`$ $`I\mathrm{Hom}(W,V),`$ $`J\mathrm{Hom}(V,W)`$, as usual, and by $`[A,B]_q`$ we mean a $`q`$-commutator:
$$[A,B]_q=qABq^1BA.$$
We claim that solutions of these “$`q`$– deformed” ADHM equations can be used as an input for the construction of instantons on $`_{q,\mathrm{}}^4`$ of rank $`r=dimW`$ and instanton charge $`k=dimV`$. Let us sketch this construction. Define an operator
$$𝒟\mathrm{Hom}_{𝒜_{q,\mathrm{}}}((VVW)_{}𝒜_{q,\mathrm{}},(VV)_{}𝒜_{q,\mathrm{}})$$
by the formula
$$𝒟=\left(\begin{array}{ccc}B_1qz_1& qB_2+qz_2& I\\ B_2^{}\overline{z}_2& qB_1^{}\overline{z}_1& J^{}\end{array}\right).$$
Now we can go through the same manipulations as in section 2: assume that $`𝒟`$ is surjective, and its kernel is a free module, and define a connection 1-form by the expression (5). The same formal computation as in section 2 shows that the curvature of this connection is anti-self-dual.
In order to ensure that $`𝒟`$ is surjective, it is probably necessary to replace the algebra $`𝒜_{q,\mathrm{}}`$ with some bigger algebra containing $`𝒜_{q,\mathrm{}}`$ as a subalgebra. This bigger algebra should play the role of the algebra of smooth functions on our noncommutative $`^4`$. For $`\mathrm{}=0`$, $`q1`$ there is even a natural candidate for this bigger algebra: it should consist of $`C^{\mathrm{}}`$ functions on $`^2`$ with some suitable growth conditions at infinity and the product defined by
(38)
$$\begin{array}{c}(fg)(z_1,z_2,\overline{z}_1,\overline{z}_2)=\hfill \\ \hfill \mathrm{exp}\left(\mathrm{ln}(q)\left(z_1\overline{z}_1^{}\frac{^2}{z_1\overline{z}_1^{}}+z_2\overline{z}_2^{}\frac{^2}{z_2\overline{z}_2^{}}z_1^{}\overline{z}_1\frac{^2}{z_1^{}\overline{z}_1}z_2^{}\overline{z}_2\frac{^2}{z_2^{}\overline{z}_2}\right)\right)\\ \hfill f(z_1,z_2,\overline{z}_1,\overline{z}_2)g(z_1^{},z_2^{},\overline{z}_1^{},\overline{z}_2^{})|_{z_1^{}=z_1,z_2^{}=z_2}.\end{array}$$
Assuming that this formal expression exists, it is easy to check that the product is associative, that polynomial functions form a subalgebra with respect to it, and that this subalgebra is isomorphic to $`𝒜_{q,\mathrm{}}`$.
It is natural to conjecture that all instantons on $`_{q,\mathrm{}}^4`$ arise from this deformed ADHM construction. Note that in this case the deformed ADHM equations are not hyperkähler moment map equations, and one cannot use the hyperkähler quotient construction to infer the existence of hyperkähler metric on the quotient space.
The algebro-geometric part of the story can also be generalized. We did not go through this carefully, but nevertheless would like to indicate one result. It appears that the $`q`$– deformed ADHM data can be interpreted in terms of sheaves on a more general noncommutative $`^2`$ than the one defined in section 3. The graded algebra corresponding to this noncommutative $`^2`$ is generated by degree one elements $`z_1,z_2,z_3`$ with the quadratic relations
$$qz_1z_2q^1z_2z_1=2\mathrm{}z_3^2,[z_i,z_3]=0,i=1,2.$$
This algebra is one of the Artin-Schelter regular algebras of dimension three . It is characterized by the fact that the corresponding noncommutative variety $`_{q,\mathrm{}}^2`$ contains as subvarieties a commutative quadric and a noncommutative line. The latter is given by the equation $`z_3=0`$. In the limit $`q1`$ the plane $`_{q,\mathrm{}}^2`$ reduces to $`_{\mathrm{}}^2`$, and the union of the quadric and the line turns into the triple commutative line $`l`$ which played such a prominent role in this paper. If $`q1`$, then in the limit $`\mathrm{}0`$ the quadric turns into a union of two intersecting commutative lines $`z_1=0`$ and $`z_2=0`$.
For any $`q`$ the line $`z_3=0`$ should be regarded as “the line at infinity” (which is noncommutative for $`q1`$). It is plausible that the $`q`$– deformed ADHM data are in one-to-one correspondence with bundles, or maybe torsion–free sheaves, on $`_{q,\mathrm{}}^2`$ with a trivialization on this line.
## 10. Appendix
In this section we define a $``$–product on the space of complex-valued $`C^{\mathrm{}}`$ functions on $`^n`$ whose derivatives of arbitrary order are polynomially bounded. The $``$–product endows this space with a structure of a $``$-algebra and reduces to the Wigner-Moyal product (3) on polynomial functions.
###### Definition 10.1.
Let $`\mathrm{\Phi }`$ be a topological vector space which is a subspace of the space of $`C^{\mathrm{}}`$ functions on $`^n,`$ and let $`\mathrm{\Phi }^{}`$ be the space of distributions on $`\mathrm{\Phi }`$. Let $`f`$ be a $``$-valued function on $`^n`$ which simultaneously is a distribution in $`\mathrm{\Phi }^{}`$. $`f`$ is called a multiplier if for any $`\varphi \mathrm{\Phi }`$ $`f\varphi \mathrm{\Phi }`$.
The set of multipliers of $`\mathrm{\Phi }^{}`$ is obviously a subspace of $`\mathrm{\Phi }^{}.`$
###### Definition 10.2.
Let $`f\mathrm{\Phi }^{}`$. $`f`$ is called a convolute if for any $`\varphi \mathrm{\Phi }`$ we have
$$(f\varphi )(x)(f(\xi ),\varphi (x+\xi ))\mathrm{\Phi },$$
and this expression depends continuously on $`\varphi `$. The above expression is called the convolution of $`f`$ with $`\varphi `$.
The set of convolutes is obviously a subspace of $`\mathrm{\Phi }^{}.`$
We will denote the Fourier duals of $`\mathrm{\Phi }`$ and $`\mathrm{\Phi }^{}`$ by $`\stackrel{~}{\mathrm{\Phi }}`$ and $`\stackrel{~}{\mathrm{\Phi }^{}},`$ respectively. If $`f\mathrm{\Phi },`$ then $`\stackrel{~}{f}\stackrel{~}{\mathrm{\Phi }}`$ will be the Fourier transform of $`f,`$ etc.
###### Definition 10.3.
The Schwartz space $`𝒮(^n)`$ is the space of $``$-valued $`C^{\mathrm{}}`$ functions on $`^n`$ such that $`\varphi 𝒮`$ if and only if all the norms
(39)
$$\underset{x}{sup}x^kD^m\varphi (x),k=0,1,2,\mathrm{},$$
are finite. Here $`m=(m_1,\mathrm{},m_n)`$ is an arbitrary polyindex.
Convergence on $`𝒮`$ is defined using the family of norms (39). Then $`𝒮`$ becomes a complete countably normed space .
###### Proposition 10.4.
A function $`f𝒮^{}`$ is a multiplier if and only if it is a $`C^{\mathrm{}}`$ function on $`^n`$ all of whose derivatives are polynomially bounded.
###### Proof.
Obvious. ∎
The following theorem proved in describes the subspace of convolutes of $`𝒮^{}`$:
###### Theorem 10.5.
A distribution $`f𝒮^{}`$ is a convolute if and only if it has the form
$$f=\underset{|\alpha |<r}{}D^\alpha f_\alpha (x),$$
where $`r`$ is a positive integer, and $`f_\alpha `$ are $`C^0`$ functions on $`^n`$ which decrease at infinity faster than any negative power of $`x`$.
The functions which decrease at infinity faster than any negative power will be called rapidly decreasing.
The following theorem is proved in , vol. 2, ch. III:
###### Theorem 10.6.
Fourier transform and its inverse act as automorphisms on both $`𝒮`$ and $`𝒮^{}`$.
From now on we identify $`𝒮\stackrel{~}{𝒮},`$ $`𝒮^{}\stackrel{~}{𝒮^{}}.`$
###### Theorem 10.7.
Fourier transform and its inverse establish an isomorphism between the space of multipliers and the space of convolutes of $`𝒮^{}`$.
###### Proof.
By the preceding theorem, it is sufficient to show that the Fourier transform of every multiplier is a convolute, and vice versa. The former fact is proved in , vol. 2, ch. III. Let us prove the converse.
By theorem 10.5, every convolute has the form
$$f(x)=\underset{|\alpha |<r}{}D^\alpha f_\alpha (x)$$
for some $`r`$ and rapidly decreasing continuous functions $`f_\alpha `$. Let
$$\stackrel{~}{f}_\alpha (p)=f_\alpha (x)e^{\sqrt{1}px}d^nx.$$
be the Fourier transform of $`f_\alpha (x).`$ Since the integrals
$$x^\beta f_\alpha (x)e^{\sqrt{1}px}d^nx$$
are absolutely convergent, the functions $`\stackrel{~}{f}_\alpha `$ are $`C^{\mathrm{}}`$ functions. Furthermore, the Fourier transform of $`f`$ is equal to
$$\stackrel{~}{f}(p)=\underset{|\alpha |<r}{}(\sqrt{1}p)^\alpha \stackrel{~}{f}_\alpha (p)$$
(see , vol. 2, ch. III), hence $`\stackrel{~}{f}`$ is also a $`C^{\mathrm{}}`$ function. Finally, since by the preceding theorem the Fourier transform of any element of $`𝒮^{}`$ is again an element of $`𝒮^{},`$ $`\stackrel{~}{f}`$ and all its derivatives are polynomially bounded. Hence $`\stackrel{~}{f}`$ is a multiplier. ∎
###### Definition 10.8.
Let $`\omega `$ be a skew-symmetric real-valued bilinear form on $`^n`$. The $``$–product on the space of convolutes of $`𝒮^{}`$ is defined by
$$(\stackrel{~}{f}\stackrel{~}{g})(p)=\stackrel{~}{f}(q)\stackrel{~}{g}(pq)e^{\sqrt{1}\omega (p,q)}\frac{d^nq}{(2\pi )^n}.$$
###### Theorem 10.9.
The $``$–product is well-defined and makes the space of convolutes of $`𝒮^{}`$ into an algebra over $``$.
###### Proof.
We will prove that the $``$–product of two convolutes of $`𝒮^{}`$ is well-defined, and is again a convolute of $`𝒮^{}`$. The rest is obvious.
It is sufficient to consider the case when
$$\stackrel{~}{f}(p)=D^\alpha \stackrel{~}{f}_0(p),\stackrel{~}{g}(p)=D^\beta \stackrel{~}{g}_0(p).$$
Then, integrating by parts, we may rewrite the $``$–product in the following form:
$$(1)^{|\alpha |}\stackrel{~}{f}_0(q)\frac{^\alpha }{q^\alpha }\left[\frac{^\beta }{p^\beta }\stackrel{~}{g}_0(pq)e^{\sqrt{1}\omega (p,q)}\right]\frac{d^nq}{(2\pi )^n}.$$
Derivatives acting on the exponential bring down powers of $`p,`$ so the integral can be rewritten as
$$P(p,\frac{}{p})\stackrel{~}{f}_0(q)\frac{^\beta }{p^\beta }\stackrel{~}{g}_0(pq)e^{\sqrt{1}\omega (p,q)}\frac{d^nq}{(2\pi )^n},$$
where $`P(u,v)`$ is a homogeneous polynomial of degree $`|\alpha |`$. We now use the Leibniz rule repeatedly to rewrite the expression above as
$$P(p,\frac{}{p})Q(q,\frac{}{p})\left[\stackrel{~}{f}_0(q)\stackrel{~}{g}_0(pq)e^{\sqrt{1}\omega (p,q)}\right]\frac{d^nq}{(2\pi )^n},$$
where $`Q(u,v)`$ is a homogeneous polynomial of degree $`|\beta |`$. Because both $`\stackrel{~}{f}_0`$ and $`\stackrel{~}{g}_0`$ are rapidly decreasing, the integral converges absolutely and defines a $`C^0`$ function of $`q`$ which is rapidly decreasing. Hence the $``$–product of $`\stackrel{~}{f}_0`$ and $`\stackrel{~}{g}_0`$ has the form
$$\underset{|m||\alpha |+|\beta |}{}D^m\stackrel{~}{h}_m(p),$$
where the functions $`\stackrel{~}{h}_m(p)`$ are continuous and rapidly decreasing. It follows that the space of convolutes is closed under the $``$-product. ∎
###### Corollary 10.10.
The space of multipliers of $`𝒮^{}`$ inherits a product from the $``$–product on the space of convolutes of $`𝒮^{}`$, and this product makes the space of multipliers into an algebra over $`.`$ Polynomials form a subalgebra of this algebra isomorphic to the Weyl algebra with generators $`x_i,i=1,\mathrm{},n,`$ and relations
$$[x_i,x_j]=2\sqrt{1}\omega _{ij}.$$
###### Proof.
The first statement is an immediate consequence of theorems 10.7 and 10.9. The second statement follows from a simple computation. ∎
It is this product on the space of multipliers that we call the Wigner-Moyal product and denote with $``$.
## Acknowledgements
We are grateful to A. Beilinson, V. Ginzburg, L. Katzarkov, N. Nekrasov, T. Pantev, and A. Yekutieli for useful discussions. We also wish to thank the Institute for Advanced Study, Princeton, NJ, for a very stimulating atmosphere. |
warning/0002/astro-ph0002315.html | ar5iv | text | # Perturbative Analysis of Adaptive Smoothing Methods in Quantifying Large-Scale Structure
## 1 Introduction
The large-scale distribution of galaxies is one of the most important sources to study formation and evolution of cosmic structures. Now there are two ongoing large-scale redshift surveys of galaxies, the Sloan Digital Sky Survey (SDSS, Gunn & Weinberg 1995) and the Anglo-Australian Telescope 2dF Survey (Colless 1998). These surveys will bring us enormous information of three-dimensional galaxy distribution and are expected to revolutionarily improve our knowledge of the large-scale structure in the universe.
The two-point correlation function of galaxies, or its Fourier transform, the power spectrum are simple as well as powerful tools, and have been widely used to quantify clustering of galaxies (e.g. Totsuji & Kihara 1969, Peebles 1974, Davis & Peebles 1983). These two quantities are based on the second-order moments of matter fluctuations. When the fluctuations are Random Gaussian distributed, the two-point correlation function or the power spectrum contains full statistical information of the fluctuations. Even though the initial seed of structure formation is often assumed to be random Gaussian distributed as predicted by the standard inflation scenarios (Guth & Pi 1982, Hawking 1982, Starobinsky 1982), this simple assumption has not been observationally established. Moreover, the cosmic structures observed today are more or less affected by nonlinear gravitational evolution. Therefore, only with these two quantities, we cannot study the large-scale structure properly.
Other statistical methods have been proposed and is expected to play complimentary roles to the traditional analyses based on the second-order moment. The skewness parameter characterizes the asymmetry of one point probability distribution function (PDF) of the density field and has been investigated in deep (Peebles 1980, Fry 1984, Juszkiewicz, Bouchet & Colombi 1993, Bernardeau 1994, Scoccimarro 1998, Seto 1999). Beside higher-order moments (such as skewness), there are other statistical approaches designed to directly measure the geometrical or morphological aspects of galaxy clustering. Connectivity of the isodensity contour is an interesting target for these approaches. For example, the genus statistics were proposed by Gott, Melott & Dickinson (1986), the area statistics (equivalently the level crossing statistics) by Ryden (1988), and percolation analysis by Klypin (1988). In addition, the Minkowski functionals recently attain much attention (e.g. Minkowski 1903, Mecke, Buchert & Wagner 1994, Schmalzing & Buchert 1997, Kerscher et al. 1997).
To analyze observed cosmic structures, smoothing operation becomes crucially important in some cases. The observed galaxies are distributed in point-like manner, but geometrical or morphological analyses, such as, the genus statistics, are usually based on continuous (smoothed) density field (see also Babul & Starkman 1992, Luo & Vishniac 1995). We have traditionally used filters with spatially fixed smoothing radius for analyzing the large-scale structure. Even though this method is the simplest from theoretical point of views, other possibilities are worth investigated. The local statistical fluctuations due to the discreteness of particles are determined by the number of particles contained in the smoothing kernel. If we use a spatially fixed filter, we can measure smoothed quantities at overdensity regions relatively more accurately than at underdense regions. As a result, quality of information becomes inhomogeneous. This inhomogeneity is caused by the simple choice to spatially fix the smoothing radius. There must be more efficient methods to resolve cosmic structure from particle distribution (Hernquist & Katz 1989). Actually, Springel et al. (1998) have pointed out that the signal to noise ratio of the genus statistics is considerably improved by using adaptive smoothing methods. Adaptive methods are based on Lagrangian description, use nearly same number of “particles” (mass elements) to construct smoothed density field (Hernquist & Katz 1989) and are expected to be less affected by discreteness of mass elements. Therefore, it seems reasonable that we can resolve cosmic structures more efficiently, using these methods.
In this article, we perturbatively analyzed quantitative effects caused by adaptive smoothing methods. We pay special attention to three representative examples, the skewness parameter, the genus statistics and the area statistics. As the skewness is basically defined by the one point PDF, we can, in principle, discuss it without making continuous density field. But its analysis is very instructive to see nonlinear effects accompanied with adaptive smoothing methods. As a first step, we mainly study the density field in real space and do not discuss the effects of biasing (Kaiser 1984, Bardeen, Bond, Kaiser & Szalay 1986, Dekel & Lahav 1999).
This article is organized as follows. In §2 we describe basic properties of the adaptive smoothing and introduce its two main approaches. Then perturbative formulas are derived for each of them. In §3 we discuss the skewness parameters of the density field smoothed by these two adaptive methods. We evaluate them using second-order perturbation theory. Some of results in this section can be straightforwardly applied to the density field in the redshift space. Then we discuss weakly nonlinear effects of the genus and the area statistics using the multidimensional Edgeworth expansion method explored by Matsubara (1994). To characterize the isodensity contour, parameterization based on the volume fraction above a given density threshold is often adopted. In §4.1 we discuss this parameterization in perturbative manner. In §4.2 we explicitly evaluate the generalized skewness which is closely related to the genus and the area statistics. In §4.3 and §4.4, we show the weakly nonlinear effects on these two statistics with various smoothing methods. We make a brief summary in §5.
## 2 Adaptive Smoothing Method
The (unsmoothed) density contrast field $`\delta (𝒙)`$ at a point $`𝒙`$ is defined in terms of the mean density of the universe $`\overline{\rho }`$ and the local density $`\rho (𝒙)`$ as
$$\delta (𝒙)=\frac{\rho (𝒙)\overline{\rho }}{\overline{\rho }}.$$
(1)
In this article we assume that the primordial density fluctuations obey Random Gaussian distribution which is completely characterized by the (linear) matter power spectrum. Unless we state explicitly, we limit our analysis in the real space density field. But some of our results are straightforwardly applied to the redshift space quantities, as shown in the next section.
Isotropic filters with spatially constant smoothing radius $`R`$ have been traditionally used to obtain continuous smoothed density field $`\delta _{FR}(𝒙)`$ as follows
$$\delta _{FR}(𝒙)=𝑑𝒙^{}\delta (𝒙^{})W(|𝒙^{}𝒙|;R).$$
(2)
Here the function $`W(𝒙;R)`$ is a spatial filter function and the subscript $`F`$ indicate the fixed smoothing. Most of theoretical analyses in the large-scale structure have been based on this fixed smoothing method. As for the functional shape of $`W(𝒙;R)`$, two kinds of functions are often used (e.g. Bardeen et al. 1986, Matsubara 1995). One is the Gaussian filter and defined as
$$W(𝒙;R)=(2\pi )^{3/2}R^3\mathrm{exp}\left(𝒙^2/2R^2\right).$$
(3)
The other one is the top-hat filter and has a compact support as
$`W(𝒙;R)=\{\begin{array}{cc}3/(4\pi R^3)\hfill & (|𝒙|R)\hfill \\ 0\hfill & (|𝒙|>R).\hfill \end{array}`$ (4)
In this article we mainly use the Gaussian filter. This filter is useful for quantifying the large-scale structure from observed noisy data sets. In addition, algebraic manipulations for the Gaussian filter are generally much simpler than for the top-hat filter.
Next let us discuss the basic properties of adaptive smoothing methods (Hernquist & Katz 1989, Thomas & Couchman 1992, Springel et al. 1998). The essence of these methods is to change the smoothing radius $`R`$ as a function of position $`𝒙`$ according to its local density contrast. With a given spherically symmetric kernel $`W`$, we determine the smoothing radius $`R(𝒙)`$ so that the total mass included within the kernel becomes constant.
$$\overline{\rho }𝑑𝒙^{}(1+\delta (𝒙^{}))R(𝒙)^3W(|𝒙^{}𝒙|;R(𝒙))=\overline{\rho }R^3.$$
(5)
The radius $`R(𝒙)`$ becomes smaller than the standard value $`R`$ in a overdense region and becomes larger in a underdense region. In a system constituted by equal mass particles as in standard N-body simulations, the smoothing radius $`R(𝒙)`$ is determined so that the total number of particles in a filter becomes constant. Thus adaptive smoothing is basically Lagrangian description and their smoothing radii are closely related to the resolution of spatial structures.
We can solve the variable smoothing radius $`R(𝒙)`$ in equation (5) by perturbatively expanding the deviation $`\delta R(𝒙)R(𝒙)R`$. In this procedure we regard the density contrast $`\delta `$ as the order parameter of the perturbative expansion. After some calculations we obtain the first-order solution as follows
$$\delta R(𝒙)=\frac{1}{3}\delta _{FR}(𝒙)R+O(\delta ^2).$$
(6)
This simple result seems quite reasonable with the relation below.
$$R(𝒙)^3(1+\delta _{FR}(𝒙))=R^3(1+O(\delta ^2)).$$
(7)
This relation roughly shows that the total mass within the smoothing radius $`R(𝒙)`$ does not depend on position $`𝒙`$.
With the variable smoothing radius $`R(𝒙)`$ (solution for eq.) we can practice adaptive smoothing. As pointed out by Hernquist & Katz (1989) for the smoothed particle hydrodynamics (SPH), there exist two different methods (gather and scatter approaches) to assign the smoothed density contrast field at each point $`𝒙`$. The gather approach is simply use the solution $`R(𝒙)`$ at the point $`𝒙`$ in interest and the smoothed field is formally written as
$$\delta _{GR}(𝒙)=𝑑𝒙^{}\delta (𝒙^{})W(|𝒙^{}𝒙|;R(𝒙))C(R),$$
(8)
the subscript $`G`$ indicates the gather approach. In this case, the volume average of the first term in the right hand side dose not vanish and we have added a term $`C(R)`$ so that the total volume average of $`\delta _{GR}(𝒙)`$ becomes zero.
In the scatter approach we use the solution $`R(𝒙^{})`$ for each point where a mass element exists. We can write down the smoothed field at $`𝒙`$ as
$$\delta _{SR}(𝒙)=𝑑𝒙^{}\delta (𝒙^{})W(|𝒙^{}𝒙|;R(𝒙^{})),$$
(9)
the subscript $`S`$ represents the scatter approach. In this case the volume average becomes zero. We only dilute the mass element at point $`𝒙^{}`$ with the density profile proportional to $`W(|𝒙^{}𝒙|;R(𝒙^{}))`$ around that point. Note that the spatial dependence of the smoothing radius $`R()`$ is different between equations (8) and (9).
Next we evaluate equations for $`\delta _{GR}(𝒙)`$ and $`\delta _{SR}(𝒙)`$ up to second-order of the density contrast $`\delta `$ using perturbative solution of the smoothing radius $`R(𝒙)=R+\delta R(𝒙)`$ given in equation (6). The results are given as
$`\delta _{GR}(𝒙)`$ $`=`$ $`\delta _{FR}(𝒙){\displaystyle \frac{1}{3}}\delta _{FR}(𝒙)R{\displaystyle \frac{}{R}}\delta _{FR}(𝒙)+{\displaystyle \frac{1}{6}}R{\displaystyle \frac{d}{dR}}\sigma _R^2+O(\delta ^3),`$ (10)
$`\delta _{SR}(𝒙)`$ $`=`$ $`\delta _{FR}(𝒙){\displaystyle \frac{R}{3}}{\displaystyle 𝑑𝒙^{}_RW(|𝒙^{}𝒙|;R)\delta (𝒙^{})\delta _{FR}(𝒙^{})}+O(\delta ^3).`$ (11)
The formula for the scatter approach is somewhat complicated, compared with the gather approach. Also in numerical analysis, the scatter approach requires higher computational costs (Springel et al. 1998). This reflects nonlocal character of the smoothing radius.
Equations (10) and (11) show apparently that the corrections due to the adaptive methods start from second-order of $`\delta `$. Therefore, their effects are expected to be comparable to second-order (nonlinear) effects predicted by cosmological gravitational perturbation theory (Peebles 1980, Fry 1984, Goroff et al. 1986). Adaptive smoothing methods modify the quantities which characterize the nonlinear mode couplings, such as the skewness parameter of density field.
If we use the Gaussian filter, the leading-order correction for the scatter approach is expressed as follows
$$\frac{d𝒌}{(2\pi )^3}\frac{d𝒍}{(2\pi )^3}\mathrm{exp}\left[\frac{(2𝒍^2+𝒌^2+2𝒌𝒍)R^2}{2}\right]\delta (𝒌)\delta (𝒍)\mathrm{exp}[i(𝒌+𝒍)𝒙]\frac{(𝒌+𝒍)^2R^2}{3},$$
(12)
where $`\delta (𝒌)`$ and $`\delta (𝒍)`$ are the Fourier coefficients of the density contrast and defined as
$$\delta (𝒌)=𝑑𝒙\delta (𝒙)\mathrm{exp}(i𝒌𝒙).$$
(13)
Formula (12) is useful to quantitatively evaluate nonlinear effects caused by the scatter approach.
## 3 Skewness
In this section we investigate modifications of the skewness parameter $`S`$ caused by the two adaptive methods. Skewness is a fundamental quantities to characterize asymmetry of the one point PDF of the density field (Peebles 1980, Fry 1984, Juszkiewicz, Bouchet & Colombi 1993, Bernardeau 1994). It is defined as
$$S=\frac{\delta ^3}{\sigma ^4},$$
(14)
where the angular bracket $``$ represents to take the ensemble average and $`\sigma (\delta ^2^{1/2})`$ is the rms fluctuation of $`\delta `$. Here, we discuss the leading-order contributions for the numerator $`\delta ^3`$ and denominator $`\sigma ^4`$. As we have already commented in §1, the skewness parameter can be discussed without making continuous density field. It can be basically defined by the count probability distribution function, and spatial relation between one region and another one is unnecessary (e.g. Gazta$`\stackrel{~}{\mathrm{n}}`$aga 1992, Bouchet et al. 1993, Kim & Strauss 1998, and references therein, see also Colombi, Szapudi & Szalay 1998). Therefore our effort in this article to resolve cosmic structures by using the adaptive smoothing might be irrelevant for observational determination of the skewness parameter. But perturbative analysis in this section is very useful to grasp nonlinear effects caused by the adaptive smoothing methods and become basis for studying statistics of isodensity contours such as the genus statistics or the area statistics discussed in the next section.
The leading-order contribution for the rms fluctuation $`\sigma `$ is written in terms of the linear (primordial) power spectrum $`\delta _1(𝒌)\delta _1(𝒍)=(2\pi )^3\delta _{Drc}(𝒌+𝒍)P(k)`$ ($`\delta _1(𝒌)`$ : linear mode, $`\delta _{Drc}()`$: Dirac’s delta function). With the Fourier transformed filter function $`w(kR)`$ we have
$$\sigma _R^2=\delta _R^2(𝒙)=\frac{d𝒌}{(2\pi )^3}P(k)w(kR)^2,$$
(15)
where the suffix $`R`$ is added to explicitly indicate the smoothing radius $`R`$. Throughout in this article, we use power-law spectra $`P(k)`$ as
$$P(k)=Ak^n,3<k1.$$
(16)
for these scale-free models the normalization factor $`A`$ becomes irrelevant and we can simply put $`A=1`$ below. Furthermore, as shown later, the skewness parameter does not depend on the smoothing radius in our leading-order analysis. From equation (15) we have the variance $`\sigma _R^2`$ for the Gaussian filter as
$$\sigma _R^2=_0^{\mathrm{}}\frac{dk}{2\pi ^2}k^{n+2}e^{k^2R^2}=\frac{R^{n3}}{(2\pi )^2}\mathrm{\Gamma }\left(\frac{3+n}{2}\right).$$
(17)
The integral (15) logarithmically diverges for $`n=3`$, but skewness $`S`$ is well-behaved in the limit $`n3`$ from above. As it shows interesting behavior at this specific spectral index, we also discuss quantities at $`n=3`$ regarding them as the limit values.
Calculation of the third-order moment $`\delta ^3`$ is more complicated than that of the variance $`\sigma ^2`$ discussed so far. When the initial fluctuation is random Gaussian distributed as assumed in this article, the linear contribution for the third-order moment becomes exactly zero due to the symmetric distribution of the density contrast $`\delta `$ around the origin $`\delta =0`$. Nonlinear mode couplings induce asymmetry in this distribution. Therefore, we resort to higher-order perturbation theory. The leading-order contribution for the skewness parameter without smoothing operation is given by Peebles (1980) in the case of Einstein de-Sitter background as
$$S=\frac{34}{7}.$$
(18)
It is convenient to use the Fourier space representation to calculate the third-order moment for the smoothed density field. Following the standard procedure, we expand a nonlinear Fourier modes of overdensity $`\delta `$ and the (irrotational) peculiar velocity field $`𝑽`$ as (Fry 1984, Goroff et al. 1986)
$`\delta (𝒙)`$ $`=`$ $`\delta _1(𝒙)+\delta _2(𝒙)+\mathrm{},`$
$`𝑽(𝒙)`$ $`=`$ $`𝑽_1(𝒙)+𝑽_2(𝒙)+\mathrm{},`$ (19)
where $`\delta _1(𝒙)`$ and $`𝑽_1(𝒙)`$ are the linear modes and $`\delta _2(𝒙)`$ and $`𝑽_2(𝒙)`$ the second-order modes. We perturbatively solve the continuity, Euler and Poisson equations,
$`{\displaystyle \frac{}{t}}\delta (𝒙)+{\displaystyle \frac{1}{a}}[𝑽(𝒙)\{1+\delta (𝒙)\}]`$ $`=`$ $`0,`$
$`{\displaystyle \frac{}{t}}𝑽(𝒙)+{\displaystyle \frac{1}{a}}[𝑽(𝒙)]𝑽(𝒙)+{\displaystyle \frac{_ta}{a}}𝑽(𝒙)+{\displaystyle \frac{1}{a}}\varphi (𝒙)`$ $`=`$ $`0,`$
$`^2\varphi (𝒙)4\pi a^2\rho (t)\delta (𝒙)`$ $`=`$ $`0,`$
where $`a`$ represents the scale factor. The second-order solution in $`𝒌`$-space is given as
$$\delta _2(𝒌)=\frac{d𝒍}{(2\pi )^3}\delta _1(𝒍)\delta _1(𝒌𝒍)J(𝒍,𝒌𝒍),$$
(20)
or in $`𝒙`$-space
$$\delta _2(𝒙)=\frac{d𝒌}{(2\pi )^3}\frac{d𝒍}{(2\pi )^3}e^{i𝒌𝒙}\delta _1(𝒍)\delta _1(𝒌𝒍)J(𝒍,𝒌𝒍),$$
(21)
where the kernel $`J`$ is defined by
$$J(𝒌,𝒍)=\frac{1}{2}(1+K)+\frac{𝒌𝒍}{2}\left(\frac{1}{k^2}+\frac{1}{l^2}\right)+\frac{1}{2}(1K)\frac{(𝒌𝒍)^2}{k^2l^2}.$$
(22)
The factor $`K(\mathrm{\Omega },\lambda )`$ weakly depends on the density parameter $`\mathrm{\Omega }`$ and cosmological constant $`\lambda `$ as shown in the fitting formula (Matsubara 1995, see also Bouchet et al. 1992)
$$K(\mathrm{\Omega },\lambda )\frac{3}{7}\mathrm{\Omega }^{1/30}\frac{\lambda }{80}\left(1\frac{3}{2}\lambda \mathrm{log}_{10}\mathrm{\Omega }\right).$$
(23)
In the ranges of two parameters $`\mathrm{\Omega }`$ and $`\lambda `$
$$0.1\mathrm{\Omega }1,0.1\lambda 1,$$
(24)
the difference of $`K(\mathrm{\Omega },\lambda )`$ from $`K=3/7`$ is within $`8\%`$. Therefore, in the following analysis we basically study the Einstein de-Sitter background and use $`K=3/7`$.
Using the second-order solution (21) we can derive the well known formula for the third-order moment as follows (Juszkiewicz et al. 1993)
$$\delta _R^3=6\frac{d𝒌}{(2\pi )^3}\frac{d𝒍}{(2\pi )^3}P(k)P(l)J(𝒌,𝒍)w(kR)w(lR)w(|𝒌+𝒍|R).$$
(25)
Let us simplify this six-dimensional integral $`d𝒌d𝒍`$. In the case of the Gaussian filter
$$w(kR)=\mathrm{exp}(k^2R^2/2),$$
(26)
we can change $`\delta _R^3`$ to the following form (Matsubara 1994)
$$\delta _R^3=\frac{3}{28\pi ^4}(5I_{220}+7I_{131}+2I_{222}),$$
(27)
where we have defined
$$I_{abc}=_0^{\mathrm{}}𝑑k_0^{\mathrm{}}𝑑l_1^1𝑑u\mathrm{exp}[R^2(k^2+l^2+ukl)]k^al^bu^cP(k)P(l).$$
(28)
For a power-law initial fluctuation $`P(k)=k^n`$, we obtain a final closed formula (Matsubara 1994, $`Ł`$okas et al. 1995)
$$S_F(n)=3F(\frac{n+3}{2},\frac{n+3}{2},\frac{3}{2};\frac{1}{4})\left(n+\frac{8}{7}\right)F(\frac{n+3}{2},\frac{n+3}{2},\frac{5}{2};\frac{1}{4}),$$
(29)
where $`F`$ is the Hypergeometric function.
In the case of the top-hat filter whose Fourier transform is given by
$$w(kR)=\frac{3}{(kR)^3}(\mathrm{sin}kRkR\mathrm{cos}kR),$$
(30)
the final form of $`S`$ becomes very simple as follows (Juszkiewicz et al. 1993, Bernardeau 1994)
$$S_F(n)=\frac{34}{7}(n+3).$$
(31)
This formula is not only valid for pure power-law initial fluctuations but also for general power spectra with effective spectral index defined at the smoothing radius $`R`$ as
$$n\frac{d\mathrm{ln}\sigma _R^2}{d\mathrm{ln}R}3.$$
(32)
Equations (29) and (31) are only the leading-order contribution and more higher-order effects might change them considerably. Thus it is quite important to compare these analytic formulas with fully nonlinear numerical simulations and clarify validity of the perturbative formulas. There are many works on this topic and the analytic predictions show surprisingly good agreement with numerical simulations, even at $`\sigma 1`$ (e.g. Baugh, Gazta$`\stackrel{~}{\mathrm{n}}`$aga & Efstathiou 1995, Hivon et al. 1995, Juszkiewicz et al. 1995, $`Ł`$okas et al. 1995).
So far we have discussed skewness $`S`$ with fixed smoothing methods. For the third-order moments $`\delta _R^3`$, the second-order effects caused by the gravitational evolution and that caused by the adaptive smoothing are decoupled, as we can see from equations (10) and (11). Thus we can write the skewness parameter for adaptive methods in the following forms
$`S_G`$ $`=`$ $`S_F+\mathrm{\Delta }S_G,`$ (33)
$`S_S`$ $`=`$ $`S_F+\mathrm{\Delta }S_S.`$ (34)
Here $`\mathrm{\Delta }S_G`$ and $`\mathrm{\Delta }S_S`$ are the correction terms caused by the adaptive smoothing methods. In the next two subsections we calculated these terms explicitly.
### 3.1 Gather Approach
First we calculate the correction term $`\mathrm{\Delta }S_G`$ for the gather approach. With equation (10) this term is easily transformed to the following equation (see Appendix A.1 for derivation)
$$\mathrm{\Delta }S_G=\frac{d\mathrm{ln}\sigma _R^2}{d\mathrm{ln}R}.$$
(35)
For a power-law spectrum we have simple equation below
$$\mathrm{\Delta }S_G(n)=(n+3).$$
(36)
In derivation of equation (35) we only use Gaussianity of the one point PDF of the linear smoothed field $`\delta _R`$. Therefore, these formulas do not depend on the shape of the smoothing filter nor the cosmological parameters $`\mathrm{\Omega }`$ or $`\lambda `$. Furthermore they are valid also in the redshift space, if we use the the distant observer approximation. Thus equation (35) has strong predictability.
Hivon et al. (1995) perturbatively examined the skewness parameters in redshift space and evaluate them both for the top-hat filter and the Gaussian filter. They also compared their analytic results with numerical results. They found that these two show agreement only in the range $`\sigma \stackrel{<}{}\text{ }0.1`$, in contrast to the skewness parameter in the real space $`\sigma \stackrel{<}{}\text{ }1.0`$. They commented that this limitation is mainly due to the finger of god effects (e.g. Davis & Peebles 1983).
Here we use their analytic results and combine our new formula with them. In figure 1 we present the skewness parameters for various spectral indexes $`n`$ both in the real and redshift spaces. For simplicities we limit our analysis for the Einstein de-Sitter background.
For the Gaussian filter, the skewness parameter by the gather method is a increasing function of spectral index $`n`$ both in real and redshift spaces. This dependence is contrast to the skewness with the fixed smoothing method. Comparing the skewness in the real and the redshift spaces, $`n`$ dependence of the gather method is somewhat weaker in real space than in redshift space, but this tendency is also different from the fixed smoothing.
For the top-hat filter, there is no spectral index $`n`$\- dependence in the real space. We have $`S=34/7`$ which is the same as the unsmoothed value (Peebles 1980). Bernardeau (1994) pointed out that the skewness $`S`$ filtered with the top-hat filter in Lagrangian space does not depend on the power spectrum and is given by $`S=34/7`$. As the adaptive smoothing is basically Lagrangian description, this fact seems reasonable. In the case of the redshift space we have a fitting formula below
$$S_G(n)=\frac{35.2}{7}0.15(n+3),(\mathrm{Einstein}\mathrm{de}\mathrm{Sitter}\mathrm{background})$$
(37)
which is based on formula (49) of Hivon et al (1995). Again $`n`$ dependence is very weak and becomes weaker for $`\mathrm{\Omega }<1`$ (see Fig.4 of Hivon et al. 1995). Finally, we comment the possibility that our perturbative treatment of the redshift space skewness becomes worse in the adaptive methods than in the fixed smoothing method. In the adaptive methods, smoothing radius of a high density region becomes smaller and the (strongly nonlinear) finger of god effects might not be suppressed well.
### 3.2 Scatter Approach
Next we calculate the correction term $`\mathrm{\Delta }S_S`$ for the scatter approach. We only discuss the real space density field smoothed with the Gaussian filter (eq.). From equation (12) we obtain the following equation (see Appendix A.1),
$$\mathrm{\Delta }S_{SR}(n)\sigma _R^4=2\frac{d𝒌}{(2\pi )^3}\frac{d𝒍}{(2\pi )^3}\mathrm{exp}\left[\frac{(3𝒍^2+2𝒌^2+2𝒌𝒍)R^2}{2}\right]P(k)P(l)(𝒌+𝒍)^2R^2.$$
(38)
The six dimensional integral $`d𝒌d𝒍`$ is simplified to a three dimensional integral $`dkdldu`$ as in equations (25) and (28). Then we have the following relation
$`\mathrm{\Delta }S_{SR}(n)\sigma _R^4`$ $`=`$ $`{\displaystyle \frac{1}{4\pi ^4}}{\displaystyle _0^{\mathrm{}}}𝑑k{\displaystyle _0^{\mathrm{}}}𝑑l{\displaystyle _1^1}𝑑u\mathrm{exp}\left[{\displaystyle \frac{(3l^2+2k^2+2klu)R^2}{2}}\right]`$ (39)
$`\times k^2l^2P(k)P(l)(k^2+l^2+2klu)R^2.`$
For a pure-power law fluctuation we obtain the following analytic formula
$`\mathrm{\Delta }S_S(n)`$ $`=`$ $`2^{(n+3)/2}3^{(n+5)/2}(n+3)\{{\displaystyle \frac{2}{3}}(n+3)F({\displaystyle \frac{5+n}{2}},{\displaystyle \frac{5+n}{2}},{\displaystyle \frac{5}{2}},{\displaystyle \frac{1}{6}})`$ (40)
$`5F({\displaystyle \frac{3+n}{2}},{\displaystyle \frac{5+n}{2}},{\displaystyle \frac{3}{2}},{\displaystyle \frac{1}{6}})\}.`$
In contrast to the previous gather approach, this result is valid only to the real space skewness with the Gaussian filter. In table 1 we present numerical values of $`\mathrm{\Delta }S_S(n)`$. In figure 2 we show $`S_S(n)`$ as a function of the spectral index $`n`$. We can see that $`n`$ dependence is similar to the gather approach but now it becomes weaker. If we change $`n`$ from $`3`$ to $`1`$, skewness $`S`$ changes $`25\%`$ for the scatter approach, $`45\%`$ for the gather approach, and $`38\%`$ for the fixed smoothing.
## 4 Statistics of Isodensity Contour
The genus number is a topological quantity and defined by the number of the homotopy classes of closed curves that may be drawn on a surface without cutting them into two pieces. This definition seems highly mathematical, but there are more intuitive methods to count the genus number. First one is to notice the number of holes and isolated regions of the surface in interest. Second one is to count stationary points of the surface along one spatial direction (Adler 1981, Bardeen et al. 1986). With these equivalent methods, we can calculate the genus density as follows
$`\mathrm{Genus}\mathrm{density}`$ $`=`$ $`{\displaystyle \frac{N(\mathrm{holes})N(\mathrm{isolated}\mathrm{regions})}{\mathrm{volume}}}`$ (41)
$`=`$ $`{\displaystyle \frac{N(\mathrm{maxima})+N(\mathrm{minima})N(\mathrm{saddle}\mathrm{points})}{2\times \mathrm{volume}}}.`$ (42)
For example, in the case of one-sphere, we have $`N(\mathrm{holes})=0`$ and $`N(\mathrm{isolated}\mathrm{regions})=1`$ and genus number becomes $`1`$. We obtain the same result with equation (42). The genus number density of isodensity contour of the large-scale structure is a powerful measure to quantify connectivity of galaxy clustering, such as, filamentary networks, sheet-like or bubble-like structures. The genus density of a high density contours is expected to be negative as the surfaces would show disconnected meatball-like structure. But the genus density for contours around the mean density $`\delta 0`$ would be positive as they would look like highly connected sponge-like structure (Gott, Melott & Dickinson 1986). The genus density as a function of the matter density threshold is called the genus statistics and has been widely investigated both numerically and observationally (Gott, Weinberg & Melott 1987, Weinberg, Gott & Melott 1987, Melott, Weinberg & Gott 1988, Gott et al. 1989, Park & Gott 1991, Park, Gott & da Costa 1992, Weinberg & Cole 1992, Moore et al. 1992, Vogeley, Park, Geller, Huchra & Gott 1994, Rhoads, Gott & Postman 1994, Matsubara & Suto 1996, Coles, Davies & Pearson 1996, Sahni et al. 1997, Protogeros & Weinberg 1997, Coles, Pearson, Borgani, Plionis & Moscardini 1998, Canavezes et al. 1998, Springel et al. 1998)
Usually we use the local expression (42) to analytically study the genus statistics. For the genus density of isodensity contour at $`\nu \delta /\sigma `$, this expression is written as follows (Doroshkevich 1970, Adler 1981, Bardeen et al. 1986, Hamilton, Gott & Weinberg 1986)
$$G(\nu )=\frac{1}{2}\delta _{Drc}[\delta (𝒙)\nu \sigma ]\delta _{Drc}[_1\delta (𝒙)]\delta _{Drc}[_2\delta (𝒙)]|_3\delta (𝒙)|(_{11}\delta (𝒙)_{22}\delta (𝒙)_{12}\delta (𝒙)^2),$$
(43)
where $`\delta _{Drc}()`$ represents the Dirac’s delta function. The first one $`\delta _{Drc}[\delta (𝒙)\nu \sigma ]`$ specifies the contour $`\nu \delta /\sigma `$. The second and third ones $`\delta _{Drc}[_1\delta (𝒙)]`$, $`\delta _{Drc}[_2\delta (𝒙)]`$ specify the stationary points along $`x_3`$ direction. The term $`(_{11}\delta (𝒙)_{22}\delta (𝒙)_{12}\delta (𝒙)^2)`$ is the determinant of the Hesse-matrix and assigns proper signatures for the stationary points corresponding to signs of equation (42). Even though equation (43) introduces a specific spatial direction ($`x_3`$-axis), Seto et al. (1997) derived a rotationarilly symmetric formula, and studied nonlinear evolution of the genus statistics using the Zeldovich approximation (Zeldovich 1970)
In the case of an isotropic random Gaussian fluctuation which is usually assumed as the initial condition of the structure formation, the complicated formula (43) is simplified to (Doroshkevich 1970, Adler 1981, Bardeen et al. 1986, Hamilton, Gott & Weinberg 1986)
$$G(\nu )=\frac{1}{(2\pi )^2}\left(\frac{\sigma _1^2}{3\sigma ^2}\right)^{3/2}e^{\nu ^2/2}(1\nu ^2),$$
(44)
where $`\sigma _1^2`$ is defined as
$$\sigma _1^2=(\delta )^2=\frac{1}{2\pi ^2}_0^{\mathrm{}}𝑑kP(k)k^4w(kR)^2.$$
(45)
Nonlinear evolution of the genus statistics had been studied using N-body simulations, but most analytical predictions for the genus statistics have been based on the linear formula (44). To compare it with observed distribution of galaxies we have to use sufficiently large smoothing radius to reduce nonlinearities. However, such a large smoothing radius is not statistically preferable for the finiteness of our survey volume.
Matsubara (1994) improved this difficulty by taking into account of weakly nonlinear effects in the genus statistics (see also Hamilton 1988, Okun 1990, Matsubara & Yokoyama 1996, Seto et al. 1997). He used the multidimensional Edgeworth expansion method and added the first-order nonlinear correction to the linear formula. His result is written as
$$G(\nu )=\frac{1}{(2\pi )^2}\left(\frac{\sigma _1^2}{3\sigma ^2}\right)^{3/2}e^{\nu ^2/2}\left[H_2(\nu )+\sigma \left(\frac{S}{6}H_5(\nu )+\frac{3T}{2}H_3(\nu )+3UH_1(\nu )\right)+O(\sigma ^2)\right].$$
(46)
This formula is valid for statistically isotropic and homogeneous weakly random Gaussian fields. Here functions $`H_n(\nu )(1)^ne^{\nu ^2/2}(d/d\nu )^ne^{\nu ^2/2}`$ are the Hermite polynomials,
$`H_0(x)`$ $`=`$ $`1,H_1(x)=x,H_2(x)=x^21,`$
$`H_3(x)`$ $`=`$ $`x^33x,H_5(x)=x^510x^3+15x.`$ (47)
In equation (46), $`S`$ is the skewness parameter discussed in the previous section. $`T`$ and $`U`$ are called the generalized skewness parameters and defined by
$`T`$ $`=`$ $`{\displaystyle \frac{1}{2\sigma ^2\sigma _1^2}}\delta ^2\mathrm{\Delta }\delta ,`$
$`U`$ $`=`$ $`{\displaystyle \frac{3}{4\sigma _1^4}}\delta \delta \mathrm{\Delta }\delta .`$ (48)
Matsubara & Suto (1996) examined the perturbative formula (46) using N-body simulations. For power-law spectra with $`n=1,0`$ and $`1`$, they found that this formula are in reasonable agreement with numerical results in the range $`0.2\stackrel{<}{}\text{ }\nu \sigma \stackrel{<}{}\text{ }0.4`$.
Next let us briefly summarize the area statistics $`N_3(\nu )`$ which were proposed by Ryden (1988) and investigated detailedly by Ryden et al. (1989). The area statistics are defined as the mean area of isodensity contour surface per unit volume. For statistically homogeneous and isotropic fluctuations, the area statistics are equal to twice the mean number of isodensity contour crossings along a straight line of unit length. These two statistics are thus equivalent (beside factor 2), but the contour crossing statistics are easier to compute numerically (Ryden 1988).
As in equation (43), the area statistics for isodensity contour $`\delta =\nu \sigma `$ is written as
$$N_3(\nu )=\delta _{Drc}[\delta (𝒙)\nu \sigma ]|\delta (𝒙)|.$$
In the case of isotropic Random Gaussian fluctuations, we have the following formula (Ryden 1988)
$$N_3(\nu )=\frac{2}{\pi }\left(\frac{\sigma _1^2}{3\sigma ^2}\right)^{1/2}e^{\nu ^2/2}.$$
(49)
Weakly nonlinear effects on the area statistics can be discussed with a similar technique used to derive equation (46) (Matsubara 1995). In this case, we need information of the density field up to its first spatial derivative and nonlinear correction is expressed in terms of two parameters $`S`$ and $`T`$ as follows
$$N_3(\nu )=\frac{2}{\pi }\left(\frac{\sigma _1^2}{3\sigma ^2}\right)^{1/2}e^{\nu ^2/2}\left[1+\sigma \left(\frac{S}{6}H_3(\nu )+\frac{T}{2}H_1(\nu )\right)+O(\sigma ^2)\right].$$
(50)
In the rest of this section we consider nonlinear effects of the adaptive smoothing methods on the genus and the area statistics. We calculate the generalized skewness parameters $`T`$ and $`U`$ both for the gather approach and the scatter approach. We limit our analysis to the real space density field smoothed by Gaussian filters.
### 4.1 Reparameterization of Isodensity Contour
Equation (46) is weakly non-Gaussian genus density for isodensity contour surfaces parameterized by the simple definition $`\nu =\delta /\sigma `$. There is another conventional method to name contour surfaces. In this method, we notice the volume fraction $`f`$ above the density threshold of the contour in interest (e.g. Gott, Melott & Dickinson 1986, Gott et al. 1989), and parameterize the contour using value $`\nu _r`$ defined by
$$\nu _r\mathrm{erf}^1(f),$$
(51)
where the suffix $`r`$ indicates “ reparameterization” and $`\mathrm{erf}(x)`$ is the error function defined by
$$\mathrm{erf}(x)\frac{1}{\sqrt{2\pi }}_x^{\mathrm{}}𝑑ye^{y^2/2}.$$
(52)
This procedure is a kind of Gaussianization. Two methods coincide $`\nu =\nu _r`$ when the one point PDF $`P(\nu )`$ is Gaussian distributed. If we use this new parameterization, the genus curve is apparently invariant under a monotonic mapping of the density contrast field $`\delta `$. Furthermore, it has been long known that the genus curve with $`\nu _r`$ parameterization (51) nearly keeps its original symmetric shape (eq.) in the course of weakly nonlinear gravitational evolution of density field (e.g. Springel et al. 1998 and references therein). Almost the same kind of tendency has been confirmed for the area statistics (Ryden et al. 1989). Weakly nonlinear area density with $`\nu _r`$ parameterization remains at its linear shape very well. Here, let us relate these two parameterization methods for weakly nonlinear regime. Using the Edgeworth expansion method, the one point PDF $`P(\nu )`$ is written in terms of the skewness $`S`$ up to the first-order nonlinear correction
$$P(\nu )=\frac{1}{\sqrt{2\pi }}e^{\nu ^2/2}\left[1+\frac{\sigma S}{6}H_3(\nu )+O(\sigma ^2)\right].$$
(53)
Juszkiewicz et al. (1995) examined this approximation using N-body simulations. They found that the above formula is accurate until $`\sigma S`$ reaches 1. Therefore the inequality
$$\sigma \stackrel{<}{}\text{ }S^10.2,$$
(54)
would be a standard for the validity of the perturbative analysis in this subsection. The volume fraction $`f(\nu )`$ above the threshold $`\delta =\nu \sigma `$ is given by
$$f(\nu )=_\nu ^{\mathrm{}}𝑑xP(x)=\mathrm{erf}(\nu )+\frac{\sigma S}{6\sqrt{2\pi }}\left\{e^{\nu ^2}(\nu ^21)\right\}+O(\sigma ^2).$$
(55)
With equations (51) and (55) we obtain correspondence between $`\nu `$ and $`\nu _r`$ as follows
$$\nu =\nu _r+\frac{\sigma S}{6}\left\{\nu _r^21\right\}+O(\sigma ^2).$$
Finally the genus density $`G_r(\nu _r)`$ in this new parameterization is given by $`G(\nu )=G_r(\nu _r)`$ and written as
$`G_r(\nu _r)`$ $`=`$ $`{\displaystyle \frac{1}{(2\pi )^2}}\left({\displaystyle \frac{\sigma _1^2}{3\sigma ^2}}\right)^{3/2}e^{\nu _r^2/2}[H_2(\nu _r)+\sigma (H_3(\nu )(S+{\displaystyle \frac{3}{2}}T)`$ (56)
$`+H_1(\nu )(S+3U))+O(\sigma ^2)].`$
Similarly we have the following result for the area statistics
$$N_{3r}(\nu _r)=\frac{2}{\pi }\left(\frac{\sigma _1^2}{3\sigma ^2}\right)^{1/2}e^{\nu _r^2/2}\left[1+\sigma \left(\frac{S}{3}+\frac{T}{2}\right)H_1(\nu _r)+O(\sigma ^2)\right].$$
(57)
In this case, the first nonlinear correction is simply proportional to $`\nu _r`$ (see eq.) and it completely vanishes when we have $`S=3/2T`$.
Later in §4.3 and §4.4, we will confirm that the nonlinear correction (proportional to $`\sigma `$) for the genus and area statistics with the fixed smoothing method are very small for $`\nu _r`$ parameterization, as experientally known in N-body simulations. In the followings, we use these two parameterizations $`\nu `$ and $`\nu _r`$.
### 4.2 Generalized Skewness
The generalized skewness $`T`$ and $`U`$ are basic ingredients to perturbatively evaluate the weakly non-Gaussian effects on the genus and the area statistics. For the Gaussian filter with a fixed smoothing radius, explicit formulas valid for the power-law initial fluctuations were derived by Matsubara (1994). They are given as follows
$`T_{FR}`$ $`=`$ $`3F({\displaystyle \frac{n+3}{2}},{\displaystyle \frac{n+5}{2}},{\displaystyle \frac{3}{2}},{\displaystyle \frac{1}{4}})\left(n+{\displaystyle \frac{18}{7}}\right)F({\displaystyle \frac{n+3}{2}},{\displaystyle \frac{n+5}{2}},{\displaystyle \frac{5}{2}},{\displaystyle \frac{1}{4}})`$ (58)
$`+{\displaystyle \frac{4(n2)}{105}}F({\displaystyle \frac{n+3}{2}},{\displaystyle \frac{n+5}{2}},{\displaystyle \frac{7}{2}},{\displaystyle \frac{1}{4}}),`$
$`U_{FR}`$ $`=`$ $`F({\displaystyle \frac{n+5}{2}},{\displaystyle \frac{n+5}{2}},{\displaystyle \frac{5}{2}},{\displaystyle \frac{1}{4}}){\displaystyle \frac{7n+16}{35}}F({\displaystyle \frac{n+5}{2}},{\displaystyle \frac{n+5}{2}},{\displaystyle \frac{7}{2}},{\displaystyle \frac{1}{4}}).`$ (59)
As shown in the case of the skewness parameter $`S`$ analyzed in §3.1 and §3.2, the second-order (first nonlinear) effects caused by the adaptive smoothing methods are decoupled from those induced by gravitational mode couplings. Thus we can calculate them separately and express the total values in forms similar to equations (33) and (34).
First, we analyze correction terms $`\mathrm{\Delta }T_{GR}`$ and $`\mathrm{\Delta }U_{GR}`$ for the gather approach. We define these terms by the following equations
$$T_{GR}=T_{FR}+\mathrm{\Delta }T_{GR},U_{GR}=U_{FR}+\mathrm{\Delta }U_{GR}.$$
(60)
After some tedious algebra using equation (10), we obtain the leading-order correction terms as follows (Appendix A.2)
$`\mathrm{\Delta }T_{GR}`$ $`=`$ $`{\displaystyle \frac{2}{3}}(n+4),`$ (61)
$`\mathrm{\Delta }U_{GR}`$ $`=`$ $`{\displaystyle \frac{1}{3}}(n+5).`$ (62)
These results are similar to the correction term for the skewness $`\mathrm{\Delta }S_{GR}=n+3`$ (eq.) which does not depend on the shape of the filter function. However, situation is not so much simple here. For the Gaussian filter we have the following relation,
$$R\frac{\delta _R(𝒙)}{R}=R^2\mathrm{\Delta }_x\delta _R(𝒙).$$
(63)
This relation plays important roles to derive equations (61) and (62). But it does not hold for general filter functions and the simple results given in equations (61) and (62) are specific to the Gaussian filter. We summarize numerical data for the parameters$`T_{GR}`$ and $`U_{GR}`$ in Table 3.
Generalized skewness for the gather approach is much more complicated. If we write down them in the form
$$T_{SR}=T_{SR}+\mathrm{\Delta }T_{SR},U_{SR}=U_{SR}+\mathrm{\Delta }U_{SR},$$
(64)
correction terms $`\mathrm{\Delta }T_{SR}`$ and $`\mathrm{\Delta }U_{SR}`$ are written in the manner similar to equation (39) as follows (Appendix A.3)
$`2\mathrm{\Delta }T_{SR}(n)\sigma _R^2\sigma _{1R}^2`$ $`=`$ $`{\displaystyle \frac{1}{6\pi ^4}}{\displaystyle _0^{\mathrm{}}}𝑑k{\displaystyle _0^{\mathrm{}}}𝑑l{\displaystyle _1^1}𝑑u\mathrm{exp}\left[{\displaystyle \frac{(3l^2+2k^2+2klu)R^2}{2}}\right]`$ (65)
$`\times k^2l^2P(k)P(l)(k^2+l^2+2klu)(k^2+l^2+klu)R^2,`$
$`{\displaystyle \frac{4}{3}}\mathrm{\Delta }U_{SR}(n)\sigma _{1R}^4`$ $`=`$ $`{\displaystyle \frac{1}{6\pi ^4}}{\displaystyle _0^{\mathrm{}}}𝑑k{\displaystyle _0^{\mathrm{}}}𝑑l{\displaystyle _1^1}𝑑u\mathrm{exp}\left[{\displaystyle \frac{(3l^2+2k^2+2klu)R^2}{2}}\right]`$ (66)
$`\times k^4l^4P(k)P(l)(k^2+l^2+klu)(1u^2)R^2.`$
We can calculate them explicitly as follows
$`\mathrm{\Delta }T_{SR}`$ $`=`$ $`2^{(n+5)/2}3^{(n+9)/2}[4(n+3)F({\displaystyle \frac{n+5}{2}},{\displaystyle \frac{n+5}{2}},{\displaystyle \frac{5}{2}},{\displaystyle \frac{1}{6}})`$ (67)
$`+{\displaystyle \frac{5}{2}}(n+3)(n+5)F({\displaystyle \frac{n+5}{2}},{\displaystyle \frac{n+7}{2}},{\displaystyle \frac{5}{2}},{\displaystyle \frac{1}{6}})`$
$`{\displaystyle \frac{13(n+5)}{2}}F({\displaystyle \frac{n+3}{2}},{\displaystyle \frac{n+7}{2}},{\displaystyle \frac{3}{2}},{\displaystyle \frac{1}{6}})`$
$`12(n+3)F({\displaystyle \frac{n+5}{2}},{\displaystyle \frac{n+5}{2}},{\displaystyle \frac{3}{2}},{\displaystyle \frac{1}{6}})],`$
$`\mathrm{\Delta }U_{SR}`$ $`=`$ $`2^{(n+5)/2}3^{(n+5)/2}[4F({\displaystyle \frac{n+5}{2}},{\displaystyle \frac{n+5}{2}},{\displaystyle \frac{5}{2}},{\displaystyle \frac{1}{6}})`$ (68)
$`+{\displaystyle \frac{5(n+5)}{9}}F({\displaystyle \frac{n+5}{2}},{\displaystyle \frac{n+7}{2}},{\displaystyle \frac{5}{2}},{\displaystyle \frac{1}{6}})`$
$`4F({\displaystyle \frac{n+5}{2}},{\displaystyle \frac{n+5}{2}},{\displaystyle \frac{3}{2}},{\displaystyle \frac{1}{6}})].`$
We present numerical data of the parameters $`T_{SR}`$ and $`U_{SR}`$ in Table 4. Note that the magnitude of generalized skewness $`T`$ and $`U`$ becomes very small in the scatter approach. This fact becomes important in the next subsection.
### 4.3 Weakly Nonlinear Genus Statistics
In figures 3 to 5, we show the weakly nonlinear genus density smoothed by three different methods (fixed, gather and scatter), using two types of parameterizations $`\nu `$ and $`\nu _r`$. All of these curves are smoothed by the Gaussian filter (eq.). We plot the normalized genus curves <sup>1</sup><sup>1</sup>1Note that the amplitude of $`G(0)`$ or $`G_r(0)`$ are not changed by the first-order correction of $`\sigma `$. $`G(\nu )/G(0)`$ or $`G_r(\nu _r)/G_r(0)`$ to see deviation from the symmetric linear genus curve $`(1\nu ^2)\mathrm{exp}(\nu ^2/2)`$. Upper panel of figure 3 is essentially same as Fig.1 of Matsubara (1994).
First, comparing upper and bottom panels of Fig.3, we can confirm the fact experientally known in N-body simulations. Weakly nonlinear genus curves for the fixed smoothing method are very close to the linear symmetric shape in the case of $`\nu _r`$ parameterization (e.g. Springel et al. 1998). For a spectral index with $`n1`$, three curves for $`\sigma =0,0.2`$ and $`0.4`$ are nearly degenerated.
In Fig.4 we present genus curves with the gather smoothing. From Figs.3 and 4 it is apparent that deviations from the linear curves become larger in the gather approach, especially in tail parts. But these deviations become smaller with using parameter $`\nu _r`$.
Bottom panel of Fig.4 is calculated under conditions (gather approach and parameter $`\nu _r`$) similar to Fig.7 of Springel et al. (1998) which is obtained from N-body simulations. However overall shapes of these two are different. The minimum value of $`G_r(\nu _r)`$ are attained around the point $`\nu _r1.5`$ in our result, but this point is $`\nu _r1.5`$ in theirs. This difference might be caused by the difference of adopted filter function. We use the Gaussian filter but a different kernel (a spline kernel that is often used in SPH simulations) is adopted in their calculation (Monaghan & Lattanzio 1985).
In figure 5 we show results for the scatter approach. Nonlinear effects are more prominent than two cases analyzed earlier. As shown in equations (46) and (56), nonlinear correction of the genus curves are written by combination of terms proportional to parameters $`S`$, $`T`$ and $`U`$. Some of their contribution cancel out, as realized in the case of the fixed or gather smoothing methods. However, amplitude of parameters $`T`$ and $`U`$ for the scatter approach becomes very small (see Table 4), and cancellation becomes weaker.
Based on the definition of the genus density (eq.), nonlinear evolution of isodensity contours is sometimes described with such terminologies as, sponge-like (connected topology) or meatball like (disconnected topology). In the case of Random Gaussian initial fluctuations, linear theory predicts symmetry of the genus statistics with respect to the sign of density contrast $`\delta `$, and geometry of both high and low density tails look meat-ball like with negative genus density. If we use the fixed or gather smoothing methods, nonlinear effects make the genus number of a high density contour (e.g. $`\nu =2`$) smaller, and topology of that region becomes more meatball-like (see figures 3 and 4). In contrast, contour of low density threshold (e.g. $`\nu =2`$) is transformed in the direction of sponge-like topology, as quantified by the increase of genus number density (see figures 3 and 4).
It is not easy to understand behaviors of Fig.5 for the scatter approach. But, changing point of views, we can discuss characters of this approach by comparing figures for various smoothing methods. To characterize nonlinear effects accompanied with the scatter approach, we apply the topological interpretation mentioned in the last paragraph. As shown in Fig.6 where various smoothing methods are compared, the scatter approach makes low density regions more meatball-like. But high density regions are transformed in the opposite direction. This trend shows remarkable contrast to the nonlinear gravitational effects traced by the simple fixed smoothing method.
### 4.4 Weakly Nonlinear Area Statistics
In Figs.7 to 9, we plot the weakly nonlinear area statistics, using equation (50) for $`\nu `$ parameterization and equation (57) for $`\nu _r`$ parameterization. As in the analysis of the genus statistics, we normalize amplitude of the area density by $`N_3(0)`$ or $`N_{3r}(0)`$. Comparing upper and bottom panels of figure 7, it is apparent that $`\nu _r`$ parameterization is very effective to keep the original linear shape against weakly nonlinear effects. This fact has been confirmed experimentally in N-body simulations (Ryden et al. 1989) and is quite similar to the situation in the genus statistics explained in previous subsection. With $`\nu _r`$ parameterization (bottom panel of Fig.7), three curves for $`\sigma =0,0.2`$ and $`0.4`$ are almost completely overlapped for all spectral indexes $`n`$. This fact seems reasonable as we have $`S_F3/2T_F`$ for the fixed smoothing method (see eq. and Tables 1, 3).
If we use the adaptive smoothing methods, weakly nonlinear correction on $`N_3(\nu )`$ ($`\nu `$ parameterization) are considerable as shown in upper panels of Figs.8 and 9. This correction becomes apparently smaller for the gather approach, but not for the scatter approach. We have already commented that the first-order nonlinear correction for the area statistics is characterized by two parameters $`S`$ and $`T`$. For the scatter approach, $`T`$ parameter is very small for spectral indexes $`n>1`$, and cancellation mentioned in §4.3 is not effective.
## 5 Summary
Observational analysis of galaxy clustering is one of the central issues in modern cosmology. Various methods have been proposed to quantify the clustering, and many of them (e.g. topological analyses of isodensity contour) are based on continuous smoothed density field. However what we can observe directly is distribution of point-like galaxies. Thus smoothing operation is crucially important in the analyses of the large-scale structure. From theoretical point of views, filters with spatially constant smoothing radii are natural choice and have been widely adopted so far. But we should notice that there are no strong convincing reasons to stick to this traditional method.
There are few galaxies in void regions even at semi-nonlinear scales. In these regions density field obtained with fixed smoothing radius might be considerably affected by the discreteness of (point-like) mass elements, and might hamper our analyses of the cosmic structures. Adaptive smoothing is basically Lagrangian description, and we use nearly same number of particles to construct smoothed density field at each point. Thus it is quite possible that the adaptive methods are more efficient than the fixed methods to analyze the large-scale structure. Actually, Springel et al. (1998) have recently pointed out that using adaptive filters, signal to noise ratio of the genus statistics is improved even at weakly nonlinear scales $`\stackrel{>}{}\text{ }10h^1\mathrm{Mpc}`$.
In this article, we have developed a perturbative analysis of adaptive smoothing methods that are applied to quantify the large-scale structure. Even though adaptive methods might be promising approaches in observational cosmology, this kind of analytic investigation has not been done so far. Our targets are weakly nonlinear effects induced by two typical adaptive approaches, the gather approach and the scatter approach (Hernquist & Katz 1989). The concept of these methods is easily understood with equations (8) and (9). The gather approach is easier to handle analytically. Numerical costs dealing with discrete particles’ systems are also lower in this approach (Springel et al. 1998). Effects caused by these two adaptive methods start from second-order of $`\delta `$ in perturbative sense. They modify quantities which characterize the nonlinear mode couplings induced by gravity (e.g. Peebles 1980).
In §3 we have investigated the skewness parameter $`S`$ which is a fundamental measure to quantify asymmetry of one point PDF. We have shown that the skewness for a gather top-hat filter does not depend on the spectral index $`n`$ in real space, and very weakly depends on it ($`S35.20.15n`$: Einstein de-Sitter background) in redshift space. In the case of Gaussian filter, the skewness parameters show similar behaviors both in the scatter and gather approaches. They are increasing functions of $`n`$, in contrast to the fixed smoothing method.
Next in §4, the genus and area statistics have been studied with Gaussian adaptive filters. Our analysis is based on the multidimensional Edgeworth expansion explored by Matsubara (1994). We use two quantities $`\nu (\delta /\sigma )`$ and $`\nu _r`$ to parameterize isodensity contours. The latter $`\nu _r`$ is defined by the volume fraction above a given density threshold (Gott, Melott, & Dickinson 1986). It is explicitly shown that using this parameterization, two statistics with the fixed smoothing method are very weakly affected by semi-nonlinear gravitational dynamics, as experientally confirmed by N-body simulations. For the gather smoothing, we found that the $`\nu _r`$\- parameterization is more effective to keep original linear shape of the area statistics than of the genus statistics.
The parameters $`S,T`$ and $`U`$ which characterize the nonlinear corrections of isodensity contour depend largely on the filtering methods. We can characterize nonlinear effects of these methods in somewhat intuitive manner, using results for the genus statistics. The scatter approach makes low density tails more meatball-like, but high density tails are transformed in the direction of sponge-like (connected) topology. This is a remarkable difference from fixed or gather smoothing methods.
Our investigation in this article has been fully analytical one, using perturbative technique of cosmological density field. Numerical analyses based on N-body simulations would play complementary roles to results obtained here, and thus are very important. Apart from numerical investigations, perturbative treatment given in this article would be also developed further in several ways. The smoothed velocity field is crucially important material in observational cosmology as it is supposed to be less contaminated by effects of biasing (Dekel 1994, Strauss & Willick 1995). But our observational information is limited to the line of sight peculiar velocities only at points where astrophysical objects exist. Thus construction of smoothed velocity field contains similar characters as discussed in this article (e.g. Bernardeau & van de Weygaert 1996). There is another (more technical) problem that has not mentioned so far. In this article we have only studied spherically symmetric filter functions. Springel et al. (1998) have shown that signal to noise ratio of the genus curves is further improved by using a triaxial kernel, taking account of tensor information of local density field. This point must be also worth studying.
The author would like to thank J. Yokoyama for discussion and an anonymous referee for useful comments. He also thanks H. Sato and N. Sugiyama for their continuous encouragement. This work was partially supported by the Japanese Grant in Aid for Science Research Fund of the Ministry of Education, Science, Sports and Culture No. 3161.
## Appendix A Derivations of Parameters
In this appendix we derive expressions for the correction terms $`\mathrm{\Delta }S`$, $`\mathrm{\Delta }T`$ and $`\mathrm{\Delta }U`$ given in the main text. First we perturbatively expand the density contrast field smoothed by an adaptive filter as
$$\delta _A(𝒙)=\delta _1(𝒙)+\delta _2(𝒙)+\delta _{2A}(𝒙)+\mathrm{},$$
(A1)
where $`\delta _1(𝒙)`$ is the linear mode, $`\delta _2(𝒙)`$ is the second-order mode induced by gravity, and $`\delta _{2A}`$ is the second-order correction term caused by an adaptive smoothing (the suffix $`A`$ represents “adaptive”). Then the third-order moment for $`\delta _A(𝒙)`$ is given as
$$\delta _A^3=3\delta _2\delta _1^2+3\delta _{2A}\delta _1^2+\mathrm{}.$$
(A2)
Thus the first-order correction term for the third-order moment is given as
$$3\delta _{2A}\delta _1^2.$$
(A3)
In the same manner we have the following correction terms for $`\delta _A^2^2\delta _A`$ and $`\delta _A\delta _a^2\delta _A`$ as
$`\delta _1^2^2\delta _{2A}+2\delta _1\delta _{2A}^2\delta _1,`$ (A4)
$`\delta _1\delta _1^2\delta _{2A}+2\delta _1\delta _{2A}^2\delta _1.`$ (A5)
We can write down the second-order correction terms $`\delta _{2A}`$ for the gather and scatter approaches with smoothing radius $`R`$ (see eqs. and )
$`\delta _{2GR}(𝒙)`$ $`=`$ $`{\displaystyle \frac{1}{3}}\delta _{1R}(𝒙)R{\displaystyle \frac{}{R}}\delta _{1R}(𝒙)+{\displaystyle \frac{1}{6}}R{\displaystyle \frac{d}{dR}}\sigma _R^2,`$ (A6)
$`\delta _{2SR}(𝒙)`$ $`=`$ $`{\displaystyle \frac{R}{3}}{\displaystyle 𝑑𝒙^{}_RW(|𝒙^{}𝒙|;R)\delta _1(𝒙^{})\delta _{1R}(𝒙^{})}.`$ (A7)
where $`\delta _{1R}(𝒙)`$ represents the smoothed linear mode, $`W(|𝒙^{}𝒙|;R)`$ is a filter function. The variance $`\sigma _R`$ of the matter fluctuations is given as
$`\sigma _R^2`$ $`=`$ $`\delta _{1R}^2+O(\delta ^4)`$ (A8)
$`=`$ $`{\displaystyle \frac{d𝒌}{(2\pi )^3}w(kR)^2P(k)}+O(\delta ^4),`$ (A9)
where $`w(kR)`$ is a Fourier transformed filter function.
For the Gaussian filter (see eq.) the above equations are given with the linear Fourier modes $`\delta _1(𝒌)`$ as
$`\delta _{2GR}(𝒙)`$ $`=`$ $`{\displaystyle \frac{d𝒌}{(2\pi )^3}\frac{d𝒍}{(2\pi )^3}\mathrm{exp}\left[\frac{(𝒍^2+𝒌^2)R^2}{2}\right]\delta _1(𝒌)\delta _1(𝒍)\frac{𝒌^2R^2}{3}\mathrm{exp}[i(𝒌+𝒍)𝒙]}`$ (A10)
$`+{\displaystyle \frac{1}{6}}R{\displaystyle \frac{d}{dR}}\sigma _R^2,`$
$`\delta _{2SR}(𝒙)`$ $`=`$ $`{\displaystyle \frac{d𝒌}{(2\pi )^3}\frac{d𝒍}{(2\pi )^3}\mathrm{exp}\left[\frac{(2𝒍^2+𝒌^2+2𝒌𝒍)R^2}{2}\right]\delta _1(𝒌)\delta _1(𝒍)\frac{(𝒌+𝒍)^2R^2}{3}}`$ (A11)
$`\times \mathrm{exp}[i(𝒌+𝒍)𝒙],`$
Next we comment on the ensemble average of variables. We assume that the linear Fourier modes of density fluctuation are random Gaussian distributed. If variables $`\{A,B,C,D\}`$ obeys multivariable Gaussian distribution, we generally have the following relation
$$ABCD=ABCD+ACBD+ADBC.$$
(A12)
For the linear Fourier modes the above equation becomes
$`\delta _1(𝒌)\delta _1(𝒍)\delta _1(𝒎)\delta _1(𝒏)`$ $`=`$ $`(2\pi )^6P(k)P(l)\delta _{Drc}(𝒌+𝒎)\delta _{Drc}(𝒍+𝒏)`$ (A13)
$`+(2\pi )^6P(k)P(m)\delta _{Drc}(𝒌+𝒍)\delta _{Drc}(𝒎+𝒏)`$
$`+(2\pi )^6P(k)P(l)\delta _{Drc}(𝒌+𝒏)\delta _{Drc}(𝒍+𝒎).`$
Here $`\delta _{Drc}()`$ is the Dirac’s delta function and $`P(k)`$ is the matter power spectrum. We evaluate expressions (A3)-(A5) using relations (A12)-(A13).
### A.1 Skewness
For the gather approach the real-space representation (A6) is more convenient. Using property (A12) we obtain the following result
$`\delta _{2SR}\delta _{1R}^2`$ $`=`$ $`\delta _{1R}^3\left(R{\displaystyle \frac{}{R}}\delta _{1R}\right)+{\displaystyle \frac{1}{2}}\delta _{1R}^2{\displaystyle \frac{d}{dR}}\sigma _R^2`$ (A14)
$`=`$ $`3\sigma _{1R}^2\delta _{1R}{\displaystyle \frac{}{R}}\delta _{1R}+{\displaystyle \frac{1}{2}}\sigma _R^2{\displaystyle \frac{d}{dR}}\sigma _R^2`$ (A15)
$`=`$ $`\sigma _R^2{\displaystyle \frac{d}{dR}}\sigma _R^2.`$ (A16)
The correction term for the skewness $`S_G`$ is written as equation (35)
$$\mathrm{\Delta }S_G=\frac{\delta _{2GR}\delta _{1R}^2}{\sigma _R^4}=\frac{1}{\sigma _R^2}\frac{d}{dR}\sigma _R^2=\frac{d\mathrm{ln}\sigma _R^2}{d\mathrm{ln}R}.$$
(A17)
For power-law models we have a simple relation $`\sigma _R^2R^{n3}`$, and the above expression becomes
$$\mathrm{\Delta }S_G=(n+3).$$
(A18)
Here we should notice that relations (A17) and (A18) do not depend on the choice of filter functions.
Next let us evaluate the correction term for the skewness parameter in the case of the scatter approach. In this case we limit our analysis for a Gaussian filter (eq.). From equation (A11) we have
$`3\delta _{2SR}(𝒙)\delta _{1R}(𝒙)^2`$ $`=`$ $`{\displaystyle \frac{d𝒌}{(2\pi )^3}\frac{d𝒍}{(2\pi )^3}\frac{d𝒎}{(2\pi )^3}\frac{d𝒏}{(2\pi )^3}(𝒌+𝒍)^2R^2}`$
$`\mathrm{exp}\left[{\displaystyle \frac{(2𝒍^2+𝒌^2+𝒎^2+𝒏^2+2𝒌𝒍)R^2}{2}}\right]`$
$`\times \delta _1(𝒌)\delta _1(𝒍)\delta _1(𝒎)\delta _1(𝒏)\mathrm{exp}[i(𝒌+𝒍+𝒎+𝒏)𝒙].`$
Using equation (A13) we can simplify the above integral as
$$2\frac{d𝒌}{(2\pi )^3}\frac{d𝒍}{(2\pi )^3}\mathrm{exp}\left[\frac{(3𝒍^2+2𝒌^2+2𝒌𝒍)R^2}{2}\right]P(k)P(l)(𝒌+𝒍)^2R^2.$$
(A20)
Note that the integrad of this expression depends only on the information of the shape of the triangle determined by two vectors $`𝒌`$ and $`𝒍`$. This triangle is characterized by three quantities, namely, two sides $`k=|𝒌|`$, $`l=|𝒍|`$ and cosine between them $`u𝒌𝒍/(kl)`$ with $`1u1`$. We change variables from $`\{𝒌,𝒍\}`$ to $`\{k,l,u\}`$. The volume element is deformed as
$$d𝒌d𝒍8\pi ^2dkdldu.$$
(A21)
Thus we obtain equation (39). This equation looks somewhat complicated. For power-law models, however, we can easily evaluate it using mathematica (Wolfram 1996) and finally obtain analytical expression (40) given in the main text.
### A.2 Generalized Skewness for the Gather Approach
For this approach we have the following relation for a Gaussian filter
$$R\frac{}{R}\delta _{1R}(𝒙)=R^2^2\delta _{1R}(𝒙).$$
(A22)
Therefore the correction terms (A4) and (A5) can be written by combinations of the following five variables <sup>2</sup><sup>2</sup>2In this subsection we denote $`\delta _{1R}(𝒙)`$ simply by $`\delta `$.
$$\{\delta ,\delta ,^2\delta ,^3\delta ,^4\delta \}.$$
(A23)
For example, equation (A4) is written as
$$\frac{1}{3}R^2\left[\delta ^2^2(\delta ^2\delta )+2\delta ^2(^2\delta )(^2\delta )2\delta ^2\delta \delta ^2\delta \right]$$
(A24)
Using property (A12), the above expression is deformed as
$$\frac{1}{3}R^2\left[3^2\delta ^2\delta \delta \delta +2\delta ^3\delta \delta \delta +3\delta ^4\delta \delta \delta +4\delta ^2\delta ^2\right]$$
(A25)
The moments appeared in the above equation can be written in terms of $`P(k)`$ as
$`\delta ^2\delta `$ $`=`$ $`\delta \delta ={\displaystyle \frac{dk}{2\pi ^2}k^4P(k)e^{k^2R^2}}`$ (A26)
$`^2\delta ^2\delta `$ $`=`$ $`^3\delta \delta =^4\delta \delta ={\displaystyle \frac{dk}{2\pi ^2}k^6P(k)e^{k^2R^2}}`$ (A27)
For a power-law models ($`P(k)k^n`$). These integrals are evaluated respectively as
$$\sigma _R^2R^2\left(\frac{n+3}{2}\right),\sigma _R^2R^4\left(\frac{n+3}{2}\right)\left(\frac{n+5}{2}\right).$$
(A28)
With the definition of $`T`$ parameter (eq.) we obtain the final result that is given in equation (61) as
$$\mathrm{\Delta }T_{GR}=\frac{2}{3}(n+4)$$
(A29)
To calculate the correction term (A5), let us use the Fourier space representation (A10).<sup>3</sup><sup>3</sup>3We obtain the same result starting from equation (A6). It is straightforward to get
$`\delta _1\delta _1^2\delta _{2A}+2\delta _1\delta _{2A}^2\delta _1`$ $`=`$ $`{\displaystyle \frac{d𝒌}{(2\pi )^3}\frac{d𝒍}{(2\pi )^3}\frac{d𝒎}{(2\pi )^3}\frac{d𝒏}{(2\pi )^3}\frac{k^2R^2}{3}}`$ (A30)
$`\times \mathrm{exp}\left[{\displaystyle \frac{(𝒍^2+𝒌^2+𝒎^2+𝒏^2)R^2}{2}}\right]`$
$`\times \delta _1(𝒌)\delta _1(𝒍)\delta _1(𝒎)\delta _1(𝒏)`$
$`\times \mathrm{exp}[i(𝒌+𝒍+𝒎+𝒏)𝒙],`$
$`\times [(𝒎𝒏)(𝒌+𝒍)^2`$
$`(𝒎(𝒌+𝒍)𝒏^2)(𝒏(𝒌+𝒍)𝒎^2)].`$
With equation (A13), the above integral becomes
$`{\displaystyle \frac{1}{3}}{\displaystyle \frac{d𝒌}{(2\pi )^3}\frac{d𝒍}{(2\pi )^3}k^2R^2\mathrm{exp}\left[(𝒍^2+𝒍^2)R^2\right]P(k)P(l)[4k^2l^24(𝒌𝒍)^2]}`$ (A31)
$`=`$ $`{\displaystyle \frac{1}{6\pi ^4}}{\displaystyle _0^{\mathrm{}}}𝑑k{\displaystyle _0^{\mathrm{}}}𝑑l{\displaystyle _1^1}𝑑u\mathrm{exp}\left[(l^2+k^2)R^2\right]k^6l^4P(k)P(l)(1u^2)R^2`$
$`=`$ $`{\displaystyle \frac{2}{9\pi ^4}}{\displaystyle _0^{\mathrm{}}}k^6P(k)\mathrm{exp}[k^2R^2]{\displaystyle 𝑑ll^4P(l)\mathrm{exp}[l^2R^2]}.`$
For power-law models this expression becomes (see eqs.\[A26\]-\[A28\])
$$\frac{8}{9}\sigma _R^4\left(\frac{n+3}{2}\right)^2\left(\frac{n+5}{2}\right).$$
(A32)
Using definition of $`U`$ parameter (eq.) we obtain equation (62) as
$$\mathrm{\Delta }U_{GR}=\frac{1}{3}(n+5).$$
(A33)
Note that results (A29) and (A33) are not valid for general filters. Equation (A22) that holds for the Gaussian filter plays crucial roles to derive them.
### A.3 Generalized Skewness for the Scatter Approach
First we evaluate the correction term given in equation (A4). With the Fourier space representation (A10) we obtain the following equation
$`\delta _1^2^2\delta _{2A}+2\delta _1\delta _{2A}^2\delta _1`$ $`=`$ $`{\displaystyle \frac{d𝒌}{(2\pi )^3}\frac{d𝒍}{(2\pi )^3}\frac{d𝒎}{(2\pi )^3}\frac{d𝒏}{(2\pi )^3}\frac{(𝒌+𝒍)^2R^2}{3}}`$ (A34)
$`\times \mathrm{exp}\left[{\displaystyle \frac{(2𝒍^2+𝒌^2+𝒎^2+𝒏^2+2𝒌𝒍)R^2}{2}}\right]`$
$`\times \delta _1(𝒌)\delta _1(𝒍)\delta _1(𝒎)\delta _1(𝒏)\mathrm{exp}[i(𝒌+𝒍+𝒎+𝒏)𝒙],`$
$`\times [(𝒌+𝒍)^2+𝒎^2+𝒏^2]`$
With the formula (A13), this twelfth-dimensional integral becomes
$$\frac{2}{3}\frac{d𝒌}{(2\pi )^3}\frac{d𝒍}{(2\pi )^3}\mathrm{exp}\left[\frac{(3𝒍^2+2𝒌^2+2𝒌𝒍)R^2}{2}\right]P(k)P(l)(𝒌+𝒍)^2R^2[(𝒌+𝒍)^2+𝒍^2+𝒌^2].$$
(A35)
changing variables from $`d𝒌d𝒍`$ to $`dkdldu`$ as shown in relation (A21), we obtain the result essentially same as equation (65) as
$`\delta _1^2^2\delta _{2A}+2\delta _1\delta _{2A}^2\delta _1`$ $`=`$ $`{\displaystyle \frac{1}{6\pi ^4}}{\displaystyle _0^{\mathrm{}}}𝑑k{\displaystyle _0^{\mathrm{}}}𝑑l{\displaystyle _1^1}𝑑u\mathrm{exp}\left[{\displaystyle \frac{(3l^2+2k^2+2klu)R^2}{2}}\right]`$ (A36)
$`\times k^2l^2P(k)P(l)(k^2+l^2+2klu)(k^2+l^2+klu)R^2,`$
As in the case of skewness parameter, we can evaluate this complicated integrals with mathematica and obtain equation (67).
In the same manner, equation (A5) is written as
$`\delta _1\delta _1^2\delta _{2A}+2\delta _1\delta _{2A}^2\delta _1`$ $`=`$ $`{\displaystyle \frac{d𝒌}{(2\pi )^3}\frac{d𝒍}{(2\pi )^3}\frac{d𝒎}{(2\pi )^3}\frac{d𝒏}{(2\pi )^3}\frac{(𝒌+𝒍)^2R^2}{3}}`$ (A37)
$`\times \mathrm{exp}\left[{\displaystyle \frac{(2𝒍^2+𝒌^2+𝒎^2+𝒏^2+2𝒌𝒍)R^2}{2}}\right]`$
$`\times \delta _1(𝒌)\delta _1(𝒍)\delta _1(𝒎)\delta _1(𝒏)`$
$`\times \mathrm{exp}[i(𝒌+𝒍+𝒎+𝒏)𝒙],`$
$`\times [(𝒎𝒏)(𝒌+𝒍)^2`$
$`(𝒎(𝒌+𝒍)𝒏^2)(𝒏(𝒌+𝒍)𝒎^2)].`$
This expression is simplified to the following form
$`{\displaystyle \frac{1}{3}}{\displaystyle \frac{d𝒌}{(2\pi )^3}\frac{d𝒍}{(2\pi )^3}(𝒌+𝒍)^2R^2\mathrm{exp}\left[\frac{(2𝒍^2+𝒌^2+𝒎^2+𝒏^2+2𝒌𝒍)R^2}{2}\right]}`$
$`\times P(k)P(l)[4k^2l^24(𝒌𝒍)^2].`$ (A38)
Again, changing variables, we obtain the expression (66) as
$`\delta _1\delta _1^2\delta _{2A}+2\delta _1\delta _{2A}^2\delta _1`$ $`=`$ $`{\displaystyle \frac{1}{6\pi ^4}}{\displaystyle _0^{\mathrm{}}}𝑑k{\displaystyle _0^{\mathrm{}}}𝑑l{\displaystyle _1^1}𝑑u(1u^2)R^2`$ (A39)
$`\times \mathrm{exp}\left[{\displaystyle \frac{(3l^2+2k^2+2klu)R^2}{2}}\right]`$
$`\times k^4l^4P(k)P(l)(k^2+l^2+2klu).`$
For power-law models we can evaluate this integral using mathematica and obtain equation (68).
TABLE 1
skewness for the gather approach (Gaussian filter)
| spectral index $`n`$ | 1 | 0 | -1 | -2 | -3 |
| --- | --- | --- | --- | --- | --- |
| $`S_F(n)`$ | 3.029 | 3.144 | 3.468 | 4.022 | 4.857 |
| $`\mathrm{\Delta }S_G(n)`$ | 4.000 | 3.000 | 2.000 | 1.000 | 0 |
| $`S_G(n)`$ | 7.029 | 6.144 | 5.468 | 5.022 | 4.857 |
TABLE 2
skewness for the scatter approach
| spectral index $`n`$ | 1 | 0 | -1 | -2 | -3 |
| --- | --- | --- | --- | --- | --- |
| $`S_F(n)`$ | 3.029 | 3.144 | 3.468 | 4.022 | 4.857 |
| $`\mathrm{\Delta }_SS(n)`$ | 3.031 | 2.576 | 2.045 | 1.277 | 0 |
| $`S_S(n)`$ | 6.060 | 5.720 | 5.513 | 5.299 | 4.857 |
TABLE 3
generalized skewness for the gather approach
| spectral index $`n`$ | 1 | 0 | -1 | -2 | -3 |
| --- | --- | --- | --- | --- | --- |
| $`T_F(n)`$ | 2.020 | 2.096 | 2.312 | 2.681 | 3.238 |
| $`\mathrm{\Delta }T_G(n)`$ | 3.333 | 2.667 | 2.000 | 1.333 | 0.667 |
| $`T_G(n)`$ | 5.353 | 4.763 | 4.312 | 4.014 | 3.905 |
| $`U_G(n)`$ | 1.431 | 1.292 | 1.227 | 1.222 | 1.272 |
| $`\mathrm{\Delta }U_G(n)`$ | 2.000 | 1.667 | 1.333 | 1.000 | 0.667 |
| $`U_G(n)`$ | 3.431 | 2.959 | 2.560 | 2.222 | 1.929 |
TABLE 4
generalized skewness for the scatter approach
| spectral index $`n`$ | 1 | 0 | -1 | -2 | -3 |
| --- | --- | --- | --- | --- | --- |
| $`T_F(n)`$ | 2.020 | 2.096 | 2.312 | 2.681 | 3.238 |
| $`\mathrm{\Delta }T_S(n)`$ | -2.082 | -1.908 | -1.723 | -1.451 | -0.963 |
| $`T_S(n)`$ | -0.0623 | 0.1882 | 0.5892 | 1.230 | 2.275 |
| $`U_F(n)`$ | 1.431 | 1.292 | 1.227 | 1.222 | 1.272 |
| $`\mathrm{\Delta }U_S(n)`$ | -1.265 | -1.145 | -1.027 | -0.8916 | -0.7105 |
| $`U_S(n)`$ | 0.1662 | 0.1474 | 0.2000 | 0.3301 | 0.5611 | |
warning/0002/astro-ph0002377.html | ar5iv | text | # Quantum corrections to the ground state energy of inhomogeneous neutron matter
## Abstract
We estimate the quantum corrections to the ground state energy in neutron matter (which could be termed as well either shell correction energy or Casimir energy) at subnuclear densities, where various types of inhomogeneities (bubbles, rods, plates) are energetically favorable. We show that the magnitude of these energy corrections are comparable to the energy differences between various types of inhomogeneous phases. We discuss the dependence of these corrections on a number of physical parameters (density, filling factor, temperature, lattice distortions).
PACS numbers: 21.10.Dr, 21.65.+f, 97.60.Jd
The investigation of the nuclear matter in the neutron star crust below the saturation density leads to the consideration of exotic shapes of the nuclei immersed in a neutron gas. It was realized long ago that when the nuclei in dense matter occupy more than half of the space it is energetically favorable to “turn the nuclei inside out” and form a bubble phase . To date a large number of calculations have been performed pertaining to the structure of the neutron star crust. The liquid drop model or the Thomas–Fermi approximation calculations predict rather small energy differences between different phases, of the order of a few $`keV/fm^3.`$ (N.B. even though we often refer to energy, we actually mean energy density.). Apparently an agreement has been reached concerning the existence of the following chain of phase changes as the density is increasing: nuclei $``$ rods $``$ plates $``$ tubes $``$ bubbles $``$ uniform matter. The density range for these phase transitions is $`0.040.1fm^3`$ . Moreover, it was established that these phases exist up to temperatures of about $`10MeV`$ . At densities of the order of several nuclear densities the quark degrees of freedom get unlocked and the formation of various quark matter droplets embedded in nuclear matter becomes then energetically favorable .
The appearance of different phases is attributed to the interplay between the Coulomb and surface energies. Since most of the published works were based on the minimization of some density functional in a single Wigner–Seitz cell, the calculation of the shell correction or Casimir energy has been omitted. In Hartree–Fock calculations these quantum corrections to the ground state energy of neutron matter are obviously automatically incorporated. The Hartree–Fock calculations performed so far were limited to “spherical Wigner–Seitz cells”, which is arguably a reasonable approximation for the “nuclei in neutron gas” phase only. To our knowledge there exist only one study on this subject where the shell effects due to the bound nucleons only however (mainly protons) have been taken into account . It was determined that the shell correction energy is smaller than the energy difference between different phases and it was thus concluded that quantum corrections to the ground state energy will not lead to any qualitative changes in the sequence of the nuclear shape transitions in the neutron star crust.
Our goal is to reach a comprehensive understanding of the so called shell correction or Casimir energy in neutron stars. There is no well established terminology for the energy corrections we are considering here, even though the problem has been addressed before to some extent by other authors. In the case of finite systems, the energy difference between the true binding energy and the liquid drop energy of a given system is typically refered to as shell correction energy. In field theory a somewhat similar energy appears, due to various fluctuation induced effects and it is generically referred to as the Casimir energy :
$$E_{Casimir}=_{\mathrm{}}^{\mathrm{}}𝑑\epsilon \epsilon [g(\epsilon ,𝐥)g_0(\epsilon )],$$
(1)
where $`g_0(\epsilon )`$ is the density of states per unit volume for the fields in the absence of any objects, $`g(\epsilon ,𝐥)`$ is the density of states per unit volume in the presence of some “foreign”objects, such as plates, spheres, etc., and $`𝐥`$ is an ensemble of geometrical parameters describing these objects and their relative geometrical arrangement. A similar formula can be written for neutron matter energy
$$E_{nm}=_{\mathrm{}}^\mu 𝑑\epsilon \epsilon g(\epsilon ,𝐥)_{\mathrm{}}^{\mu _0}𝑑\epsilon \epsilon g_0(\epsilon ,𝐥),$$
(2)
with the notable difference in the upper integration limit. In the above equation $`g_0(\epsilon ,𝐥)`$ stands for the Thomas–Fermi or liquid drop density of states of the inhomogeneous phase and $`g(\epsilon ,𝐥)`$ for the true quantum density of states in the presence of inhomogeneities. The parameters: $`\mu `$ and $`\mu _0`$ are determined from the requirement that the system has a given average density
$$\rho =_{\mathrm{}}^\mu 𝑑\epsilon g(\epsilon ,𝐥)=_{\mathrm{}}^{\mu _0}𝑑\epsilon g_0(\epsilon ,𝐥).$$
(3)
Since in infinite matter the presence of various inhomogeneities does not lead to the formation of discrete levels, one might expect to refer to corresponding energy correction for neutron matter as the Casimir energy. In Ref. the authors computed a somewhat different quantity however, than the one we are interested in this work, the correction to the ground state energy arising from existence of almost discrete levels inside a nucleus in an infinite medium. Strictly speaking these levels are not discrete, but form narrow energy bands due to the tunneling between neighboring nuclei. The effects we shall consider here arise from the “outside” states, which is in complete analogy with the procedure for computing the Casimir energy. As we shall show, these energy corrections, arising from the existence of these truly delocalized states, are larger than those considered in Ref. . We have considered similar issues earlier in finite systems and to some extent in infinite 2–dimensional systems as well in Refs. .
In order to better appreciate the nature of the problem we are addressing in this work, let us consider the following situation. Let us imagine that two spherical identical bubbles have been formed in an otherwise homogeneous neutron matter. For the sake of simplicity, we shall assume that the bubbles are completely hollow. We shall sidestep the question of stability of each bubble for the moment and assume that they are stable and rigid as well. We shall ignore the role of long range forces, namely the Coulomb interaction in the case of neutron stars, as their main contribution is to the smooth, liquid drop or Thomas–Fermi part of the total energy. Under such circumstances one can ask the following apparently innocuous question: “What determines the most energetically favorable arrangement of the two bubbles?” According to a liquid drop model approach (completely neglecting for the moment the possible stabilizing role of the Coulomb forces) the energy of the system should be insensitive to the relative positioning of the two bubbles. A similar question was raised in condensed matter studies, concerning the interaction between two impurities in an electron gas. In the case of two “weak” and point–like impurities the dependence of the energy of the system as a function of the relative distance between the two impurities $`𝐚`$ is given by (spin coordinates are not displayed)
$$E(𝐚)=\frac{1}{2}𝑑𝐫_1𝑑𝐫_2V_1(𝐫_1)\chi (𝐫_1𝐫_2𝐚)V_2(𝐫_2),$$
(4)
where $`\chi (𝐫_1𝐫_2𝐚)`$ is the Lindhard response function of a homogeneous Fermi gas and $`V_1(𝐫_1)`$ and $`V_2(𝐫_2)`$ are the potentials describing the interaction between impurities and the surrounding electron gas. At large distances $`k_Fa1`$ this expression leads to the interaction first derived by Ruderman and Kittel :
$$E(𝐚)\frac{\mathrm{}^2}{2mk_Fa^3}\mathrm{cos}(2k_Fa),$$
(5)
where $`k_F`$ is the Fermi wave vector and $`m`$ is the fermion mass
$$\mu =\frac{\mathrm{}^2k_F^2}{2m}.$$
(6)
This asymptotic behavior is valid only for point–like impurities, when $`k_FR1`$, and where $`R`$ stands for the radius of the two impurities. This condition is typically violated for either nuclei embedded in a neutron gas or bubbles, when typically $`k_FR1`$. As we shall show, in the case of large “impurities” (when $`k_FR1`$) the interaction energy changes in a rather dramatic manner. If one replaces the “weak” impurities with “strong” point–like impurities, only the magnitude of the interaction changes at large distances, but not the form . The interaction (5) has a pure quantum character, and any “noise” (e.g. temperature) leads to a quick disappearance of the oscillatory behavior and with it of the power law character, and the regular Debye screening (which is exponential in character) sets in instead.
The lesson one can learn from this analysis however, is that quantum corrections are most likely responsible for the interaction of two bubbles/nuclei embedded in a Fermi gas and the form of the interaction (5) suggests the most natural way to proceed. The argument of the cosine is nothing else but the classical action in units of $`\mathrm{}`$ of the bouncing periodic orbit between the two impurities. The exact form and magnitude of the coefficient in front of the cosine can be obtained in a semiclassical approximation only after a careful estimation of the leading order correction to the leading semiclassical result. Using the 3–dimensional extension of the semiclassical approximation to the so called small disks problem , we were able to obtain a significantly simpler and more transparent derivation of this interaction than the original derivation as follows. The correction to the single–particle propagator, which depends on the presence of the two weak widely separated point–like impurities ($`R/a1`$ and $`k_FR1`$) is
$`\delta G(𝐫,𝐫^{},k)`$ $``$ $`G_0(𝐫,𝐫_1,k)G_0(𝐫_1,𝐫_2,k)G_0(𝐫_2,𝐫^{},k)`$ (7)
$`+`$ $`G_0(𝐫,𝐫_2,k)G_0(𝐫_2,𝐫_1,k)G_0(𝐫_1,𝐫^{},k),`$ (8)
where
$$G_0(𝐫_1,𝐫_2,k)=\frac{m\mathrm{exp}(ik|𝐫_1𝐫_2|)}{2\pi \mathrm{}^2|𝐫_1𝐫_2|}$$
(9)
is the free single–particle propagator. Since only “periodic orbits” contribute to the density of states, the correction to the density of states, due to the presence of the two impurities and which depends on their relative separation only is given by
$$\delta g(k,|𝐫_1𝐫_2|)\mathrm{Im}\left[\frac{i\mathrm{exp}(2i|𝐫_1𝐫_2|)}{k|𝐫_1𝐫_2|}\right]$$
(10)
and the corresponding correction to the ground state energy is given by the obvious formula
$`\delta (|𝐫_1𝐫_2|)`$ $``$ $`{\displaystyle _{kk_F}}k𝑑k\left({\displaystyle \frac{\mathrm{}^2k^2}{2m}}\mu \right)\delta g(k,|𝐫_1𝐫_2|)`$ (11)
$``$ $`{\displaystyle \frac{\mathrm{cos}(2k_F|𝐫_1𝐫_2|)}{|𝐫_1𝐫_2|^3}}.`$ (12)
Only the leading term in the limit $`|𝐫_1𝐫_2|\mathrm{}`$ is explicitly shown here. The proportionality coefficient is naturally determined by the impurity strength.
The formation of various inhomogeneities in an otherwise uniform Fermi gas can be characterized by several natural dimensionless parameters, $`k_Fa1`$, where as above $`a`$ is a characteristic separation distance between two such inhomogeneities, $`k_FR1`$, where $`R`$ is a characteristic size of such an inhomogeneity, and $`k_Fs1`$, where $`s`$ is a typical “skin” thickness of such objects. The fact that the first two parameters, $`k_Fa`$ and $`k_FR`$, are both very large makes a semiclassical approach natural. Since the third parameter, $`k_Fs`$, is never too large or too small, one might be tempted to discard a semiclassical treatment of the entire problem altogether. However, there is a large body of evidence pointing towards the fact that even though this parameter in real systems is of order unity, the approximation $`k_Fs1`$, which we shall adopt in this work, is surprisingly accurate . Moreover the corrections arising from considering $`k_Fs=𝒪(1)`$ should lead to an overall energy shift mainly, which is largely independent of the separation among various objects embedded in a Fermi gas. On one hand, this type of shift can be accounted for in principle in a correctly implemented liquid drop model or Thomas–Fermi approximation. On the other hand, the semiclassical corrections to the ground state energy arising from the relative arrangement of various inhomogeneities have to be computed separately, as they have a different physical nature. We are thus lead to the natural assumption that a simple hard–wall potential model for various types of inhomogeneities appearing in a neutron Fermi gas is a reasonable starting point to estimate quantum corrections to the ground state energy, see Refs. and earlier references therein. We shall refer to these quantum corrections to the energy as shell effects in the rest of the paper. One might expect that such simplifications will result in an overestimation of the magnitude of shell effects, but the qualitative pattern should remain the same.
We shall consider spherical bubble–like, rod–like and plate–like phases only here and we shall estimate the shell correction or Casimir energy arising due to a regular arrangement of such inhomogeneities in an otherwise homogeneous neutron gas. One can distinguish two types of “bubbles”: i) nuclei–like structures embedded in a neutron gas and ii) void–like structures. By voids we mean the regions in which the nuclear density is significantly lower than in the surrounding space. In the first case i), the single particle wave functions can be separated into roughly two classes, those localized mostly inside the nuclei–like structures and those which are completely delocalized. A fermion in a delocalized state will spend some time inside the “nuclei” too, but since the potential experienced by a nucleon is deeper there, the local momentum is larger and thus the relative time and relative probability to find a nucleon in this region is smaller. One can approximately replace then the “nuclei” with an effective repulsive potential of roughly the same shape. In the case of a “bubble”, when the probability to find a nucleon inside a “bubble” is reduced, again such an approximation appears as reasonable. The “nuclei” and “bubbles” we are refering to here, are not necessarily spherical, but could have the shape of a rod or plate as well. There are of course a number of “resonant” delocalized states, whose amplitude behaves in a manner just opposite to the one we have described here. However, the number of such “resonant” states is typically small and we thus do not expect large effects due to them. Moreover, since such states are concentrated mostly inside a “nucleus” or a “bubble” one does not expect them to affect in a major way the relative positioning of two “nuclei” or two “bubbles”. Nevertheless, these are some issues, which certainly deserve more scrutiny in the future, even though we hardly expect that a more comprehensive analysis will lead to qualitative changes of our conclusions. In all these phases the shell effects depend on the structure and stability of periodic orbits in the system . Except for the plate–like nuclei phase, where the shell energy can be computed exactly, for other geometries one should calculate the contribution from all periodic orbits. This is rather tedious task, since they proliferate exponentially as a function of their length , and moreover this is not really necessary to perform. If one is interested in the gross structure only of the shell effects, the contribution of the shortest periodic orbits should suffice for defining the gross shell structure. (We remind the reader, in an infinite medium there are really no shells as in a finite system, but we refer to the corresponding effects in this manner only, due to their similar origin because of the appearance of periodic orbits.) Since the contribution of any given periodic orbit leads to an oscillatory contribution to the density of states any suitably chosen energy averaging over the spectrum, and in particular a finite temperature as well, will leave only the contributions due to the shortest periodic orbits. Since a periodic orbit of length $``$ will lead to detail on an energy scale of the order $`\mathrm{\Delta }E=\mathrm{}^2\pi ^2/2m^2`$, performing an averaging over an energy interval $`\mathrm{\Delta }E`$ will effectively mask the contribution of orbits of length $``$ or larger. Moreover, since the geometry of the rod–like and spherical phases admit only unstable (hyperbolic) orbits, the longer the orbits, the lesser their contribution is, due their decreased stability.
The simplest system consist of plate–like nuclei with the neutron gas filling the space between slabs. The shell energy for this system per unit volume can be easily evaluated:
$`{\displaystyle \frac{E_{shell}}{L^3}}`$ $`=`$ $`{\displaystyle \frac{EE_{Weyl}+\mathrm{\Delta }E}{L^3}},`$ (13)
$`\mathrm{\Delta }E`$ $`=`$ $`\mu (\rho _0\rho _{Weyl})L^3`$ (14)
where the exact and the Weyl (smooth) energy per unit volume are given by
$`{\displaystyle \frac{E}{L^3}}`$ $`=`$ $`{\displaystyle \frac{2}{L^3}}{\displaystyle \frac{\mathrm{}^2}{2ma^2}}{\displaystyle \frac{\pi ^3}{2}}\left({\displaystyle \frac{L}{a}}\right)^2[{\displaystyle \frac{1}{4}}\left({\displaystyle \frac{k_Fa}{\pi }}\right)^4N`$ (15)
$``$ $`{\displaystyle \frac{N(N+1)(2N+1)(3N^2+3N1)}{120}}],`$ (16)
$`{\displaystyle \frac{E_{Weyl}}{L^3}}`$ $`=`$ $`{\displaystyle \frac{2}{L^3}}{\displaystyle \frac{\mathrm{}^2}{2ma^2}}{\displaystyle \frac{\pi ^3}{2}}\left({\displaystyle \frac{L}{a}}\right)^2[{\displaystyle \frac{1}{5}}\left({\displaystyle \frac{k_Fa}{\pi }}\right)^5`$ (17)
$``$ $`{\displaystyle \frac{1}{8}}\left({\displaystyle \frac{k_Fa}{\pi }}\right)^4]`$ (18)
In the above formula
$$N=\mathrm{Int}\left[\frac{k_Fa}{\pi }\right]$$
(19)
stands for the integer part of the argument in the square brackets, and $`a=L2R`$ is the distance between slabs and $`R`$ is the half of the width of the slab. Here $`L^3`$ is the volume of an elementary (cubic) cell and the factor $`{}_{}{}^{}2_{}^{}`$ in front stands for the two spin states. The average matter density (the number of neutrons per unit volume) $`\rho _0`$ and the smoothed density $`\rho _{Weyl}`$ are determined by relations
$`\rho _0=2{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{d^2k}{(2\pi )^2}\mathrm{\Theta }\left(\mu \frac{\mathrm{}^2k^2}{2m}\frac{\mathrm{}^2n^2\pi ^2}{2ma^2}\right)}`$ (20)
$`={\displaystyle \frac{2}{L^3}}{\displaystyle \frac{\pi }{4}}\left({\displaystyle \frac{L}{a}}\right)^2\left[\left({\displaystyle \frac{k_Fa}{\pi }}\right)^2N{\displaystyle \frac{N(N+1)(2N+1)}{6}}\right].`$ (21)
$`\rho _{Weyl}={\displaystyle \frac{2}{L^3}}{\displaystyle \frac{\pi }{2}}\left({\displaystyle \frac{L}{a}}\right)^2\left[{\displaystyle \frac{1}{3}}\left({\displaystyle \frac{k_Fa}{\pi }}\right)^3{\displaystyle \frac{1}{4}}\left({\displaystyle \frac{k_Fa}{\pi }}\right)^2\right].`$ (22)
Using these formulas one can show that the shell correction energy has the behavior
$$\frac{E_{shell}}{L^3}=\frac{\mathrm{}^2k_F^2}{40a^2Lm}G\left(\frac{k_Fa}{\pi }\right),$$
(23)
where $`G(x)`$ is an approximate periodic function of its argument, for $`x1`$), $`G(x+1)G(x)`$, with properties $`G(x=n/2)0`$ and approximately $`1G(x)1`$. Furthermore
$$\rho _{out}=\frac{\rho _0}{v}$$
(24)
is the actual density of the neutron gas between the two slabs and
$$v=1u=\frac{L2R}{L}$$
(25)
is the filling factor, which is the ratio of the occupied volume to the volume of the cell. One can show also that
$$\rho _0=\rho _{Weyl}+\frac{k_F}{12La}F\left(\frac{k_fa}{\pi }\right),$$
(26)
where $`\rho _{Weyl}`$ is the Weyl approximation to the density and $`F(x+1)F(x)`$ is an approximate periodic function of its argument too, for $`x1`$, with properties $`1F(x)0.5`$ and $`F(x=n)=1`$. This periodicity leads to the clear pattern of “valleys” ($`k_Fa=(n+3/4)\pi `$) and “ridges” ($`k_Fa(n+1/4)\pi `$) in the profile of the shell energy shown in Fig. 1a. These features of the energy and density are naturally related to fact that these quantities are almost periodic functions in the classical action along the only periodic orbit in the system, i.e. in the variable $`S=2k_Fa`$.
In the case of rod–like and spherical voids we shall use the semiclassical theory in order to compute the shell energy. Since we are interested only in the “gross shell structure” we have to take into account a few of the shortest periodic orbits among the nearest neighbors only. The lengths of the shortest periodic orbits depend on the lattice type. In the following we will assume the simple cubic and simple square lattices for spherical and rod–like phases respectively. The expression for the shell energy density and the neutron density reads:
$`{\displaystyle \frac{E_{shell}}{L^3}}`$ $`=`$ $`{\displaystyle \frac{1}{L^3}}{\displaystyle _0^\mu }(\epsilon \mu ){\displaystyle \underset{i}{}}g_{shell}(\epsilon ,L_i)d\epsilon `$ (27)
$`\rho _0`$ $`=`$ $`{\displaystyle \frac{1}{L^3}}{\displaystyle _0^\mu }\left[g_{Weyl}(\epsilon )+{\displaystyle \underset{i}{}}g_{shell}(\epsilon ,L_i)\right]𝑑\epsilon ,`$ (28)
where $`g_{shell}(\epsilon ,L_i)`$ denotes the contribution to the level density due to the orbit $`L_i`$ and $`g_{Weyl}`$ is the smooth level density determined using the Weyl prescription .
For the rod–like phase we took into account four orbits of the length $`2L_1=2(L2R)`$ and four orbits of the length $`2L_2=2(L\sqrt{2}2R)`$. Introducing longer orbits did not lead to noticeable changes in the patterns presented here. Hence the shell energy per volume is equal to:
$$\frac{E_{shell}}{L^3}=\frac{1}{L^3}_0^\mu (\epsilon \mu )\underset{i=1}{\overset{2}{}}A_ig_{shell}(\epsilon ,L_i)d\epsilon ,$$
(29)
where $`A_1=A_2=4`$ and the chemical potential $`\mu `$ is determined by the condition:
$$\rho _0=\rho _{Weyl}+\frac{1}{L^3}_0^\mu \underset{i=1}{\overset{2}{}}A_ig_{shell}(\epsilon ,L_i)d\epsilon .$$
(30)
A periodic orbit of the type considered by us gives actually a contribution with a factor 1/2, since only half of it belongs to a particular elementary cell. Because there are two spin states, and thus eight orbits in total, each type of orbit eventually is weighted by four. The density of states was evaluated using the convolution of the exact 1–dimensional density of states and the density of states given by Gutzwiller trace formula for the 2–dimensional system of disks, which is the cross section of the rod–like system we are interested in. In some cases such a procedure can lead to spurious contributions, which are however rather easy to single out, see Refs. . For a given periodic orbit of length $`2L_i`$, the shell correction to the density of states is given by the following expression:
$`g_{shell}(\epsilon ,L_i)={\displaystyle \frac{mLL_i}{2\pi \mathrm{}^2}}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\text{J}_0(2nkL_i)}{\mathrm{sinh}n\kappa _i}},`$ (31)
where the summation is over repetitions of this orbit and
$$\epsilon =\frac{\mathrm{}^2k^2}{2m}.$$
(32)
When one is interested in the gross shell structure then the contribution of long orbits as well as the contributions due to repetitions of short primitive orbits vanish under energy averaging.
The explicit form of the shell energy and of the fluctuating part of the density reads:
$`{\displaystyle \frac{E_{shell}}{L^3}}={\displaystyle \frac{1}{L^3}}{\displaystyle \frac{\mathrm{}^2k_F^2}{2m\pi }}{\displaystyle \frac{1}{4}}{\displaystyle \underset{i=1}{\overset{2}{}}}A_i{\displaystyle \frac{L}{L_i}}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\text{J}_2(2nk_FL_i)}{n^2\mathrm{sinh}(n\kappa _i)}},`$ (33)
$`\rho _0=\rho _{Weyl}+{\displaystyle \frac{k_F}{4\pi L^2}}{\displaystyle \underset{i=1}{\overset{2}{}}}A_i{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\text{J}_1(2nk_FL_i)}{n\mathrm{sinh}(n\kappa _i)}}.`$ (34)
The parameter $`\kappa _i`$ determines the stability of the orbit $`L_i`$:
$$\kappa _i=\mathrm{ln}\left[1+\frac{L_i}{R}+\sqrt{\frac{L_i}{R}\left(\frac{L_i}{R}+2\right)}\right].$$
(35)
The shell energy as a function of the anti–filling factor (relative void volume) $`u={\displaystyle \frac{\pi R^2}{L^2}}`$ and $`\rho _0`$ is shown in Fig 1b. The shell energy has a smaller amplitude then in the case of the plate–like phase. This comes about because the periodic orbits are now hyperbolic in the plane perpendicular to the rods. Note however that the pattern of “valleys” and “ridges” looks very similar to the one for the slabs. This is because the main contribution due to the classical orbit of length $`2L_1`$ is the same. There are small interference effects caused by the orbit of length $`2L_2`$ however. Since it is longer, this second trajectory contributes with a smaller weight.
For the case of spherical voids there are $`26`$ periodic orbits between nearest neighbors of three different lengths $`2L_1=2(L2R)`$, $`2L_2=2(L\sqrt{2}R)`$ and $`2L_3=2(L\sqrt{3}R)`$. Thus the shell energy and density are equal to:
$`{\displaystyle \frac{E_{shell}}{L^3}}`$ $`=`$ $`{\displaystyle \frac{1}{L^3}}{\displaystyle _0^\mu }(\epsilon \mu ){\displaystyle \underset{i=1}{\overset{3}{}}}A_ig_{shell}(\epsilon ,L_i)d\epsilon ,`$ (36)
$`\rho _0`$ $`=`$ $`\rho _{Weyl}+{\displaystyle \frac{1}{L^3}}{\displaystyle _0^\mu }{\displaystyle \underset{i=1}{\overset{3}{}}}A_ig_{shell}(\epsilon ,L_i)d\epsilon .`$ (37)
The contribution due to one periodic orbit to the fluctuating part of the level density reads:
$$g_{shell}(\epsilon ,L_i)=\frac{mL_i}{2\pi \mathrm{}^2k}\underset{n=1}{\overset{\mathrm{}}{}}\frac{\mathrm{cos}(2nkL_i)}{\mathrm{sinh}^2(n\kappa _i)}.$$
(38)
Hence we get:
$`{\displaystyle \frac{E_{shell}}{L^3}}={\displaystyle \frac{1}{L^3}}{\displaystyle \frac{\mathrm{}^2k_F^2}{2m}}{\displaystyle \underset{i=1}{\overset{3}{}}}{\displaystyle \frac{A_i}{8\pi (k_FL_i)^2}}\times `$ (39)
$`{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{[2nk_FL_i\mathrm{cos}(2nk_FL_i)\mathrm{sin}(2nk_FL_i)]}{n^3\mathrm{sinh}^2(n\kappa _i)}},`$ (40)
$`\rho _0=\rho _{Weyl}+{\displaystyle \frac{1}{L^3}}{\displaystyle \frac{1}{4\pi }}{\displaystyle \underset{i=1}{\overset{3}{}}}A_i{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\mathrm{sin}(2nk_FL_i)}{n\mathrm{sinh}^2(n\kappa _i)}},`$ (41)
where $`A_1=6,A_2=12,A_3=8`$ respectively. The shell energy for the spherical phase is shown in Fig 1c. In this case the anti–filling factor is given by $`u={\displaystyle \frac{4}{3}}{\displaystyle \frac{\pi R^3}{L^3}}`$. A stronger interference pattern due to the orbits $`L_2`$ and $`L_3`$ can be seen. The amplitude of the shell effects is also lower due to the greater instability of the orbit on one hand, and due to the smaller relative volume of the scatterers on the other hand.
In a similar manner one can obtain the interaction energy between two isolated bubbles at large separations ($`a=L2RR`$)
$$E_{}\frac{\mathrm{}^2k_FR^2}{8\pi m}\frac{\mathrm{cos}(2k_Fa)}{a^3}.$$
(42)
When compared with the interaction (5) one observes a similar behaviour, even though now the two “impurities” are large $`k_FR1`$. It can be shown however that if one computes instead the same energy for fixed chemical potential, instead of particle number as was done here, the bubble–bubble interaction will decay inversely proportional to the square of the separation . In a recent paper the Casimir energy for similar arrangements has been calculated using the semiclassical approximation. In the case of Casimir energy the situation is somewhat simple, since instead of two independent dimensionless parameters, $`k_FR`$ and $`k_FL`$, only one dimensionless parameter exists, $`R/L`$. Thus the Casimir energy for two spheres has naturally the form
$$E_{}^{Cas}=\frac{\mathrm{}c}{L}F\left(\frac{R}{L}\right),$$
(43)
with an unknown function $`F(x)`$. A similar, but much stronger result can be obtained for the critical Casimir energy , where one can show that the theory is conformal invariant. The authors of Refs. provide also a very compelling argument why the semiclassical approximation should be particularly accurate for the calculation of the Casimir energy in case of ideal metallic boundaries and they show that using only the single periodic orbit the Casimir energy for two spheres is given by
$$E_{}^{Cas}=\frac{\pi ^3\mathrm{}cR}{720L^2},$$
(44)
and dismiss this result as being valid for large separations, since it contradicts their expectations that it should agree with the Casimir–Polder interaction
$$E_{}^{CP}\frac{\mathrm{}cR^6}{L^7}.$$
(45)
The authors of Ref. argue that the contributions arising from the diffractive paths discussed in Refs. , should eventually lead to additional contributions, which will cancel exactly this longer range interaction and in the end, the authors hope that the Casimir–Polder result will be retrieved. The difference between these two results for the Casimir energy is very similar to the difference between the interaction (5) between two point–like impurities ($`k_FR1`$) and the interaction between two “fat” bubbles ($`k_FR1`$) . In the case of “fat” bubbles, the contribution of diffractive orbits are exponentially small ($`\mathrm{exp}(\alpha k_FR)`$, where $`\alpha `$ is of order unity) , as one would naturally expect in the case when rays are a very good approximation to the wave phenomena. The resolution of this apparent conundrum lies in resolving the clear clash of limits. When the size of the scatterer $`R`$ decreases the contribution of the diffractive paths (creeping orbits) increases and an increasingly larger number of them contribute significantly to the scattering and thus to the propagator. In the limit $`k_FR0`$ the standard geometric orbit approach has to be modified, see Ref. and our discussion around Eqs. (7–10). It is notable that in the case of the critical Casimir effect, even longer range interactions ($`1/a^{1+ϵ}`$, with very small $`ϵ`$) between two spheres are possible .
The structure of the shell energies shown in the Fig. 1 indicates the existence of the optimal void sizes (with respect to the shell effects) for a given outside nucleon density. Note that for all phases and for $`\rho _0>0.05fm^3`$ the shell energy exhibits a remarkable softness toward adding additional neutrons to the system (the “valleys” and “ridges” are almost horizontal in the Fig. 1). Hence one can conclude that once the size of the voids have been determined by minimization of the total energy of the system, an increase in the number of neutrons outside the voids will not affect much the shell energy of the system. However, the surface energy will be affected.
In the Fig. 2 we show the shell energies as a function of $`\rho _0`$ for the optimal filling factors and nuclear radii determined in Ref. . One can see that the amplitudes of the shell energies in the region $`\rho _00.040.07`$ are of the order of $`10keV/fm^3`$, $`3keV/fm^3`$ and $`0.05keV/fm^3`$ for plate–like, rod–like and bubble–like phases, respectively. There are usually one or two shallow shell energy minima for the density range $`\rho _0>0.03fm^3`$. The minima are more pronounced in the case of spherical bubble–like phase mainly due to the stronger interference effects caused by longer orbits.
Once a phase is formed there is a positional order maintained by the Coulomb repulsion between spherical nuclei, rods or slabs . Although the Coulomb energy is a smooth function of the void displacement , the shell energy is not. Since several different orbits contribute to the shell effects (except for slab–like phase) the displacement of a single bubble–like or rod–like void from its equilibrium position in the lattice will give rise to the interference effects. The interference pattern will depend on the type lattice. For the simple cubic and simple square lattices for spherical nuclei and rods, respectively, we show in the Fig. 3 the changes in the energy due to such “defects”. For the plate–like system there is only one direction of displacement (we do not consider the shear mode) denoted by $`x`$ perpendicular to the slab (Fig. 3a). Since the rod–like phase is a two-dimensional system, in the Fig 3b we have shown the shell energy as a function of two perpendicular displacements $`x`$ and $`y`$. They are perpendicular to the rods and point in the direction of the nearest neighbor. The same axes have been chosen for the spherical system although it will not exhaust all possible directions in the system. The behavior of the shell energy in this case is shown in Fig 3c.
The structure of the shell energy surface as a function of a displacement depends on the lengths of the shortest periodic orbits. Except for the trivial plate–like phase, in the rod–like and spherical phase there exist directions into which is easier to locally deform the lattice.
In Fig. 4 we show the pattern of the energy changes induced by deforming the rod–like lattice. We considered only volume conserving deformations. The square lattice was stretched by a factor $`\alpha `$ in the $`x`$–direction, by a factor $`\beta `$ in the $`y`$–direction and also the angle between the two axes has been changed to $`\gamma `$. In order to preserve the volume all these three parameters should satisfy the condition
$$\alpha \beta \mathrm{sin}\gamma =1.$$
(46)
The case $`\alpha =\beta `$ and $`\gamma =\pi /3`$ correspond to a perfect triangular lattice.
Increasing the temperature will weaken the shell effects. At sufficiently high temperatures the nuclear lattice will disappear. At smaller temperatures however, when the lattice can be regarded as frozen, the rise of the temperature will affect mainly shell effects in the neutron gas. In order to wash out completely the shell effects the temperature $`T`$ should be of the order of half of the distance between shells The spacing between two consecutive shells is determined by the length of the shortest orbit, $`a=L2R`$. Thus the energy distance between shells can be determined from the requirement: $`2ka=2\pi `$ and is given by the expression:
$$\mathrm{\Delta }E=\frac{\mathrm{}^2\pi ^2}{2m(L2R)^2}.$$
(47)
For the optimal filling factors and lattice constants of various phases obtained in Ref. one obtains the following estimates for the critical temperature:
$`T_c`$ $``$ $`32MeV\text{ for plate–like system},`$ (48)
$`T_c`$ $``$ $`19MeV\text{ for rod–like system},`$ (49)
$`T_c`$ $``$ $`12MeV\text{ for spherical system}.`$ (50)
An accurate description of the shell effects as a function of the temperature can be obtained using the temperature averaged level density :
$$g_{shell}(\epsilon ,T)=\underset{p.o.}{}\frac{A_i\tau _{i,n}(T)}{\mathrm{sinh}\tau _{i,n}(T)}g_{shell}(\epsilon ,L_i),$$
(51)
where the sum is taken over all periodic orbits including the number of repetitions $`n`$ of the orbit and
$$\tau _{i,n}=\frac{2\pi TmnL_i}{\mathrm{}^2k}.$$
(52)
Consequently the oscillating part of the free energy density is given by the formula:
$$\frac{F_{shell}}{L^3}=\frac{1}{L^3}_0^\mu (\epsilon \mu )g_{shell}(\epsilon ,T)𝑑\epsilon .$$
(53)
The estimates for different phases are shown in Fig. 5. For simplicity we have retained in these calculations the value of the Fermi momentum $`k_F`$ equal to its zero temperature limit, and therefore the neutron matter density is not temperature independent in these figures. One can see that thermal effects will wash out the shell correction energy at temperatures of the order of 10 MeV and higher.
Now at the end of this analysis we suspect that there are a lot of other effects, which might be relevant. We did not consider periodic orbits bouncing between three or more objects. An orbit bouncing between two bodies leads to a pairwise interaction. Orbits bouncing between three or more bodies would lead to genuine many body interactions. We have also considered only perfectly smooth objects. If one allows for some degree of corrugation of these surfaces, many more periodic orbits are likely to appear and that would lead to even more complicated interactions and more complicated interference patterns. The fact that corrugation can influence in a significant, perhaps major way, the Casimir energy, has already been predicted and measured experimentally . The long range character of the interaction together with its oscillatory nature could very easily be at the origin of disorder, even at zero temperature. At finite temperature disorder is more likely to occur, due to entropic effects . We did not consider here the role of pairing, which we expect however to lead to a certain flattening of the shell effects , which, however, should not be interpreted as disappearance of shell effects. Especially at subnuclear densities neutron pairing should be rather strong . A completely different type of softness of these structures has been argued in Ref. , according to which the mantles of neutron stars resembles more liquid crystals than solids.
In the paper we have studied the shell effects in the neutron medium filled by different nuclear phases. To our knowledge this is the first approach which considers specifically the shell effects in the outside neutron gas and we aimed at discussing its basic features. Even though in principle Hartree–Fock calculations include in principle such effects already, the calculations performed so far were too narrow in scope and did not address this issue specifically. Using semiclassical methods, we have analyzed the structure of the shell energy as a function of the density, filling factor, lattice distortions and temperature. We expect that our result overestimate somewhat the amplitude of the shell effects. However, the emerging qualitative overall picture should remain valid and further microscopic studies are highly desirable. The main lesson one should remember from this work is that the amplitude of the shell energy effects is comparable with the energy differences between various phases determined in simpler liquid drop type models. The magnitude of the quantum corrections to the ground state energy of the inhomogeneous neutron matter we have found is significantly larger than that determined in Ref.. The analysis of Ref. was limited however to the motion of nucleons inside nuclei embedded in a lower density neutron gas.
Our results suggest that the inhomogeneous phase has perhaps an extremely complicated structure, maybe even completely disordered, with several types of shapes present at the same time.
The DOE financial support is gratefully acknowledged. This research was supported in part by the Polish Committee for Scientific Research (KBN) under Contract No. 2 P03B 040 14. AB thanks N.D. Whelan, O. Agam, T. Guhr and S.C. Creagh for discussions and correspondence concerning various aspects of the semiclassical approximation. PM thanks the Nuclear Theory Group for hosting his visit in Seattle. AB thanks the members of the Nuclear Theory Group at the Institute of Theoretical Physics in Warsaw for their hospitality. And last but not least, AB thanks W.A. Weidenmüller for being such a gracious host. |
warning/0002/quant-ph0002025.html | ar5iv | text | # REFERENCES
First experimental test of Bell inequalities performed
using a non-maximally entangled state
M.Genovese <sup>*</sup><sup>*</sup>* genovese@ien.it. Tel. 39 011 3919234, fax 39 011 3919259,G.Brida, C.Novero
Istituto Elettrotecnico Nazionale Galileo Ferraris, Str. delle Cacce
91,I-10135 Torino
E. Predazzi
Dip. Fisica Teorica Univ. Torino e INFN, via P. Giuria 1, I-10125 Torino
Abstract
We describe the realisation of a new test of Bell inequalities using a new scheme obtained by the superposition of type I parametric down conversion produced in two different non-linear crystals pumped by the same laser, but with different polarisations. This experiment is the first test of Bell inequalities using a non-maximally entangled state and thus represents an important step in the direction of eliminating the detection loophole.
The idea that Quantum Mechanics (QM) could be an incomplete theory, representing a statistical approximation of a complete deterministic theory (where observable values are fixed by some hidden variable) appeared already in 1935 thank to the celebrate Einstein-Podolsky-Rosen paper .
Historically, the quest for hidden variable theories stopped when Von Neumann published a theorem asserting the impossibility of constructing a hidden variable theory reproducing all the results of QM. For long time, the prestige of Von Neumann led to an acritical acceptation of this theorem, but it was then discovered that one of his hypotheses was too restrictive so that the program of a hidden variable theory was still possible.
The next fundamental progress in discussing possible extensions of QM was the discovery of Bell that any realistic Local Hidden Variable LHV theory must satisfy certain inequalities which can be violated in QM leading in principle to a possible experimental test of the validity of QM as compared to LHV.
Since then, many interesting experiments have been devoted to a test of Bell inequalities , leading to a substantial agreement with quantum mechanics and disfavouring LHV theories, but, so far, no experiment has yet been able to exclude definitively such theories. In fact, so far, one has always been forced to introduce a further additional hypothesis , due to the low total detection efficiency, stating that the observed sample of particle pairs is a faithful subsample of the whole. This problem is known as detection or efficiency loophole. The research for new experimental configurations able to overcome the detection loophole is of course of the greatest interest.
In the 90’s big progresses in this direction have been obtained by using parametric down conversion (PDC) processes.
This technique has been largely employed for generating ”entangled” photon pairs, i.e. pairs of photons described by a common wave function which cannot be factorized into the product of two distinct wave functions pertaining to separated photons.
The generation of entangled states by parametric down conversion (PDC) has replaced other techniques, such as the radiative decay of excited atomic states, as it was in the celebrated experiment of A. Aspect et al. , for it overcomes some former limitations. In particular, having angular correlations better than 1 mrad, it overcomes the poor angular correlation of atomic cascade photons, that was at the origin of the small total efficiency of this type of experiments in which one is forced to select a small subsample of the produced photons, leading inevitably to the detection loophole.
The first experiments using this technique, where performed with type I PDC, which gives phase and momentum entanglement and can be used for a test of Bell inequalities using two spatially separated interferometers , as realised by . The use of beam splitters, however, strongly reduces the total quantum efficiency.
In alternative, one can generate a polarisation entangled state . It appears, however, that the creation of couples of photons entangled from the point of view of polarisation, which is by far the most diffuse case due to the easy experimental implementation, still suffers severe limitations, as it was pointed out recently in the literature . The essence of the problem is that in generating this state, half of the initial photon flux is lost (in most of the used configurations), and one is, of necessity, led to assume that the photon’s population actually involved in the experiment is a faithful sample of the original one, without eliminating the efficiency loophole.
Recently, an experiment where a polarisation entangled state is directly generated, has been realised using Type II PDC . This scheme has permitted, at the price of delicate compensations for having identical arrival time of the ordinary and extraordinary photon, a much higher total efficiency than the previous ones, which is, however, still far from the required value of $`0.81`$. Also, some recent experiments studying equalities among correlations functions rather than Bell inequalities are far from solving these problems . A large interest remains therefore for new experiments increasing total quantum efficiency in order to reduce and finally overcome the efficiency loophole.
Some years ago, a very important theoretical step in this direction has been performed recognising that, while for maximally entangled pairs a total efficiency larger than to 0.81 is required to obtain an efficiency-loophole free experiment, for non maximally entangled pairs this limit is reduced to 0.67 (in the case of no background). However, it must be noticed that, for non-maximally entangled states, the largest discrepancy between quantum mechanics and local hidden variable theories is reduced: thus a compromise between a lower total efficiency and a still sufficiently large value of this difference will be necessary when realising of an experiment addressed to overcome the detection loophole.
Considering a polarization entangled state of photons of the form
$$|\psi =\frac{|H|H+f|V|V}{\sqrt{(1+|f|^2)}}$$
(1)
where $`H`$ and $`V`$ indicate horizontal and vertical polarisations respectively, the parameter $`f`$ describes how much the state 1 differs from a maximally entangled one.
The region corresponding to the one where the detection loophole is eliminated in the plane $`f`$ and $`\eta `$ (total detection efficiency) is shown in fig.1.
Let us now describe more in detail our experimental set-up. It derives from developing a proposal made by Hardy and is based on the creation of a polarisation (non maximally-) entangled states of the form 1 via the superposition of the spontaneous fluorescence emitted by two non-linear crystals driven by the same pumping laser. The crystals are put in cascade along the propagation direction of the pumping laser and the superposition is obtained by using an appropriate optics.
More in details (see fig. 2), two crystals of $`LiIO_3`$ (10x10x10 mm) are placed along the pump laser propagation, 250 mm apart, a distance smaller than the coherence length of the pumping laser. This guarantees indistinguishibility in the creation of a couple of photons in the first or in the second crystal. A couple of planoconvex lenses of 120 mm focal length centred in between, focalises the spontaneous emission from the first crystal into the second one maintaining exactly, in principle, the angular spread. A hole of 4 mm diameter is drilled into the centre of the lenses in order to allow transmission of the pump radiation without absorption and, even more important, without adding stray-light, because of fluorescence and diffusion of the UV radiation. This configuration, which realises what is known as an ”optical condenser”, was chosen among others, using an optical simulation program, as a compromise between minimisation of aberrations (mainly spherical and chromatic) and losses due to the number of optical components. The pumping beam at the exit of the first crystal is displaced from its input direction by birefringence: the small quartz plate (5 x5 x5 mm) in front of the first lens of the condensers compensates this displacement, so that the input conditions are prepared to be the same for the two crystals, apart from alignment errors. Finally, a half-wavelength plate immediately after the condenser rotates the polarisation of the Argon beam and excites in the second crystal a spontaneous emission cross-polarised with respect to the first one. With a phase matching angle of $`51^o`$, the spontaneous emissions at 633 and 789 nm (which are the wave lengths to be used for the test) are emitted at $`3.5^o`$ and $`4^o`$ respectively. The dimensions and positions of both plates are carefully chosen in order that they do not intersect this two conjugated emissions.
We have used as photo-detectors two avalanche photodiodes with active quenching (EG&G SPCM-AQ) with a sensitive area of 0.025 $`mm^2`$ and dark count below 50 counts/s. The PDC signal was coupled to an optical fiber (carrying the light to the detectors) by means of a microscope objective with magnification 20. The quantum efficiency, included the fiber coupling, has been measured to be $`0.535\pm 0.008`$ at 633 nm .
The output signals from the detectors are pulses of Transistor-Transistor-Logic (TTL) like amplitude levels which are routed to a two channel counter, in order to have the number of events on single channel, and to a Time to Amplitude Converter (TAC) circuit, followed by a single channel analyzer, for selecting and counting coincidence events.
A very interesting degree of freedom of this configuration is given by the fact that by tuning the pump intensity between the two crystals, one can easily tune the value of $`f`$, which determines how far from a maximally entangled state ($`f=1`$) the produced state is. This is a fundamental property, which permits to select the most appropriate state for the experiment.
The main problem of this configuration is the alignment, which is of the utmost relevance for a high visibility. The solution of this problem lies in a technique, that had been already applied in our laboratory for metrological studies, namely the use of an optical amplifier scheme, where a solid state laser is injected into the crystals together with the pumping laser, an argon laser at 351 nm wavelength (see fig.1). If the angle of injection is selected appropriately, a stimulated emission along the correlated direction appears, allowing an easy identification of the two correlated directions. Then, stopping the stimulated emission of the first crystal, and rotating the polarisation of the diode laser one obtaines the stimulated emission in the second crystal and can check the superposition with the former.
We think that the proposed scheme will lead to a further step towards a conclusive experimental test of non-locality in quantum mechanics. The analysis of the experiments realised up to now shows in fact that visibility of the wanted effect (essentially visibility of interference fringes) and overall quantum detection efficiency are the main parameters in such experiments. One first advantage of the proposed configuration with respect to most of the previous experimental set-ups is that all the entangled pairs are selected (and not just $`<50\%`$ as with beams splitters); furthermore, it does not require delicate compensations for the optical paths of the ordinary and extraordinary rays emerging from the crystal.
For this time being, the results which we are going to present are still far from a definite solution of the detection loophole; nevertheless, being the first test of Bell inequalities using a non-maximally entangled state, they represents an important step in this direction. Furthermore, this configuration allows to use any pair of correlated frequencies and not only the degenerate ones. We have thus realised this test using for a first time two different wave lengths (at $`633`$ and $`789`$ nm).
An experiment which presents analogies with our, has been realised recently in ref. . The main difference between the two experiments is that in the two crystals are very thin and in contact with orthogonal optical axes: this allows a ”partial” superposition of the two emissions with opposite polarisation. This overlapping is mainly due to the finite dimension of the pump laser beam, which reflects into a finite width of each wavelength emission. A much better superposition can be obtained with the present scheme, by fine tuning the crystals’ and optics’ positions and using the parametric amplifier trick. Furthermore, in the experiment of ref. the value of $`f`$ is in principle tunable by rotating the polarisation of the pump laser; however, this reduces the power of the pump producing PDC already in the first crystal, while in our case the whole pump power can always be used in the first crystal, tuning the PDC produced in the second one.
As a first check of our apparatus, we have measured the interference fringes, varying the setting of one of the polarisers, while leaving the other fixed. We have found a high visibility, $`V=0.973\pm 0.038`$.
Our results are summarised by the value obtained for the Clauser-Horne sum,
$$CH=N(\theta _1,\theta _2)N(\theta _1,\theta _2^{})+N(\theta _1^{},\theta _2)+N(\theta _1^{},\theta _2^{})N(\theta _1^{},\mathrm{})N(\mathrm{},\theta _2)$$
(2)
which is strictly negative for local realistic theory. In (2), $`N(\theta _1,\theta _2)`$ is the number of coincidences between channels 1 and 2 when the two polarisers are rotated to an angle $`\theta _1`$ and $`\theta _2`$ respectively ($`\mathrm{}`$ denotes the absence of selection of polarisation for that channel)
On the other hand, quantum mechanics predictions for $`CH`$ can be larger than zero: for a maximally entangled state the largest value is obtained for $`\theta _1=67^o.5`$ , $`\theta _2=45^o`$, $`\theta _1^{}=22^o.5`$ , $`\theta _2^{}=0^o`$ and corresponds to a ratio
$$R=[N(\theta _1,\theta _2)N(\theta _1,\theta _2^{})+N(\theta _1^{},\theta _2)+N(\theta _1^{},\theta _2^{})]/[N(\theta _1^{},\mathrm{})+N(\mathrm{},\theta _2)]$$
(3)
equal to 1.207.
For non-maximally entangled states the angles for which CH is maximal are somehow different and the maximum is reduced to a smaller value. The angles corresponding to the maximum can be evaluated maximising Eq. 2 with
$`\begin{array}{c}N[\theta _1,\theta _2]=[ϵ_1^{||}ϵ_2^{||}(Sin[\theta _1]^2Sin[\theta _2]^2)+\hfill \\ ϵ_1^{}ϵ_2^{}(Cos[\theta _1]^2Cos[\theta _2]^2)\hfill \\ (ϵ_1^{}ϵ_2^{||}Sin[\theta _1]^2Cos[\theta _2]^2+ϵ_1^{||}ϵ_2^{}Cos[\theta _1]^2Sin[\theta _2]^2)\hfill \\ +|f|^2(ϵ_1^{}ϵ_2^{}(Sin[\theta _1]^2Sin[\theta _2]^2)+ϵ_1^{||}ϵ_2^{||}(Cos[\theta _1]^2Cos[\theta _2]^2)+\hfill \\ (ϵ_1^{||}ϵ_2^{}Sin[\theta _1]^2Cos[\theta _2]^2+\hfill \\ ϵ_1^{}ϵ_2^{||}Cos[\theta _1]^2Sin[\theta _2]^2)\hfill \\ +(f+f^{})((ϵ_1^{||}ϵ_2^{||}+ϵ_1^{}ϵ_2^{}ϵ_1^{||}ϵ_2^{}ϵ_1^{}ϵ_2^{||})(Sin[\theta _1]Sin[\theta _2]Cos[\theta _1]Cos[\theta _2])]/(1+|f|^2)\hfill \end{array}.`$ (11)
where (for the case of non-ideal polariser) $`ϵ_i^{||}`$ and $`ϵ_i^{}`$ correspond to the transmission when the polariser (on the branch $`i`$) axis is aligned or normal to the polarisation axis respectively.
In our case, we have generated a state with $`f0.4`$ which corresponds, for $`\theta _1=72^o.24`$ , $`\theta _2=45^o`$, $`\theta _1^{}=17^o.76`$ and $`\theta _2^{}=0^o`$, to $`R=1.16`$.
Our experimental result is $`CH=512\pm 135`$ coincidences per second, which is almost four standard deviations different from zero and compatible with the theoretical value predicted by quantum mechanics. In terms of the ratio (3), our result is $`1.082\pm 0.031`$.
For the sake of comparison, one can consider the value obtained with the angles which optimise Bell inequalities violation for a maximally entangled state (given before Eq. 3). The result is $`CH=92\pm 89`$, which, as expected, shows a smaller violation than the value obtained with the correct angles setting.
In summary, this is the first measurement of the violation of the Clauser-Horne inequality (or other Bell inequalities) using a non-maximally entangled state and thus represents and interesting result as a first step in the direction of eliminating the detection loophole. Further developments in this sense are the purpose of this collaboration.
Acknowledgements We would like to acknowledge the support of the Italian Space Agency under contract LONO 500172 and of MURST via special programs ”giovani ricercatori”, Dip. Fisica Teorica Univ. Torino. |
warning/0002/cond-mat0002425.html | ar5iv | text | # Electronic transport through a local state linearly coupled to phonon modes
## I Introduction
Recently, we addressed the problem of the correspondence between two different descriptions of the phenomenon of inelastic electronic transitions between two distinct set of continuum electron states, through an intermediate resonant state coupled to an oscillator . Even though the problem of interest there was the influence of vibration damping on the spectroscopy of adsorbates with the scanning tunneling microscope, it was shown that the formalism adopted, taken over from a description of resonant tunneling in semiconductor structures with the necessary modifications, was in fact reducible to a simpler resolvent method, in the no-damping limit, but with the restriction of just one vibration coordinate. Our intent at that point was to make a connection with a previous paper on STM, but an obvious possibility then was the eventual extension of that correspondence between the two approaches, from just one vibration coordinate to a Brioullin zone of modes. The present paper addresses this issue, showing that indeed a simple resolvent model, without appealing to many-body field theory, seems to be able to describe inelastic tunneling in small semiconductor structures, at least with the restrictions of a $`0`$ K temperature and an Einstein band of optical phonons. Even so, we believe it is of some interest to present this analysis, as it may provide a simpler means of assessing the essential physics of the situation and also as a starting point for further useful extensions.
On the other hand, even though this model can be applied to many real distinct situations in physics and in chemistry, as pointed out by Gadzuk when drawing similar comparisons, we will be referring specifically to inelastic resonant electronic transport in semiconductors (for a recent exposition based on nonequilibrium Green functions, see Ref. ), so that our results can be compared directly with the more general, and more powerful, many-body approach used in Ref. , which employs a two-particle many-body Green’s function and which was subsequently simplified to a singe-particle many-body Green’s function treatment . Other methods have been published, such as those of Ref. , for example. But with the present, simpler, method even the inclusion of more than one intermediate, resonant state should not introduce major difficulties. The essentials of this method were used many years ago by Domcke and Cederbaum , to describe inelastic electron scattering by molecules in the gas phase.
## II Resonance Green’s Phonon Operator
As usual, the total spinless Hamiltonian is a sum of three parts, the electron Hamiltonian, the phonon Hamiltonian and the electron-phonon interaction, $`H=H_{el}+H_{ph}+H_{int}`$:
$`H_{el}=`$ $`ϵ_a^0c_a^{}c_a+{\displaystyle \underset{l}{}}ϵ_lc_l^{}c_l+{\displaystyle \underset{k}{}}ϵ_kc_k^{}c_k`$ (1)
$`+`$ $`{\displaystyle \underset{l}{}}(V_{al}c_a^{}c_l+V_{la}c_l^{}c_a)`$ (2)
$`+`$ $`{\displaystyle \underset{k}{}}(V_{ak}c_a^{}c_k+V_{ka}c_k^{}c_a),`$ (3)
$$H_{ph}=\underset{q}{}\mathrm{}\omega _q(b_q^{}b_q+1/2),$$
(4)
and
$$H_{int}=c_a^{}c_a\underset{q}{}M_q(b_q^{}+b_q).$$
(5)
In $`H_{el}`$ the first term describes the localized resonant electron state, the second and third terms describe the electronic states in the left and right (isolated) leads, and the last two terms represent the hybridization between the localized electronic resonance and the two leads. $`H_{ph}`$ describes the harmonic phonon states and $`H_{int}`$ is the linear electron-phonon interaction, whose strength is given by the coupling constant $`M_q`$, acting only when an electron occupies the intermediate, resonant state.
We proceed to compute the resonance Green’s function from the corresponding resolvent operator $`G=(ϵH)^1`$, using the equations-of-motion method, as in the original Anderson description of localized magnetic states associated with impurities in metals . In doing this, we assume the validity of the Born-Oppenheimer approximation and include the interaction Hamiltonian in the resonant electron state energy, writing it as a phonon operator
$$ϵ_a=ϵ_a(b_q^{()})=ϵ_a^0+\underset{q}{}M_q(b_q^{}+b_q),$$
(6)
thereby rendering the electron Hamiltonian $`H_{el}`$ adiabatically dependent on the state of the lattice oscillators.
From the equations of motion for the operators
$$i\mathrm{}\frac{c_\alpha }{t}=[c_\alpha ,H]$$
(7)
we obtain, taking into account the commutation of electron and phonon operators under the adiabatic hypothesis, the set of equations
$`i\mathrm{}{\displaystyle \frac{c_a}{t}}=`$ $`ϵ_ac_a(t)+{\displaystyle \underset{k}{}}V_{ak}c_k(t)+{\displaystyle \underset{l}{}}V_{al}c_l(t)`$ (8)
$`i\mathrm{}{\displaystyle \frac{c_k}{t}}=`$ $`ϵ_kc_k(t)+V_{ka}c_a(t)`$ (9)
$`i\mathrm{}{\displaystyle \frac{c_l}{t}}=`$ $`ϵ_lc_l(t)+V_{la}c_a(t).`$ (10)
Fourier transforming to the energy domain, we get the corresponding algebraic system
$`ϵc_a(ϵ)=`$ $`ϵ_ac_a(ϵ)+{\displaystyle \underset{k}{}}V_{ak}c_k(ϵ)+{\displaystyle \underset{l}{}}V_{al}c_l(ϵ)`$ (11)
$`ϵc_k(ϵ)=`$ $`ϵ_kc_k(ϵ)+V_{ka}c_a(ϵ),k,`$ (12)
$`ϵc_l(ϵ)=`$ $`ϵ_lc_l(ϵ)+V_{la}c_a(ϵ),l,`$ (13)
which can be rewritten in matrix form
$`\left[ϵ\mathrm{𝟏}\left(\begin{array}{c}ϵ_a\\ 𝐕_{\mathrm{𝐤𝐚}}\\ 𝐕_{\mathrm{𝐥𝐚}}\end{array}\begin{array}{c}𝐕_{\mathrm{𝐚𝐤}}^𝐓\\ ϵ_𝐤\\ \mathrm{𝟎}\end{array}\begin{array}{c}𝐕_{\mathrm{𝐚𝐥}}^𝐓\\ \mathrm{𝟎}\\ ϵ_𝐥\end{array}\right)\right]\left(\begin{array}{c}c_a\\ 𝐜_𝐤\\ 𝐜_𝐥\end{array}\right)=\left(\begin{array}{c}0\\ \mathrm{𝟎}\\ \mathrm{𝟎}\end{array}\right),`$ (29)
where $`𝐕_{\mathrm{𝐤𝐚}}`$ is a column matrix, $`𝐕_{\mathrm{𝐚𝐤}}^𝐓`$ the row matrix transpose of $`𝐕_{\mathrm{𝐚𝐤}}`$, $`ϵ_𝐤`$ and $`ϵ_𝐥`$ are diagonal matrices with elements $`ϵ_{k_1},ϵ_{k_2},\mathrm{}`$ and $`ϵ_{l_1},ϵ_{l_2},\mathrm{}`$, respectively, $`\mathrm{𝟏}`$ is the identity matrix of the appropriate dimension and $`\mathrm{𝟎}`$ are null matrices of the appropriate dimension also.
This matrix equality can be put in the form
$$\left(ϵ\mathrm{𝟏}𝐅_{\mathrm{𝐞𝐥}}\right)𝐜=\mathrm{𝟎},$$
(30)
writing the previous matrices more compactly, in an obvious notation. Introducing now another diagonal phonon operator matrix
$$𝐅_{\mathrm{𝐩𝐡}}=\left[\mathrm{}\omega _q(b_q^{}b_q+1/2)\right]\mathrm{𝟏},$$
(31)
we define an electron-phonon resolvent operator
$`𝐆^{\mathrm{𝐭𝐨𝐭𝐚𝐥}}(\zeta )`$ $`=`$ $`\left[(\zeta +i\eta )\mathrm{𝟏}(𝐅_{\mathrm{𝐞𝐥}}+𝐅_{\mathrm{𝐩𝐡}})\right]^1`$ (32)
$`=`$ $`\left[(\zeta +i\eta )\mathrm{𝟏}𝐅^{\mathrm{𝐭𝐨𝐭𝐚𝐥}}\right]^1,`$ (33)
$`\zeta `$ being the total energy parameter and $`\eta `$ a positive infinitesimal. Then $`\left[(\zeta +i\eta )\mathrm{𝟏}𝐅^{\mathrm{𝐭𝐨𝐭𝐚𝐥}}\right]𝐆^{\mathrm{𝐭𝐨𝐭𝐚𝐥}}(\zeta )=\mathrm{𝟏}`$ which, expanded, means
$`\left(\begin{array}{c}(\zeta +i\eta )ϵ_aH_{ph}\mathrm{𝟏}\\ 𝐕_{\mathrm{𝐤𝐚}}\\ 𝐕_{\mathrm{𝐥𝐚}}\end{array}\begin{array}{c}𝐕_{\mathrm{𝐚𝐤}}^𝐓\\ (\zeta +i\eta )\mathrm{𝟏}ϵ_𝐤H_{ph}\mathrm{𝟏}\\ \mathrm{𝟎}\end{array}\begin{array}{c}𝐕_{\mathrm{𝐚𝐥}}^𝐓\\ \mathrm{𝟎}\\ (\zeta +i\eta )\mathrm{𝟏}ϵ_𝐥H_{ph}\mathrm{𝟏}\end{array}\right)`$ (43)
$`\left(\begin{array}{c}G_{aa}\\ 𝐆_{\mathrm{𝐤𝐚}}\\ 𝐆_{\mathrm{𝐥𝐚}}\end{array}\begin{array}{c}𝐆_{\mathrm{𝐚𝐤}}\\ 𝐆_{\mathrm{𝐤𝐤}}\\ 𝐆_{\mathrm{𝐥𝐤}}\end{array}\begin{array}{c}𝐆_{\mathrm{𝐚𝐥}}\\ 𝐆_{\mathrm{𝐤𝐥}}\\ 𝐆_{\mathrm{𝐥𝐥}}\end{array}\right)=\left(\begin{array}{c}1\\ \mathrm{𝟎}\\ \mathrm{𝟎}\end{array}\begin{array}{c}\mathrm{𝟎}\\ \mathrm{𝟏}\\ \mathrm{𝟎}\end{array}\begin{array}{c}\mathrm{𝟎}\\ \mathrm{𝟎}\\ \mathrm{𝟏}\end{array}\right).`$ (62)
Introducing a common index $`m`$ denoting both $`k`$ and $`l`$, we have the system of equations
$`[(\zeta +i\eta )ϵ_aH_{ph}]G_{aa}{\displaystyle \underset{m=k,l}{}}V_{am}G_{ma}=1`$ (63)
$`V_{ma}G_{aa}+[(\zeta +i\eta )ϵ_mH_{ph}]G_{ma}=0.`$ (64)
In the second equalitiy we solve for $`G_{ma}`$ (where $`m=k,l`$) and substitute it in the first, obtaining the Green’s phonon operator for the localized electronic resonance
$`G_{aa}(\zeta )=`$ $`\{[(\zeta +i\eta )ϵ_aH_{ph}]`$ (65)
$``$ $`{\displaystyle \underset{m=k,l}{}}|V_{am}|^2[(\zeta +i\eta )ϵ_mH_{ph}]^1\}^1`$ (66)
as a function of the total energy parameter, $`\zeta `$. At zero K temperature, the initial energy in the transition process, equal to the total energy, will comprise the incoming electron energy, $`ϵ_{k_i}`$ plus the ground state lattice energy, $`1/2_q\mathrm{}\omega _q`$. With this last quantity fixed, we use as the energy parameter the incoming electron energy and write
$`G_{aa}(ϵ_{k_i})=\{[(ϵ_{k_i}+1/2{\displaystyle \underset{q}{}}\mathrm{}\omega _q+i\eta )ϵ_a^0{\displaystyle \underset{q}{}}M_q(b_q^{}+b_q){\displaystyle \underset{q}{}}\mathrm{}\omega _q(b_q^{}b_q+1/2)]`$ (67)
$`{\displaystyle \underset{m}{}}|V_{am}|^2[ϵ_{k_i}+1/2{\displaystyle \underset{q}{}}\mathrm{}\omega _q+i\eta ϵ_m{\displaystyle \underset{q}{}}\mathrm{}\omega _q(b_q^{}b_q+1/2)]^1\}.`$ (68)
However, since in general $`(x+i\eta )^1=P(x^1)i\pi \delta (x)`$, $`P`$ designating the principal part, we get
$`[ϵ_{k_i}+i\eta ϵ_m{\displaystyle \underset{q}{}}\mathrm{}\omega _qb_q^{}b_q]^1`$ (69)
$`=P[(ϵ_{k_i}ϵ_m{\displaystyle \underset{q}{}}\mathrm{}\omega _qb_q^{}b_q)^1]`$ (70)
$`i\pi \delta (ϵ_{k_i}ϵ_m{\displaystyle \underset{q}{}}\mathrm{}\omega _qb_q^{}b_q)`$ (71)
and defining an electronic shift $`\mathrm{\Delta }`$ and width $`\mathrm{\Gamma }`$ phonon operators by
$`\mathrm{\Delta }(ϵ_{k_i})=P{\displaystyle \underset{m=k,l}{}}\left|V_{am}\right|^2(ϵ_{k_i}ϵ_m{\displaystyle \underset{q}{}}\mathrm{}b_q^{}b_q)^1`$ (72)
$`\mathrm{\Gamma }(ϵ_{k_i})=2\pi {\displaystyle \underset{m=k,l}{}}\left|V_{am}\right|^2\delta (ϵ_{k_i}ϵ_m{\displaystyle \underset{q}{}}\mathrm{}b_q^{}b_q),`$ (73)
we can rewrite the Green’s operator for the resonant state as
$`G_{aa}(ϵ_{k_i})=\{[ϵ_{k_i}ϵ_a^0{\displaystyle \underset{q}{}}M_q(b_q^{}+b_q){\displaystyle \underset{q}{}}\mathrm{}\omega _qb_q^{}b_q]`$ (74)
$`\mathrm{\Delta }(ϵ_{k_i})+{\displaystyle \frac{i}{2}}\mathrm{\Gamma }(ϵ_{k_i})\}^1.`$ (75)
At this point, and for our present purposes, we neglect the operator nature of both $`\mathrm{\Delta }`$ and $`\mathrm{\Gamma }`$ as well as their dependence on $`ϵ`$, obtaining a simplified Green’s operator
$$G_{aa}(ϵ_{k_i})=[ϵ_{k_i}\overline{ϵ}_a\underset{q}{}M_q(b_q^{}+b_q)\underset{q}{}\mathrm{}\omega _qb_q^{}b_q]^1$$
(76)
with a complex resonance energy $`\overline{ϵ}_a=ϵ_a^0+\mathrm{\Delta }i\mathrm{\Gamma }/2`$, corresponding to a lorentzian profile.
## III The Transition Amplitude
Having obtained the resonance Green’s operator, we now proceed to compute the transition amplitude, from the $`T`$ matrix operator, $`T=V+VGV`$, between an initial electron-phonon state $`|k_i;n_q=0,q`$ with zero phonons and a final state $`|l_f;n_1,n_2,\mathrm{},n_N`$, with $`n_q`$ excited phonons in mode $`q`$, in a total of $`N`$ lattice vibration modes. Since we assume no direct transitions between electron $`k`$ states and $`l`$ states, we have
$`l_f|T|k_i=`$ $`l_f|({\displaystyle \underset{l}{}}V_{la}c_l^{}c_a)G_{aa}({\displaystyle \underset{k}{}}V_{ak}c_a^{}c_k)|k_i`$ (77)
$`=`$ $`V_{l_f,a}V_{a,k_i}\left[ϵ_{k_i}\overline{ϵ}_a{\displaystyle \underset{q}{}}M_q(b_q^{}+b_q){\displaystyle \underset{q}{}}\mathrm{}\omega _qb_q^{}b_q\right]^1,`$ (78)
still a phonon operator, to be inserted between the vibrational ground state of the lattice and all possible excited states:
$`l_f;n_1,n_2,\mathrm{},n_N|T|k_i;0_1,0_2,\mathrm{},0_N=n_1,n_2,\mathrm{},n_N|l_f|T|k_i|0_1,0_2,\mathrm{},0_N`$ (79)
$`=V_{l_f,a}V_{a,k_i}n_1,n_2,\mathrm{},n_N|\left[ϵ_{k_i}\overline{ϵ}_a{\displaystyle \underset{q}{}}M_q(b_q^{}+b_q){\displaystyle \underset{q}{}}\mathrm{}\omega _qb_q^{}b_q\right]^1|\mathrm{𝟎},`$ (80)
where we wrote the lattice initial, ground state as $`|\mathrm{𝟎}`$.
Next, we proceed to diagonalize this last denominator by appealing to the operator $`U`$ (as in Refs. and )
$$U=\mathrm{exp}\left[\underset{q}{}\frac{M_q}{\mathrm{}\omega _q}(b_q^{}b_q)\right]$$
(81)
and inserting $`U^1U=\mathrm{𝟏}`$ in the matrix element (3.2) above, after noting its effect on the phonon operators, $`b_q^{}`$ and $`b_q`$,
$`Ub_q^{}U^1=b_q^{}{\displaystyle \frac{M_q}{\mathrm{}\omega _q}}`$ (82)
$`Ub_qU^1=b_q{\displaystyle \frac{M_q}{\mathrm{}\omega _q}},`$ (83)
results that assume symmetry of the phonon bands, specifically that $`M_q=M_q`$ and $`\omega _q=\omega _q`$. We have also employed the well known operator theorem that $`\mathrm{exp}\{A\}B\mathrm{exp}\{A\}=B+[A,B]`$ if $`[A,B]`$ is a c-number. Then, the matrix element in (3.2) can be written
$`n_1,n_2,\mathrm{},n_N|U^1U\left[ϵ_{k_i}\overline{ϵ}_a{\displaystyle \underset{q}{}}M_q(b_q^{}+b_q){\displaystyle \underset{q}{}}\mathrm{}\omega _qb_q^{}b_q\right]^1U^1U|\mathrm{𝟎}`$ (84)
and diagonalizes to
$`n_1,n_2,\mathrm{},n_N|U^1\left[ϵ_{k_i}\overline{ϵ}_a+{\displaystyle \underset{q}{}}{\displaystyle \frac{M_q^2}{\mathrm{}\omega _q}}{\displaystyle \underset{q}{}}\mathrm{}\omega _qb_q^{}b_q\right]^1U|\mathrm{𝟎}.`$ (85)
We proceed to calculate this matrix element. First we note that, since $`M_q=M_q`$ and $`\omega _q=\omega _q`$, we can write
$`U|\mathrm{𝟎}`$ $`=`$ $`e^{_q\frac{M_q}{\mathrm{}\omega _q}(b_q^{}b_q)}|0_1,0_2,\mathrm{},0_N`$ (86)
$`=`$ $`e^{\frac{M_{q_1}}{\mathrm{}\omega _{q_1}}(b_{q_1}^{}b_{q_1})}|0_1e^{\frac{M_{q_2}}{\mathrm{}\omega _{q_2}}(b_{q_2}^{}b_{q_2})}|0_2\mathrm{}`$ (88)
$`\mathrm{}e^{\frac{M_{q_N}}{\mathrm{}\omega _{q_N}}(b_{q_N}^{}b_{q_N})}|0_N`$
from which, taking advantage of the operator rule $`\mathrm{exp}\{A+B\}=\mathrm{exp}\{A\}\mathrm{exp}\{B\}\mathrm{exp}\{1/2[A,B]\}`$, we obtain
$`e^{\frac{M_q}{\mathrm{}\omega _q}(b_q^{}b_q)}|0_q`$ (89)
$`=`$ $`e^{\frac{1}{2}\left(\frac{M_q}{\mathrm{}\omega _q}\right)^2}e^{\left(\frac{M_q}{\mathrm{}\omega _q}\right)b^{}}e^{\left(\frac{M_q}{\mathrm{}\omega _q}\right)b}|0_q`$ (90)
$`=`$ $`e^{\frac{1}{2}\left(\frac{M_q}{\mathrm{}\omega _q}\right)^2}{\displaystyle \underset{m_q=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(M_q/\mathrm{}\omega _q)^{m_q}}{\sqrt{m_q!}}}|m_q,`$ (91)
for each vibration mode $`q`$. We then obtain the product of similar terms for all $`N`$ phonon modes. Operating on that product with the diagonalized denominator phonon operator, we have
$`\left[ϵ_{k_i}\overline{ϵ}_a{\displaystyle \underset{q}{}}\mathrm{}\omega _qb_q^{}b_q\right]^1U|0_1,0_2,\mathrm{},0_N`$ (92)
$`=`$ $`e^{\frac{1}{2}_q\left(\frac{M_q}{\mathrm{}\omega _q}\right)^2}{\displaystyle \underset{m_1=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{m_2=0}{\overset{\mathrm{}}{}}}\mathrm{}{\displaystyle \underset{m_N=0}{\overset{\mathrm{}}{}}}()^{\left(_qm_q\right)}`$ (93)
$`\times `$ $`{\displaystyle \frac{(M_1/\mathrm{}\omega _1)^{m_1}(M_2/\mathrm{}\omega _2)^{m_2}\mathrm{}(M_N/\mathrm{}\omega _N)^{m_N}}{\sqrt{m_1!m_2!\mathrm{}m_N!}}}`$ (94)
$`\times `$ $`\left[ϵ_{k_i}\overline{ϵ}_a+{\displaystyle \underset{q}{}}{\displaystyle \frac{M_q^2}{\mathrm{}\omega _q}}{\displaystyle \underset{q}{}}\mathrm{}\omega _qm_q\right]^1|m_1,m_2,\mathrm{},m_N.`$ (95)
It remains to calculate the action of the operator $`U^1`$ on the intermediate vibrational state ket $`|m_1,m_2,\mathrm{},m_N`$:
$`U^1`$ $`|m_1,m_2,\mathrm{},m_N`$ (96)
$`=`$ $`e^{_q\frac{M_q}{\mathrm{}\omega _q}(b_q^{}b_q)}|m_1,m_2,\mathrm{},m_N,`$ (97)
a product of terms of the type
$`e^{\frac{M_q}{\mathrm{}\omega _q}(b_q^{}b_q)}|m_q=`$ (98)
$`e^{\frac{1}{2}(\frac{M_q}{\mathrm{}\omega _q})^2}{\displaystyle \underset{s_q=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{r_q=0}{\overset{\mathrm{}}{}}}`$ $`()^r{\displaystyle \frac{(M_q/\mathrm{}\omega _q)^{r_q+s_q}}{r_q!s_q!}}{\displaystyle \frac{\sqrt{m_q!(m_qr_q+s_q)!}}{(m_qr_q)!}}|m_qr_q+s_q.`$ (99)
Performing the internal product with the corresponding bra $`n_q|`$, we use the orthogonality condition, $`n_q|m_qr_q+s_q=\delta (n_qm_q+r_qs_q)`$, and get, for the $`q^{th}`$ mode, a contribution (including the pre-factor in Eq.(3.8) above)
$`e^{\frac{1}{2}\left(\frac{M_q}{\mathrm{}\omega _q}\right)^2}{\displaystyle \underset{m_q=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\left(M_q/\mathrm{}\omega _q\right)^{m_q}}{\sqrt{m_q!}}}n_q|U^1|m_q`$ (100)
$`=e^{(\frac{M_q}{\mathrm{}\omega _q})^2}\left({\displaystyle \frac{M_q}{\mathrm{}\omega _q}}\right)^{n_q}(n_q!)^{1/2}`$ (101)
$`\times {\displaystyle \underset{m_q=0}{\overset{\mathrm{}}{}}}()^{m_q}{\displaystyle \underset{r_q=0}{\overset{m_q}{}}}()^{r_q}{\displaystyle \frac{(n_q!)[(M_q/\mathrm{}\omega _q)^2]^{r_q}}{r_q!(m_qr_q)!(n_q+r_qm_q)!}}.`$ (102)
However, this last summation is a generalized Laguerre polynomial :
$`{\displaystyle \underset{r_q=0}{\overset{m_q}{}}}()^{r_q}{\displaystyle \frac{n_q!}{r_q!(m_qr_q)!(n_q+r_qm_q)!}}[(M_q/\mathrm{}\omega _q)^2]^{r_q}`$ (103)
$`=L_{m_q}^{n_qm_q}[(M_q/\mathrm{}\omega _q)^2]`$ (104)
and the whole transition amplitude matrix becomes
$`l_f;n_1,n_2,\mathrm{},n_N|T|k_i;0_1,0_2,\mathrm{},0_N`$ (105)
$`=V_{l_f,a}V_{a,k_i}e^{_q\left(\frac{M_q}{\mathrm{}\omega _q}\right)^2}\left[\left({\displaystyle \frac{M_1}{\mathrm{}\omega _1}}\right)^{n_1}(n_1!)^{1/2}\left({\displaystyle \frac{M_2}{\mathrm{}\omega _2}}\right)^{n_2}(n_2!)^{1/2}\mathrm{}\left({\displaystyle \frac{M_N}{\mathrm{}\omega _N}}\right)^{n_N}(n_N!)^{1/2}\right]`$ (106)
$`\times {\displaystyle \underset{m_1=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{m_2=0}{\overset{\mathrm{}}{}}}\mathrm{}{\displaystyle \underset{m_N=0}{\overset{\mathrm{}}{}}}()^{m_1}L_{m_1}^{n_1m_1}[(M_1/\mathrm{}\omega _1)^2]()^{m_2}L_{m_2}^{n_2m_2}[(M_2/\mathrm{}\omega _2)^2]\mathrm{}()^{m_N}L_{m_N}^{n_Nm_N}[(M_N/\mathrm{}\omega _N)^2]`$ (107)
$`\times \left[ϵ_{k_i}\overline{ϵ}_a+{\displaystyle \underset{q}{}}\left({\displaystyle \frac{M_q^2}{\mathrm{}\omega _q}}\right)(m_1\mathrm{}\omega _1+m_2\mathrm{}\omega _2+\mathrm{}+m_N\mathrm{}\omega _N)\right]^1`$ (108)
It is not necessary to distinguish the two cases, $`n_qm_q`$ and $`m_qn_q`$, if we take into consideration the symmetry properties of the generalized Laguerre polynomials,
$$()^\nu \nu !x^\nu L_\nu ^{\mu \nu }(x)=()^\mu \mu !x^\mu L_\mu ^{\nu \mu }(x)$$
(109)
which, applied to the present case, allows us to write, when $`n_qm_q`$
$`L_{m_q}^{n_qm_q}[(M_q/\mathrm{}\omega _q)^2]`$ $`=`$ $`()^{n_qm_q}\left({\displaystyle \frac{n_q!}{m_q!}}\right)\left[\left({\displaystyle \frac{M_q}{\mathrm{}\omega _q}}\right)^2\right]^{(m_qn_q)}`$ (111)
$`\times L_{n_q}^{m_qn_q}[(M_q/\mathrm{}\omega _q)^2].`$
From now on, we will always write the matrix element as displayed in Eq.(3.14) above. Up to this point what we have is basically the generalization, for $`N`$ vibration modes, of the results for just one vibration coordinate. But the similarities between these and the results of Wingreen et al. in Ref. lead us to believe that the $`N`$-mode situation could somehow be reduced to a single vibration mode inelastic scattering event. This is in fact possible, if we impose the restriction of an Einstein band of phonons, that is to say $`\omega _q=\omega _0`$ for all modes q. In order to perform this reduction from $`N`$ to just one “mode”, we begin by considering only the last two phonon modes in Eq.(3.14) above, with the complete energy denominator:
$`{\displaystyle \underset{m_{N1}=0}{\overset{\mathrm{}}{}}}()^{m_{N1}}L_{m_{N1}}^{n_{N1}m_{N1}}(g_{N1}){\displaystyle \underset{m_N=0}{\overset{\mathrm{}}{}}}()^{m_N}L_{m_N}^{n_Nm_N}(g_N)`$ (112)
$`\left[(zA_{N2})(m_{N1}\mathrm{}\omega _{N1}+m_N\mathrm{}\omega _N)\right]^1`$ (113)
where we introduced the short notations $`z=ϵ_{k_i}\overline{ϵ}_a+_q(M_q^2/\mathrm{}\omega _q)`$, $`A_{N2}=m_1\mathrm{}\omega _1+m_2\mathrm{}\omega _2+\mathrm{}+m_{N2}\mathrm{}\omega _{N2}`$ and $`g_q=(M_q/\mathrm{}\omega _q)^2`$. Having isolated these factors, we keep all summation indices $`m_q`$ fixed except for $`m_{N1}`$ and $`m_N`$, impose $`\omega _{N1}=\omega _N=\omega _0`$, and rewrite the double summation of these last two modes as
$$\underset{m_{N1}=0}{\overset{\mathrm{}}{}}\underset{m_N=0}{\overset{\mathrm{}}{}}_{m_{N1},m_N}$$
(114)
which, by way of the general algebraic rule
$$\underset{k=0}{\overset{\mathrm{}}{}}\underset{l=0}{\overset{\mathrm{}}{}}_{k,l}=\underset{m=0}{\overset{\mathrm{}}{}}\underset{p=0}{\overset{m}{}}_{p,mp}$$
(115)
can now be rewritten as
$$\underset{m=0}{\overset{\mathrm{}}{}}\frac{()^m}{(zA_{N2})m\mathrm{}\omega _0}\underset{p=0}{\overset{m}{}}L_p^{n_{N1}p}(g_{N1})L_{mp}^{n_Nm+p}(g_N).$$
(116)
Appealing now to the sum rule (A7) for the generalized Laguerre polynomials, deduced in the Appendix below, the summation over the index $`p`$ obeys the equality
$`{\displaystyle \underset{p=0}{\overset{m}{}}}L_p^{n_{N1}p}(g_{N1})L_{mp}^{n_Nm+p}(g_N)`$ (117)
$`=L_m^{(n_{N1}+n_N)m}(g_{N1}+g_N),`$ (118)
and the contraction of the last two vibration modes summations originates the single summation
$$\underset{m=0}{\overset{\mathrm{}}{}}\frac{()^m}{(zA_{N2})m\mathrm{}\omega _0}L_m^{(n_{N1}+n_N)m}(g_{N1}+g_N).$$
(119)
Keeping on contracting all Laguerre polynomials over the phonon modes, we end up with just one Laguerre polynomial and then the transition amplitude matrix element can be written in a “one-mode” form:
$`l_f;n_1,n_2,\mathrm{},n_N|T|k_i;0_1,0_2,\mathrm{},0_N=`$ (120)
$`V_{l_f,a}V_{a,k_i}e^g\left[{\displaystyle \underset{q}{}}(g_q)^{n_q/2}(n_q!)^{1/2}\right]`$ (121)
$`\times {\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{()^mL_m^{nm}(g)}{\left[ϵ_{k_i}\overline{ϵ}_a+\lambda m\mathrm{}\omega _0\right]}},`$ (122)
$`n=_qn_q`$ being the total number of final excited phonons into all the lattice vibration modes, $`g=_q(M_q/\mathrm{}\omega _0)^2=_qg_q`$ and $`\lambda =_q(M_q^2/\mathrm{}\omega _0)`$.
## IV Summing over Final Phonon States
The transition probability will be given by the modulus square of amplitude above. Summing over all possible final phonon states $`n_1,n_2,\mathrm{},n_N|`$ such that the total number of excited phonons is $`n`$, we have
$`\left|l_f;n|T|k_i;\mathrm{𝟎}\right|^2`$ $`=`$ $`|V_{l_f,a}|^2|V_{a,k_i}|^2e^{2g}`$ (125)
$`\times {\displaystyle \underset{n_1}{}}{\displaystyle \underset{n_2}{}}\mathrm{}{\displaystyle \underset{n_N}{}}{\displaystyle \frac{g_1^{n_1}}{n_1!}}{\displaystyle \frac{g_2^{n_2}}{n_2!}}\mathrm{}{\displaystyle \frac{g_N^{n_N}}{n_N!}}`$
$`\times \left|{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{()^mL_m^{nm}(g)}{\left[ϵ_{k_i}\overline{ϵ}_a+\lambda m\mathrm{}\omega _0\right]}}\right|^2.`$
However, in general $`x^n/n!=()^nL_n^n(x)`$ and, taking advantage once more of the algebraic properties of the generalized Laguerre polynomials, expressed in this instance by the sum rule (A10) in the Appendix below, we obtain, keeping in mind the restriction $`n_1+n_2+\mathrm{}+n_N=n`$,
$`{\displaystyle \underset{n_1}{}}{\displaystyle \underset{n_2}{}}\mathrm{}`$ $`{\displaystyle \underset{n_N}{}}{\displaystyle \frac{g_1^{n_1}}{n_1!}}{\displaystyle \frac{g_2^{n_2}}{n_2!}}\mathrm{}{\displaystyle \frac{g_N^{n_N}}{n_N!}}`$ (126)
$`=`$ $`{\displaystyle \underset{n_1}{}}{\displaystyle \underset{n_2}{}}\mathrm{}{\displaystyle \underset{n_N}{}}()^{n_1}L_{n_1}^{n_1}(g_2)()^{n_2}L_{n_2}^{n_2}(g_2)`$ (128)
$`\mathrm{}()^{n_N}L_{n_N}^{n_N}(g_N)`$
$`=`$ $`()^nL_n^n(g)={\displaystyle \frac{g^n}{n!}}.`$ (129)
Consequently, the transition probability from zero to all possible number $`n`$ of final phonons (and from electron state $`k_i`$ to electron state $`l_f`$), imposing energy conservation in the overall process, will be given by
$`{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}\left|l_f;n|T|k_i;\mathrm{𝟎}\right|^2=`$ (130)
$`|V_{l_f,a}|^2|V_{a,k_i}|^2e^{2g}`$ $`{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{g^n}{n!}}\delta (ϵ_{k_i}ϵ_{l_f}n\mathrm{}\omega _0)`$ (132)
$`\times \left|{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{()^mL_m^{nm}(g)}{\left[ϵ_{k_i}\overline{ϵ}_a+\lambda m\mathrm{}\omega _0\right]}}\right|^2,`$
for all the $`N`$ phonon modes, but in a form that reproduces a one-mode situation, as sought. And if we now go from the domain of wavevectors $`k_i`$ and $`l_f`$, to the domain of incoming and outgoing electron energies, $`ϵ_i`$ and $`ϵ_f`$, respectively, and taking matrix elements $`V_{l_f,a}`$ and $`V_{a,k_i}`$ independent of the energy in the range of interest, we get a transmission matrix $`T(ϵ_i,ϵ_f)`$ given by
$`T(ϵ_i,ϵ_f)`$ (135)
$`=\mathrm{\Gamma }_l\mathrm{\Gamma }_ke^{2g}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{g^n}{n!}}\delta (ϵ_iϵ_fn\mathrm{}\omega _0)`$
$`\times \left|{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{()^mL_m^{nm}(g)}{\left[ϵ_i\overline{ϵ}_a+\lambda m\mathrm{}\omega _0\right]}}\right|^2,`$
where $`\mathrm{\Gamma }_k=2\pi _{k_i}|V_{a,k_i}|^2\delta (ϵϵ_{k_i})`$ and $`\mathrm{\Gamma }_l=2\pi _{l_f}|V_{a,l_f}|^2\delta (ϵϵ_{l_f})`$ are the partial widths of the intermediate resonant state due to the coupling to the continuum of electron states, initial and final.
In this final result we may represent the generalized Laguerre polynomial as
$$L_m^{nm}(g)=\underset{j=0}{\overset{m}{}}\frac{\mathrm{\Gamma }(n+1)}{\mathrm{\Gamma }(j+nm+1)(mj)!}\frac{(g)^j}{j!}$$
(136)
and then, using the algebraic rule in Eq.(3.19), the factor inside the modulus square of Eq.(4.3) above can be rewritten as
$`{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{j=0}{\overset{m}{}}}{\displaystyle \frac{()^m\mathrm{\Gamma }(n+1)}{\mathrm{\Gamma }(jm+n+1)(mj)!}}{\displaystyle \frac{(g)^j}{j!}}\left[ϵ_i\overline{ϵ}_a+\lambda m\mathrm{}\omega _0\right]^1`$ (137)
$`={\displaystyle \underset{m_1=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{m_2=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{()^{(m_1+m_2)}\mathrm{\Gamma }(n+1)}{\mathrm{\Gamma }(m_2+n+1)(m_2!)}}{\displaystyle \frac{(g)^{m_1}}{(m_1)!}}\left[ϵ_i\overline{ϵ}_a+\lambda (m_1+m_2)\mathrm{}\omega _0\right]^1.`$ (138)
However, if we take into consideration the presence of the poles of the gamma function $`\mathrm{\Gamma }(m_2+n+1)`$ at the points $`m_2=n+1,n+2,\mathrm{}`$, we conclude that the upper limit in the $`m_2`$ summation is, in fact, not $`\mathrm{}`$ but $`n`$. Consequently, we can write the transmission matrix above as (changing labels, from $`m_1`$ to $`m`$, and $`m_2`$ to $`j`$)
$`T(ϵ_i,ϵ_f)=`$ $`\mathrm{\Gamma }_l\mathrm{\Gamma }_ke^{2g}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{g^n}{n!}}\delta (ϵ_iϵ_fn\mathrm{}\omega _0)`$ (140)
$`\times \left|{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{j=0}{\overset{n}{}}}{\displaystyle \frac{()^jn!}{(nj)!j!}}{\displaystyle \frac{g^m}{m!}}\left[ϵ_i\overline{ϵ}_a+\lambda (m+j)\mathrm{}\omega _0\right]^1\right|^2,`$
which is the result expressed by Eq. (32) in Ref. and by Eq. (10a) in Ref. . We may conclude, therefore, that the relationship between the end results of the two approaches goes beyond a numerical indistinguishability, and that a consistent description of inelastic resonant tunneling in heterostructures can also be achieved, without the need to use many-body Green’s functions.
## V Conclusions
Our aim in the present paper was to clarify the possible validity of the antecipated connection between two different approaches to the same basic phenomenon, not only for one vibrational degree of freedom as done before, but for a whole set of vibration modes, even if with the restriction of equal frequencies. Fundamental to the whole exposition were the algebraic properties of the generalized Laguerre polynomials. Furthermore, we believe that it may become possible, with this simpler method, to treat more general situations in a direct fashion, in particular the consideration of non-Lorentzian lineshapes, of non-linear electron-phonon coupling, and of more than one intermediate resonance.
###### Acknowledgements.
The author thanks Professor P.R. Antoniewicz for a critical reading of the manuscript and gratefully acknowledges financial support from F.L.A.D. (Luso-American Development Foundation) and from the University of the Azores, Portugal.
## Sum rules for the generalized Laguerre polynomials
We give here proof of the sum rules between generalized Laguerre polynomials used above. A set of useful references regarding orthogonal polynomials is given below (Refs. ). Starting from the generating function
$$w(x,z)=(1+z)^\alpha e^{xz}=\underset{n=0}{\overset{\mathrm{}}{}}c_n(x)z^n$$
(141)
expanded in powers of the complex variable $`z`$ (such that $`|z|<1`$), the coefficients $`c_n(x)`$ are given by
$`c_n(x)=`$ $`{\displaystyle \frac{1}{2\pi i}}{\displaystyle _\mathrm{\Gamma }}{\displaystyle \frac{w(x,z)}{z^{n+1}}}𝑑z={\displaystyle \frac{1}{n!}}{\displaystyle \frac{d^nw(x,z)}{dz^n}}_{z=0}`$ (142)
$`=`$ $`{\displaystyle \frac{1}{n!}}{\displaystyle \frac{d^n[(1+z)^\alpha e^{xz}]}{dz^n}}_{z=0},`$ (143)
where the contour $`\mathrm{\Gamma }`$ encloses the origin. Using Leibniz’s rule for the n-th derivative of a product of two functions, we get
$`c_n(x)=`$ $`{\displaystyle \frac{1}{n!}}{\displaystyle \underset{r=0}{\overset{n}{}}}\left(\begin{array}{c}n\\ r\end{array}\right){\displaystyle \frac{\alpha !}{(\alpha r)!}}(x)^{nr}`$ (146)
$`=`$ $`{\displaystyle \underset{r=0}{\overset{n}{}}}{\displaystyle \frac{\mathrm{\Gamma }(\alpha +1)}{\mathrm{\Gamma }(\alpha n+r+1)(nr)!}}{\displaystyle \frac{(x)^r}{r!}}`$ (147)
$`=`$ $`L_n^{(\alpha n)}(x).`$ (148)
Hence, we can write
$$w(x,z)=(1+z)^\alpha e^{xz}=\underset{n=0}{\overset{\mathrm{}}{}}L_n^{(\alpha n)}(x)z^n,$$
(149)
using the definition of generalized (as opposed to “associated”) Laguerre polynomials, where the upper index can be any complex number.
Some convenient representations of the generalized Laguerre polynomials, which also remain valid for $`\alpha `$ a negative integer, are
$`L_n^{(\alpha )}=`$ $`{\displaystyle \underset{r=0}{\overset{n}{}}}{\displaystyle \frac{\mathrm{\Gamma }(\alpha +n+1)}{\mathrm{\Gamma }(\alpha +r+1)(nr)!}}{\displaystyle \frac{(x)^r}{r!}}`$ (150)
$`=`$ $`{\displaystyle \frac{1}{n!}}{\displaystyle \underset{r=0}{\overset{n}{}}}{\displaystyle \frac{(n)_r}{r!}}(\alpha +r+1)_{nr}x^r.`$ (151)
In this last expression we introduced the Pochhammer symbols, $`(a)_r=\mathrm{\Gamma }(a+r)/\mathrm{\Gamma }(a)`$.
But now, with the help of the generating function above, we can prove the sum rule required. In fact,
$`(1+z)^\alpha e^{xz}(1+z)^\beta e^{yz}`$ (152)
$`=`$ $`{\displaystyle \underset{r=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{s=0}{\overset{\mathrm{}}{}}}L_r^{(\alpha r)}(x)L_s^{(\beta s)}(y)z^{r+s}`$ (153)
$`=`$ $`(1+z)^{\alpha +\beta }e^{(x+y)z}`$ (154)
$`=`$ $`{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}L_n^{(\alpha +\beta n)}(x+y)z^n`$ (155)
Equating equal powers of the variable $`z`$, we conclude that $`r+s=n`$ and consequently
$$\underset{r=0}{\overset{n}{}}L_r^{(\alpha r)}(x)L_{nr}^{(\beta +rn)}(y)=L_n^{(\alpha +\beta n)}(x+y),$$
(156)
our first desired result. The summation index $`r`$ can not attain values larger than $`n`$, since the lower index of a generalized Laguerre polynomial is always nonnegative.
Next we prove the second sum rule used above. It is, in fact, a consequence of the first sum rule, for which the particular case $`\alpha =\beta =0`$ gives
$$\underset{r=0}{\overset{n}{}}L_r^{(r)}(x_1)L_{nr}^{(rn)}(x_2)=L_n^{(n)}(x_1+x_2).$$
(157)
If we now consider only the first two factors in the product (4.2) above, we may write, for a fixed value of $`n`$ (and of $`n_3,\mathrm{},n_N`$),
$`{\displaystyle \underset{n_1+n_2=n(n_3+\mathrm{}+n_N)}{}}`$ $`L_{n_1}^{n_1}(g_1)L_{n_2}^{n_2}(g_2)`$ (158)
$`=`$ $`{\displaystyle \underset{n_1=0}{\overset{n(n_3+\mathrm{}+n_N)}{}}}L_{n_1}^{n_1}(g_1)L_{[n(n_3+\mathrm{}+n_N)n_1]}^{[n(n_3+\mathrm{}+n_N)n_1]}(g_2)`$ (159)
$`=`$ $`L_{[n(n_3+\mathrm{}+n_N)]}^{[n(n_3+\mathrm{}+n_N)]}(g_1+g_2).`$ (160)
We repeat this contracting process over all $`N`$ factors, ending up with just one final Laguerre polynomial, $`L_n^n(g)`$. In general, then,
$`{\displaystyle \underset{n_1+n_2+\mathrm{}+n_N=n}{}}`$ $`L_{n_1}^{n_1}(x_1)L_{n_2}^{n_2}(x_2)\mathrm{}L_{n_N}^{n_N}(x_N)`$ (162)
$`=L_n^n(x),`$
with $`x=x_1+x_2+\mathrm{}+x_N`$ and $`n=n_1+n_2+\mathrm{}+n_N`$. |
warning/0002/hep-ph0002296.html | ar5iv | text | # Manifestaion of SUSY in B decays aafootnote a Talk given at the Third International Conference on B Physics and CP Violation, December 3 -7, 1999, Taipei.
## 1 Introduction
In order to explore supersymetry (SUSY) indirect search experiments can play a complementary role to direct search for SUSY particles at collider experiments. Since SUSY particles may affect flavor changing neutral current (FCNC) processes and CP violation in B and K meson decays, it is possible that new experiments in B decay at both $`e^+e^{}`$ colliders and hadron machines reveal new physics signals which can be interpreted as indirect evidence of SUSY.
In the context of SUSY models flavor physics has important implications. Because the squark and the slepton mass matrices become new sources of flavor mixings generic mass matrices would induce too large FCNC and lepton flavor violation (LFV) effects if the superpartners’ masses are in a few-hundred-GeV region. For example, if we assume that the SUSY contribution to the $`K^0\overline{K}^0`$ mixing is suppressed because of the cancellation among the squark contributions of different generations, the squarks with the same gauge quantum numbers must be highly degenerate in masses at least for the first two generations.
There are several scenarios to solve this flavor problem. In the minimal supergravity model flavor problem are avoided by taking SUSY soft-breaking terms as flavor-blind structure. The scalar mass terms are assumed to be common for all scalar fields at the Planck scale and therefore there are no FCNC effects nor LFV from the squark and slepton sectors at this scale. The physical squark and slepton masses are determined taking account of renormalization effects from the Planck to the weak scale. This induces sizable SUSY contributions to various FCNC and LFV processes.
In this talk we consider two types of SUSY models and discuss FCNC processes. The first one is the minimal supersymmetric standard model (MSSM) with a universal SUSY breaking terms at the Planck scale which is realized in the minimal supergravity model. The other is the SU(5) grand unified theory with right-handed neutrino supermultiplets. This model incorporates the see-saw mechanism for neutrino mass generation. In the latter case the neutrino Yukawa coupling constants can be a source of the flavor mixing in the right-handed-down-type-squark sector and due to this mixing the time-dependent CP asymmetry of radiative $`B`$ decay can be as large as 30% and the ratio of $`B_s`$$`\overline{B}_s`$ mixing and $`B_d`$$`\overline{B}_d`$ mixing deviates from the prediction for the standard model (SM) and the MSSM without the neutrino interaction.
## 2 Update of FCNC Processes in the Supergravity Model
In the minimal SM various FCNC processes and CP violation in B and K decays are determined by the Cabibbo-Kobayashi-Maskawa (CKM) matrix. Constraints on the parameters of the CKM matrix can be conveniently expressed in terms of the unitarity triangle. With CP violation at B factory as well as rare K decay experiments we will be able to check consistency of the unitarity triangle and at the same time search for effects of physics beyond the SM. In order to distinguish possible new physics effects it is important to identify how various models can modify the SM predictions.
Although general SUSY models can change the lengths and the angles of the unitarity triangle in variety ways, the supergravity model predicts a specific pattern of the deviation from the SM.$`^\mathrm{?}`$ Namely, we can show that the SUSY loop contributions to FCNC amplitudes approximately have the same dependence on the CKM elements as the SM contributions. In particular, if we assume that there are no CP violating phases from SUSY breaking sectors, the complex phase of the $`B^0\overline{B^0}`$ mixing amplitude does not change even if we take into account the SUSY contributions. The case with supersymmetric CP phases was also studied within the minimal supergravity model and it was shown that effects of new CP phases on the $`B^0\overline{B^0}`$ mixing amplitude and the direct CP asymmetry in the $`bs\gamma `$ process are small once constraints from neutron and electron EDMs are included.$`^\mathrm{?}`$
We calculate various FCNC processes in the supergravity model with universal soft breaking terms at a high energy scale. The results are summarized as follows.
1. The amplitude for $`bs\gamma `$ can receive a large contribution from the SUSY and the charged-Higgs-top-quark loop diagrams. The experimental branching ratio puts a strong constraint on SUSY parameter space. Since the SUSY contribution can interfere with other contributions either constructively or destructively we cannot exclude the light charged Higgs boson region unlike the non-SUSY type II two Higgs doublet model.$`^\mathrm{?}`$
2. When the sign of the $`bs\gamma `$ amplitude is opposite to that of the SM, B($`bsl^+l^{}`$) can be twice larger than the SM prediction. This can occur for a large $`\mathrm{tan}\beta `$ region where $`\mathrm{tan}\beta `$ is the ratio of two vacuum expectation values of Higgs fields. The deviation is also evident in the differential branching ratio and the lepton forward-backward asymmetry.$`^\mathrm{?}`$
3. In terms of consistency check of the unitarity triangle the supergravity model has the following features.$`^\mathrm{?}`$ (i)$`\mathrm{\Delta }M_{B_d}`$ and $`ϵ_K`$ are enhanced by the SUSY and charged-Higgs loop effects. When these quantities are normalized by the corresponding quantities in the SM they are almost independent of the CKM matrix element, and the enhancement factors for $`\mathrm{\Delta }M_{B_d}`$ and $`ϵ_K`$ are almost equal. (ii) The branching ratios for $`K^+\pi ^+\nu \overline{\nu }`$ and $`K_L\pi ^0\nu \overline{\nu }`$ processes are suppressed compared to the SM prediction. Again the suppression factor are almost the same for two branching ratios and does not depend strongly on the CKM matrix element. (iii) CP asymmetries in various B decay modes such as $`BJ/\psi K_S`$ and the ratio of $`\mathrm{\Delta }M_{B_s}`$ and $`\mathrm{\Delta }M_{B_d}`$ are the same as the SM prediction.
In Fig. 1 we present the correlation between $`\mathrm{\Delta }M_{B_d}`$ and B($`K_L\pi ^0\nu \overline{\nu }`$) normalized by the corresponding quantities in the SM for $`\mathrm{tan}\beta `$ = 3. The constraint on the SUSY parameter space from the recent improved SUSY Higgs search is implemented.$`^\mathrm{?}`$ We have calculated the SUSY particle spectrum based on two different assumptions on the initial conditions of renormalization group equations. The minimal case corresponds to the minimal supergravity where all scalar fields have a common SUSY breaking mass at the GUT scale. For “nonminimal” we enlarge the SUSY parameter space by relaxing the initial conditions for the SUSY breaking parameters, namely all squarks and sleptons have a common SUSY breaking mass whereas an independent SUSY breaking parameter is assigned for Higgs fields. The square(dot) points correspond to the minimal (enlarged) parameter space of the supergravity model. We can see that the $`\mathrm{\Delta }M_{B_d}`$ (and $`ϵ_K`$) can deviated from the SM by 20% whereas the deviation in $`K_L\pi ^0\nu \overline{\nu }`$ and $`K^+\pi ^+\nu \overline{\nu }`$ processes are small. These deviations may be evident in future when B factory experiments provide additional information on the CKM parameters.
## 3 FCNC in SUSY GUT with Right-handed Neutrino
In this section we consider FCNC and LFV of charged lepton decays in the model of a SU(5) SUSY GUT which incorporates the see-saw mechanism for the neutrino mass generation.$`^\mathrm{?}`$ In this model sources of the flavor mixing are Yukawa coupling constant matrices for quarks and leptons as well as that for the right-handed neutrinos. Because the quark and lepton sectors are related by GUT interactions, the flavor mixing relevant to the CKM matrix can generate LFV such as $`\mu e\gamma `$ and $`\tau \mu \gamma `$ processes $`^\mathrm{?}`$ in addition to FCNC in hadronic observables.$`^\mathrm{?}`$ In the SUSY model with right-handed neutrinos, branching ratios of the LFV processes can become large enough to be measured in near-future experiments. $`^\mathrm{?}`$ When we consider the right-handed neutrinos in the context of GUT, the flavor mixing related to the neutrino oscillation can be a source of the flavor mixing in the squark sector. We show that due to the large mixing of the second and third generations suggested by the atmospheric neutrino anomaly, $`B_s`$$`\overline{B}_s`$mixing, the time-dependent CP asymmetry of the $`BM_s\gamma `$ process, where $`M_s`$ is a CP eigenstate including the strange quark, can have a large deviation from the SM prediction.$`^\mathrm{?}`$
The Yukawa coupling part and the Majorana mass term of the superpotential for the SU(5) SUSY GUT with right-handed neutrino supermultiplets is given by $`W=\frac{1}{8}f_U^{ij}\mathrm{\Psi }_i\mathrm{\Psi }_jH_5+f_D^{ij}\mathrm{\Psi }_i\mathrm{\Phi }_jH_{\overline{5}}+f_N^{ij}N_i\mathrm{\Phi }_jH_5+\frac{1}{2}M_\nu ^{ij}N_iN_j`$ ,where $`\mathrm{\Psi }_i`$, $`\mathrm{\Phi }_i`$ and $`N_i`$ are $`\mathrm{𝟏𝟎}`$, $`\overline{\mathrm{𝟓}}`$ and $`\mathrm{𝟏}`$ representations of SU(5) gauge group. $`i,j=1,2,3`$ are the generation indices. $`H_5`$ and $`H_{\overline{5}}`$ are Higgs superfields with $`\mathrm{𝟓}`$ and $`\overline{\mathrm{𝟓}}`$ representations.
The renormalization effects due to the Yukawa coupling constants induce various FCNC and LFV effects from the mismatch between the quark/lepton and squark/slepton diagonalization matrices. In particular the large top Yukawa coupling constant is responsible for the renormalization of the $`\stackrel{~}{q}_L`$ and $`\stackrel{~}{u}_R`$ mass matrices. At the same time the $`\stackrel{~}{e}_R`$ mass matrix receives sizable corrections between the Planck and the GUT scales and various LFV processes are induced. In a similar way, if the neutrino Yukawa coupling constant $`f_N^{ij}`$ is large enough, the $`\stackrel{~}{l}_L`$ mass matrix and the $`\stackrel{~}{d}_R`$ mass matrix receive sizable flavor changing effects due to renormalization between the Planck and the right-handed neutrino mass scales and the Planck and the GUT scales, respectively. These are sources of extra contributions to LFV processes and various FCNC processes such as $`bs\gamma `$, the $`B^0`$$`\overline{B}^0`$ mixing and the $`K^0`$$`\overline{K}^0`$ mixing.
The flavor mixing in the $`\stackrel{~}{d}_R`$ sector can induce large time-dependent CP asymmetry in the $`BM_s\gamma `$ process. Using the Wilson coefficients $`c_7`$ and $`c_7^{}`$ in the effective Lagrangian for the $`bs\gamma `$ decay $`=c_7\overline{s}\sigma ^{\mu \nu }b_RF_{\mu \nu }+c_7^{}\overline{s}\sigma ^{\mu \nu }b_LF_{\mu \nu }+H.c.`$, the asymmetry is written as $`\frac{\mathrm{\Gamma }(t)\overline{\mathrm{\Gamma }}(t)}{\mathrm{\Gamma }(t)+\overline{\mathrm{\Gamma }}(t)}=\xi A_t\mathrm{sin}\mathrm{\Delta }m_dt,A_t=\frac{2\mathrm{I}\mathrm{m}(\mathrm{e}^{\mathrm{i}\theta _\mathrm{B}}\mathrm{c}_7\mathrm{c}_7^{})}{|c_7|^2+|c_7^{}|^2},`$ where $`\mathrm{\Gamma }(t)`$ ($`\overline{\mathrm{\Gamma }}(t)`$) is the decay width of $`B^0(t)M_s\gamma `$ ($`\overline{B}^0(t)M_s\gamma `$) and $`M_s`$ is some CP eigenstate ($`\xi =+1(1)`$ for a CP even (odd) state) such as $`K_S\pi ^0`$.$`^\mathrm{?}`$ $`\mathrm{\Delta }m_d=2|M_{12}(B_d)|`$ and $`\theta _B=\mathrm{arg}M_{12}(B_d)`$ where $`M_{12}(B_d)`$ is the $`B_d`$$`\overline{B}_d`$ mixing amplitude. Because the asymmetry can be only a few percent in the SM, a sizable asymmetry is a clear signal of new physics beyond the SM.
We calculated various FCNC and LFV observables in this model under the assumption of the universal soft breaking terms at the Planck scale. As typical examples of the neutrino parameters, we consider the following parameter set corresponding to the Mikheyev-Smirnov-Wolfenstein (MSW) small mixing case. $`m_\nu =(2.236\times 10^3,\mathrm{\hspace{0.17em}3.16}\times 10^3,\mathrm{\hspace{0.17em}5.92}\times 10^2)\mathrm{eV}`$ and the Maki-Nakagawa-Sakata (MNS) matrix is given by
$`V_{\mathrm{MNS}}`$ $`=`$ $`\left(\begin{array}{ccc}0.999& 0.0371& 0\\ 0.0262& 0.707& 0.707\\ 0.0262& 0.707& 0.707\end{array}\right).`$ (4)
We also take $`M_\nu `$ to be proportional to a unit matrix with a diagonal element of $`M_R=4\times 10^{14}`$ GeV. We fix CKM parameters as $`V_{cb}=0.04`$, $`|V_{ub}/V_{cb}|=0.08`$ and $`\delta _{13}=60^{}`$. We take $`\mathrm{tan}\beta =5`$ and vary other SUSY parameters. We take account of various phenomenological constraints on SUSY parameters including B($`bs\gamma `$). We also calculated B($`\mu e\gamma `$) and $`ϵ_K`$ and imposed constraints from these quantities.
The upper part of Fig.2 shows a correlation between $`\mathrm{\Delta }m_s/\mathrm{\Delta }m_d`$ (ratio of $`B_s`$$`\overline{B}_s`$ and $`B_d`$$`\overline{B}_d`$ mass splittings) and B($`\tau \mu \gamma `$). We can see that $`\mathrm{\Delta }m_s/\mathrm{\Delta }m_d`$ can be enhanced up to 30% compared to the SM prediction. This feature is quite different from the minimal supergravity model without the GUT and right-handed neutrino interactions $`^\mathrm{?}`$ where $`\mathrm{\Delta }m_s/\mathrm{\Delta }m_d`$ is almost the same as the SM value. $`A_t`$ for the same parameter set is shown as a function of B($`\tau \mu \gamma `$) in the lower part of Fig. 2. We can see that $`|A_t|`$ can be close to 30% when B($`\tau \mu \gamma `$)is larger than $`10^8`$. The large asymmetry arises because the renormalization effect due to $`f_N`$ induces sizable contribution to $`c_7^{}`$ through gluino–$`\stackrel{~}{d}_R`$ loop diagrams. In this example possible new physics signals in B($`\tau \mu \gamma `$), $`B_s`$$`\overline{B}_s`$ mixing and $`A_t`$ all come from the renormalization effect on squark and slepton mass matrices from the large neutrino Yukawa coupling constant. Because these signals provide quite different signatures compared to the SM and the minimal supergravity model without GUT and right-handed neutrino interactions, future experiments in $`B`$ physics and LFV can provide us important clues on the interactions at very high energy scale.
The work was supported in part by the Grant-in-Aid of the Ministry of Education, Science, Sports and Culture, Government of Japan (No.09640381), Priority area “Supersymmetry and Unified Theory of Elementary Particles” (No. 707), and “Physics of CP Violation” (No.09246105).
## References |
warning/0002/math0002011.html | ar5iv | text | # Stability Properties of the Riemann Ellipsoids
## 1 Introduction
A. The Dirichlet problem consists of the “linear” motions of an ideal, incompressible, homogeneous, self–gravitating fluid mass having an ellipsoidal shape at any instant. This problem originated in attempts to determine the shapes of the planets, and has a long and distinguished history (see e.g. chapter 1 of and references therein). Let us only remark here that Newton and Maclaurin established the existence of rigidly rotating oblate spheroids and that Jacobi described rigidly rotating asymmetric ellipsoids. Dirichlet proved the existence of “linear” fluid motions, such that the position $`x(t,y)`$ at time $`t`$ of any material particle $`y`$ is given by
$$x(t,y)=F(t)y,F(t)\mathrm{SL}(3).$$
(1.1)
If, as we assume, the reference configuration of the fluid mass is a ball of radius $`\rho `$, then the free surface of the fluid determined by the motion (1.1) is at every instant an ellipsoid with semiaxes $`\rho a_1`$, $`\rho a_2`$, $`\rho a_3`$, where $`a_1`$, $`a_2`$ and $`a_3`$ are the singular values of the matrix $`F`$ (namely, the eigenvalues of $`\sqrt{FF^T}`$).
The Maclaurin spheroids and Jacobi ellipsoids are steady rigid motions of the Dirichlet problem. Dedekind identified a new class of steady, but non–rigid, motions of the problem; they were constructed by transposition of the Jacobi ellipsoids, thus bringing the particle relabelling symmetry of the system into play. Riemann reformulated the equations of motion of the Dirichlet problem in a particularly convenient form for studying asymmetric configurations (as we shall show, this reformulation corresponds to the passage to a four–to–one covering space). Riemann also classified all steady motions of the system, showing that they are generated by combinations of spatial and body symmetries, and studied their Lyapunov (or nonlinear) orbital stability.
The steady motions found by Riemann, which are called Riemann ellipsoids, are motions (1.1) such that
$$F(t)=\mathrm{exp}(t\mathrm{\Omega }_l)A\mathrm{exp}(t\mathrm{\Omega }_r)$$
for some constant diagonal matrix $`A=\mathrm{diag}(a_1,a_2,a_3)`$ and constant antisymmetric matrices $`\mathrm{\Omega }_l`$ and $`\mathrm{\Omega }_r`$. Except for the special case of the Jacobi ellipsoids, for which $`\mathrm{\Omega }_r=0`$, Riemann ellipsoids do not perform rigid motions: the motion is the composition of an internal rotation, a stretch along the principal axes, and a spatial rotation; the free surface retains a rotating ellipsoidal shape, but the fluid particles describe rosette–shaped motions which are either closed (periodic) or open (quasi–periodic) depending on the two angular frequencies $`\omega _l`$ and $`\omega _r`$ (that is, the vectors corresponding to the antisymmetric matrices $`\mathrm{\Omega }_l`$ and $`\mathrm{\Omega }_r`$). Riemann proved that there are five types of Riemann ellipsoids with distinct semiaxes, characterized by different relations between the angular velocities and the semiaxes. For two types (hereafter referred to as $`\mathrm{S}_2`$ and $`\mathrm{S}_3`$), both $`\omega _l`$ and $`\omega _r`$ are parallel to the same principal axis of the ellipsoid; this axis is either the shortest ($`\mathrm{S}_3`$) or the middle ($`\mathrm{S}_2`$) one. For the other three types, called I, II, and III, both $`\omega _l`$ and $`\omega _r`$ belong to one of the two principal planes containing the longest axis of the ellipsoid.<sup>1</sup><sup>1</sup>1The distinction between $`\mathrm{S}_2`$ and $`\mathrm{S}_3`$ ellipsoids that we make here is not standard; Chandrasekhar calls them collectively ellipsoids of type S.
Since the Riemann ellipsoids are reduced equilibria, namely equilibria modulo symmetries, of the Dirichlet problem, it is meaningful to study their stability modulo symmetries, which we will frequently refer to simply as stability. Riemann found that all ellipsoids of type $`\mathrm{S}_3`$ and a subfamily of the $`\mathrm{S}_2`$ ellipsoids are stable. Riemann’s analysis does not lead to any definite conclusion about the (in)stability of the remaining ellipsoids, because they are saddle points of his Lyapunov function.
More recently, Chandrasekhar studied the ‘ellipticity’, or ‘spectral stability’, of the Riemann ellipsoids. We say here that a (relative) equilibrium is elliptic if the eigenvalues of the linearization at the equilibrium are purely imaginary; note that, accordingly, we treat the case of zero eigenvalues as elliptic, even though this is not customary. Chandrasekhar numerically found a region of ellipticity of the ellipsoids of type $`\mathrm{S}_2`$ larger than the region of Lyapunov stability found by Riemann. He also found nonempty regions of elliptic ellipsoids of type I and III, but none of type II. As we shall see, some of these conclusions are not entirely correct, and in fact require substantial modification.
B. For a Hamiltonian system such as the Dirichlet problem, ellipticity is a necessary, but not sufficient, condition for Lyapunov stability. Therefore, the ellipticity analysis is of twofold interest: on one hand, the absence of ellipticity implies Lyapunov instability; on the other hand, it opens up the problem of the stability properties of those Riemann ellipsoids that are elliptic, but of undetermined Lyapunov stability. The present article is mainly devoted to a study of the latter question.
First of all, as a prerequisite for this analysis, we repeat the Lyapunov stability and the ellipticity analysis of the Riemann ellipsoids. While our analysis confirms Riemann’s Lyapunov stability results, we find significant discrepancies with Chandrasekhar’s results: there is a nonempty region of elliptic ellipsoids of type II and a region of elliptic ellipsoids of type III that is substantially larger than the one reported by Chandrasekhar. Therefore, the regions of known (Lyapunov) instability of the ellipsoids of types II and III are substantially smaller than those found by Chandrasekhar. Moreover, the ellipticity regions of the ellipsoids of type I, II and III have a finer and more complex structure than was previously realized.
The main goal of the article is to use the methods of Hamiltonian perturbation theory to investigate the long–time stability properties of the elliptic Riemann ellipsoids of undetermined Lyapunov stability. In fact, the Dirichlet problem is a Hamiltonian system with symmetry and the Riemann ellipsoids are the relative equilibria of this system. They correspond to the equilibria of the reduced system, which has four degrees of freedom for generic values of the momenta and three degrees of freedom in special cases. (The reduced phase space is in fact an orbifold, with a singular set containing the Riemann ellipsoids of type S, but this difficulty is overcome by passing to Riemann’s covering space.) Since the reduced system has more than two degrees of freedom, KAM theory cannot be used to show the stability of its equilibria. Therefore, we shall base our investigation on Nekhoroshev theory, which allows one to obtain “practical” stability results, namely, to prove that motions starting sufficiently near the equilibrium remain near it for times growing exponentially with (some power) of the initial distance; specifically, we shall use some recent results obtained in (see also , , ). The results of our investigation indicate that all elliptic Riemann ellipsoids are Nekhoroshev–stable, with the possible exception of a finite number of codimension one families of ellipsoids satisfying certain low–order resonance conditions.
In view of the complexity of the system, we resort to numerical methods at certain steps in the ellipticity analysis and in the construction of the normal forms for the reduced Hamiltonian. We numerically compute the eigenvalues of the linearization and numerically construct the normal forms for a finite (though rather large and, we think, significant) number of values of the momentum. Therefore, the forementioned results are not rigorous, but we think that our analysis provides strong evidence for our conclusions.
C. Our analysis is based on a Hamiltonian formulation (on a covering manifold) of the Dirichlet problem, which is briefly developed in Section 2. For the sake of brevity, we skip most of the computational details and focus on the conceptual features. The Riemann ellipsoids are described in Section 3. Overall, Sections 2 and 3 contain a treatment of the Dirichlet problem which, although concise, is complete and detailed; this seems necessary, because the existing literature does not seem to include a comparably comprehensive treatment.
Sections 4–6 contain the ellipticity and the Nekhoroshev–stability analyses. In Section 4 we review the results on Lyapunov stability and instability and describe the results of our ellipticity analysis; a detailed comparison with Chandrasekhar’s results is included in Appendix B. We briefly review the necessary notions for the Nekhoroshev–stability analysis in Section 5 and describe the results for the Riemann ellipsoids in Section 6. A Conclusion follows.
All of the numerical computations have been performed using the software package Mathematica.
## 2 The Dirichlet Problem
A. Dirichlet’s equations. Dirichlet showed that the linear motions (1.1) form an invariant submanifold of the phase space of an ideal, incompressible, homogeneous, self–gravitating fluid, with the boundary condition that the pressure is constant at the free surface. He expressed the restriction of the hydrodynamic equations of motion to this submanifold as equations for the curve $`F(t)\mathrm{SL}(3)`$, obtaining a second order system on $`\mathrm{SL}(3)`$.
To describe these results, we introduce the following notation, which will be used throughout the article. We denote by $``$ and $``$ the Euclidean scalar product and norm on $`^3`$ and by $`D_\rho `$ the ball $`\{y^3:y\rho \}`$. We use the standard inner product and norm on $`L(3)`$ given by
$$A,B=\mathrm{tr}(AB^T)\mathrm{and}\left|A\right|^2=A,A.$$
(2.1)
(Note that the subspaces of symmetric and antisymmetric matrices are $`,`$–orthogonal complements of one another in $`L(3)`$.) We shall frequently identify $`^3`$ and $`\mathrm{so}(3)`$ by means of the isomorphism $`\omega \widehat{\omega }`$, where $`\widehat{\omega }\xi =\omega \times \xi `$ for all $`\xi ^3`$. Note that $`\widehat{\omega },\widehat{\xi }=\mathrm{\hspace{0.33em}2}\omega \xi `$ for all $`\omega ,\xi ^3`$. Following Riemann, we utilize the “singular value decomposition” of matrices, nowadays a standard tool in numerical linear algebra: If $`F\mathrm{SL}(3)`$ and $`a_1,a_2,a_3`$ are its singular values, then there exist matrices $`U_l,U_r\mathrm{SO}(3)`$ such that
$$F=U_lAU_r^T,A=\mathrm{diag}(a_1,a_2,a_3).$$
We shall refer to the triplet $`(A,U_l,U_r)`$ as a “SVD” of the matrix $`F`$. The matrices $`A`$, $`U_l`$ and $`U_r`$ are not unique, but one can clearly fix $`A`$ by ordering its entries, e.g. $`a_1a_2a_3`$ (see also Proposition 1 below). Finally, we set
$$𝒱(a)=2\pi g_0^{\mathrm{}}\left[(s+a_1^2)(s+a_2^2)(s+a_3^2)\right]^{1/2}𝑑s,$$
where $`g`$ denotes the gravitational constant; by a slight abuse of notation, we shall regard $`𝒱`$ and its derivatives as functions either of $`a=(a_1,a_2,a_3)`$ or of the diagonal matrix $`A=\mathrm{diag}(a_1,a_2,a_3)`$.
Dirichlet’s result can now be expressed as follows: Consider a curve $`tF(t)\mathrm{SL}(3)`$ defined on some open interval $``$ and let $`(A(t),U_l(t),U_r(t))`$ be any singular value decomposition of $`F(t)`$. Then the map
$$x:\times D_\rho ^3,(t,y)F(t)y$$
is a solution of the hydrodynamical equation for an ideal, incompressible, homogeneous, self–gravitating fluid with constant pressure at the boundary iff
$$_F\left[\ddot{F}+U_l𝒱^{}(A)U_r^T\right]=0,$$
(2.2)
where $`𝒱^{}=\mathrm{diag}(\frac{𝒱}{a_1},\frac{𝒱}{a_2},\frac{𝒱}{a_3})`$ and $`_FG=G\frac{1}{3}G,FF^T`$ for any $`GL(3)`$. Note that the matrix $`U_l𝒱^{}(A)U_r^T`$ is independent of the particular SVD of $`F`$ used; this is proven by observing that, if $`U,V\mathrm{SO}(3)`$ and $`A=UA^{}V`$ with $`A`$ and $`A^{}`$ diagonal, then $`𝒱^{}(A)=𝒱^{}(UA^{}V)=U𝒱^{}(A^{})V`$. In fact, Dirichlet used the polar decomposition, not the SVD, when writing the force term in the equations of motion.
Remark: In his work on Dirichlet problem, Chandrasekhar required that the relative pressure $`pp_o`$, where $`p_o`$ is the constant pressure on the surface of the fluid, be nonnegative at any point of the fluid. Following Riemann, we do not make this assumption. (There are many fluids that support even very high tensile states.)
Since the gravitational forces are conservative and invariant under spatial rotations and particle relabelling, equations (2.2) inherit the conservation of energy, angular momentum and circulation from the hydrodynamic equations. For future reference, we observe that if we set $`c=\frac{m\rho ^2}{5}`$, where $`m`$ denotes the mass of the fluid, then the energy of any motion (1.1) is $`c\left[\frac{1}{2}\left|\dot{F}\right|^2+𝒱(a)\right]`$ and the angular momentum is the vector $`k`$ satisfying $`\widehat{k}=\mathrm{\hspace{0.33em}2}c\text{skew}\left(\dot{F}F^T\right)`$; here $`\text{skew}(G)=\frac{1}{2}(GG^T)`$. The conservation of circulation is equivalent to the constancy of $`\text{skew}\left(F^T\dot{F}\right)`$.
B. Riemann’s equation. Shortly after Dirichlet’s work appeared, Riemann rewrote equation (2.2) using what is nowadays called the singular value decomposition of matrices and left–trivializing the factors $`T\mathrm{SO}(3)`$, namely expressing the tangent vectors $`\dot{U}_l`$ and $`\dot{U}_r`$ in terms of the antisymmetric matrices $`\mathrm{\Omega }_l=U_l^T\dot{U}_l`$ and $`\mathrm{\Omega }_r=U_r^T\dot{U}_r`$. Once these substitutions have been made, equation (2.2) takes the form
$$_A\left[\ddot{A}+\mathrm{\hspace{0.33em}2}(\mathrm{\Omega }_l\dot{A}\dot{A}\mathrm{\Omega }_r)+(\dot{\mathrm{\Omega }}_lAA\dot{\mathrm{\Omega }}_r)+(\mathrm{\Omega }_l^2A2\mathrm{\Omega }_lA\mathrm{\Omega }_r+A\mathrm{\Omega }_r^2)+𝒱^{}(A)\right]=0$$
(2.3)
(which should be compared with equations ($`\alpha `$) in ). This equation determines a second order system on $`𝒜\times \mathrm{SO}(3)\times \mathrm{SO}(3)`$, where
$$𝒜=\{\mathrm{diag}(a_1,a_2,(a_1a_2)^1):a_1>a_2>a_1^{1/2}\}$$
(if $`A`$ has repeated eigenvalues, $`\dot{\mathrm{\Omega }}_l`$ and $`\dot{\mathrm{\Omega }}_r`$ are not uniquely determined by (2.3)). This system is equivalent to the restriction of Dirichlet’s equation (2.2) to the submanifold
$$Q=\{F\mathrm{SL}(3):F\mathrm{has}\mathrm{distinct}\mathrm{singular}\mathrm{values}\}.$$
More precisely, because of the nonuniqueness of the SVD, the two systems are related not by a diffeomorphism, but by a four–to–one covering map:
###### Proposition 1
The map
$$𝒞:𝒜\times \mathrm{SO}(3)\times \mathrm{SO}(3)Q,(A,U_l,U_r)U_lAU_r^T$$
is a 4–fold covering map which maps solutions of equation (2.3) into solutions of equation (2.2).
Proof. If $`(A,U_l,U_r)`$ and $`(A,V_l,V_r)`$ are any two SVDs of $`FQ`$, with $`A=\mathrm{diag}(a_1,a_2,a_3)`$, $`a_1>a_2>a_3`$, then $`V_l=R_iU_l`$ and $`V_r=R_iU_r`$ for some $`i=0,\mathrm{},3`$, where
$$R_0=\mathrm{𝟏},R_1=\mathrm{diag}(1,1,1),R_2=\mathrm{diag}(1,1,1),\text{and}R_3=\mathrm{diag}(1,1,1).$$
(2.4)
Thus, the map $`𝒞`$ is four–to–one. That it is a submersion is proven, for instance, by showing that the determinant of its Jacobian is nonzero. (Using $`a_1,a_2`$ as coordinates on $`𝒜`$ and left–trivializing the factors $`T\mathrm{SO}(3)`$, one obtains the determinant $`(a_1^2a_2^2)(a_1^2a_3^2)(a_2^2a_3^2)`$.)
In Riemann’s formulation, the use of the SVD has the advantage that the absence of the matrices $`U_l`$ and $`U_r`$ from the equation (2.3) makes manifest the $`\mathrm{SO}(3)\times \mathrm{SO}(3)`$ invariance of the Dirichlet problem. From our perspective, the formulation on the covering also has the advantage that, once we pass to the Hamiltonian formulation, the symplectically reduced phase space associated to the cotangent bundle of the covering is a manifold, while reduction of the system on $`T^{}Q`$ produces singularities.
Remarks: (i) The submanifold $`Q`$ is not invariant under the dynamics of the Dirichlet problem (axisymmetric configurations can evolve into asymmetric ones, and vice versa) but this is not a problem for our analysis: the Riemann ellipsoids with asymmetric configurations clearly belong to $`Q`$ for all time and, since $`Q`$ is open, it will be a consequence of our stability estimates that nearby motions also remain in $`Q`$ for the time under consideration.
(ii) The restriction to the case of distinct singular values excludes all axisymmetric Riemann ellipsoids. The stability properties of the Maclaurin spheroids, which rotate rigidly about the axis of symmetry, are completely known (see ); the remaining axisymmetric ellipsoids merit further investigation.
C. Hamiltonian formulation. We now describe the Hamiltonian formulation of Riemann’s equation (2.3) on the cotangent bundle of the covering manifold $`𝒜\times \mathrm{SO}(3)\times \mathrm{SO}(3)`$. We shall use the first two diagonal entries $`a_1`$ and $`a_2`$ as coordinates on the set $`𝒜`$, by means of the diffeomorphism $`\alpha :𝒜`$, where
$$=\{(b_1,b_2)^2:b_1>b_2>b_1^{1/2}\}\text{and}\alpha (b_1,b_2):=\mathrm{diag}(b_1,b_2,\frac{1}{b_1b_2}).$$
(We denote by $`b=(b_1,b_2)`$ the first two singular values, so as to distinguish $`b^2`$ from $`a^3`$; wherever we write $`b_3`$, it stands for $`(b_1b_2)^1`$, not for an independent variable.) Furthermore we identify the cotangent bundle of each factor $`\mathrm{SO}(3)`$ with $`\mathrm{SO}(3)\times ^3`$ as follows: we first identify the cotangent space $`T_U^{}\mathrm{SO}(3)`$ with $`T_U\mathrm{SO}(3)`$ by means of the inner product $`\frac{1}{2},`$ (see (2.1)) and then identify $`T_U\mathrm{SO}(3)`$ with $`^3`$ by means of the $`\widehat{}`$ isomorphism, thus associating to each $`\mu T_U^{}\mathrm{SO}(3)`$ the unique $`m^3`$ satisfying
$$\mu U\widehat{\omega }=\frac{1}{2}\widehat{m},\widehat{\omega }=m\omega \omega ^3.$$
After these identifications, we work on the sixteen–dimensional manifold
$$=\times ^2\times (\mathrm{SO}(3))^2\times (^3)^2,$$
which is diffeomorphic to $`T^{}\left(𝒜\times \mathrm{SO}(3)\times \mathrm{SO}(3)\right)`$. We shall denote by $`(b,c,U,m)`$ the elements of $``$, with $`b`$, $`c^2`$, $`U=(U_l,U_r)(\mathrm{SO}(3)^2)`$, and $`m=(m_l,m_r)(^3)^2`$ and use the induced scalar and vector products on $`(^3)^2`$ given by
$$(\omega _l,\omega _r)(\xi _l,\xi _r)=\omega _l\xi _l+\omega _r\xi _r,(\omega _l,\omega _r)\times (\xi _l,\xi _r)=(\omega _l\times \xi _l,\omega _r\times \xi _r).$$
When writing tangent vectors to $``$ we shall again left–trivialize and identify the tangent spaces to the two factors $`\mathrm{SO}(3)`$ with $`^3`$; thus, $`(v,w,\omega ,n)`$ denotes the vector $`(v,w,(U_l\widehat{\omega }_l,U_r\widehat{\omega }_r),n)T_{(b,c,U,m)}`$. Finally, in order to avoid the introduction of a new symbol, we write $`𝒱(b)`$ to denote $`𝒱(b_1,b_2,\frac{1}{b_1b_2})`$.
###### Proposition 2
Riemann’s equation (2.3) on $`𝒜\times \mathrm{SO}(3)\times \mathrm{SO}(3)`$ is equivalent to the Hamiltonian system on the manifold $``$ defined by the Hamiltonian
$$H(b,c,U,m)=\frac{1}{2}c𝕂(b)c+\frac{1}{2}m𝕁(b)m+𝒱(b),$$
(2.5)
where
$`𝕂(b)`$ $`=`$ $`{\displaystyle \frac{1}{b_2^2b_3^2+b_1^2b_3^2+b_1^2b_2^2}}\left(\begin{array}{cc}b_1^2\left(b_2^2+b_3^2\right)& b_3\\ b_3& b_2^2\left(b_1^2+b_3^2\right)\end{array}\right)`$
$`𝕁(b)`$ $`=`$ $`\left(\begin{array}{cc}J_1(b)& J_2(b)\\ J_2(b)& J_1(b)\end{array}\right)\mathrm{with}\{\begin{array}{cc}& J_1(b)=\mathrm{diag}(\frac{b_2^2+b_3^2}{(b_2^2b_3^2)^2},\frac{b_1^2+b_3^2}{(b_1^2b_3^2)^2},\frac{b_1^2+b_2^2}{(b_1^2b_2^2)^2})\hfill \\ & J_2(b)=\mathrm{diag}(\frac{2b_2b_3}{(b_2^2b_3^2)^2},\frac{2b_1b_3}{(b_1^2b_3^2)^2},\frac{2b_1b_2}{(b_1^2b_2^2)^2})\text{ ,}\hfill \end{array}`$
and the “left–trivialized” (in the sense just specified) symplectic structure $`\sigma `$ is defined by
$$\sigma (b,c,U,m)((v,w,\omega ,n),(v^{},w^{},\omega ^{},n^{}))=vw^{}v^{}w+\omega n^{}n\omega ^{}+m(\omega \times \omega ^{}).$$
(2.6)
This proposition can be proven by a direct computation showing that Hamilton’s equations for the Hamiltonian system defined by (2.5) and (2.6) are equivalent to Riemann’s equations. Alternatively, one can observe that Riemann’s equations are equivalent to the Lagrangian system defined on $`T𝒜\times T\mathrm{SO}(3)\times T\mathrm{SO}(3)`$ by the Lagrangian
$$((A,\dot{A}),U_l\widehat{\omega }_l,U_r\widehat{\omega }_r)=\frac{1}{2}\left|\dot{A}+\widehat{\omega }_lAA\widehat{\omega }_r\right|^2𝒱(A)$$
(2.7)
and construct the Hamiltonian system using the Legendre transform. For the sake of brevity, we skip all details, remarking only that the relation between ‘angular momenta’ and ‘angular velocities’ is $`\omega =𝕁(b)m`$.
Remark: The Lagrangian nature of Dirichlet’s problem, with Lagrangian equal to the difference of the kinetic and the potential energies, is implicit in Riemann’s work and has been repeatedly used since then. One can verify by a direct computation that the Euler–Lagrange equations for the Lagrangian (2.7) are equivalent to Riemann’s equations (2.3).
The Hamiltonian system $`(,H,\sigma )`$ of Proposition 2 is invariant under the symplectic action $`\mathrm{\Phi }`$ of $`\mathrm{SO}(3)\times \mathrm{SO}(3)`$ on $``$ given by
$$\mathrm{\Phi }_{(R_l,R_r)}(b,c,(U_l,U_r),m)=(b,c,(R_lU_l,R_rU_r),m).$$
Since the group does not act on the factor $`\times ^2`$, while it acts on each factor $`\mathrm{SO}(3)\times ^3`$ as the left–trivialized cotangent lift of the left action of $`\mathrm{SO}(3)`$ on itself, all results from the elementary case of the left action of $`\mathrm{SO}(3)`$ on $`T^{}\mathrm{SO}(3)`$ apply (see, e.g., or ). Thus, one concludes that the action $`\mathrm{\Phi }`$ has the momentum map
$$𝐉:(^3)^2,(a,c,U,m)(U_lm_l,U_rm_r).$$
For any $`\eta (^3)^2`$, one may identify $`𝐉^1(\eta )`$ and $`\times ^2\times (\mathrm{SO}(3))^2`$, in which case the immersion $`i_\eta :𝐉^1(\eta )`$ is given by $`i_\eta (a,c,U)=(a,c,U,(U_l^T\eta _l,U_r^T\eta _r))`$. The two components $`U_lm_l`$ and $`U_rm_r`$ of the momentum map coincide, up to constant rescalings, with the total angular momentum and the circulation; see Section 2.A.
The next Proposition, which describes the reduced system $`(P_\eta ,\sigma _\eta ,H_\eta )`$, also follows immediately from standard results for the reduction of an $`SO(3)`$ invariant Hamiltonian system on $`T^{}\mathrm{SO}(3)`$. We denote by $`S_\rho ^2`$ the sphere of radius $`\rho `$. As usual, $`G_\eta `$ denotes the isotropy subgroup of $`\eta `$ with respect to the coadjoint action; in the case at hand, $`G_\eta =\{\text{rotations about }\eta _l\}\times \{\text{rotations about }\eta _r\}`$.
###### Proposition 3
(i) For any $`\eta =(\eta _l,\eta _r)(^3)^2`$ with $`\eta _l0\eta _r`$, the reduced phase space $`𝐉^1(\eta )/G_\eta `$ can be identified with
$$P_\eta =\times ^2\times \left(S_{\eta _l}^2\times S_{\eta _r}^2\right).$$
After this identification, the canonical projection $`\pi _\eta :𝐉^1(\eta )P_\eta `$ is $`\pi _\eta (b,c,U)=(b,c,(U_l^T\eta _l,U_r^T\eta _r))`$. The symplectic structure $`\sigma _\eta `$ on $`P_\eta `$ defined by the equality $`i_\eta ^{}\sigma =\pi _\eta ^{}\sigma _\eta `$ is
$$\sigma _\eta (b,c,m)((v,w,m\times \omega ),(v^{},w^{},m\times \omega ^{}))=vw^{}v^{}wm\omega \times \omega ^{}.$$
The reduced Hamiltonian $`H_\eta :P_\eta `$, defined by $`H_\eta \pi _\eta =Hi_\eta `$, is
$$H_\eta (b,c,m)=\frac{1}{2}c𝕂(b)c+\frac{1}{2}m𝕁(b)m+𝒱(b).$$
(2.8)
(ii) If $`\eta =(\eta _l,0)`$ or $`\eta =(0,\eta _l)`$ with $`\eta _l0`$, then the reduced phase space can be identified with $`\times ^2\times S_{\eta _l}^2`$ and one has $`\pi _\eta (b,c,U_l)=(b,c,U_l^T\eta _l)`$ and
$`\sigma _\eta (b,c,m_l)((v,w,m_l\times \omega _l),(v^{},w^{},m_l\times \omega _l^{}))=vw^{}v^{}wm_l\omega _l\times \omega _l^{}`$
$`H_\eta (b,c,m_l)=\frac{1}{2}c𝕂(b)c+\frac{1}{2}m_lJ_1(b)m_l+𝒱(b).`$
The reduced system generically has four degrees of freedom, but only three in the special case (ii), which arises in the case of the so–called “irrotational” Riemann ellipsoids. The case $`\eta _l=\eta _r=0`$ yields a reduced system with two degrees of freedom, the only equilibrium of which is the stationary sphere.
In view of the application of Nekhoroshev theory, we note that the reduced Hamiltonian $`H_\eta `$ is an analytic function. This is seen by observing that the self–gravitational potential can be written in the form
$$𝒱(b)=\frac{4\pi g}{\sqrt{b_1^2b_3^2}}F\left(\mathrm{arccos}\left(\frac{b_3}{b_1}\right)|\sqrt{\frac{b_1^2b_2^2}{b_1^2b_3^2}}\right)$$
where $`F(\phi |k)`$ is the incomplete elliptic integral of the first kind, which is analytic for $`\phi (0,\frac{\pi }{2})`$ and $`k(0,1)`$ (see , pg. 4–5, 299).
D. Equivalence and symmetry of reduced systems. As we now discuss, (i) certain classes of reduced systems are equivalent and (ii) the reduced systems have certain discrete symmetries. For the sake of conciseness, we consider here only the generic case, in which both $`\eta _l`$ and $`\eta _r`$ are nonzero. (The case in which one of the momenta vanishes is recovered with obvious modifications.) There are two distinct reasons for equivalence of different reduced systems:
* The reduced system $`(P_\eta ,H_\eta ,\sigma _\eta )`$ depends on $`\eta =(\eta _l,\eta _r)`$ only through the norms of $`\eta _l`$ and $`\eta _r`$. This equivalence reflects the invariance of the Dirichlet problem under spatial rotation.
* Dedekind proved that if $`F(t)`$ is a solution of the Dirichlet problem, then $`F(t)^T`$ is also a solution . (The proof follows immediately from the observations that $`_{F^T}G^T=(_FG)^T`$ for any matrix $`G`$ and that $`U_lAU_r^T`$ is a singular value decomposition of $`F`$ iff $`U_rAU_l^T`$ is a singular value decomposition of $`F^T`$.) Chandrasekhar calls ‘adjoint’ the two motions determined by $`F(t)`$ and $`F(t)^T`$. This invariance of the Dirichlet problem has a counterpart in the Hamiltonian formulation on the covering space $``$, which is invariant under the symplectomorphism exchanging $`U_l`$ with $`U_r`$ and $`m_l`$ with $`m_r`$. This map takes solutions to solutions, but exchanges the level sets of $`𝐉`$. Correspondingly, the reduced systems for the values $`(\eta _l,\eta _r)`$ and $`(\eta _r,\eta _l)`$ of the momentum map are conjugated to one another by the diffeomorphism $`(b,c,(m_l,m_r))(b,c,(m_r,m_l))`$.
Furthermore, each reduced system is invariant under two discrete actions:
* The $`_2`$–action $`(b,c,m)(b,c,m)`$, which reflects the invariance of the original system with respect to the combination of spatial reflection and time reversal.
* The $`_4`$–action $`(b,c,m)(b,c,R_im)`$, $`i=0,\mathrm{},3`$, where $`R_0,\mathrm{},R_3`$ are given by (2.4) and $`R_i(m_l,m_r)=(R_im_l,R_im_r)`$. The points $`(b,c,R_im)`$ are obtained from one another by a reflection about one and the same coordinate axis on the two spheres; see figure 1. Hence the $`_4`$–orbit of the point $`(b,c,m)`$ consists of two distinct points if $`m`$ is parallel to a basis vector $`(e_j,e_j)`$, $`j=1,2,3`$, and of four distinct points otherwise.
In the stability analysis, we need not distinguish equivalent reduced systems, nor distinct points within the $`_4\times _2`$ orbits, in each reduced system. A $`_4\times _2`$–orbit consists of eight, four, or two points, depending on the number of nonzero components of $`m_l`$ and $`m_r`$. We will say that all points in such an orbit are $`_4\times _2`$–equivalent and regard two points in the reduced space as adjoint if any two points in their $`_4\times _2`$–orbits are adjoint.
By a Riemann ellipsoid we will mean an equivalence class, under spatial rotations (E1) and transposition (E2), of motions of the Dirichlet problem corresponding to an equilibrium of the reduced system. Thus each Riemann ellipsoid can be identified with the $`_4\times _2`$–orbit of an equilibrium of some reduced system.<sup>2</sup><sup>2</sup>2Note that, as was previously stated, we exclude from our consideration all axisymmetric ellipsoids.
Remark: The $`_4`$–symmetry is not present in Dirichlet’s system (2.2), but is clearly inherited from the passage to the fourfold covering manifold $``$ of $`T^{}Q`$. (The existence of this symmetry of Riemann’s equation is mentioned by Chandrasekhar (, pg. 72), but its origin is not identified there.) Since all points in a $`_4`$–orbit correspond to the same state of the fluid, the passage to the covering introduces some redundancy, but allows us to avoid singular reduction. In fact, Dirichlet’s equations (2.2) determine a $`\mathrm{SO}(3)\times \mathrm{SO}(3)`$–invariant Hamiltonian system on $`T^{}Q`$; however, the action is not free and the reduced phase space is the orbifold $`\times ^2\times (S^2\times S^2)/^4`$, rather than a smooth manifold. The singular set consists of all states in which $`m_l`$ and $`m_r`$ are both parallel to the same axis of the ellipsoid, and thus contains the Riemann ellipsoids of type S.
## 3 Riemann Ellipsoids
We now study the equilibria of the reduced systems of Proposition 3, and thus the Riemann ellipsoids. The results of this analysis are due to Riemann. As we mentioned in the Introduction, there are five classes of Riemann ellipsoids and the momenta $`m_l`$ and $`m_r`$ of the ellipsoids of each type are completely determined by their semiaxes. Correspondingly, in the set $``$ there are “existence regions” of the Riemann ellipsoids of each type. In order to specify these regions, and the momenta of the corresponding Riemann ellipsoids, we introduce some notation. First we set
$$C_n(x,y,z)=\mathrm{\hspace{0.33em}2}\pi g_0^{\mathrm{}}\left[(s+x^2)(s+y^2)(s+z^2)\right]^{3/2}s^n𝑑s,n=0,1,2.$$
The functions
$$G_\pm ^S(x,y,z)=\frac{(xz)^4}{xz}\left[\left[xy^2z\pm (x^2y^2x^2z^2+y^2z^2)\right]C_1(x,y,z)+(xz\pm y^2)C_2(x,y,z)\right]$$
determine the momenta of the ellipsoids of type S, while the functions
$`G(x,y,z)`$ $`=`$ $`x^2\left(y^2z^2\right)C_1(x,y,z)+\left(y^24z^2\right)\left[z^2C_1(x,y,z)+C_2(x,y,z)\right]`$
$`D(x,y,z)`$ $`=`$ $`x^2(y^2z^2)+z^2(4z^2y^2)`$
$`G_\pm ^R(x,y,z)`$ $`=`$ $`(yz)^4\left(x^2(y\pm 2z)^2\right){\displaystyle \frac{x^2z^2}{x^2y^2}}{\displaystyle \frac{G(x,y,z)}{D(x,y,z)}}`$
determine the momenta of the ellipsoids of type I, II and III; we refer to the latter three classes of ellipsoids collectively as type R ellipsoids.<sup>3</sup><sup>3</sup>3In Riemann’s article, the study of the existence of these ellipsoids precedes that of the ellipsoids of type S. Using these functions we define the five sets
$`_{\mathrm{S}_2}`$ $`=`$ $`\{b:G_{}^S(b_1,b_2,b_3)0\}`$
$`_{\mathrm{S}_3}`$ $`=`$ $`\{b:G_+^S(b_1,b_3,b_2)0\}`$
$`_\mathrm{I}`$ $`=`$ $`\{b:b_12b_2b_3\}`$ (3.1)
$`_{\mathrm{II}}`$ $`=`$ $`\{b:b_12b_2+b_3,D(b_1,b_3,b_2)<0\}`$
$`_{\mathrm{III}}`$ $`=`$ $`\{b:b_1b_2+2b_3,G(b_1,b_2,b_3)>0\}`$
and five pairs of real–valued maps $`\mu _\alpha ^\pm `$, one for each type $`\alpha `$ of ellipsoids; these maps and their domains are described in Table 1, where $`\{e_1,e_2,e_3\}`$ denotes the standard basis vector of $`^3`$ and
$$N_\pm ^\tau (x,y,z)=\frac{1}{2}\left[\sqrt{G_+^\tau (x,y,z)}\pm \sqrt{G_{}^\tau (x,y,z)}\right]\tau =R,S.$$
| Type ($`\alpha `$) | Domain | $`\mu _\alpha ^\pm (b)`$ |
| --- | --- | --- |
| $`\mathrm{S}_2`$ | $`_{\mathrm{S}_2}`$ | $`N_\pm ^S(b_1,b_2,b_3)e_2`$ |
| $`\mathrm{S}_3`$ | $`_{\mathrm{S}_3}`$ | $`N_\pm ^S(b_1,b_3,b_2)e_3`$ |
| I | $`_\mathrm{I}`$ | $`N_\pm ^R(b_1,b_3,b_2)e_1+N_\pm ^R(b_3,b_1,b_2)e_3`$ |
| II | $`_{\mathrm{II}}`$ | $`N_\pm ^R(b_1,b_3,b_2)e_1+N_{}^R(b_3,b_1,b_2)e_3`$ |
| III | $`_{\mathrm{III}}`$ | $`N_\pm ^R(b_1,b_2,b_3)e_1+N_{}^R(b_2,b_1,b_3)e_2`$ |
Table 1: The momenta of the Riemann Ellipsoids
###### Proposition 4
(i) A point $`(b,c,m)P_\eta `$ is an equilibrium of the reduced system $`(P_\eta ,H_\eta ,\sigma _\eta )`$, $`\eta _l0\eta _r`$, iff it is $`_4\times _2`$–equivalent to either the point $`(b,0,(\mu _\alpha ^+(b),\mu _\alpha ^{}(b)))`$ or the point $`(b,0,(\mu _\alpha ^{}(b),\mu _\alpha ^+(b)))`$ for some $`b_\alpha `$ and some $`\alpha =\mathrm{S}_2,\mathrm{S}_3,\mathrm{I},\mathrm{II},\mathrm{III}`$.
(ii) If one component of $`\eta `$ equals zero, then the equilibria of the six–dimensional reduced system are the points $`_4\times _2`$–equivalent to $`(b,0,\mu _\alpha ^+(b))`$. This situation is met only at the equilibria of type $`\mathrm{S}_2`$ satisfying $`G_+^S(b_1,b_2,b_3)=G_{}^S(b_1,b_2,b_3)`$ and those of type I satisfying $`b_1^2b_2^2+b_1^2b_3^2+b_2^2b_3^23b_2^4=0`$; in both cases, $`\mu _\alpha ^+(b)0=\mu _\alpha ^{}(b)`$.
The proof of this proposition can be reconstructed from Riemann’s and Chandrasekhar’s work. A sketch of the argument is given in Appendix A.
The existence regions (3.1) are drawn in figures 2 using the coordinates $`(x,y)=(b_2/b_1,b_3/b_1)`$, which have the advantage that the unbounded region $``$ is mapped onto the bounded triangle $`\{(x,y):\mathrm{\hspace{0.17em}0}<y<x<1\}`$. The borders of the various regions are determined by the definitions (3.1). Note that, since the existence regions of ellipsoids of different types have nonempty intersections, there exist Riemann ellipsoids of up to four different types having equal semiaxes.
The Riemann ellipsoids belonging to a six–dimensional reduced space have either zero angular momentum or zero circulation and are called “irrotational” by Chandrasekhar. The existence of irrotational ellipsoids of type $`\mathrm{S}_2`$ was known to Riemann, while the existence of those of type I seems not to have been noticed before. Note that the curve of irrotational ellipsoids of type $`\mathrm{S}_2`$ divides the region $`_{\mathrm{S}_2}`$ into two parts, in one of which the two momenta are coparallel, while in the other they are counterparallel.<sup>4</sup><sup>4</sup>4We say that two vectors $`u`$ and $`v`$ are coparallel (counterparallel) if $`uv=uv`$ ($`=uv`$).
Remark: Since Chandrasekhar requires that the relative pressure be positive inside the ellipsoids, his region of existence of the ellipsoids of type II is smaller than ours, lying below the “zero–pressure curve” $`D(b_1,b_3,b_2)C_0+6b_2^2C_1+3C_2=0`$ shown in figure 2.c (see , chapter 7, formula (195)).
## 4 Stability and Ellipticity of the Riemann Ellipsoids
A. Results. We now review the known results about the Lyapunov stability of the Riemann ellipsoids, regarded here as equilibria in the reduced phase space, and we describe the ellipticity analysis. From now on, we restrict ourselves to the generic case in which the reduced system has four degrees of freedom; the irrotational Riemann ellipsoids will be studied elsewhere.
###### Proposition 5
(i) All $`\mathrm{S}_3`$–ellipsoids and all counterparallel $`\mathrm{S}_2`$–ellipsoids are Lyapunov stable.
(ii) The coparallel $`\mathrm{S}_2`$–ellipsoids for $`b_{\mathrm{S}_2}\mathrm{Int}(_{\mathrm{III}})`$ (i.e., $`G(b_1,b_2,b_3)>0`$) are not elliptic and hence are Lyapunov unstable.
The first statement is due to Riemann , who proved it using arguments essentially equivalent to using the reduced Hamiltonian as a Lyapunov function. Riemann also showed that the reduced Hamiltonian has a saddle point, and thus cannot be used as a Lyapunov function, at any other ellipsoid. Riemann interpreted this fact as an indication that these ellipsoids are unstable (“labil”). However, as best we know this fact is unproven, except in the case of the $`\mathrm{S}_2`$ ellipsoids considered in statement (ii), the instability of which follows from Chandrasekhar’s work, since he showed that the linearization of the equations of motion has an eigenvalue with nonzero real part at these points; see section 4.B.
Remark: Given the structure of the problem, the Lyapunov stability of the equilibria of the reduced system corresponds to orbital stability of the solutions of the Dirichlet problem (see e.g. , , ).
The ellipticity of the Riemann ellipsoids has been studied by Chandrasekhar; most of his results are based on numeric calculations. We have repeated the numerical analysis of the ellipticity, finding the same results as Chandrasekhar for the $`\mathrm{S}_2`$ ellipsoids and essentially the same results for the ellipsoids of type I, but significantly different results for the ellipsoids of types II and III. The conclusions of our investigation are the following:
Numerical Conclusions 1: (i) The coparallel $`\mathrm{S}_2`$–ellipsoids for $`b`$ in the complement of $`_{\mathrm{S}_2}\mathrm{Int}(_{\mathrm{III}})`$, i.e. $`b`$ satisfying $`G(b_1,b_2,b_3)0`$, are elliptic.
(ii) There are nonempty subregions of $`_\mathrm{I}`$, $`_{\mathrm{II}}`$ and $`_{\mathrm{III}}`$ consisting of elliptic ellipsoids; these regions are the shaded regions in figures 3 and 4. The ellipsoids of types I, II and III in the complement of these subregions are not elliptic and hence are Lyapunov unstable.
This analysis has been carried out by analytically constructing the linearization of the reduced Hamiltonian vector field at the reduced equilibria and by numerically evaluating this matrix at some tens of thousands of ellipsoids of each type. Details about this analysis are given in the next subsection. (As explained there, Riemann sketches an argument that would imply statement (i).) Here we comment on these results.
$``$ The elliptic $`\mathrm{S}_2`$ ellipsoids of unknown Lyapunov stability lie between the “irrotational” curve and the border of region $`_{\mathrm{III}}`$. This is the dark shaded region in figure 3.b.
$``$ Our computations confirm the existence of two ellipticity regions for the ellipsoids of type I, as shown in figure 3 of , but disclose a finer structure of the lower region than that found by Chandrasekhar: the upper border of this region appears to transversally intersect the border $`b_1=2b_2b_3`$ of the existence region at about $`b_2/b_1=.515`$, but there is a very narrow crescent of nonelliptic ellipsoids along this line, approximately in the range $`.5<b_2/b_1<.503`$, which is shown in figure 4.a. We could not determine with certainty if this crescent actually touches the line $`b_1=2b_2b_3`$.
$``$ The region of elliptic ellipsoids of type II consists of three narrow fringes, shown enlarged in figure 4.b,c,d. None of these regions were detected by Chandrasekhar, who stated that no ellipsoid of type II is elliptic (, pg. 169). Note that two of these fringes belong (for the most part) to the region of positive relative pressure considered by Chandrasekhar.
$``$ The region of ellipticity of the ellipsoids of type III has a three–lobed structure, the three lobes being separated by two very thin gaps; see figure 4.e, f. Chandrasekhar detected only the upper of these three lobes. To within the precision of our numerical computations, the two gaps close at approximately $`b_2/b_1=.01`$, but it is possible that, in fact, they persist to the origin.
In an attempt to understand the discrepancies between our results and Chandrasekhar’s, we tested our calculations using his linearization of the equations of motion, finding complete agreement. It is therefore possible that Chandrasekhar’s conclusions were simply based on too few sample cases to obtain a detailed global picture of the structure of the regions of ellipticity. A more detailed comparison with Chandrasekhar’s result is given in Appendix B.
We conclude this analysis by observing that the sources of the richness and intricacy of the structure of the ellipticity regions resulting from these numerical computations are not yet understood. Of course, we cannot be sure that our numerical computations, even though accurate, ultimately detected the finest details of this structure, particularly in view of the nontrivial structure of the ellipticity regions for ellipsoids approaching a degenerate (flat or rodlike) configuration. Many features remain to be understood. For instance, figures 3.a and 4.f suggest that there are (nontrivial) relations between the ellipticity regions of the ellipsoids I (resp. III) and the boundary of the existence region of the ellipsoids $`\mathrm{S}_2`$ (resp. $`\mathrm{S}_3`$), which could be the object of a bifurcation analysis. Additional analytical insight into the problem is certainly needed.
Remarks: (i) In his articles, Chandrasekhar refers to ellipticity as “stability”, while Riemann designated by this name what is now called Lyapunov stability. Because of this, and the previously mentioned fact that Riemann regarded as unstable the ellipsoids which are saddle points of the Lyapunov function, Riemann and Chandrasekhar found different “stability” and “instability” regions. Chandrasekhar interpreted this difference as being due to some errors in Riemann’s work (see page 875 in and pages 185–7 in ). Riemann’s results are fully confirmed by our numerical computations (with the exception of the characterization of the saddle points as unstable, of course) and are, in fact, compatible with Chandrasekhar’s results once the distinction between Lyapunov stability and ellipticity is taken into account.
(ii) In the quoted references, Chandrasekhar refers to for an explanation of “the origin of his \[Riemann’s\] errors”. In that article, to the best of our understanding, the source of these supposed errors is ascribed to the fact that Riemann treated the Dirichlet problem as a Lagrangian system (which is correct, as we have noted in Section 2), while according to the system is Lagrangian only in the case of the ellipsoids of type S; remark (vi) in page 185 of seems to indicate that Chandrasekhar shared this opinion.
B. Numerical Computations. We now briefly describe the numerical procedures used in the verification of Numerical Conclusions 1 and in the Nekhoroshev analysis. We also provide the key components of the proof of statement (ii) of Proposition 5. We use canonical coordinates obtained by introducing a pair of local “Poincaré coordinates” $`(q,p)`$ on each sphere in the reduced phase space; these coordinates coincide with the Poincaré elements used in the Kepler problem. The use of coordinates is not strictly necessary for the ellipticity analysis, but is very useful for the construction of the normal forms used in the analysis of the Nekhoroshev stability.
###### Lemma 1
For any $`\rho >0`$, the map
$$(q,p)\stackrel{~}{m}_\rho (q,p):=(p\sqrt{\rho \frac{q^2+p^2}{4}},q\sqrt{\rho \frac{q^2+p^2}{4}},\rho \frac{q^2+p^2}{2})$$
(4.1)
is a diffeomorphism from the disk $`\{(q,p)^2:q^2+p^2<2\rho \}`$ onto the half–sphere $`\{m^3:m=\rho ,m_3>0\}`$. It is symplectic in the sense that $`(\stackrel{~}{m}_\rho )_{}\left(dpdq\right)(m\times \omega ,m\times \omega ^{})=m\omega \times \omega ^{}`$ for all $`m`$ in the half sphere and all $`\omega ,\omega ^{}^3`$.
Proof. The map (4.1) is smooth and onto; its inverse $`m\sqrt{2/(\rho +m_3)}(m_2,m_1)`$ is also smooth. If $`v=(v_1,v_2,v_3)`$ is tangent to the sphere at the point $`m=\stackrel{~}{m}_\rho (q,p)`$, then $`v=v_q\frac{\stackrel{~}{m}}{q}+v_p\frac{\stackrel{~}{m}}{p}`$ for some constants $`v_q`$ and $`v_p`$; specifically, if we set $`k_1=2\rho q^2p^2`$ and $`k_2=k_1+2\rho `$, then $`k_1\sqrt{k_2}v_q=qpv_1(k_2p^2)v_2`$ and $`k_1\sqrt{k_2}v_p=(k_2q^2)v_1qpv_2`$. Thus $`(\stackrel{~}{m}_\rho )_{}\left(dpdq\right)(m\times \omega ,m\times \omega ^{})=(m\times \omega )_p(m\times \omega ^{})_q(m\times \omega )_q(m\times \omega ^{})_p=m(\omega \times \omega ^{})`$ for any $`\omega ,\omega ^{}^3`$.
When studying a given equilibrium $`(b^{},0,M^{})`$, with $`M^{}=(m_l^{},m_r^{})`$, we shall translate the origin of the coordinates $`b`$ to the equilibrium value $`b^{}`$; for simplicity, we will retain the notation $`b`$ for the translated coordinates. Moreover, we shall use ‘rotated’ Poincaré coordinates $`(q_1,p_1)`$ and $`(q_2,p_2)`$ centered at $`m_l^{}`$ and $`m_r^{}`$ on the spheres; specifically, for given $`m^{}^3`$ we define
$$(q,p)\stackrel{~}{m}(q,p;m^{}):=R_m^{}\stackrel{~}{m}_m^{}(q,p)$$
where $`R_m^{}`$ is an orthogonal matrix such that $`R_m^{}e_3=\frac{m^{}}{m^{}}`$. Since rotations are symplectic, the coordinates $`\stackrel{~}{m}(;m^{})`$ are symplectic. These coordinates are not uniquely defined: right–multiplication of $`R_m^{}`$ by a rotation about $`e_3`$ produces a new system of rotated Poincaré coordinates centered at $`m^{}`$. To compress the notation, we write $`(q,p)`$ for $`(q_1,q_2,p_1,p_2)`$ and<sup>5</sup><sup>5</sup>5In what follows, to keep the notation simple, we will write all vectors as row vectors, even though they must be regarded as column vectors in all matrix formulas.
$$M(q,p;M^{})=(\stackrel{~}{m}(q_1,p_1;m_l^{}),\stackrel{~}{m}(q_2,p_2;m_r^{})).$$
Thus, the reduced Hamiltonian (2.8) takes the form
$$(b,c,q,p;b^{},M^{})=\frac{1}{2}c𝕂(b^{}+b)c+\frac{1}{2}M(q,p;M^{})𝕁(b^{}+b)M(q,p;M^{})+𝒱(b^{}+b).$$
We now study the block structure of the Hessian matrix $`^{\prime \prime }(0;b^{},M^{})`$. To obtain a Hessian with as few nontrivial blocks as possible, we construct the coordinates using matrices such that $`R_{m_l^{}}e_1=R_{m_r^{}}e_1`$ is parallel to a principal axis orthogonal to $`\text{span}[m_l,m_r]`$. Using this choice of coordinates and denoting by $`_{bb}^{\prime \prime }`$ the two–by–two block $`\frac{^2}{b_ib_j}`$, etc, we have
###### Lemma 2
For all equilibria, the blocks $`_{cb}^{\prime \prime }`$, $`_{cq}^{\prime \prime }`$, $`_{cp}^{\prime \prime }`$, $`_{bp}^{\prime \prime }`$, and $`_{qp}^{\prime \prime }`$ equal zero and $`_{cc}^{\prime \prime }=𝕂`$. Moreover, $`_{bq}^{\prime \prime }=0`$ for the equilibria of type S.
Proof. The equalities $`_{cb}^{\prime \prime }=_{cq}^{\prime \prime }=_{cp}^{\prime \prime }=0`$ and $`_{cc}^{\prime \prime }=𝕂`$ are obvious. Note that the ‘standard’ Poincaré parameterization (4.1) of the sphere satisfies
$$\frac{\stackrel{~}{m}}{q}(0,0;\rho )=\sqrt{\rho }e_2,\frac{\stackrel{~}{m}}{p}(0,0;\rho )=\sqrt{\rho }e_1,\frac{^2\stackrel{~}{m}}{qp}(0,0;\rho )=\mathrm{\hspace{0.33em}0}.$$
(4.2)
Using the latter equality and denoting by $`\frac{\stackrel{~}{M}}{q_j}`$ (resp. $`\frac{𝕁}{q_j}`$) the vector (resp. matrix) with entries given by the derivatives with respect to $`q_j`$ of the entries of $`\stackrel{~}{M}`$ (resp. $`𝕁`$), etc., we have
$`{\displaystyle \frac{^2}{b_ip_j}}(0;b^{},M^{})`$ $`=`$ $`R{\displaystyle \frac{\stackrel{~}{M}}{p_j}}(0,0){\displaystyle \frac{𝕁}{b_i}}(b^{})M^{}`$
$`{\displaystyle \frac{^2}{q_ip_j}}(0;b^{},M^{})`$ $`=`$ $`R{\displaystyle \frac{\stackrel{~}{M}}{q_i}}(0,0)𝕁(b^{})R{\displaystyle \frac{\stackrel{~}{M}}{p_j}}(0,0)`$ (4.3)
$`{\displaystyle \frac{^2}{b_iq_j}}(0;b^{},M^{})`$ $`=`$ $`R{\displaystyle \frac{\stackrel{~}{M}}{q_j}}(0,0){\displaystyle \frac{𝕁}{b_i}}(b^{})M^{},`$
where $`R`$ is the six–by–six block diagonal matrix with diagonal blocks $`R_{m_l^{}}`$ and $`R_{m_r^{}}`$.
Due to the block structure of $`𝕁`$, the two ‘components’ of $`\frac{𝕁}{b_i}(b^{})M^{}`$ lie in $`\text{span}[m_l,m_r]`$. (Here we refer to $`m_l`$ and $`m_r`$ as the ‘components’ of $`(m_l,m_r)`$.) On the other hand, (4.2) and the condition that $`R_{m_l^{}}e_1=R_{m_r^{}}e_1`$ be orthogonal to $`\text{span}[m_l,m_r]`$ imply that the components of $`R\frac{\stackrel{~}{M}}{p_j}(0,0)`$ are orthogonal to $`\text{span}[m_l,m_r]`$ for $`j=1,2`$; hence (4.3) implies that $`\frac{^2}{b_ip_j}(0;b^{},M^{})=0`$. To see that $`\frac{^2}{q_ip_j}(0;b^{},M^{})`$ equals zero, note that $`R_{m_l^{}}e_1=R_{m_r^{}}e_1`$ is, by construction, an eigenvector of $`J_1(b^{})`$ and $`J_2(b^{})`$; hence (4.2) implies that the components of $`R^T𝕁(b^{})R\frac{\stackrel{~}{M}}{p_j}(0,0)`$ are orthogonal to the components of $`\frac{\stackrel{~}{M}}{q_j}(0,0)`$.
In the case of the S ellipsoids, the two components of $`M^{}`$ are both parallel to the same coordinate axis. Due to the block structure of $`𝕁`$, the components of $`\frac{𝕁}{b_i}(b^{})M^{}`$ are parallel to the components of $`M^{}`$; thus, both components of $`R^T\frac{𝕁}{b_i}M^{}`$ are parallel to $`e_3`$. By (4.2) and (4.3), this implies $`_{bq}^{\prime \prime }=0`$.
Let us now consider the ellipsoids of type $`\mathrm{S}_2`$ . With a suitable choice of the coordinates, the Hessian has the diagonal block structure
$$^{\prime \prime }(0;b^{},M^{})=\mathrm{diag}[_{bb}^{\prime \prime },𝕂,_{qq}^{\prime \prime },_{pp}^{\prime \prime }].$$
We limit ourselves to stating without proof the following facts, which can be proven using elementary, though sometimes laborious, computations:
* The blocks $`𝕂`$ and $`_{qq}^{\prime \prime }`$ are positive definite for all the $`\mathrm{S}_2`$–ellipsoids.
* The block $`_{pp}^{\prime \prime }`$ is positive definite for all the counterparallel $`\mathrm{S}_2`$–ellipsoids and is indefinite for all the coparallel $`\mathrm{S}_2`$–ellipsoid.
The block $`_{bb}^{\prime \prime }`$ has a relatively simple expression, but the analytical study of its definiteness is quite laborious. We verified numerically that it is positive definite at all ellipsoids of type $`\mathrm{S}_2`$ . (Chandrasekhar’s study of this block in is also numerical.) The linearization at the equilibrium of the Hamiltonian vector field of $``$ is $`X(b^{},M^{})(b,c,q,p)`$, where the matrix $`X(b^{},M^{})`$ is the product of the symplectic matrix and $`^{\prime \prime }(0;b^{},M^{})`$. Thus, at a type $`\mathrm{S}_2`$ ellipsoid, the matrix $`X(b^{},M^{})`$ is block diagonal, with the two $`4\times 4`$ blocks
$$X_{bc}=\left(\begin{array}{cc}0& 𝕂\\ _{bb}^{\prime \prime }& 0\end{array}\right)\mathrm{and}X_{qp}=\left(\begin{array}{cc}0& _{pp}^{\prime \prime }\\ _{qq}^{\prime \prime }& 0\end{array}\right).$$
Since the eigenvalues of each of these two matrices are purely imaginary if its two nonzero blocks are positive semi–definite, we arrive to statement (i) of Numerical Conclusions 1. This conclusion is not rigorous because it is based on the numerical verification of the definiteness of $`_{bb}^{\prime \prime }`$. Statement (ii) of Proposition 5 follows from the definiteness of $`_{qq}^{\prime \prime }`$ and the indefiniteness of $`_{pp}^{\prime \prime }`$.
Remark: In , Riemann suggests an analytic argument that would show that for all ellipsoids of type S satisfying $`m_lm_r`$, i.e. $`G_{}^S(b_1,b_2,b_3)>0`$, the function $`b(b,0,0,0;b^{},M^{})`$ has a minimum at $`b=0`$. Since the block $`_{bb}^{\prime \prime }`$ is positive semi–definite at such a minimum, this result could be used to construct an analytic proof of statement (i) of Numerical Conclusions 1. However, as Riemann explicitly states, he does not actually carry out the minimality argument in .
We now consider the ellipsoids of types I, II and III. The Hessian has the form
$$^{\prime \prime }(0;b^{},M^{})=\left(\begin{array}{cccc}_{bb}^{\prime \prime }& 0& _{bq}^{\prime \prime }& 0\\ 0& 𝕂& 0& 0\\ _{qb}^{\prime \prime }& 0& _{qq}^{\prime \prime }& 0\\ 0& 0& 0& _{pp}^{\prime \prime }\end{array}\right).$$
The eigenvalue problem of this matrix factorizes into the eigenvalue problems for the three matrices
$$𝕂,_{pp}^{\prime \prime },\left(\begin{array}{cc}_{bb}^{\prime \prime }& _{bq}^{\prime \prime }\\ _{qb}^{\prime \prime }& _{qq}^{\prime \prime }\end{array}\right).$$
One can prove that the two–by–two matrix $`_{pp}^{\prime \prime }`$ is indefinite for any ellipsoid of type I, II or III. (This implies that these ellipsoids are saddle points of the reduced Hamiltonian, as stated by Riemann.) The eigenvalue problem for the linearization of the Hamiltonian vector field determined by this Hessian does not factorize. Therefore the analysis of the ellipticity is particularly difficult. We know of no rigorous ellipticity results for these ellipsoids.
## 5 Nekhoroshev–stability of the Riemann–ellipsoids: Generalities
We base our study of the Nekhoroshev–stability of the Riemann ellipsoids on some recent results in Nekhoroshev theory, which we describe in this section with reference to the problem at hand; complete proofs can be found in (see also , , ). As is common in Hamiltonian perturbation theory, the procedure consists of the construction of the Birkhoff normal form of order four for the reduced Hamiltonian and testing of the normal form for certain properties that we now specify.
A. Birkhoff Normal Form. Consider an elliptic equilibrium $`(b^{},M^{})`$ with $`M^{}=(m_l^{},m_r^{})`$, $`m_l^{}0m_r^{}`$. Use the canonical coordinates $`\xi =(b,c,q,p)`$ defined as in Section 4.B. Let $`_j(\xi ;b^{},M^{})`$ denote the $`j`$–th term in the Taylor series at $`\xi =0`$ for the reduced Hamiltonian $``$, so that $`_j`$ is a homogeneous polynomial of order $`j`$ in $`\xi `$ and
$$(\xi ;b^{},M^{})=_2(\xi ;b^{},M^{})+_3(\xi ;b^{},M^{})+_4(\xi ;b^{},M^{})+\mathrm{}.$$
Let $`X(b^{},M^{})\xi `$ be the Hamiltonian vector field of $`_2`$. Since the equilibrium is elliptic, the matrix $`X(b^{},M^{})`$ has purely imaginary eigenvalues
$$\pm i\omega _j(b^{},M^{}),j=1,\mathrm{},4,$$
with the convention $`\omega _j0`$.
The first step in the construction of the Birkhoff normal form for the Hamiltonian $``$ is the “symplectic diagonalization” of the quadratic term $`_2`$. It is known that if the eigenvalues of $`X(b^{},M^{})`$ are all distinct, then there is a linear canonical change of coordinates
$$\xi \mathrm{\Xi }=(B,C,Q,P)=T(b^{},M^{})^1\xi $$
such that $`\widehat{}_2(\mathrm{\Xi };b^{},M^{}):=_2(T(b^{},M^{})\mathrm{\Xi };b^{},M^{})`$ has the form
$$\widehat{}_2=s_1\omega _1\frac{B_1^2+C_1^2}{2}+s_2\omega _1\frac{B_2^2+C_2^2}{2}+s_3\omega _3\frac{Q_1^2+P_1^2}{2}+s_4\omega _4\frac{Q_2^2+P_2^2}{2},$$
where $`s_j=\pm 1`$. The matrix $`T`$ and the numbers $`s_j`$ are constructed as follows: For $`j=1,..,4`$, let $`x_j^\pm =x_j^{}\pm ix_j^{\prime \prime }^8`$, with $`x_j^{},x_j^{\prime \prime }^8`$, denote any eigenvectors of the matrix $`X(b^{},M^{})`$ associated to the eigenvalues $`\pm i\omega _j`$. One can show that
$$\mathrm{\Gamma }_j(b^{},M^{}):=x_j^{}J_8x_j^{\prime \prime }\mathrm{\hspace{0.33em}0},j=1,..,4,$$
where $`J_8=\mathrm{diag}[J_4,J_4]`$; here $`J_4`$ denotes the standard four–by–four symplectic matrix $`\left(\begin{array}{cc}\mathrm{𝟎}& \mathrm{𝟏}\\ \mathrm{𝟏}& \mathrm{𝟎}\end{array}\right)`$. We set $`(u_j,v_j)=(x_j^{},x_j^{\prime \prime })`$ if $`\mathrm{\Gamma }_j>0`$ and $`(u_j,v_j)=(x_j^{\prime \prime },x_j^{})`$ if $`\mathrm{\Gamma }_j<0`$, and define
$$T=(\frac{u_1}{\gamma _1},\frac{u_2}{\gamma _2},\frac{v_1}{\gamma _1},\frac{v_2}{\gamma _2},\frac{u_3}{\gamma _3},\frac{u_4}{\gamma _4},\frac{v_3}{\gamma _3},\frac{v_4}{\gamma _4}),s_j=\mathrm{Sign}(\mathrm{\Gamma }_j),$$
where $`\gamma _j=\sqrt{|\mathrm{\Gamma }^j|}`$; the notation means that the eight vectors are the rows of the matrix. Note that the condition that all the eigenvalues of $`X`$ be distinct (and hence nonzero) can be regarded as the condition that the “frequency vector”
$$\mathrm{\Omega }=(s_1\omega _1,\mathrm{},s_4\omega _4)$$
has no resonances of order one or two:
$$\mathrm{\Omega }\nu 0\mathrm{for}\mathrm{all}\nu ^2\mathrm{such}\mathrm{that}|\nu |=1,2,$$
(5.1)
where, for integer vectors, $`|\nu |=|\nu _1|+\mathrm{}+|\nu _4|`$.
Remark: In the case of the ellipsoids of type S, the eigenvalue problem for the matrix $`X(b^{},M^{})`$ factorizes into that for two $`4\times 4`$ blocks. Since the construction of the matrix $`T`$ also factorizes, it is necessary to check the nonresonance condition (5.1) for all resonances of order one, but only for the four resonances $`(1,\pm 1,0,0)`$, $`(0,0,1,\pm 1)`$ of order two.
For the construction of the Birkhoff normal forms we use the complex coordinates $`U=(W_1,W_2,W_3,W_4,Z_1,Z_2,Z_3,Z_4)`$ defined by
$$W_j=\frac{iB_j+C_j}{\sqrt{2}},Z_j=\frac{B_j+iC_j}{\sqrt{2}},W_{j+2}=\frac{iQ_j+P_j}{\sqrt{2}},Z_{j+2}=\frac{Q_j+iP_j}{\sqrt{2}}(j=1,2)$$
(note that $`W`$ are the coordinates, $`Z`$ the momenta). Since $`\mathrm{\Xi }=\mathrm{\Sigma }U`$, with
$$\mathrm{\Sigma }=\left(\begin{array}{cc}\mathrm{\Sigma }_4& 0\\ 0& \mathrm{\Sigma }_4\end{array}\right)\text{and}\mathrm{\Sigma }_4=\frac{1}{\sqrt{2}}\left(\begin{array}{cccc}i& 0& 1& 0\\ 0& i& 0& 1\\ 1& 0& i& 0\\ 0& 1& 0& i\end{array}\right),$$
the Hamiltonian in terms of the complex coordinates $`U`$ is $`H(U;b^{},M^{}):=(T\mathrm{\Sigma }U;b^{},M^{})`$ and its quadratic part is
$$H_2(U;b^{},M^{})=\underset{j=1}{\overset{4}{}}i\mathrm{\Omega }_j(b^{},M^{})Z_jW_j.$$
We now describe the construction of the Birkhoff normal form of order four of the Hamiltonian $`H`$. Given a function $`f`$ of $`W=(W_1,W_2,W_3,W_4)`$ and $`Z=(Z_1,Z_2,Z_3,Z_4)`$, we denote its Taylor series by $`f(W,Z)=_{j,k^4}f_{jk}W^jZ^k`$, where $`W^j=W_1^{j_1}W_2^{j_2}W_3^{j_3}W_4^{j_4}`$ etc. For any integer vector $`\nu ^4`$, we define the $`\nu `$–th harmonic of $`f`$ by
$$f_\nu (W,Z)=\underset{\genfrac{}{}{0pt}{}{j,k^4}{jk=\nu }}{}f_{jk}W^jZ^k.$$
The average of $`f`$ is its harmonic $`f_0`$. The spectrum of $`f`$ is $`\mathrm{Sp}(f)=\{\nu ^4:f_\nu 0\}`$.
We construct the Birkhoff normal form by means of the so–called Lie method. Let $`\mathrm{\Phi }^\chi `$ denote the time–one–map of the flow of the Hamiltonian vector field of a function $`\chi `$; the Lie transform generated by $`\chi `$ of a function $`f`$ is $`f\mathrm{\Phi }^\chi `$. If $`f`$ and $`\chi `$ are analytic in an open neighbourhood of a point (the equilibrium point, in our case), then $`f\mathrm{\Phi }^\chi `$ is also analytic in some sufficiently small, but nonempty, open neighbourhood of such a point and there one has $`f\mathrm{\Phi }^\chi =_{j=0}^{\mathrm{}}\frac{1}{j!}L_\chi ^jf`$. (We adopt the sign convention $`L_fg=\{f,g\}=\frac{f}{Z}\frac{g}{W}\frac{g}{Z}\frac{f}{W}`$ for the Poisson brackets.)
The generating function $`\chi `$ of the Lie transform is constructed by solving the so–called homological equation, which has the form $`\{H_2,\chi \}=ff_0`$ for some function $`f`$. The solution of this equation is formally given by $`\chi =𝒮_\mathrm{\Omega }(f)`$, where
$$𝒮_\mathrm{\Omega }(f)=\underset{\nu \mathrm{Sp}(f)\{0\}}{}\frac{f_\nu }{i\mathrm{\Omega }\nu }.$$
Obviously, this is well defined if $`\mathrm{\Omega }`$ does not resonate with any $`\nu \mathrm{Sp}(f)`$, and if the series converges. In our case $`f`$ will always be a polynomial, either $`f=H_3`$ or $`f=H_4^{}`$ (see below), so no convergence problems exist.
The fourth–order Birkhoff normal form of $`H=H_2+H_3+H_4+\mathrm{}`$, if it exists, is constructed using two Lie transforms, which average the terms of degree three and four of $`H`$, respectively. The first Lie transform is generated by the solution
$$\chi _1=𝒮_\mathrm{\Omega }(H_3)$$
(5.2)
of the homological equation $`\{H_2,\chi _1\}=H_3`$ and is well defined if
$$\mathrm{\Omega }(b^{},M^{})\nu 0\nu \mathrm{Sp}(H_3).$$
(5.3)
Since $`H_3_0=0`$ and $`L_{\chi _1}^2H_2=L_{\chi _1}H_3`$ one sees that this Lie transform conjugates $`H=H_2+H_3+H_4+\mathrm{}.`$ to
$$H^{}=H_2+H_4^{}+\mathrm{},\text{where}H_4^{}=\frac{1}{2}L_{\chi _1}H_3+H_4.$$
The second Lie transform is generated by the solution $`\chi _2=𝒮_\mathrm{\Omega }(H_4^{})`$ of $`\{H_2,\chi _2\}=H_4^{}H_4^{}_0`$. It is well–defined if
$$\mathrm{\Omega }(b^{},M^{})\nu 0\nu \mathrm{Sp}(H_4^{})\{0\}$$
(5.4)
and leads to the fourth–order Birkhoff normal form
$$H^{\prime \prime }=H_2+H_4^{\prime \prime }+\mathrm{}\text{with}H_4^{\prime \prime }=\frac{1}{2}L_{\chi _1}H_3+H_4_0.$$
(5.5)
For our purposes, the only important points are controlling the existence of the fourth–order Birkhoff normal form and testing that it possesses the convexity properties discussed above. Neither of these tasks require the actual computation of the generator $`\chi _2`$: one needs only test (5.4) and compute $`H_4^{\prime \prime }`$.
B. Nekhoroshev stability of elliptic equilibria. Let us introduce the action functions $`I=(I_1,\mathrm{},I_4)`$, $`I_j=iW_jZ_j`$, and regard $`H_4^{\prime \prime }`$ as a quadratic form on $`^4`$ by writing
$$H_4^{\prime \prime }(I)=\frac{1}{2}IA(b^{},M^{})I$$
where $`A(b^{},M^{})`$ is a $`4\times 4`$ symmetric matrix. The Birkhoff normal form of order four (5.5) is said to be
* convex if the quadratic form $`H_4^{\prime \prime }`$ is definite, i.e. the eigenvalues of $`A`$ are either all positive or all negative.
* quasi–convex if the restriction of $`H_4^{\prime \prime }`$ to the subspace orthogonal to $`\mathrm{\Omega }`$ is definite, i.e.
$$\mathrm{\Omega }I=0\mathrm{and}H_4^{\prime \prime }(I)=0I=0.$$
* directionally quasi–convex if the restriction of the quadratic form $`H_4^{\prime \prime }`$ to the plane orthogonal to $`\mathrm{\Omega }`$ is nonvanishing in the “first 16–ant”, i.e.
$$\mathrm{\Omega }I=\mathrm{\hspace{0.33em}0},H_4^{\prime \prime }(I)=\mathrm{\hspace{0.33em}0}\mathrm{and}I_1,\mathrm{},I_40I=0.$$
Clearly, each notion generalizes the previous one.
If the Hamiltonian $`H(W,Z)`$ and its Birkhoff normal form of order four are analytic, then any of these conditions implies the Nekhoroshev–stability of the equilibrium, in the precise sense that for any small $`ϵ>0`$ one has
$$I(0)ϵI(t)ϵ^\alpha \text{for}|t|\mathrm{exp}ϵ^\beta $$
for some positive constants $`\alpha `$ and $`\beta `$. There are different possible values for the constants $`\alpha `$ and $`\beta `$, which we report here with reference to the case under consideration of a system with four degrees of freedom:
* If the Birkhoff normal form is directionally quasi–convex, then one can take $`\alpha =\beta =\frac{1}{4}`$ and it is also possible to prove a stricter confinement of motions on correspondingly shorter times, namely $`\alpha =\frac{1+k}{4+k}`$, $`\beta =\frac{1}{4+k}`$ for any $`k>0`$ (see ).
* If the Birkhoff normal form is quasi–convex (directional quasi–convexity does not seem to be sufficient), then it is possible to prove the above estimates with $`\alpha =1`$ and $`\beta =\frac{1}{16}`$. (See .)
* If the Birkhoff normal form of order four is quasi–convex and it is possible to construct the Birkhoff normal form of order $`s>4`$, then one also has $`\alpha =1`$, $`\beta =\frac{s3}{16}`$ (see ).
Remark: The original motivation for the quoted articles , on the Nekhoroshev stability of elliptic equilibria was precisely the present study of the Riemann ellipsoids. Until a few years ago, all Nekhoroshev stability results assumed that the frequency vector would satisfy a strong nonresonance condition, typically a diophantine one. This would be insufficient for the study of a problem such as the present one, in which the frequencies of the equilibrium, and therefore its nonresonance properties, depend on continuously varying parameters. The notion of directional quasi–convexity was introduced in in connection with the study of the Nekhoroshev–stability of the triangular points of the restricted three–body problem. As will be seen below, this notion also plays a central role in the study of the Riemann ellipsoids, since the majority of the ellipsoids are not quasi–convex. (Nekhoroshev estimates also hold if a sufficiently high order Birkhoff normal form exists and satisfies some “steepness” condition; directional quasi–convexity allows us to avoid these generalizations, the hypotheses of which would be rather difficult to verify and which would lead to significantly worse values for the constants $`\alpha `$ and $`\beta `$.)
C. KAM theory. If the Birkhoff normal form of order four is nondegenerate, in the sense that $`detA(b^{},M^{})0`$, then KAM theory applies, ensuring that in any sufficiently small neighbourhood of the equilibrium, the majority of the initial data gives rise to motions which are quasi–periodic with four frequencies. (See e.g. .)
## 6 Nekhoroshev–stability of the Riemann ellipsoids: Numerical results
A. Results. We have numerically constructed the Birkhoff normal forms for a (quite large) number of sample Riemann ellipsoids. Specifically, within the existence region of each type of ellipsoid, we considered a mesh in the plane $`(\frac{b_2}{b_1},\frac{b_3}{b_1})`$ determined by vertical lines uniformly spaced at a distance $`.0025`$; the number of mesh points on each vertical lines typically varied between twenty and one hundred, depending on the length of the line (however, in some cases we used as many as five hundred points). This analysis leads to the following
Numerical Conclusions 2. The Birkhoff normal form of order four exists and is directionally quasi–convex and nondegenerate for all the Riemann ellipsoids except for those lying on a finite number of curves in the set $``$ corresponding to resonances of order up to four.
We regard the fact that all of the computed nonresonant Birkhoff normal forms are directionally quasi–convex as indicating that all nonresonant Riemann ellipsoids are Nekhoroshev stable. We investigated the existence of low–order resonances in greater detail, without restriction to the points of the considered mesh (see subsection 6.B for some details), finding that there are a finite number of curves in the plane $`(\frac{b_2}{b_1},\frac{b_3}{b_1})`$ at which at least one of the nonresonance conditions (5.1), (5.3), (5.4) is violated. Specifically, we found 8 different resonances for the type $`\mathrm{S}_2`$, 52 for the type I, 33 for the type II, and 47 for the type III. For example, figure 5.a reports the resonant curves for the ellipsoids of type I.
We found that the Birkhoff normal form is quasi–convex only for a few ellipsoids of each type; in turn, very few of these ellipsoids are actually convex. For example, figure 5.b shows the quasi–convex ellipsoids in the lower region of elliptic type I ellipsoids. The determination of the quasi–convex ellipsoids has some interest because, as discussed in the previous Section, such equilibria may have stronger stability properties if they are nonresonant up to sufficiently high order. (From this perspective, the finer distinction between convexity and quasi–convexity does not instead seem to be as significant.)
Nondegeneracy of the Birkhoff normal forms implies that the KAM theorem applies, so that the majority of the motions near each equilibrium in the reduced phase space are quasi–periodic with four frequencies. “Reconstruction” of these motions then gives quasi–periodic motions of the Dirichlet problem, with up to (and in fact, typically) eight frequencies.
B. Numerical Procedures. We give now some information about the numerical procedures adopted here. The fourth–order Taylor expansion of the Hamiltonian in coordinates $`(b,c,q,p)`$, and hence the Hamiltonian matrix $`X(b^{},M^{})`$, has been constructed analytically for each type of ellipsoid. The coefficients of these polynomials have been evaluated for any ellipsoid in our mesh, and all other operations have then been performed on polynomials with numeric coefficients. These include:
* Computing the eigenvalues and the eigenvectors of $`X(b^{},M^{})`$ and testing for the presence of resonances of order one and two. Note that once the eigenvectors are known, the matrix $`S(b^{},M^{})`$ determining the coordinates $`\mathrm{\Xi }`$ and $`U`$ is also known.
* Constructing the normal forms by means of formulas (5.5) and (5.2). This requires computing Poisson brackets, i.e. derivatives, and averages of polynomials with numeric coefficients; these are easily implemented operations.
* Checking the nonresonance conditions (5.3) and (5.4), which assure that the normal forms can be constructed. This involves determining the spectrum of a polynomial.
* Testing the normal forms, namely the numeric vector $`\mathrm{\Omega }`$ and matrix $`A`$, for convexity, quasi–convexity, or directional quasi–convexity.
The chosen implementation of the latter two points requires some explanation.
Spectra and Resonances. We computed the spectra of $`H_3`$ and $`H_4^{\prime \prime }`$ at every point $`(b^{},M^{})`$ in the considered mesh and took the unions $`\mathrm{\Sigma }_3`$ and $`\mathrm{\Sigma }_4`$ of all these sets. (The union is taken because the numerically computed spectrum at a given point may lack very small harmonics.) We identified as resonances those $`\nu \mathrm{\Sigma }_3\mathrm{\Sigma }_4`$ for which the scalar product $`\mathrm{\Omega }(b^{},M^{})\nu `$ either vanishes or changes sign on our grid; of course, use of a finer grid could reveal additional resonances. (The resonant curves of figure 5.a were constructed by computing the approximate zeros of $`\mathrm{\Omega }(b^{},M^{})\nu `$, testing for changes of sign on a finer grid.)
Quasi and Directional Quasi–convexity. If the matrix $`A`$ is not convex, then the tests for quasi–convexity and directional quasi–convexity can be carried out as follows. Let $`R\mathrm{SO}(3)`$ be such that $`R\mathrm{\Omega }=(1,0,0,0)`$, so that the Hessian of the restriction of the quadratic form $`H_4^{\prime \prime }`$ to the subspace orthogonal to $`\mathrm{\Omega }`$ is the lower right hand $`3\times 3`$ block $`\stackrel{~}{A}`$ of the matrix $`RAR^T`$. Quasi–convexity of $`H_4^{\prime \prime }`$ is equivalent to convexity of the matrix $`\stackrel{~}{A}`$.
Under the additional hypothesis that the restriction of $`H_4^{\prime \prime }(I)`$ to the subspace orthogonal to $`\mathrm{\Omega }`$ is nondegenerate,<sup>6</sup><sup>6</sup>6This hypothesis is not crucial, but it is satisfied by all the computed ellipsoids. the test for directional quasi–convexity can be performed as follows. The quadratic form $`H_4^{\prime \prime }`$ is directionally quasi–convex iff any “asymptotic vector”, namely any nonzero vector $`I`$ satisfying $`I\stackrel{~}{A}I=0`$, points out of the first 16–ant, i.e. has at least one negative and one positive entry. Since these vectors form a cone (of dimension two, if $`\stackrel{~}{A}`$ is nondegenerate) we need only test the asymptotic vectors in the intersection of the cone and the unit sphere, namely two ellipses. In order to determine these ellipses, note that in the absence of quasi–convexity, the matrix $`\stackrel{~}{A}`$ has at least one eigenvalue of either sign. We may assume that the eigenvalues of $`\stackrel{~}{A}`$ satisfy $`\alpha _1>0`$ and $`\alpha _2,\alpha _3<0`$, replacing $`\stackrel{~}{A}`$ with $`\stackrel{~}{A}`$ if necessary. (This is possible because $`\stackrel{~}{A}`$ and $`\stackrel{~}{A}`$ have the same asymptotic vectors. Also, the nondegeneracy hypothesis assures that all of the eigenvalues are nonzero.) Let $`\{\stackrel{~}{v}_1,\stackrel{~}{v}_2,\stackrel{~}{v}_3\}`$ be an orthonormal eigenbasis associated to $`\stackrel{~}{A}`$, so that $`S=(\stackrel{~}{v}_1,\stackrel{~}{v}_2,\stackrel{~}{v}_3)\mathrm{SO}(3)`$ satisfies $`S^TAS=\mathrm{diag}(|\alpha _1|,|\alpha _2|,|\alpha _3|)`$. In this basis, the cone of asymptotic vectors is given by
$$\{x=(x_1,x_2,x_3)^3:x_1^2=\frac{|\alpha _2|}{\alpha _1}x_2^2+\frac{|\alpha _3|}{\alpha _1}x_3^2\}.$$
The two ellipses of unit vectors on this cone consist of the vectors
$$\widehat{x}^\pm (\theta )=(\pm \sqrt{\frac{|\alpha _2|}{\alpha _1+|\alpha _2|}\mathrm{cos}^2\theta +\frac{|\alpha _3|}{\alpha _1+|\alpha _3|}\mathrm{sin}^2\theta },\sqrt{\frac{\alpha _1}{\alpha _1+|\alpha _2|}}\mathrm{cos}\theta ,\sqrt{\frac{\alpha _1}{\alpha _1+|\alpha _3|}}\mathrm{sin}\theta ),\theta S^1.$$
In the original basis in the subspace orthogonal to $`\mathrm{\Omega }`$, these vectors are
$$\widehat{I}^\pm (\theta )=(\widehat{I}_1,\widehat{I}_2,\widehat{I}_3)=S\widehat{x}^\pm (\theta ).$$
These are three–dimensional vectors belonging to the subspace orthogonal to $`(1,0,0,0)=R\mathrm{\Omega }`$, so the corresponding four–dimensional vectors are $`I^\pm (\theta )=R^T(0,\widehat{I}_1,\widehat{I}_2,\widehat{I}_3)`$. The normal form is directionally quasi–convex iff for every $`\theta S^1`$ the vector $`I^+(\theta )`$ has at least one negative and one positive entry. (It is obviously sufficient to consider only one of the two circles.) In the computations, we have tested this condition by checking that, at each of a number (typically $`101`$) of uniformly spaced values of $`\theta `$ in the interval $`0\theta 2\pi `$, there are two eigenvalues whose product is negative.
## 7 Conclusions
In this work, we have posed the problem of the Nekhoroshev stability of those Riemann ellipsoids that are elliptic but of unknown Lyapunov stability. We have provided numerical evidence that all Riemann ellipsoids that are not resonant up to order four are Nekhoroshev stable. We have developed a consistent and rigorous Hamiltonian formulation of the problem on a four–to–one covering manifold, which validates the procedures (and clarifies the presence of a previously not understood discrete symmetry in the problem). Within this formulation we provided explicit, precise descriptions of the existence regions and of the momenta of the Riemann ellipsoids of the various types and we reviewed the existent results on the stability and ellipticity of the Riemann ellipsoids, showing (numerically) that some of the latter regions are significantly larger than was previously thought.
At the conclusion of this work, it is our opinion that several key features of the problem are not yet understood. Some open problems include:
* Properties of axisymmetric Riemann ellipsoids other than the Maclaurin spheroids.
* Analytic determination of the regions of ellipticity. Understanding of their subtle structure.
* Lyapunov stability or instability of those Riemann ellipsoids that are saddle points of the reduced Hamiltonian.
Of course, it would also be interesting to prove the Nekhoroshev–stability rigorously, that is, without recourse to numerical calculations, but this appears to be a formidable task at present.
In conclusion, we would like to mention a point of historic interest. Riemann’s treatment of the Dirichlet problem is exceptionally advanced and uses several geometric constructions, including Poisson reduction, the amended potential, and the trivialization of the tangent and cotangent bundles of Lie groups, that are now considered to be central to modern geometric mechanics.
Riemann’s apparent lack of concern with proving the actual instability of the saddle points of the reduced Hamiltonian is probably typical of that era — this error permeated the mechanical literature long after Lyapunov clarified the subject at the beginning of the twentieth century. At the dawn of a new century, the Lyapunov (in)stability of equilibria that are both elliptic and saddles of the ‘obvious’ Lyapunov function is still a formidable problem. Nekhoroshev stability provides a weaker, but very practical, notion of stability, which has the ultimate advantage of being compatible with Lyapunov instability.
## 8 Appendix A: On the proof of Proposition 4
The proof of Proposition 4, even though elementary, is too lengthy to be reported here. Hence we indicate the main steps in the form of Lemmas that we do not prove. The relative equilibria are the critical points of the reduced Hamiltonian $`H_\eta :P_\eta `$ and hence are given by the conditions
$$_bH_\eta =0,_cH_\eta =0,m\times _mH_\eta =0.$$
(8.1)
Since $`_cH_\eta =0`$ iff $`c=0`$, we can equivalently search for critical points $`(b,m)`$ of
$$\stackrel{~}{H}_\eta (b,m):=\frac{1}{2}m𝕁(b)m+𝒱(b).$$
Note that only the first equation in (8.1) depends on the potential $`𝒱`$.
###### Lemma 3
A nonzero vector $`m=(m_l,m_r)(^3)^2`$ satisfies $`m\times _m\stackrel{~}{H}_\eta (b,m)=0`$ iff either of the following alternative conditions is verified:
* $`m_l`$ and $`m_r`$ are both parallel to the same principal axis;
* $`m_l`$ and $`m_r`$ both belong to the same principal plane, say $`e_ie_j`$ with $`b_i>b_j`$, but are not colinear with either principal axis, and, together with $`b`$, they satisfy any of the following three conditions (where $`(i,j,k)`$ is a permutation<sup>7</sup><sup>7</sup>7By a permutation, we mean either an even or an odd permutation. of $`(1,2,3)`$):
+ $`b_j<|b_i2b_k|`$ and
$`\left({\displaystyle \frac{m_l^im_r^i}{m_l^i+m_r^i}}\right)^2`$ $`=`$ $`\left({\displaystyle \frac{b_j+b_k}{b_jb_k}}\right)^4{\displaystyle \frac{b_i^2(b_j2b_k)^2}{b_i^2(b_j+2b_k)^2}}`$ (8.2)
$`{\displaystyle \frac{m_l^jm_r^j}{m_l^j+m_r^j}}`$ $`=`$ $`{\displaystyle \frac{m_l^i+m_r^i}{m_l^im_r^i}}\left({\displaystyle \frac{b_j+b_k}{b_jb_k}}\right)^2\left({\displaystyle \frac{b_i+b_k}{b_ib_k}}\right)^2{\displaystyle \frac{b_i+b_j2b_k}{b_i+b_j+2b_k}}.`$ (8.3)
+ $`b_i=2b_k+b_j`$, $`m_l^i=m_r^i`$, and $`m_l^j=m_r^j`$.
+ $`b_i=2b_kb_j`$ and $`m_l=m_r`$.
Equations (8.2) and (8.3) imply that (8.2) is also valid with the index $`i`$ and $`j`$ exchanged. If either $`m_l`$ or $`m_r`$ vanish (the case of an irrotational ellipsoid), then equations (8.2) and (8.3) are equivalent to $`k=2`$ and $`b_1^2b_2^2+b_2^2b_3^2+b_1^2b_3^2=3b_2^4`$.
We now consider the equation $`_bH_\eta =0`$. Set
$$\stackrel{~}{}_j=\{b:G_{}^S(b_1,b_j,b_k)0,G_+^S(b_1,b_j,b_k)0\},j=2,3$$
where $`(j,k)`$ is any permutation of $`(2,3)`$ and
$$_{ij}^\pm =\{b:b_j\pm (b_i2b_k),D(b_i,b_j,b_k)0,G_{}^R(b_i,b_j,b_k)>0,G_\pm ^R(b_j,b_i,b_k)>0\}$$
where $`(i,j,k)`$ is any permutation of $`(1,2,3)`$ with $`i<j`$. (We shall see below that some of these sets are empty.)
###### Lemma 4
Assume $`m=(m_l,m_r)(^3)^2`$ is nonzero and satisfies $`m\times _m\stackrel{~}{H}_\eta =0`$. Then $`_b\stackrel{~}{H}_\eta =0`$ iff one of the following two conditions is satisfied:
* $`b\stackrel{~}{}_{ij}`$, $`m_l=m_l^je_j`$, $`m_r=m_r^je_j`$ and $`(m_l^j\pm m_r^j)^2=G_\pm ^S(b_1,b_j,b_k)`$ for either $`j=2`$, $`k=3`$ or $`j=3`$, $`k=2`$.
* $`b_{ij}^{}_{ij}^+`$, $`m_l=m_l^ie_i+m_l^je_j`$, $`m_r=m_r^ie_i+m_r^je_j`$, $`(m_l^i\pm m_r^i)^2=G_\pm ^R(b_i,b_j,b_k)`$, $`(m_l^j\pm m_r^j)^2=G_\pm ^R(b_j,b_i,b_k)`$ and, if $`b_i2b_k\pm b_j`$,
$$m_l^jm_r^j=\mathrm{sign}\left((2b_kb_ib_j)\frac{(m_l^j+m_r^j)(m_l^i+m_r^i)}{m_l^im_r^i}\right)\sqrt{G_{}^R(b_j,b_i,b_k)},$$
for any permutation $`(i,j,k)`$ of $`(1,2,3)`$ with $`i<j`$.
The proof of part (i) of Proposition 4 is completed by showing that the “existence regions” of the previous Lemma coincide with those of the Proposition and that all the equilibria described in that Lemma are $`_4\times _2`$–equivalent or adjoint to those indicated in the Proposition:
###### Lemma 5
The sets $`_{23}^+`$, $`_{23}^{}`$ and $`_{12}^{}`$ are empty, while $`\stackrel{~}{}_2=_{\mathrm{S}_2}`$, $`\stackrel{~}{}_3=_{\mathrm{S}_3}`$, $`_{13}^{}=_\mathrm{I}`$, $`_{13}^+=_{\mathrm{II}}`$, and $`_{12}^+=_{\mathrm{III}}`$.
###### Lemma 6
Under the assumptions of Lemma 4, a point $`(b,0,m)`$ is a relative equilibrium for some reduced system iff it is $`_4\times _2`$–equivalent to one of the points $`(b,0,(\mu _\alpha ^\pm (b),\mu _\alpha ^{}(b)))`$.
Part (ii) of the Proposition is a consequence of the following
###### Lemma 7
The maps $`\mu _\alpha ^\pm `$, $`\alpha =\mathrm{S}_3,\mathrm{II},\mathrm{III}`$, are never zero. For $`\alpha =\mathrm{S}_2`$ and $`\alpha =\mathrm{I}`$, they vanish on the sets indicated in Proposition 4.
## 9 Appendix B: On Chandrasekhar’s linearization
To facilitate the comparison of our results with Chandrasekhar’s, in figure 6 we display the regions of elliptic ellipsoids of types II and III using the coordinates employed by Chandrasekhar, namely $`(\frac{b_1}{b_2},\frac{b_3}{b_2})`$ for the ellipsoids of type II and $`(\frac{b_1}{b_3},\frac{b_2}{b_3})`$ for the ellipsoids of type III. (For the latter, Chandrasekhar also uses the coordinates $`(\frac{b_2}{b_3},\frac{b_2}{b_1})`$; the corresponding figure should be compared to figure 3.b.) Among all these ellipsoids, Chandrasekhar identified as elliptic (“stable”, in his terminology) only the ellipsoids in the lowest ellipticity region in figure 6.d.
As we have already remarked, Chandrasekhar’s linearization of Riemann’s equations yields the same results as ours (after correcting a few misprints that we report in the Remark below for the convenience of the interested reader). Hence, it seems probable that the meshes used by Chandrasekhar in his numerical study were too coarse, or too small, to capture all of the ellipticity regions shown in figure 6. Certainly, the very few ellipsoids whose frequencies are reported in (three of type II and four of type III, which are indicated by large points in figure 6) are consistent with our result.
Chandrasekhar did not exploit (nor, apparently, notice) the Lagrangian and Hamiltonian structure of the Dirichlet problem and worked with the full sixteen–dimensional Riemann’s system, without exploiting the symmetry of the problem. The linearization of the reduced and unreduced systems are, of course, equivalent: the four additional frequencies appearing in the full system (two of which are identically zero) do not alter the ellipticity.
Remark: In , the frequencies of the linearized equations of motion are determined by the system of nine equations (116). The nine-by-nine determinant of the coefficients of this system is given in formula (133) of , where there are two misprints: the entry (7,2) should be multiplied by $`2+(\beta 2)\gamma a_1^2/a_2^2`$, and the entry (8,3) should be multiplied by $`2+(\gamma 2)\beta a_1^2/a_3^2`$. This determinant is reduced to the eight-by-eight determinant (134), in which the entry (1,8) should be changed to $`\mathrm{\Omega }_2\mathrm{\Omega }_3+2\frac{a_1}{a_2}\mathrm{\Omega }_2\mathrm{\Omega }_3^{}`$, the entry (2,7) to $`\mathrm{\Omega }_2\mathrm{\Omega }_3+2\frac{a_1}{a_3}\mathrm{\Omega }_2^{}\mathrm{\Omega }_3`$ and the entry (3,5) to $`\mathrm{\Omega }_2\mathrm{\Omega }_3+\frac{a_2}{a_3}\mathrm{\Omega }_2^{}\mathrm{\Omega }_3^{}`$; the entries (8,1), (7,2) and (5,3) should be changed accordingly: the entry $`(i,j)`$ is obtained from the entry $`(j,i)`$ by exchanging $`\mathrm{\Omega }_k`$ and $`\mathrm{\Omega }_k^{}`$ ($`k=2,3`$). Also, the symbol $`B_{123}`$ appearing in equation (135) of , should be defined as in equation (104), Chapter 3, of . Let us explicitly note that these misprints are not the sources of the discrepancies between Chandrasekhar’s and our conclusions: it is only after they have been corrected that equation (134) of detects as elliptic the regions that Chandrasekhar describes as “stable”. |
warning/0002/hep-th0002083.html | ar5iv | text | # Bidyon as an electromagnetic model for charged particle with spin
(Received October 16, 2002; Revised February, 2003)
## Abstract
A general model of nonlinear electrodynamics with dyon singularities is considered. We consider the field configuration having two dyon singularities with identical electric and opposite magnetic charges and we name it bidyon. We investigate the sum of two dyon solutions as an initial approximation to the bidyon solution. We consider the case when the velocities of the dyons have equal modules and opposite directions on a common line. It is shown that the associated field configuration has a constant full angular momentum which is independent of distance between the dyons and their speed. This property permits a consideration of this bidyon configuration as an electromagnetic model for charged particle with spin. We discuss the possible electrodynamic world with oscillating bidyons as particles.
A hypothetical particle having both electric and magnetic charges is said to be dyon . An electromagnetic field configuration with $`N`$ point dyons satisfies the following two differential conditions:
$$\{\begin{array}{ccc}\hfill \mathrm{Div}𝐃& =& 4\pi ȷ^0\hfill \\ \hfill \mathrm{Div}𝐁& =& 4\pi \overline{ȷ}^0\hfill \end{array},$$
(1)
where $`\mathrm{Div}𝐃\overline{}_iD^i`$ , $`\overline{}_\mu {\displaystyle \frac{1}{\sqrt{|g|}}}{\displaystyle \frac{}{x^\mu }}\sqrt{|g|}`$ , $`gdet(g_{\mu \nu })`$ ,
the Latin indices take the values $`1,2,3`$, the Greek ones take the value $`0,1,2,3`$, $`g_{\mu \nu }`$ is a metric of space-time coordinate system,
$$ȷ^0\frac{1}{\sqrt{|g|}}\underset{n=1}{\overset{N}{}}\stackrel{n}{d}\delta (𝐱\stackrel{n}{𝐚}),\overline{ȷ}^0\frac{1}{\sqrt{|g|}}\underset{n=1}{\overset{N}{}}\stackrel{n}{b}\delta (𝐱\stackrel{n}{𝐚}),$$
(2)
$`\stackrel{n}{d}`$ is an electric charge and $`\stackrel{n}{b}`$ is a magnetic one for $`n`$-th singular point, $`\stackrel{n}{𝐚}=\stackrel{n}{𝐚}(x^0)`$ is a trajectory of it.
Here we use the definition for three-dimensional $`\delta `$-function which is suitable for discontinuous functions $`f(𝐱)`$:
$`{\displaystyle \underset{\stackrel{n}{\mathrm{\Omega }}}{}}f(𝐱)\delta (𝐱\stackrel{n}{𝐚})(\mathrm{d}x)^3\underset{\stackrel{n}{\sigma }0}{lim}\left[{\displaystyle \frac{1}{|\stackrel{n}{\sigma }|}}{\displaystyle \underset{\stackrel{n}{\sigma }}{}}f(𝐱)\mathrm{d}\stackrel{n}{\sigma }\right],|\stackrel{n}{\sigma }|{\displaystyle \underset{\stackrel{n}{\sigma }}{}}\mathrm{d}\stackrel{n}{\sigma },`$ (3)
where $`\stackrel{n}{\mathrm{\Omega }}`$ is a region of three-dimensional space including the point $`𝐱=\stackrel{n}{𝐚}`$, $`\stackrel{n}{\sigma }`$ is a closed surface enclosing this point, $`\mathrm{d}\stackrel{n}{\sigma }`$ is an area element of the surface $`\stackrel{n}{\sigma }`$ , $`|\stackrel{n}{\sigma }|`$ is an area of the whole surface $`\stackrel{n}{\sigma }`$ .
Eqs. (1) are the part of Maxwell system of equations in any space-time coordinate system with a metric $`g_{\mu \nu }`$ (see also ).
To have a natural interaction between the dyons (see ) we must take the fields $`𝐃`$ and $`𝐁`$ satisfying some nonlinear Maxwell equations that may be written in the following general form:
$$\{\begin{array}{ccc}\hfill \overline{}_0𝐃\mathrm{Rot}𝐇& =& 4\pi \mathit{ȷ}\hfill \\ \hfill \overline{}_0𝐁+\mathrm{Rot}𝐄& =& 4\pi \overline{\mathit{ȷ}}\hfill \end{array},$$
(4)
where $`(\mathrm{Rot}𝐄)^i\epsilon ^{0ijk}E_k/x^j`$, $`\epsilon ^{0123}=|g|^{1/2}`$, $`\epsilon _{0123}=|g|^{1/2}`$,
$`E_i={\displaystyle \frac{}{D^i}},H_i={\displaystyle \frac{}{B^i}},=(𝐃,𝐁),`$ (5)
$`\mathit{ȷ}{\displaystyle \frac{1}{\sqrt{|g|}}}{\displaystyle \underset{n=1}{\overset{N}{}}}\stackrel{n}{d}\stackrel{n}{𝐕}\delta (𝐱\stackrel{n}{𝐚}),\overline{\mathit{ȷ}}{\displaystyle \frac{1}{\sqrt{|g|}}}{\displaystyle \underset{n=1}{\overset{N}{}}}\stackrel{n}{b}\stackrel{n}{𝐕}\delta (𝐱\stackrel{n}{𝐚}),\stackrel{n}{𝐕}{\displaystyle \frac{\mathrm{d}\stackrel{n}{𝐚}}{\mathrm{d}x^0}}.`$ (6)
According to Eqs. (5) we have some dependencies $`𝐄=𝐄(𝐃,𝐁)`$ and $`𝐇=𝐇(𝐃,𝐁)`$ (see also ). If $`𝐃`$, $`𝐁`$ appears only quadratically in $``$ then we have a linear electrodynamics but in general case the function $`(𝐃,𝐁)`$ defines a nonlinear electrodynamic model. In this approach the vector fields $`𝐃`$, $`𝐁`$ play the role of unknown functions for system of equation (4) with additional differential conditions (1). This representation is best suitable for an investigation of the interaction between the dyons. From the fields $`𝐄`$, $`𝐁`$ satisfying equations (1), (4) we can obtain an appropriate electromagnetic potential. In the case of the dyon singularity of electromagnetic field a space part of the four-potential has a line singularity .
The singular currents (2), (6) must satisfy to the following condition :
$`F_{\mu \nu }ȷ^\nu {\displaystyle \frac{1}{2}}\epsilon _{\mu \nu \sigma \rho }f^{\sigma \rho }\overline{ȷ}^\nu `$ $`=`$ $`0,`$ (7)
where $`F_{i0}=E_i`$, $`F_{ij}=\epsilon _{0ijk}B^k`$, $`f^{0i}=D^i`$, $`f^{ij}=\epsilon ^{0ijk}H_k`$.
Using Eqs. (1), (4), (5), (7) we can check directly the following differential conservation laws for energy-momentum tensor (in Cartesian coordinate systems):
$`{\displaystyle \frac{}{x^0}}`$ $`=`$ $`{\displaystyle \frac{}{x^j}}\left(\epsilon ^{jpq}E_pH_q\right),`$ (8)
$`{\displaystyle \frac{𝒫_i}{x^0}}`$ $`=`$ $`{\displaystyle \frac{}{x^j}}\left[\delta _i^j\left(𝐃𝐄+𝐁𝐇\right)\left(D^jE_i+B^jH_i\right)\right],`$ (9)
where $`𝒫_i\epsilon _{ipq}D^pB^q`$ or $`𝓟𝐃\times 𝐁`$ ($`\epsilon _{123}=\epsilon ^{123}=1`$).
From (8) and (9) we easily obtain that the full energy-momentum<sup>1</sup><sup>1</sup>1Note, here we take the function $``$ such that $`=0`$ for $`𝐃=𝐁=0`$. This is distinction from the designation which is used in the article for Born-Infeld electrodynamics.
$`\begin{array}{ccccccc}\hfill & =& {\displaystyle \frac{1}{4\pi }}{\displaystyle (\mathrm{d}x)^3}\hfill & ,& \hfill 𝐏& =& {\displaystyle \frac{1}{4\pi }}{\displaystyle 𝓟(\mathrm{d}x)^3},\hfill \end{array}`$ (11)
and the vector of full angular momentum
$`𝐌`$ $`=`$ $`{\displaystyle \frac{1}{4\pi }}{\displaystyle \left(𝐱\times 𝓟\right)(\mathrm{d}x)^3}`$ (12)
are conserved on time, i.e. $`\mathrm{d}/\mathrm{d}x^0=\mathrm{d}𝐏/\mathrm{d}x^0=\mathrm{d}𝐌/\mathrm{d}x^0=0`$.
Let us consider a solution of system (4), (1) having two dyon singularities with identical electric and opposite magnetic charges: $`\stackrel{1}{d}=\stackrel{2}{d}`$, $`\stackrel{1}{b}=\stackrel{2}{b}`$. We name this solution bidyon. Let us consider the case when the velocities of the singularities have equal absolute values and opposite directions on a common line. At first we use a cylindrical coordinate system $`\{z,\rho ,\phi \}`$ such that the dyon singularities are on the axis $`z`$. This configuration is shown in Fig. 1, where $`d=\pm \overline{d}`$, $`b=\pm \overline{b}`$ and $`\overline{d},\overline{b}`$ are some positive constants.
We can search the solution by some iterative procedure and we can take a sum of two moving dyon solutions as initial approximation. That is we consider the following initial approximation to the bidyon solution:
$$𝐃^{(0)}=\stackrel{1}{𝐃}+\stackrel{2}{𝐃},𝐁^{(0)}=\stackrel{1}{𝐁}+\stackrel{2}{𝐁}.$$
(13)
Here we consider the dyon solutions with constant velocity. For $`z`$\- and $`\rho `$-components of these solutions (see ) we have the following expressions:
$`\begin{array}{c}\begin{array}{ccc}\{\begin{array}{ccccc}\hfill {\displaystyle \frac{\stackrel{1}{D}_z}{d}}& =& \hfill {\displaystyle \frac{\stackrel{1}{B}_z}{b}}& =& {\displaystyle \frac{1}{\sqrt{1V^2}}}{\displaystyle \frac{z+a}{\stackrel{1}{r}^3}}\hfill \\ \hfill {\displaystyle \frac{\stackrel{1}{D}_\rho }{d}}& =& \hfill {\displaystyle \frac{\stackrel{1}{B}_\rho }{b}}& =& {\displaystyle \frac{1}{\sqrt{1V^2}}}{\displaystyle \frac{\rho }{\stackrel{1}{r}^3}}\hfill \end{array},\hfill & & \\ \{\begin{array}{ccccc}\hfill {\displaystyle \frac{\stackrel{2}{D}_z}{d}}& =& \hfill {\displaystyle \frac{\stackrel{2}{B}_z}{b}}& =& {\displaystyle \frac{1}{\sqrt{1V^2}}}{\displaystyle \frac{za}{\stackrel{2}{r}^3}}\hfill \\ \hfill {\displaystyle \frac{\stackrel{2}{D}_\rho }{d}}& =& \hfill {\displaystyle \frac{\stackrel{2}{B}_\rho }{b}}& =& {\displaystyle \frac{1}{\sqrt{1V^2}}}{\displaystyle \frac{\rho }{\stackrel{2}{r}^3}}\hfill \end{array},\hfill & & \end{array}\hfill \end{array}`$ (21)
| where | $`V{\displaystyle \frac{\mathrm{d}a}{\mathrm{d}x^0}}`$ , $`\stackrel{1}{r}=\sqrt{\left(z^{}+a^{}\right)^2+\rho ^2}`$ , $`\stackrel{2}{r}=\sqrt{\left(z^{}a^{}\right)^2+\rho ^2}`$ , |
| --- | --- |
| | $`z^{}{\displaystyle \frac{z}{\sqrt{1V^2}}},a^{}{\displaystyle \frac{a}{\sqrt{1V^2}}}`$ . |
Forms of $`\phi `$-components for the vector fields $`\stackrel{1}{𝐃}`$ , $`\stackrel{1}{𝐁}`$ , $`\stackrel{2}{𝐃}`$ , $`\stackrel{2}{𝐁}`$ depend on forms of the functions $`𝐄=𝐄(𝐃,𝐁)`$, $`𝐇=𝐇(𝐃,𝐁)`$ but forms (21) for $`\rho `$\- and $`z`$-components are independent of the specific model’s nonlinearity. The $`\phi `$-components equals zero when $`V=0`$. The lines of the vector fields $`𝐃^{(0)}`$ and $`𝐁^{(0)}`$ in $`z\rho `$-plane for $`V=0`$ are shown in Fig. 2.
Now let us calculate the vector of full angular momentum $`𝐌`$ (12) for field configuration (13) with (21). Because of a symmetry property of the element of integration into (12), for our case we have $`M_\rho =M_\phi =0`$ and
$`M_z`$ $`=`$ $`{\displaystyle \frac{1}{4\pi }}{\displaystyle 𝒫_\phi ^{(0)}\rho (\rho \mathrm{d}z\mathrm{d}\rho \mathrm{d}\phi )}.`$ (22)
Using (21) we can easily obtain the following expression:
$`𝒫_\phi ^{(0)}`$ $`=`$ $`{\displaystyle \frac{4abd\rho }{\stackrel{1}{r}^3\stackrel{2}{r}^3(1V^2)}}.`$ (23)
Substituting (23) into (22) and introducing new variables of integration we obtain
$`M_z`$ $`=`$ $`{\displaystyle \frac{bd}{\pi \left(1V^2\right)}}{\displaystyle \frac{a\rho ^3}{\stackrel{1}{r}^3\stackrel{2}{r}^3}dzd\rho d\phi }`$ (24)
$`=`$ $`{\displaystyle \frac{bd}{\pi }}{\displaystyle \frac{a^{}\rho ^3}{\stackrel{1}{r}^3\stackrel{2}{r}^3}dz^{}d\rho d\phi }`$ (25)
$`=`$ $`{\displaystyle \frac{bd}{\pi }}{\displaystyle \frac{(\rho ^{\prime \prime })^3}{(\stackrel{1}{r}^{\prime \prime })^3(\stackrel{2}{r}^{\prime \prime })^3}dz^{\prime \prime }d\rho ^{\prime \prime }d\phi },`$ (26)
where $`z^{\prime \prime }z^{}/a^{}`$ , $`\rho ^{\prime \prime }\rho /a^{}`$ ,
where $`\stackrel{\prime \prime }{\stackrel{1}{r}}\sqrt{\left(z^{\prime \prime }+1\right)^2+\left(\rho ^{\prime \prime }\right)^2}`$ , $`\stackrel{\prime \prime }{\stackrel{2}{r}}\sqrt{\left(z^{\prime \prime }1\right)^2+\left(\rho ^{\prime \prime }\right)^2}`$ .
As we see, in first change (24)$``$(25) the dependence on speed is canceled. In second change (25)$``$(26) the dependence on $`a^{}`$ is canceled. Thus we obtain that the full angular momentum for field configuration (13), (21) is independent of dyon’s speed ($`V`$) and distance between the dyons ($`2a`$)!
For calculation (26) we introduce the variables of integration $`\{\xi ,\zeta ,\phi \}`$ that appropriate to the bispherical coordinate system with unit parameter characterizing positions of focal points:
$`z^{\prime \prime }={\displaystyle \frac{\mathrm{sinh}\xi }{\mathrm{cosh}\xi \mathrm{cos}\zeta }},\rho ^{\prime \prime }={\displaystyle \frac{\mathrm{sin}\zeta }{\mathrm{cosh}\xi \mathrm{cos}\zeta }}.`$ (27)
The bispherical element of value has the form
$`(\rho ^{\prime \prime }\mathrm{d}z^{\prime \prime }\mathrm{d}\rho ^{\prime \prime }\mathrm{d}\phi )`$ $`=`$ $`{\displaystyle \frac{\mathrm{sin}\zeta \mathrm{d}\xi \mathrm{d}\zeta \mathrm{d}\phi }{\left(\mathrm{cosh}\xi \mathrm{cos}\zeta \right)^3}}.`$ (28)
We have also that
$`\stackrel{\prime \prime }{\stackrel{1}{r}}={\displaystyle \frac{\sqrt{2}\mathrm{exp}\left(\xi /2\right)}{\sqrt{\mathrm{cosh}\xi \mathrm{cos}\zeta }}},\stackrel{\prime \prime }{\stackrel{2}{r}}={\displaystyle \frac{\sqrt{2}\mathrm{exp}\left(\xi /2\right)}{\sqrt{\mathrm{cosh}\xi \mathrm{cos}\zeta }}}.`$ (29)
Substituting (27), (28), (29) into (26) and introducing the variable $`\underset{¯}{z}=\mathrm{cos}\zeta `$, we obtain
$`M_z`$ $`=`$ $`{\displaystyle \frac{bd}{8\pi }}{\displaystyle \underset{\pi }{\overset{\pi }{}}}d\phi {\displaystyle \underset{\mathrm{}}{\overset{\mathrm{}}{}}}d\xi {\displaystyle \underset{1}{\overset{1}{}}}{\displaystyle \frac{\left(1\underset{¯}{z}^2\right)}{\left(\mathrm{cosh}\xi \underset{¯}{z}\right)^2}}d\underset{¯}{z}.`$ (30)
Making firstly the integration over $`\underset{¯}{z}`$ and $`\xi `$ in the finite limits $`[\overline{\xi },\overline{\xi }]`$ we obtain
$$\underset{\mathrm{}}{\overset{\mathrm{}}{}}d\xi \underset{1}{\overset{1}{}}\frac{\left(1\underset{¯}{z}^2\right)}{\left(\mathrm{cosh}\xi \underset{¯}{z}\right)^2}d\underset{¯}{z}=\underset{\overline{\xi }\mathrm{}}{lim}\left[4\left(\mathrm{ln}\frac{\mathrm{cosh}\overline{\xi }+1}{\mathrm{cosh}\overline{\xi }1}\right)\mathrm{sinh}\overline{\xi }\right]=\mathrm{\hspace{0.17em}8}.$$
(31)
Thus we have
$`M_z`$ $`=`$ $`2bd.`$ (32)
Of course, the full angular momentum for an appropriate exact solution is conserved. This implies that we may have internal movements of the singularities, which don’t change the full angular momentum. Thus here we have verified that our choice of two moving dyons (13), (21) as the initial approximation is appropriate.
It is evident that the full momentum for field configuration (13), (21) is zero. To verify satisfaction of the conservation law for full energy, we must take a concrete function $`(𝐃,𝐁)`$. That is, we must investigate the concrete nonlinear electrodynamic model. In this case the condition $`=\mathrm{const}`$ can be used for defining a trajectory $`a(x^0,)`$ of the dyons in the initial approximation. This problem was investigated for Born-Infeld electrodynamics and it was shown that the initial field configuration may behave as nonlinear oscillator. (A wave part of the dyon solutions connected with acceleration of the singular points was not included to the initial approximation.)
The field configuration (13), (21) looks like charged particle with spin. The charge of this particle is $`2d`$ and its spin is equal to $`|M_z|`$. We may set $`2\overline{d}=e`$, where $`e`$ is the absolute value of the electron charge, and $`|M_z|=\mathrm{}/2`$. In this case we have
$`\overline{b}e={\displaystyle \frac{\mathrm{}}{2}}\overline{b}={\displaystyle \frac{e}{2}}{\displaystyle \frac{\mathrm{}}{e^2}}={\displaystyle \frac{e}{2}}{\displaystyle \frac{1}{\overline{\alpha }}}{\displaystyle \frac{\overline{d}}{\overline{b}}}=\overline{\alpha },`$ (33)
where $`\overline{\alpha }=e^2/\mathrm{}1/137`$ is the fine structure constant.
The field configuration (21) is considered here as initial approximation to unknown exact bidyon solution. This initial approximation does not include a wave part of the bidyon solution. For periodical bidyon solution this wave part must have the form of some standing wave localized near dyon singularities. Because this problem has the boundary conditions (which follow from (7), see also ) in two (moving) points (in which the dyon singularities are at a current instant of time), it is possible that the bidyon solution has some discrete set of allowable frequencies.
With the help of Lorentz transformation, from the rest oscillating bidyon we can obtain an appropriate moving bidyon solution. It is evident that the moving oscillating bidyon has both particle and wave properties.
The field model under consideration allows existence for great number of the dyon singularities with charges $`\stackrel{n}{d}`$ and $`\stackrel{n}{b}`$. We may assume that there is some kind invariance of the theory, such that
$$\stackrel{n}{d}=\pm \overline{d},\stackrel{n}{b}=\pm \overline{b}.$$
(34)
For the suitable dimensional system we can take $`\overline{d}=1`$ or $`\overline{b}=1`$. Thus we have the relation $`\overline{d}/\overline{b}`$ as the single dimensionless constant of the theory. We can set that this relation equals the fine structure constant $`\overline{\alpha }`$ (33).
We can fancy a world constructed from the great number of the bidyon-type field configurations as particles. Particles with full angular momentum divisible by $`\mathrm{}/2`$ may be constructed from some number of the bidyons. Because of $`\overline{d}/\overline{b}=\overline{\alpha }1`$, we may build a perturbation theory with the fine structure constant $`\overline{\alpha }`$ as small parameter, for some aspects of the mathematical model of this world. In this case we will have an analogy with the procedure of perturbation theory in quantum electrodynamics.
As a result we can assume that there is some correlation between the bidyon solution of a nonlinear electrodynamic model and leptons. |
warning/0002/hep-th0002209.html | ar5iv | text | # 1 Introduction
## 1 Introduction
The quest for a string representation of the $`D=4`$ continuous Yang-Mills ($`YM_D`$) gauge theory shaped, to a large extent, many branches of the contemporary mathematical physics. Currently, there are two (qualitatively overlapping) candidates for the stringy systems conjecturaly dual to $`YM_D`$. The first one, due to Polyakov, puts forward certain Ansatz for the world-sheet action which is to ensure the invariance of the Wilson loop averages $`<W_C>`$ with respect to the zig-zag backtrackings of the contour $`C`$. The complementary approach has recently sprung to life after daring conjecture of Maldacena (further elaborated in ) about the so-called AdF/CFT correspondence concerning the $`𝒩=4`$ SUSY $`YM_4`$. Being motivated by the latter correspondence, Witten has made a speculative proposal advocating that the ordinary nonsupersymmetric $`YM_4`$ exhibits confinement at least when considered for large $`N`$ in the specific strong coupling (SC) regime. Namely, employing certain SUSY $`YM_{\overline{D}},\overline{D}>D,`$ properly broken to a nonsupersymmetric $`YM_D`$ system, one is to fix an effective finite ultraviolet (UV) cut off $`\mathrm{\Lambda }`$ while keeping the relevant $`YM_D`$ coupling constant(s) sufficiently large.
At present, the support of the conjectured duality-mappings is fairly limited to a few indirect though reasonably compelling arguments. On the other hand, final justification of that or another stringy representation calls either to reproduce the $`YM_D`$ loop-equations in stringy terms or to provide with an explicit transformation of the gauge theory into a sort of string theory. It is our goal to make a step in the second direction approaching (with new tools) the old challenge: the exact reformulation of the continuous $`D3`$ $`YM_D`$ theory in terms of the microscopic colour-electric YM-fluxes. More specifically, we focus on the pattern of the smooth Gauge String inherent in the presumable large $`N`$ SC expansion (as opposed to the standard weak coupling (WC) series) valid in the SC regime similar to that of the Witten’s proposal . Complementary, the considered generic strongly coupled continuous $`YM_D`$ models in $`D=4`$ might be viewed as the prototypes of the effective low-energy $`YM_4`$ theory. The latter is supposed to result via the Wilsonian renormgroup (RG) flow of the effective actions (starting from the standard $`YM_4`$ in the WC phase) up to the confinement scale. In this perspective, the ’built in’ UV cut off $`\mathrm{\Lambda }`$ in $`D=4`$ is to be qualitatively identified with a physical scale $`\mathrm{\Lambda }_{YM_4}`$ that is of order of the lowest glueball mass.
The short-cut way to a continuous model of smooth YM-fluxes is suggested by the two-dimensional analysis. Here, the pattern of the flux-theory proposed by Gross and Taylor is encoded not only in the continuous $`YM_2`$ systems on a $`2d`$ manifold. It is also inherent in the associated RG invariant $`2d`$ lattice gauge models introduced via the plaquette-factor
$$\begin{array}{c}Z(\{\stackrel{~}{b}_k\}|U)=_Re^{F(R)}\chi _R(U);e^{F(R)}=dimRe^{\mathrm{\Gamma }(\{\stackrel{~}{b}_k\},N,\{C_p(R)\})},\end{array}$$
(1.1)
where $`dimR,\chi _R(U),`$ and $`C_p(R),p=1,\mathrm{},N,`$ stand respectively for the dimension, character and $`p`$th order Casimir operator associated to a given $`SU(N)`$ irreducible representation (irrep) $`R`$ (while $`\{\stackrel{~}{b}_kN^0\}`$ denotes a set of the dimensionless coupling constants). The key observation is that, in the $`D3`$ lattice $`YM_D`$ systems (1.1), the pattern of the appropriately constructed flux-theory as well refers (owing to a subtle $`D3`$ ’descedant’ of the $`2d`$ RG invariance) to the properly associated continuous $`YM_D`$ models.
To support this assertion, I first reformulate the strongly coupled $`D3`$ lattice $`YM_D`$ models (1.1) in terms of the Gauge String which does appropriately extend the $`D=2`$ Gross-Taylor stringy pattern into higher dimensions. This lattice theory of the YM-flux is endowed with certain continuous (rather than discrete as one might expect from the lattice formulation) group of the area-preserving homeomorphisms. In turn, the latter symmetry ensures that the considered $`D2`$ pattern of the lattice flux-theory can be employed to unambiguously define the associated $`D2`$ smooth Gauge String invariant under the area-preserving diffeomorphisms. The remarkable thing is that, in the latter continuous $`D3`$ flux-theory, one can identify (see Section 8) such SC conglomerates of the (piecewise) smooth flux-worldsheets which are in one-to-one correspondence with the judiciously associated varieties of the WC Feynman diagrams on the side of the properly specified continuous $`YM_D`$ model. Therefore, the proposed smooth Gauge String provides with the concrete realization of the old expectation (for a recent discussion see ) that certain nonperturbative effects ’close up’ the windows of the Feynman diagrams trading the latter for the string worldsheets.
More specifically, building on the nonabelian duality transformation recently proposed by the author, we show that in the $`D2`$ lattice systems (1.1) the free energy and the Wilson (multi)loop observables can be rewritten in terms of the following statistics of strings. The lattice weight $`w[\stackrel{~}{M}(T)]`$ of a given connected worldsheet $`\stackrel{~}{M}(T)`$ (with the support on a subspace, represented by certain $`2d`$ cell-complex T, of the $`2d`$ skeleton of the $`D`$-dimensional base-lattice)
$$\begin{array}{c}w[\stackrel{~}{M}(T)]=exp[\overline{A}\mathrm{\Lambda }^2\stackrel{~}{\sigma }_0(\{\stackrel{~}{b}_k\})]N^{22hb}J[\stackrel{~}{M}(T)|\{\stackrel{~}{b}_k\}],\end{array}$$
(1.2)
is composed of the three different blocks which altogether conspire so that the contribution of the strings with any backtrackings (i.e. foldings) is zero. In eq. (1.2), the first factor is the exponent of the Nambu-Goto term proportional to the total area $`\overline{A}`$ of $`\stackrel{~}{M}`$. The D-independent bare string tension $`\sigma _0=\mathrm{\Lambda }^2\stackrel{~}{\sigma }_0`$ (to be defined by eq. (1.6) below) is measured in the units of the UV cut off $`\mathrm{\Lambda }`$ squared (which, in $`D3`$, regularizes the ’transverse’ string fluctuations). Next, there appears the standard ’t Hooft factor, where h and b are respectively the genus and the number of the boundary contours $`C_k,k=1,\mathrm{},b,`$ of $`\stackrel{~}{M}`$.
Finally, given a particular model (1.1), the third term $`J[\stackrel{~}{M}(T)|\{\stackrel{~}{b}_k\}]`$ (being equal to unity for a nonselfintersecting surface $`\stackrel{~}{M}`$) is sensitive only to the topology (but not to the geometry) of selfintersections of $`\stackrel{~}{M}`$. In particular, the dependence of $`J[\stackrel{~}{M}(T)|\{\stackrel{~}{b}_k\}]`$ on the coupling constants $`\{\stackrel{~}{b}_k\}`$ is collected from the elementary weights assigned to the admissible ’movable’ singularities (to be specified later on) of the associated map $`\phi :\stackrel{~}{M}T`$. As a result, the last term (similarly to the remaining ones) is invariant under the required continuous group of the area-preserving homeomorphisms so that the pattern (1.2) directly applicable to a generic smooth worldsheet $`\stackrel{~}{M}`$.
Next, consider the $`D3`$ continuous flux-theory defined as the statistics of the smooth worldsheets $`\stackrel{~}{M}`$ which are postulated to be endowed with the weights (1.2) corresponding to the smooth mappings of $`\stackrel{~}{M}`$ into the Euclidean space $`𝐑^𝐃`$. In compliance with the above asserted SC/WC correspondence, a given specification (1.1) the smooth Gauge String refers to the following unique continuous $`YM_D`$ model. The local lagrangian of the latter is to be reconstructed as the $`D`$-dimensional ’pull-back’
$$\begin{array}{c}L_2(F)L_D(F):\{F_{\mu \nu }^a;\mu ,\nu =1,2\}\{F_{\mu \nu }^a;\mu ,\nu =1,\mathrm{},D\}\end{array}$$
(1.3)
of the $`D=2`$ lagrangian (composed of the $`O(D)`$-invariant combinations of the $`F_{\mu \nu }^aF`$ tensor)
$$\begin{array}{c}D=2:L_D(F)=\mathrm{\Lambda }^{D4}_{n2}_{rY_n}\frac{[tr(F^k)]^{p_k}}{g_r(\{\stackrel{~}{b}_k\})},\end{array}$$
(1.4)
associated to such continuous $`YM_2`$ theory that its partition function on a $`2d`$ disc (of unit area with the free boundary conditions) is equal to the plaquette-factor (1.1). (In $`D3`$, within each trace of (1.4), the involved $`(F_{\mu \nu })_j^i`$-factors are prescribed to be totally symmetrized to exclude the $`[F_{\mu \nu },F_{\rho \sigma }]`$-commutator dependent terms identically vanishing in $`D=2`$.) To separate out the ’kinematical’ D-dependent rescaling of the coupling constants, we have introduced the parameter $`\mathrm{\Lambda }`$. The latter is to be identified with the effective UV cut off, for the considered strongly coupled $`YM_D`$ systems, predetermined by the necessary regularization of the $`D3`$ YM-flux transverse fluctuations. Remark also that the N\- and $`\mathrm{\Lambda }`$-dependent coupling constants $`g_r(\{\stackrel{~}{b}_k\})g_r(\{\stackrel{~}{b}_k\},N,\mathrm{\Lambda })`$ are canonically labelled by the $`S(n)`$ irreps $`rY_n`$ (parametrized by the partitions of $`n`$: $`_{k=1}^nkp_k=n`$).
In sum, our proposal is that the particular theory of the smooth Gauge String, thus induced through (1.1), reproduces (to all orders in $`1/N`$) the corresponding local continuous $`D3`$ $`YM_D`$ model (1.4) for sufficiently large coupling constants $`\{g_rN^{2_{k=1}^np_k}O(N^0)\}`$. (The latter constants are measured in the units of $`1/N`$ and $`\mathrm{\Lambda }`$ akin to eq. (9.1).) Note that in $`D=2`$, after proper identification of the coupling constants, this correspondence is justified by the matching with the $`2d`$ pattern of Gross and Taylor. One may also expect that the asserted $`YM_D/String`$ correspondence is unique to the extent that the stringy degrees of freedom are directly identified with the worldsheets of the microscopic, conserved YM-flux immersed (by mappings with certain admissible singularities to be specified after eq. (1.11) below) into a $`D`$-dimensional base-space.
The above constriant on $`\{g_r\}`$ (selecting the $`SC`$ regime) accounts for the fact that in $`D3`$ the proposed $`YM_D/String`$ duality is limited by the stability of the YM-flux. In particular, the physical string tension $`\sigma _{ph}`$ (entering e.g. the asymptotics of the Wilson loop averages) must be positive. As we will see, despite the extra $`J[..]`$-factor in eq. (1.2), the ’bare’ and the physical string tensions (defined within the $`1/N`$ expansion) are conventionally related so that the latter condition in the large $`N`$ limit reads
$$\begin{array}{c}\{g_r(\{\stackrel{~}{b}_k\},..)\}:\sigma _{ph}=(\stackrel{~}{\sigma }_0(\{\stackrel{~}{b}_k\})\zeta _D)\mathrm{\Lambda }^2>0,\end{array}$$
(1.5)
i.e. the ’bare’ part $`\sigma _0=\mathrm{\Lambda }^2\stackrel{~}{\sigma }_0`$ should be larger than the $`D3`$ entropy contribution $`\delta \sigma _{ent}=\zeta _D\mathrm{\Lambda }^2`$ due to the transverse string fluctuations. Among our results, the central one is the explicit $`\{\stackrel{~}{b}_k\}`$-dependence of $`\sigma _0=\mathrm{\Lambda }^2\stackrel{~}{\sigma }_0(\{\stackrel{~}{b}_k\})`$ entering the worldsheet weights (1.2). It is uniquely reconstructed
$$\begin{array}{c}\mathrm{\Gamma }(\{\stackrel{~}{b}_k\},N,\{C_p(R)\})=n\stackrel{~}{\sigma }_0(\{\stackrel{~}{b}_k\})(1+O(N^1));\sigma _0=\mathrm{\Lambda }^2\stackrel{~}{\sigma }_0O(N^0),\end{array}$$
(1.6)
from the formal $`1/N`$ asymptotics (with $`RY_n^{(N)},nO(N^0)`$) of the admissible function $`\mathrm{\Gamma }(..)`$ defining the associated lattice model (1.1). For example, the smooth Gauge String induced from the Heat-Kernal action
$$\begin{array}{c}Z(g^2|U)=_{n0}_{RY_n^{(N)}}dimR\chi _R(U)exp[\stackrel{~}{g}^2C_2(R)/2],\end{array}$$
(1.7)
being presumably dual to the standard $`\mathrm{𝑌𝑀}_D`$ theory with $`g^2=\mathrm{\Lambda }^{4D}\stackrel{~}{g}^21/N`$, refers to the following bare string tension (with $`g^21/N`$)
$$\begin{array}{c}L_D(F)=tr(F_{\mu \nu }^2)/4g^2\sigma _0=g^2N\mathrm{\Lambda }^{D2}/2.\end{array}$$
(1.8)
Summarizing, the predicted pattern (1.6) of $`\sigma _0`$ implies that the continuous $`D3`$ $`YM_D`$ systems are confining at least in the large $`N`$ SC regime belonging to the domain (1.5). Although in this regime the standard WC series are expected to fail, the latter SC phenomenon suggests a mechanism of confinement for $`YM_4`$ in the WC phase (see Conclusions). Remark also that the physical string tension in eq. (1.5) generically is not adjusted to be infinitely less than the squared UV cut off $`\mathrm{\Lambda }`$ (which in $`D=4`$ matches with the above identification $`\mathrm{\Lambda }\mathrm{\Lambda }_{YM_4}`$). Actually, in the considered $`YM_D`$ models the introduced smooth flux-worlsheets might become unstable degrees of freedom even prior to the saturation of (1.5) owing to the presumable large $`N`$ phase transition(s) (generalizing the $`D=2`$ situation ).
Finally, the distinguished regime of the smooth Gauge String is the extreme large N SC limit where $`\sigma _0(\{\stackrel{~}{b}_k\})>>\mathrm{\Lambda }^2`$. As in eq. (1.5) the entropy-constant $`\zeta _D`$ is $`\{\stackrel{~}{b}_k\}`$-independent, in this limit the leading order of the Wilson loop averages $`<W_C>`$ is given by the minimal area contiribution that allows for a number of nontrivial predictions. First of all, the physical string tension in this regime merges with the ’bare’ one, $`\sigma _{ph}\sigma _0`$, determined by eq. (1.6). (We will show in Section 9 that the Witten’s SC asymptotics of $`\sigma _{ph}`$ in $`D=4`$ is semiquantitatively consistent with the latter prediction). Also, provided the minimal area worldsheet $`\stackrel{~}{M}_{min}(C)`$ has the support $`T_{min}(C)`$ on a $`2d`$ manifold (rather than on a $`2d`$ cell-complex), in this limit the pattern of the $`D3`$ averages $`<W_C>`$ is reduced to the one in the corresponding continuous $`D=2`$ $`YM_2`$ theory (1.4) conventionally defined on $`T_{min}(C)`$.
### 1.1 The $`D2`$ pattern of the Gauge String weights.
$$𝐀.\mathrm{𝐓𝐡𝐞}𝐃=\mathrm{𝟐}\mathrm{𝐜𝐚𝐬𝐞}.$$
Recall that, in the $`2d`$ framework initiated by Gross and Taylor, the partition function $`\stackrel{~}{X}_M`$ of a given continuous $`YM_2`$ theory on a $`2d`$ manifold $`M`$ of the area $`A`$ (without boundaries) is rewritten as the sum
$$\begin{array}{c}\stackrel{~}{X}_M=_{\stackrel{~}{M}}df(1)^{P_f}\frac{N^{\chi _f}}{|C_f|}exp[\sigma _0n_fA];f:\stackrel{~}{M}M,\end{array}$$
(1.9)
over the topologically distinct branched covering spaces $`\stackrel{~}{M}`$ of $`M`$ specified by the mappings f. The latter maps locally satisfy the definition of the immersion $`\stackrel{~}{M}M`$ everywhere on $`M`$ except for a set of isolated points where the singularities, corresponding to a branch-point or/and to a collapsed (to a point) subsurface connecting a few sheets, are allowed. To visualize the pattern of $`\stackrel{~}{M}`$ as the associated Riemann surface (without foldings but with the singularities to be viewed as certain collapsed $`2d`$ subsurfaces) identified with a particular string worldsheet $`\stackrel{~}{M}`$ of the total Euler characteristic $`\chi _f`$, I propose the cutting-gluing rules which readily generalize to the $`D3`$ case. As for $`|C_f|,P_f`$, and $`n_f`$, in eq. (1.9) they denote respectively the symmetry factor (i.e. the number of distinct automorphisms $`\kappa `$ of $`f:f\kappa =f`$), the ’parity’, and the degree (i.e. the number of the covering sheets) of the map $`f`$. The sum (1.9) includes all admissible disconnected contributions in such a way that the free energy $`ln[\stackrel{~}{X}_M]`$ is supposed to be deduced from (1.9) constraining the worldsheets $`\stackrel{~}{M}`$ to be connected.
To make contact with the $`T=M,\overline{A}=n_fA`$ case of (1.2), observe first that the summation (1.9) over the maps $`f`$ effectively enumerates the admissible selfintersections of $`\stackrel{~}{M}`$. In particular, the sum over the singularities of $`f`$ implies the one over the ’nonmovable’ and (the positions of) the ’movable’ branch points parametrized by the corresponding cyclic decomposition $`\{p\}`$: $`_{k=1}^nkp_k=n_f`$ of $`n_f`$. These two types of the points are respectively assigned with the nontrivial weights $`\stackrel{~}{w}_{\{p\}}`$ and $`w_{\{p\}}(\{\stackrel{~}{b}_k\})`$ (included into the measure $`df`$ of (0.5)) which are $`\{\stackrel{~}{b}_k\}`$-dependent only in the ’movable’ case. (The $`w_{\{p\}}`$-, $`\stackrel{~}{w}_{\{p\}}`$-factors, together with $`(1)^{P_f}/|C_f|`$, are collected into $`J[\stackrel{~}{M}(T)|\{\stackrel{~}{b}_k\}]`$ of (1.2).) Therefore, the summation $`df`$ implies in particular the multiple integrals
$$\begin{array}{c}_{\{p\}}\frac{[w_{\{p\}}(\{\stackrel{~}{b}_k\})]^{m_{\{p\}}}}{[m_{\{p\}}]!}_{k=1}^{m_{\{p\}}}_Md^2X_{\{p\}}^{(k)}\end{array}$$
(1.10)
over all positions $`X_{\{p\}}^{(k)}`$ on M of (a given number $`m_{\{p\}}`$ of) the movable branch points that introduces additional dependence on the area $`A=_Md^2X_{\{p\}}^{(k)}`$ of M. (’Nonmovable’ singularities of the $`f`$-map can be placed anywhere on $`M`$ but do not carry any area- or $`\{\stackrel{~}{b}_k\}`$-dependent factors.) Finally, the bare string tension $`\sigma _0=\mathrm{\Lambda }^2\stackrel{~}{\sigma }_0(\{\stackrel{~}{b}_k\})`$ (defined via eq. (1.6)) enters the exponent of (1.9) as the no-fold variant of the Nambu-Goto term. Our attention is mostly confined to the option (1.8) corresponding to the $`D2`$ $`SU(N)`$ Heat-Kernal lattice gauge theory (1.7). In this case, $`\sigma _0=\mathrm{\Lambda }^2\stackrel{~}{g}^2N/2`$ and there are only the simple transposition branch points (i.e. the ones connecting a pair of the sheets) weighted by $`w_2=\stackrel{~}{g}^2N`$.
Finally, the $`D=2`$ representation (1.9) of the continuous $`YM_2`$ theories (1.4) is valid also in the corresponding $`2d`$ lattice gauge systems (1.1) on a discretized surface $`M`$. The origin of this two-fold interpretation of (1.9) resides in the RG invariance of the latter lattice models. Complementary, owing to the symmetry of the continuous $`YM_2`$ models (1.4) under the area-preserving diffeomorphisms, the RG invariance of (1.1) results in the invariance of the lattice weights $`w[\stackrel{~}{M}]`$ under the continuous group of the area-preserving worldsheet homeomorphisms. These homeomorphisms continuously (rather than discretely) translate the positions of the singularities of the map $`f`$ everywhere on $`M`$ including interiors of the plaquettes.
$$𝐁.\mathrm{𝐄𝐱𝐭𝐞𝐧𝐬𝐢𝐨𝐧}\mathrm{𝐭𝐨}\mathrm{𝐭𝐡𝐞}𝐃\mathrm{𝟑}\mathrm{𝐜𝐚𝐬𝐞}.$$
A given $`D3`$ lattice $`YM_D`$ system (1.1) (having some lattice $`𝐋^𝐃`$ as the base-space $`B`$) can be equally viewed as the $`YM`$ model (1.1) defined on the $`2d`$ skeleton $`𝐓^𝐃`$ of $`𝐋^𝐃`$ represented by the associated $`2d`$ cell-complex. As we will discuss in Section 7, this reformulation allows to reveal certain $`D3`$ ’descedant’ of the $`2d`$ RG invariance of (1.1) that in turn foreshadows the existence of the $`D3`$ extension of the pattern (1.9).
Consider the partition function $`\stackrel{~}{X}_{L^D}`$ of the $`YM_D`$ system (1.1) defined on $`B=𝐋^𝐃`$. Combining the nonabelian duality transformation with the methods of algebraic topology , we derive that the pattern (1.9) remains to be valid in $`D3`$ with the only modification. It concerns the structure of the relevant mappings (to be summed over) specifying the admissible topology of the worldsheets $`\stackrel{~}{M}`$ of the YM-flux. The $`D=2`$ maps $`f`$ are superseded by the mappings $`\phi `$
$$\begin{array}{c}\phi :\stackrel{~}{M}T𝐓^𝐃,\end{array}$$
(1.11)
of a $`2d`$ surface $`\stackrel{~}{M}`$ onto a given subspace T (represented by a $`2d`$ cell-complex) of the $`2d`$ skeleton $`𝐓^𝐃`$ of the $`D`$-dimensional base-lattice $`𝐋^𝐃`$. Akin to the $`D=2`$ case, the maps $`\phi `$ locally comply with the definition of the immersion $`\stackrel{~}{M}T`$ anywhere except for a countable set of points on $`T`$ where certain singularities are allowed. The relevant (for the smooth implementation of the Gauge String) ones are either of the same type as in the $`2d`$ $`f`$-mappings (1.9), i.e. the branch points or/and the collapsed subsurfaces, or of the type corresponding to the homotopy retraction (see e.g. Appendix C) of the latter irregularities. In particular, the included into $`d\phi `$-measure $`\{\stackrel{~}{b}_k\}`$-dependent weights $`w_{\{p\}}(\{\stackrel{~}{b}_k\})`$ (of the movable branch points) in $`D3`$ are equal to their counterparts in the $`2d`$ case (1.10) associated to the same continuous $`YM_2`$ theory (1.4). (Akin to eq. (1.9), one can argue that the free energy $`ln[\stackrel{~}{X}_{L^D}]`$ is provided by the restriciton of the sum for partition function $`\stackrel{~}{X}_{L^D}`$ to the one over the connected worldsheets $`\stackrel{~}{M}`$.)
The $`D3`$ nature of the mappings (1.11) is reflected by the fact that, owing to possible ’higher-dimensional’ selfintersections of $`\stackrel{~}{M}(T)`$, the taget-space T generically is represented by the $`2d`$ cell-complex $`T=_kE_k`$ rather than by a $`2d`$ surface. (Recall that any $`2d`$ cell-complex, after cutting out along the links shared by more than two plaquettes, becomes a disjoint union of $`2d`$ surfaces $`E_k`$ of the areas $`A_k`$ with certain boundaries. Conversely, the set $`\{E_k\}`$ can be combined back into $`T=_kE_k`$ according to the incidence numbers corresponding to the entries of the associated incidence-matrix.) Therefore, the construction of $`\stackrel{~}{M}(T)`$ (via our cutting-gluing algorithm developed for the canonical branched coverings (1.9)) entails the appropriate generalization of the notion of the Riemann surface (with the total area $`\overline{A}=_kn_\phi ^{(k)}A_k`$). Alternatively, $`\stackrel{~}{M}(T)`$ can be viewed as the specific generalization of the branched covering space (wrapped around T) suitable for the analysis of the Gauge String.
Next, as it is shown in Section 8, the total contribution of the worldsheets $`\stackrel{~}{M}`$ with arbitrary backtrackings (that, generalizing the $`D=2`$ case, bound any zero 3-volume) vanishes. The simplest $`J[\stackrel{~}{M}(T)|\{\stackrel{~}{b}_k\}]=1`$ pattern (1.2) of $`w[\stackrel{~}{M}]`$ in $`D2`$ arises when the genus $`h`$ connected string worldsheet $`\stackrel{~}{M}(T)`$ is represented by an embedding $`T=\stackrel{~}{M}(T)`$ (i.e. $`\stackrel{~}{M}(T)`$ does not selfintersect). Nevertheless, the set of all possible $`J[\stackrel{~}{M}(T)|\{\stackrel{~}{b}_k\}]=1`$ embedding-weights (1.2) does not discriminate between those distinct models (1.1)/(1.4) which provide with one and the same bare string tension (1.6). To distinguish between the different models, the specification of the remaining immersion-weights (assigned to the selfintersecting string worlsheets) is indispensable. In this way, one reconstructs the data encoded in a given continuous $`YM_2`$ theory (1.4) that originally served as the bridge (1.3) between the Gauge String and the corresponding continuous $`YM_D`$ model.
Owing to the asserted pattern, the $`D3`$ lattice weights $`w[\stackrel{~}{M}(T)]`$ are invariant under certain continuous group of the area-preserving homeomorphisms extending the $`D=2`$ ones. As a result, the set $`\{w[\stackrel{~}{M}(T)]\}`$ can be unambiguously used to introduce the statistics of the (piecewise) smooth YM flux-worldsheets $`\stackrel{~}{M}(T)`$, with the latter homeomorphisms being traded for the corresponding diffeomorphisms. As for the sum over the worldsheets, it is specified by the one running over the $`\phi `$-mappings (1.11). This time, one is to consider the piecewise smooth immersions (with the admissible singularities which, in $`D4`$, can be restricted to the ones listed after eq. (1.11)) of the $`2d`$ manifolds $`M`$ into the $`D`$-dimensional space-time $`B=𝐑^𝐃`$ that results in the worldsheet $`\stackrel{~}{M}`$ with the support on $`T𝐑^𝐃`$. In turn, it constitutes the proper class of the $`2d`$ cell-complexes T which can play the role of the taget-spaces for the considered (piecewise) smooth maps $`\phi `$.
As we will discuss in Section 8, the $`D3`$ dynamics of the smooth Gauge String in certain sense is dramatically simpler compared to its lattice counterpart. First of all, it is convenient to make certain redefinition of the bare string tension (that amounts to the substitution $`\stackrel{~}{\sigma }_0=\lambda /2\lambda (11/N^2)/2`$ in the case of (1.7)) to get rid of a redundant subset of the ’movable’ singularities present in the $`SU(N)`$ mappings (1.11) even for a 1-sheet covering of $`T`$. Given this modification, consider the subset of the smooth worldsheets $`\stackrel{~}{M}(T)`$ (without boundaries) which are constrained to be strictly nonselfintersecting in $`D5`$ and allowed to selfintersect at an arbitrary union of isolated points in $`D=4`$. (Actually the latter condition can weakened e.g. to include non(self)intersecting boundary contour(s)). They are assigned with the simplest $`J[\stackrel{~}{M}(T)|\{\stackrel{~}{b}_k\}]=1`$ weight-pattern (1.2). The point is that in $`D4`$ the latter subset is dense in the set of all worldsheets $`\stackrel{~}{M}`$ parametrized by the (piecewise) smooth mappings (1.11) into $`𝐑^𝐃`$, provided the latter redefinition of $`\stackrel{~}{\sigma }_0`$. In particular, it justifies (at least in $`D4`$) the validity of the simple relation (1.5). As for the more complicated pattern (1.2) of the weights, in $`D4`$ it is observable for example within such Wilson loop averages $`<W_C>`$ where the corresponding minimal surface $`S_{min}(C)`$ selfintersects (or when the boundary contour $`C`$ has some zig-zag backtrackings).
## 2 Outline of the further content.
To make the analysis of the $`D2`$ lattice Gauge String more concise, we employ the Twisted Eguchi-Kawai (TEK) representation of the large $`N`$ ’infinite-lattice’ SU(N) gauge systems. Recall that in the limit $`N\mathrm{}`$ the partition function (PF) $`\stackrel{~}{X}_{L^D}`$ of a lattice $`YM_D`$ theory like (1.1) in a $`D=2p`$ volume $`L^D=N^2`$ can be reproduced (at least within both the SC and the WC series)
$$\begin{array}{c}lim_N\mathrm{}\stackrel{~}{X}_{L^D}=lim_N\mathrm{}\left(_{\rho =1}^DdU_\rho _{\mu \nu =1}^{D(D1)/2}Z(g^2|tU_\mu U_\nu U_\mu ^+U_\nu ^+)\right)^{L^D},\end{array}$$
(2.1)
through the PF $`\stackrel{~}{X}_D`$ of the associated $`D`$-matrix SU(N) TEK model with the reduced space-time dependence (where $`t=\mathrm{exp}[i2\pi /N^{\frac{2}{D}}]𝐙_𝐍`$). Therefore, in the $`N\mathrm{}`$ SC phase, the free energy of (2.1) yields the generating functional for the $`YM_D`$ string-weights assigned to the worldsheets corresponding to the $`2d`$ spheres immersed into $`𝐋^𝐃`$. On the side of the TEK model, the surfaces are ’wrapped around’ the EK base-lattice $`T_{EK}`$ which has the topology of the $`2d`$-skeleton of the $`D`$-dimensional cube with periodic boundary conditions. Thus, $`T_{EK}`$ is homeomorphic to the $`2d`$ cell-complex visualized as the union
$$\begin{array}{c}T_{EK}=_{\mu \nu =1}^{D(D1)/2}E_{\mu \nu },E_{\rho \mu }E_{\rho \nu }=l_\rho ,\end{array}$$
(2.2)
of the $`D(D1)/2`$ mutually intertwined 2-tora $`E_{\mu \nu }`$ sharing D uncontractible cycles (i.e. compactified links of the $`D`$-cube) $`l_\rho `$ in common.
To proceed further, let us first label the SU(N) irreps summed up in each $`\mu \nu `$-species of the Z-factor (1.1) entering (2.1) by
$$\begin{array}{c}R_{\mu \nu }Y_{n_{\mu \nu }}^{(N)},n_+=\underset{\mu \nu =1}{\overset{D(D1)/2}{}}n_{\mu \nu }.\end{array}$$
(2.3)
It is convenient to rewrite the SU(N) TEK PF (2.1) in the following form
$$\begin{array}{c}\stackrel{~}{X}_D=_{\{n_{\mu \nu }\}}_{\{R_{\mu \nu }Y_{n_{\mu \nu }}^{(N)}\}}t^{n_+}e^{S(\{R_{\mu \nu }\})}B(\{R_{\mu \nu }\})=_{\{n_{\mu \nu }\}}t^{n_+}B(\{n_{\mu \nu }\}),\end{array}$$
(2.4)
introducing the elementary master-integrals $`B(\{R_{\mu \nu }\})`$ (which are then composed into $`B(\{n_{\mu \nu }\})`$)
$$\begin{array}{c}B(\{R_{\mu \nu }\})=_{\{\rho \}}dU_\rho _{\{\mu \nu \}}\chi _{R_{\mu \nu }}(U_{\mu \nu });e^{S(\{R_{\mu \nu }\})}=_{\{\mu \nu \}}e^{F(R_{\mu \nu })},\end{array}$$
(2.5)
where $`U_{\mu \nu }U_\mu U_\nu U_\mu ^+U_\nu ^+`$, $`F(R)`$ is defined by eq. (1.1), and we have used the identity $`\chi _R(tV)=t^{n(R)}\chi _R(V),tZ_N,RY_{n(R)}^{(N)}`$. In the SC phase, we are concerned below, the twist-factor t in eq. (2.4) is irrelevant in the limit $`N\mathrm{}`$ so that the TEK model is reduced to the original Eguchi-Kawai one. In this regime, the $`t=1`$ correspondence (2.1) is supposed to be valid in any (not necessarily even) $`D2`$.
Next, making use of the nonabelian duality transformation , in Section 3 we rewrite the TEK master-integral $`B(\{n_{\mu \nu }\})`$ as the weighted sum of the $`Tr_{n_+}`$-characters (i.e. traces)
$$\begin{array}{c}B(\{n_{\mu \nu }\})=_{\{R_{\mu \nu }Y_{n_{\mu \nu }}^{(N)}\}}e^{S(\{R_{\mu \nu }\})}Tr_{n_+}[𝐃(A_{n_+}(\{R_{\mu \nu }\}))],\end{array}$$
(2.6)
of certain master-elements $`A_{n_+}(\{R_{\mu \nu }\})=_{\sigma S(n_+)}a(\sigma |\{R_{\mu \nu }\})\sigma `$. Being defined by eqs. (3.34),(3.35), they take values in the tensor representation of the $`S(n_+)`$ algebra (with $`n_+`$ being given by eq. (2.3)). The latter is deduced by linearity from the canonical representation for $`S(n)`$-group elements $`\sigma `$
$$\begin{array}{c}𝐃(\sigma )_{\{j^n\}}^{\{i^n\}}=\delta _{j_1}^{i_{\sigma (1)}}\delta _{j_2}^{i_{\sigma (2)}}\mathrm{}\delta _{j_n}^{i_{\sigma (n)}};\widehat{\sigma }:k\sigma (k),k=1,\mathrm{},n,\end{array}$$
(2.7)
where $`\delta _j^i`$ denotes the ’N-dimensional’ Kronecker delta function.
To relate (2.6) with the stringy pattern like (1.9), in Section 4 we represent this equation in the form suitable for the algebraic definition of the topological data. First of all, one observes (see Appendix B) that the $`Tr_{n_+}`$-trace in eq. (2.6) is simply related to the associated character of the regular $`S(n_+)`$-representation. As a result, $`B(\{n_{\mu \nu }\})`$ can be rewritten as the delta-function on the $`S(n_+)`$-algebra (that selects the contribution of the weight of the $`S(n_+)`$ unity-permutation $`\widehat{1}_{[n_+]}`$)
$$\begin{array}{c}B(\{n_{\mu \nu }\})=Tr_{n_+}[𝐃(\stackrel{~}{A}_{n_+})]=\delta _{n_+}(\mathrm{\Lambda }_{n_+}^{(1)}\stackrel{~}{A}_{n_+})\end{array}$$
(2.8)
that already played the important role in the $`D=2`$ analysis . Next, employing the Schur-Weyl duality, both the operator $`\mathrm{\Lambda }_{n_+}^{(1)}`$ (defined by the U(N) variant of eq. (3.9), see Section 3) and
$$\begin{array}{c}\stackrel{~}{A}_{n_+}=_{\{R_{\mu \nu }Y_{n_{\mu \nu }}^{(N)}\}}e^{S(\{R_{\mu \nu }\})}A_{n_+}(\{R_{\mu \nu }\})\end{array}$$
(2.9)
are reformulated entirely in terms of the symmetric group elements (i.e. no SU(N) irreps $`R_\varphi `$ are left). The resulting expression (4.8) is suitable to specify the mappings (1.11). The $`D3`$ nature of (1.11) is reflected by the outer-product structure of $`\stackrel{~}{A}_{n_+}`$ which is represented as certain combination of the $`S(n_\varphi )`$-blocks embedded to act in the common enveloping space of the $`S(n_+)=_\varphi S(n_\varphi )`$ algebra (with $`\varphi \{\mu \nu \},\{\rho \}`$). The technique, dealing with such compositions, is naturally inherited from the nonabelian duality transformation .
The remaining material is organized as following. In Section 5, we rederive by our methods the Baez-Taylor reformulation of the Gross-Taylor $`2d`$ pattern (1.9) for the Heat-Kernal model (1.7). We also briefly sketch how our method generalizes for any admissible model (1.1). Being more compact than the Gross-Taylor one, the $`D=2`$ representation of the Baez-Taylor type has the structure where the full stringy pattern is not entirely manifest. To circumvent this problem we formulate a simple prescription how to transform the latter into the former.
The $`D3`$ generalization of the $`2d`$ Baez-Taylor representation is derived in Section 6 for the Eguchi-Kawai models (2.1) which fully justifies the announced pattern (1.11) of the $`D3`$ mappings. (In particular, we make certain conjecture concerning the concise topological reinterpretation, generalizing the $`D=2`$ one , of the $`D3`$ sums like (1.9) in the formal topological limit when $`\mathrm{\Gamma }(..)0`$.) As well as in the $`D=2`$ case, to relate the obtained $`D3`$ representation with the manifest stringy pattern (1.2), a natural extension of the $`D=2`$ prescription (formulated in Section 5) is suggested. In Section 7, we discuss the structure of the asserted continuous group of the area-preserving homeomorphisms (the generic lattice Gauge String weights $`w[\stackrel{~}{M}]`$ are endowed with) and make a brief comparison with the earlier large N variants of the Wilson’s SC expansion devoid of the latter invariance. In Section 8, we discuss the major qualitative features of the Gauge String accentuating the similarities and differences between the continuous flux-theory and the conventional paradigm of the $`D3`$ ’fundamental’ strings. In the last section, we put forward a speculative proposal for the mechanism of confinement in the standard weakly-coupled continuous gauge theory (1.8) at large N. Also a preliminary contact with the two existing stringy proposals is made. Finally, the Appendices contain technical pieces of some derivations used in the main text.
## 3 The Dual Representation of $`\stackrel{~}{X}_D`$.
To derive (2.6), we apply the nonabelian duality transformation to the partition function of the SU(N) TEK model (2.1)/(1.1) following the general algorithm formulated in for a generic SU(N) $`D`$-matrix system. To begin with, the $`SU(N)`$ character is to be represented in the form (see e.g. ) reminiscent of the one of eq. (2.6). Let $`C_R`$ denote the canonical Young idempotent proportional, $`P_R=d_RC_R`$, to the Young projector
$$\begin{array}{c}P_R=\frac{d_R}{n!}_{\sigma S(n)}\chi _R(\sigma )\sigma ,RY_n,\end{array}$$
(3.1)
where $`\chi _R(\sigma ),d_R`$ are the character and the dimension associated to the $`S(n)`$ irrep $`R`$ (while $`[P_R,\sigma ]=0,\sigma S(n)`$). Then, $`\chi _R(U)`$ assumes (akin to (2.6)) the form of the trace: $`\chi _R(U)=Tr_n[𝐃(C_R)U^n],RY_n^{(N)}`$, where
$$\begin{array}{c}Tr_n[𝐃(\sigma )U^n]=_{i_1i_2..i_n=1}^NU_{i_1}^{i_{\sigma (1)}}U_{i_2}^{i_{\sigma (2)}}\mathrm{}U_{i_n}^{i_{\sigma (n)}},\end{array}$$
(3.2)
and the tensor $`𝐃(\sigma )`$ is defined in eq. (2.7).
Altogether, it implies that the master-integral (2.5) can be rewritten in the synthetic form
$$\begin{array}{c}B(\{R_{\mu \nu }\})=Tr_{4n_+}[𝐃(\mathrm{\Xi }_{4n_+}(\{R_{\mu \nu }\}))𝐃(\{U_\rho U_\rho ^+\})]_{\stackrel{~}{\rho }=1}^DdU_{\stackrel{~}{\rho }}\end{array}$$
(3.3)
where the $`S(4n_+)`$-algebra valued tensor $`𝐃(\mathrm{\Xi }_{4n_+}(\{R_{\mu \nu }\}))`$ (to be defined by eq. (3.19) below) multiplies the complementary tensor given by the ordered direct product of the $`4n_+`$ elementary $`N\times N`$ matrices $`(U_\rho )_j^i,(U_\rho ^+)_l^k`$:
$$\begin{array}{c}𝐃(\{U_\rho U_\rho ^+\})_{\rho =1}^D((U_\rho )^{n_\rho }(U_\rho ^+)^{n_\rho }),2n_+=_{\rho =1}^Dn_\rho .\end{array}$$
(3.4)
Given $`\mathrm{\Xi }_{4n_+}(\{R_{\mu \nu }\})`$, we first derive the intermediate $`S(4n_+)`$ representation
$$\begin{array}{c}B(\{R_{\mu \nu }\})=Tr_{4n_+}[𝐃(J_{4n_+}(\{R_{\mu \nu }\}))]\end{array}$$
(3.5)
which in Section 3.3 will be transformed into the final $`S(n_+)`$ form of eq. (2.6). To determine the operator $`J_{4n_+}S(4n_+)`$, in (3.3) one is to substitute the dual form of the SU(N) measure, i.e. to represent the result of the $`D`$ different $`U_\rho `$-integrations as an $`S(4n_+)`$-tensor akin to $`𝐃(\mathrm{\Xi }_{4n_+}(\{R_{\mu \nu }\}))`$.
Before we explain the latter procedure, let us make the pattern of the element $`\mathrm{\Xi }_{4n_+}(\{R_{\mu \nu }\})=_{\sigma S(4n_+)}\xi (\sigma |\{R_{\mu \nu }\})\sigma `$ more transparent relating the explicit form of the $`S(4n_+)`$-group tensor $`𝐃(\sigma )`$ with the one of (3.4). To this aim, we introduce first a particular $`S(4n_+)`$ permutation $`\alpha _{\{n_\rho \}}`$ via the mapping $`m\alpha (m),m=1,\mathrm{},4n_+`$. Then, generalizing the tensor representation (2.7), the tensor $`𝐃(\alpha _{\{n_\rho \}})`$ stands for
$$\begin{array}{c}\delta _{j_1}^{i_{\alpha (1)}}\mathrm{}\delta _{j_{n_1}}^{i_{\alpha (n_1)}}\delta _{l_{n_1+1}}^{k_{\alpha (n_1+1)}}\mathrm{}\delta _{l_{2n_1}}^{k_{\alpha (2n_1)}}\mathrm{}\delta _{l_{4n_+2n_D}}^{k_{\alpha (4n_+2n_D)}}\mathrm{}\delta _{l_{4n_+n_D}}^{k_{\alpha (4n_+n_D)}}\mathrm{}\delta _{l_{4n_+}}^{k_{\alpha (4n_+)}}.\end{array}$$
(3.6)
where, to each individual $`U_\rho `$-,$`U_\rho ^+`$-factor in the product (3.4), we associate one copy of the Kronecker delta-function. In this way, both $`𝐃(\alpha _{\{n_\rho \}})`$ and the block (3.4) are defined to act on one and the same $`S(4n_+)`$ space (to be manifestly constructed, see eqs. (3.13),(3.14) below). Complementary, the pattern of the trace (3.5) of $`J_{4n_+}(\{R_{\mu \nu }\})`$ (defined through the structures like (3.6)) naturally generalizes the above convention (3.2) to the case of (3.4).
### 3.1 The Dual form of the $`SU(N)`$ measure.
On a given base-lattice, a generic multilink integral (like in (2.1)) evidently can be expressed in terms of the 1-link integrals $`M^G(n,m)_{j_1\mathrm{}l_m}^{p_1\mathrm{}q_m}`$:
$$\begin{array}{c}(U)_{j_1}^{p_1}\mathrm{}(U)_{j_n}^{p_n}(U^+)_{l_1}^{q_1}\mathrm{}(U^+)_{l_m}^{q_m}𝑑U𝐃(U)_{\{j^n\}}^{\{p^n\}}𝐃(U^+)_{\{l^m\}}^{\{q^m\}}𝑑U\end{array}$$
(3.7)
composed of the $`N\times N`$ matrices $`(U)_{j_k}^{p_k},(U^+)_{l_k}^{q_k}`$ in the (anti)fundamental representation of the considered Lie group $`G`$. As we will show in a moment, the $`SU(N)`$ TEK partition function (PF) (2.1) is invariant under the substitution of the $`SU(N)`$ link-variables by the $`U(N)=[SU(N)U(1)]/𝐙_𝐍`$ ones. In the U(N) case, where the dual form of $`M^{U(N)}(n,m)`$ in terms of the $`S(n)`$-valued tensors (2.7) reads
$$\begin{array}{c}M^{U(N)}(n,m)_{j_1\mathrm{}l_m}^{p_1\mathrm{}q_m}=\delta [n,m]_{\sigma S(n)}𝐃(\sigma ^1\mathrm{\Lambda }_n^{(1)})_{\{j^n\}}^{\{q^n\}}𝐃(\sigma )_{\{l^n\}}^{\{p^n\}}.\end{array}$$
(3.8)
that renders manifest the previously known interrelation between (3.7) and the symmetric group’s structures. The operator $`\mathrm{\Lambda }_n^{(1)}S(n)`$ belongs to the (U(N) option of the) family
$$\begin{array}{c}\mathrm{\Lambda }_n^{(m)}=_{RY_n^{(N)}}d_R(n!dimR/d_R)^mC_R,m𝐙,\end{array}$$
(3.9)
where $`C_R=P_R/d_R`$ is defined by eq. (3.1). In eq. (3.9), $`d_R`$ and $`dimR`$ are respectively the dimension of the $`S(n)`$-irrep and chiral $`U(N)`$-irrep both described by the same Young tableau $`Y_n^{(N)}`$ containing not more than N rows. (In eqs. (4.4),(4.8) below, we will consider the SU(N) variant of $`\mathrm{\Lambda }_n^{(m)}`$ where the sum is traded for the one over the SU(N) irreps.)
As for the possibility to substitute the $`SU(N)`$ TEK link-variables by the $`U(N)=[SU(N)U(1)]/𝐙_𝐍`$ ones, the pattern of eq. (2.1) ensures that the TEK action is invariant under the $`D`$ copies of the extended transformations $`[U(1)]^D:U_\rho t_\rho U_\rho `$, where $`t_\rho U(1)`$ rather than taking value in the center-subgroup $`𝐙_𝐍`$ of $`SU(N)`$. As a result, the nondiagonal moments $`M^{SU(N)}(n,m)`$, $`nm`$, do $`not`$ contribute into the TEK PF $`\stackrel{~}{X}_D`$. As for the remaining diagonal integrals $`M^{SU(N)}(n,n)`$, the latter extended invariance justifies the required substitution: $`M^{SU(N)}(n,n)=M^{U(N)}(n,n),n𝐙_0`$. Moreover, in the context of the large N SC expansion, in eq. (3.9) one can insert the SU(N) variant of $`\mathrm{\Lambda }_n^{(1)}`$. (The difference can be traced back to the contributions of the SU(N) string-junctions which are supposed to be irrelevant to all orders in 1/N.)
Actually, the derivation of (3.5) calls for the alternative $`S(2n)`$ reformulation of the $`S(n)S(n)`$ formula (3.8) that will require the explicit form of an $`S(2n)`$-basis. For this purpose, we first recall that each individual matrix $`U_\rho ^+`$ or $`U_\rho `$ can be viewed as the operator acting on the associated elementary N-dimensional subspace $`|i_\pm (\rho )>`$ according to the pattern: $`\widehat{U}|i_{}>=_{j_{}=1}^NU_i_{}^j_{}|j_{}>`$ and similarly for $`\widehat{U}^+|i_+>`$. Complementary, for a given $`\sigma S(n)`$, the operator (2.7) acts as the corresponding permutation of the elementary subspaces $`|i_k>`$
$$\begin{array}{c}\widehat{\sigma }|i_1>|i_2>\mathrm{}|i_n>=|i_{\sigma ^1(1)}>|i_{\sigma ^1(2)}>\mathrm{}|i_{\sigma ^1(n)}>𝐃(\sigma )_{\{i^n\}}^{\{j^n\}}|j>^n\end{array}$$
(3.10)
where in the r.h.s. the summation $`_{\{j_k\}}`$ is implied. As a given $`S(4n_+)`$-basis is constructed as the outer product of the $`4n_+`$ building blocks $`|i_\pm (\rho )>`$ (ordered according to a particular prescription), the elementary subspaces $`|i_k>`$ of eq. (3.10) are represented by $`|i_\pm (\rho )>`$.
Returning to the $`S(2n)`$-reformulation of eq. (3.8), in the basis $`|I_{2n}>=|I_n^{(+)}>|I_n^{()}>`$ (with $`|I_n^{(\pm )}>=|i_\pm >^n`$) it reads
$$\begin{array}{c}𝑑U𝐃(U)_{i_1\mathrm{}i_n}^{j_1\mathrm{}j_n}𝐃(U^+)_{i_{n+1}\mathrm{}i_{2n}}^{j_{n+1}\mathrm{}j_{2n}}=𝐃(\mathrm{\Phi }_{2n}\mathrm{\Gamma }(2n)(\mathrm{\Lambda }_n^{(1)}\widehat{1}_{[n]}))_{\{i^{2n}\}}^{\{j^{2n}\}},\end{array}$$
(3.11)
where $`\widehat{1}_{[n]}`$ denotes the ’unity’-permutation of the $`S(n)`$ group, $`\mathrm{\Lambda }_n^{(1)}S(n)`$ is defined by eq. (3.9), while $`\mathrm{\Gamma }(2n)=_{\sigma S(n)}(\sigma ^1\sigma )`$, i.e.
$$\begin{array}{c}𝐃(\mathrm{\Gamma }(2n))_{\{i^{2n}\}}^{\{j^{2n}\}}=_{\sigma S(n)}𝐃(\sigma ^1)_{i_1\mathrm{}i_n}^{j_1\mathrm{}j_n}𝐃(\sigma )_{i_{n+1}\mathrm{}i_{2n}}^{j_{n+1}\mathrm{}j_{2n}}S(2n).\end{array}$$
(3.12)
To restore the $`\rho `$-labels, $`nn_\rho `$, observe first that the ordering of the $`\{U_\rho \}`$-factors in eq. (3.4) is associated to the following basis
$$\begin{array}{c}|\stackrel{~}{I}_{4n(+)}>=_{\rho =1}^D|I_{2n(\rho )}>;|I_{2n(\rho )}>=|I_{n(\rho )}^{(+)}>|I_{n(\rho )}^{()}>,\end{array}$$
(3.13)
$$\begin{array}{c}|I_{n(\rho )}^{(\pm )}>=_{\nu \rho }^{D1}|I_{n(\rho \nu )}^{(\pm )}>;|I_{n(\rho \nu )}^{(\pm )}>=|i_\pm (\rho )>^{n(\rho \nu )}|i_\pm (\nu )>^{n(\rho \nu )},\end{array}$$
(3.14)
where $`2n_+=_{\rho =1}^Dn_\rho `$, and $`|I_n^{(\pm )}>=|i_\pm >^n`$ (used in eq. (3.11)) matches with $`|I_{n(\rho )}^{(\pm )}>`$. Therefore, for a given link $`\rho `$, the left and the right $`S(n_\rho )`$-subblocks of $`\mathrm{\Gamma }(2n_\rho )`$ in eq. (3.12) act respectively on $`|I_{n(\rho )}^{(+)}>`$ and on $`|I_{n(\rho )}^{()}>`$. The same convention is used for the $`S(n_\rho )`$-subblocks in the direct product $`(\mathrm{\Lambda }_{n_\rho }^{(1)}\widehat{1}_{[n_\rho ]})`$ entering eq. (3.11).
The remaining $`S(2n_\rho )`$-operator $`\mathrm{\Phi }_{2n_\rho }`$, being considered in the alternatively ordered basis $`|\stackrel{~}{I}_{2n(\rho )}>`$ for each $`|I_{2n(\rho )}>`$-subsector,
$$\begin{array}{c}|I_{2n(\rho )}>|\stackrel{~}{I}_{2n(\rho )}>=(|i_+(\rho )>|i_{}(\rho )>)^{n_\rho }\end{array}$$
(3.15)
(with $`|i_\pm (\rho )>^{n(\rho )}=_{\nu \rho }^{D1}|i_\pm (\rho )>^{n(\rho \nu )}`$), takes the simple form of the outer product of the 2-cycle permutations $`c_2C(2)`$
$$\begin{array}{c}\mathrm{\Phi }_{2n_\rho }=(c_2)^{n_\rho }S(2n_\rho );c_2:\{12\}\{21\},\end{array}$$
(3.16)
where each $`c_2S(2)`$ acts on the ’elementary’ sector $`|i_+(\rho )>|i_{}(\rho )>`$. It completes the embedding of the $`S(2n_\rho )`$ operators (3.11), representing the individual 1-link integrals, to act in the common ’enveloping’ $`S(4n_+)`$-space. (It is noteworthy that the three $`S(2n_\rho )`$ subblocks of the inner-product in the r.h.s. of eq. (3.11) commute with each other.)
### 3.2 The Dual form of the TEK action.
Let us now turn to the derivation of the master-element $`\mathrm{\Xi }_{4n_+}(\{R_{\mu \nu }\})`$ entering the synthetic representation (3.3) of the master-integral (2.5). For this purpose, we first specify an alternative, more suitable $`S(4n_+)`$ basis
$$\begin{array}{c}|I_{4n(+)}>=_{\mu \nu =1}^{D(D1)/2}|I_{4n(\mu \nu )}>,\end{array}$$
(3.17)
$$\begin{array}{c}|I_{4n(\mu \nu )}>=(|i_+(\mu )>|i_+(\nu )>|i_{}(\mu )>|i_{}(\nu )>)^{n_{\mu \nu }}\end{array}$$
(3.18)
where in eq. (3.18) the product of the four elementary blocks $`|i_\pm (\rho )>`$ is associated to the elementary $`\mu \nu `$-holonomy $`U_{\mu \nu }`$ entering eq. (2.5). As it is derived in Appendix A, in this basis one obtains
$$\begin{array}{c}\mathrm{\Xi }_{4n_+}(\{R_{\mu \nu }\})=_{\mu \nu =1}^{D(D1)/2}P_{4n_{\mu \nu }}(R_{\mu \nu })\mathrm{\Psi }_{4n_{\mu \nu }};n_+=_{\{\mu \nu \}}n_{\mu \nu },\end{array}$$
(3.19)
where each of the operators $`P_{4n_{\mu \nu }}(R_{\mu \nu }),\mathrm{\Psi }_{4n_{\mu \nu }}S(4n_{\mu \nu })`$ is supposed to act on the corresponding $`|I_{4n(\mu \nu )}>`$ subspace of $`|I_{4n(+)}>`$. Given (3.18), $`\mathrm{\Psi }_{4n_{\mu \nu }}`$ assumes the simple form of the outer product
$$\begin{array}{c}\mathrm{\Psi }_{4n_{\mu \nu }}=(c_4)^{n_{\mu \nu }}S(4n_{\mu \nu });c_4:\{1234\}\{4123\},\end{array}$$
(3.20)
with each individual 4-cycle permutation $`c_4C(4)`$ acting on the elementary plaquette subspace
$$\begin{array}{c}|i_+(\mu )>|i_+(\nu )>|i_{}(\mu )>|i_{}(\nu )>\end{array}$$
(3.21)
ordered in accordance with the original pattern (2.5) of the plaquette-labels of $`U_{\mu \nu }`$.
As for $`P_{4n_{\mu \nu }}(R_{\mu \nu })`$, making use of the alternative basis $`|I_{4n(\mu \nu )}>|\stackrel{~}{I}_{4n(\mu \nu )}>`$ of the $`S(4n_{\mu \nu })`$-subspace in (3.18):
$$\begin{array}{c}|\stackrel{~}{I}_{4n(\mu \nu )}>=|i_+(\mu )>^{n_{\mu \nu }}|i_+(\nu )>^{n_{\mu \nu }}|i_{}(\mu )>^{n_{\mu \nu }}|i_{}(\nu )>^{n_{\mu \nu }},\end{array}$$
(3.22)
we employ (see Appendix A) the following representation
$$\begin{array}{c}P_{4n_{\mu \nu }}(R_{\mu \nu })=\widehat{1}_{[n_{\mu \nu }]}\widehat{1}_{[n_{\mu \nu }]}(C_{R_{\mu \nu }}\sqrt{d_{R_{\mu \nu }}})(C_{R_{\mu \nu }}\sqrt{d_{R_{\mu \nu }}}),\end{array}$$
(3.23)
where $`\widehat{1}_{[n_{\mu \nu }]}`$ denotes the $`S(n_{\mu \nu })`$-unity. More explicitly, the four (ordered) $`S(n_{\mu \nu })`$-factors in eq. (3.23) are postulated to act on the corresponding four (ordered) $`n_{\mu \nu }`$-dimensional subspaces (3.22) of $`|\stackrel{~}{I}_{4n_{\mu \nu }}>`$. Altogether, we have formulated the required (for the derivation of eq. (3.5)) $`embedding`$ of the $`S(4n_{\mu \nu })`$ operators $`\mathrm{\Psi }_{4n_{\mu \nu }},P_{4n_{\mu \nu }}(R_{\mu \nu })`$ into the enveloping $`S(4n_+)`$-space.
### 3.3 $`B(\{R_{\mu \nu }\})`$ as the $`Tr_{n_+}`$-character.
Combining together eqs. (3.19) and (3.11) in compliance with the pattern of eq. (3.3), we arrive at the dual representation of the master-integral (2.5) as the $`S(4n_+)`$ character (3.5) with the master-element
$$\begin{array}{c}J_{4n_+}(\{R_{\mu \nu }\})=\left(_{\mu \nu =1}^{D(D1)/2}\mathrm{\Psi }_{4n_{\mu \nu }}\right)\left(_{\rho =1}^D\mathrm{\Delta }_{2n_\rho }(\{R_{\rho \nu }\})\right),\end{array}$$
(3.24)
$$\begin{array}{c}\mathrm{\Delta }_{2n_\rho }(\{R_{\nu \rho }\})=\mathrm{\Phi }_{2n_\rho }\mathrm{\Gamma }(2n_\rho )K_{2n_\rho }(\{R_{\nu \rho }\})S(2n_\rho ).\end{array}$$
(3.25)
For later convenience, we have introduced the $`S(n_\rho )S(n_\rho )`$ combination that in the $`|I_{2n(\rho )}>`$ basis (3.13) reads
$$\begin{array}{c}K_{2n_\rho }(\{R_{\nu \rho }\})=\mathrm{\Lambda }_{n_\rho }^{(1)}\left(_{\nu \rho }^{D1}C_{R_{\rho \nu }}\sqrt{d_{R_{\rho \nu }}}\right),\end{array}$$
(3.26)
where $`\mathrm{\Lambda }_{n_\rho }^{(1)}`$ acts onto $`|I_{n(\rho )}^{(+)}>`$, while each of the $`C_{R_{\rho \nu }}`$ factors acts on the $`|I_{n(\rho \nu )}^{()}>`$-subspace of $`|I_{n(\rho )}^{()}>`$.
Next, let us complete the duality transformation trading the intermediate $`S(4n_+)`$ representation (3.5) for its final $`S(n_+)`$ pattern (2.6). To this aim, it is convenient to start with the alternative following alternative form of the master-element (3.24). As it is demonstrated in Appendix A, the dual representation (3.5) does not alter when in the element (3.24)/(3.25) one makes the substitution
$$\begin{array}{c}_{\rho =1}^D\mathrm{\Phi }_{2n_\rho }_{\rho =1}^D(\mathrm{\Phi }_{2n(\rho )})^2_{\rho =1}^D\widehat{1}_{[2n_\rho ]}=\widehat{1}_{[4n_+]}.\end{array}$$
(3.27)
so that (inside the $`Tr_{4n_+}`$-character) in eq. (3.25) all the operators $`\mathrm{\Phi }_{2n(\rho )}`$ can be omitted.
#### 3.3.1 The $`D=2`$ case.
Next, it is appropriate to proceed with the simplest $`D=2`$ case of (3.5) corresponding to the well studied continuous SU(N) gauge theory on a 2-torus. It will provide not only with a cross-check of our $`D2`$ formalism but also with the motivation for the announced reduction, $`S(4n_+)S(n_+)`$, of the enveloping space. In $`D=2`$, the this reduction is encoded in the basic property (following from eqs. (A.1),(A.3) in Appendix A) of the $`\mathrm{\Psi }_{4n}`$ operator (3.20)
$$\begin{array}{c}Tr_n[U_{\mu \nu }^n]=Tr_{4n}[𝐃(\mathrm{\Psi }_{4n})\stackrel{~}{U}_{\mu \nu }^n];\stackrel{~}{U}_{\mu \nu }=U_\mu U_\mu ^+U_\nu U_\nu ^+,\end{array}$$
(3.28)
while $`U_{\mu \nu }=U_\mu U_\nu U_\mu ^+U_\nu ^+`$. Making the substitution $`U_\rho \sigma _\rho ^{(+)},U_\rho ^+\sigma _\rho ^{()}`$, one obtains
$$\begin{array}{c}Tr_{4n}[𝐃(\left(_{\rho =1}^2(\sigma _\rho ^{(+)}\sigma _\rho ^{()})\right)\mathrm{\Psi }_{4n})]=Tr_n[𝐃(_{\rho =1}^D\sigma _\rho ^{(+)}_{\mu =1}^D\sigma _\mu ^{()})],\end{array}$$
(3.29)
where in the l.h.s. the operators $`\sigma _\rho ^{(\pm )}S(n)`$ (combined into the outer product) act on the associated $`|I_n^{(\pm )}>`$ subspaces of the $`S(4n)`$ basis (3.13). As for the ordering inside the two inner $`\rho `$-products in the r.h.s. of eq. (3.29), in both products it complies with the ordering of the $`|i_+(\rho )>`$ (or, equally, $`|i_{}(\rho )>`$) elementary blocks in eq. (3.21).
Let us apply the identity (3.29) to the $`D=2`$ option of (3.5), (3.24). Combining eq. (3.29) with the orthonormality $`P_{R_1}P_{R_2}=\delta _{R_1,R_2}P_{R_1}`$ of $`P_R=d_RC_R`$ and taking into account the standard relation $`Tr_n[C_R\sigma ]=dimR\chi _R(\sigma )/d_R`$ (between the projected $`Tr_n`$-trace and the canonical $`S(n)`$ character $`\chi _R`$, see e.g. ), one easily obtains
$$\begin{array}{c}B(R)=\frac{1}{dimR}\{\left(\frac{d_R}{n!}\right)^2_{\{\sigma _\rho S(n)\}}\frac{\chi _R([\sigma _1,\sigma _2])}{d_R}\}=\frac{1}{dimR},\end{array}$$
(3.30)
where $`RY_n^{(N)}`$, and $`[\sigma _1,\sigma _2]`$ conventionally stands for $`(\sigma _1\sigma _2\sigma _1^1\sigma _2^1)`$. We have also used that the block in the curly brakets of (3.30) is equal to unity according to the identity derived in . Together with eq. (2.4), the expression (3.30) for $`B(R)`$ precisely matches with the (genus one) result of derived by the combinatorial method of .
#### 3.3.2 The $`D3`$ case.
Returning to the generic $`D3`$ $`Tr_{4n_+}`$-character (2.5) of the master-element (3.24), the $`S(4n_+)S(n_+)`$ reduction of the enveloping space is performed with the help of the following generalization of the $`D=2`$ identity (3.29)
$$Tr_{4n_+}[𝐃(\left(\underset{\mu \nu =1}{\overset{D(D1)/2}{}}\mathrm{\Psi }_{4n_{\mu \nu }}\right)\left(\underset{\rho =1}{\overset{D}{}}(\sigma _\rho ^{(+)}\sigma _\rho ^{()})\right))]=$$
$$\begin{array}{c}=Tr_{n_+}[𝐃(\left(_{\rho =1}^D(\sigma _\rho ^{(+)}\widehat{1}_{[\frac{n_+}{n_\rho }]})\right)\left(_{\rho =1}^D(\sigma _\rho ^{()}\widehat{1}_{[\frac{n_+}{n_\rho }]})\right))].\end{array}$$
(3.31)
Let us simply explain the meaning of the above pattern, while for more details see Appendix A. In the l.h.s. of (3.31), the outer $`\mu \nu `$-product is defined in the same way as in eq. (3.24), and the operators $`\sigma _\rho ^{(\pm )}S(n_\rho )`$ (composed into the outer $`\rho `$-product) act on the associated $`|I_{n(\rho )}^{(\pm )}>`$ subspace of the $`S(4n_+)`$ basis (3.13). As for the r.h.s. of (3.31), we first construct the following $`S(n_+)`$ basis. To begin with , one is to introduce the $`N`$-dimensional spaces $`|i(\mu \nu )>,i=1,\mathrm{},N,`$ parametrized by the plaquette label $`\mu \nu =1,\mathrm{},D(D1)/2`$. Then, the $`S(n_+)`$ operators (recall that $`n_+=_{\{\mu \nu \}}n_{\mu \nu }`$) can be viewed as acting on
$$\begin{array}{c}|I_{n(+)}>=_{\mu \nu =1}^{D(D1)/2}|I_{n(\mu \nu )}>,|I_{n(\mu \nu )}>=(|i(\mu \nu )>)^{n_{\mu \nu }},\end{array}$$
(3.32)
according to the same rule (3.10) that has been already used for the $`S(4n_+)`$ operators. Given this convention, each operator $`\sigma _\rho ^{(\pm )}`$ is postulated to act on the associated $`S(n_\rho )`$ subspace $`|\stackrel{~}{I}_{n(\rho )}>`$ of $`|I_{n(+)}>`$,
$$\begin{array}{c}|\stackrel{~}{I}_{n(\rho )}>=_{\nu \rho }^{D1}(|i(\rho \nu )>)^{n_{\rho \nu }},n_\rho =_{\nu \rho }^{D1}n_{\rho \nu },\end{array}$$
(3.33)
where the ordering of the $`|i(\rho \nu )>`$ blocks matches with the one in eq. (3.32). As for $`\widehat{1}_{[n_+/n_\rho ]}`$, in eq. (3.31) it denotes the unity permutation on the $`S(n_+n_\rho )`$ subspace of (3.32) complementary to (3.33). In the $`D=2`$ case, where $`n_+=n_1=n_2`$, the general eq. (3.31) readily reduces to its degenerate variant (3.29).
To complete the construction (3.31), one should specify the ordering inside the two inner $`\rho `$-products of its r.h. side. We defer this task till the end of the section and now apply the identity (3.31) to the $`Tr_{4n_+}`$ character (3.5) of the master-element (3.24). For this purpose, all what we need is the proper identification of the composed into $`J_{4n_+}(\{R_{\mu \nu }\})`$ permutations with $`\sigma _\rho ^{(\pm )}`$. To this aim, let denote by $`\lambda _{\mu \nu }`$ and $`\lambda _\rho `$ the permutations which enter the definition (3.1) of the relevant operators $`C_{R_\varphi }=P_{R_\varphi }/d_{R_\varphi }S(n_\varphi )`$ (combined into $`K_{2n_\rho }`$ of eq. (3.26), with $`\mathrm{\Lambda }_{n_\rho }^{(1)}S(n_\rho )`$ being given by eq. (3.9)) assigned with the associated labels $`\varphi \{\mu \nu \},\{\rho \}`$. Complementary, let $`\sigma _\rho `$ stands for the permutations entering the definition (3.12) of $`\mathrm{\Gamma }(2n_\rho )`$. Identifying $`(_{\mu \rho }^{D1}\lambda _{\mu \rho })\sigma _\rho \sigma _\rho ^{()},\lambda _\rho \sigma _\rho ^1\sigma _\rho ^{(+)}`$ (without numb summation over $`\rho `$), after some routine machinery one derives for the master-element
$$\begin{array}{c}A_{n_+}(\{R_{\mu \nu }\})=_{\rho =1}^D_{R_{n_\rho }Y_{n_\rho }^{(N)}}\frac{d_{R_\rho }^2}{n_\rho !dimR_\rho }_{\sigma _\rho S(n_\rho )}F(\{\sigma _\rho \};\{R_\varphi \}),\end{array}$$
(3.34)
$$\begin{array}{c}F=\left(_{\{\mu \nu \}}C_{R_{\mu \nu }}\right)\left(_{\{\rho \}}(\sigma _\rho \widehat{1}_{[\frac{n_+}{n_\rho }]})\right)\left(_{\{\lambda \}}(\sigma _\lambda ^1C_{R_\lambda }\widehat{1}_{[\frac{n_+}{n_\lambda }]})\right),\end{array}$$
(3.35)
so that its trace (2.6) determines the $`B(\{R_{\mu \nu }\})`$-block (2.5) of the TEK partition function (PF).
Altogether, it establishes the exact duality transformation of the TEK PF (2.1) which is one of the main results of the paper. In certain sense (modulo the explicit presence of the $`\sigma _\rho `$-twists), it provides with the $`D3`$ generalization of the representation for the PF of the continuous gauge theory on a $`2d`$ manifold. Also, it can be compared with the considerably simpler pattern of the PF of the judiciously constructed solvable $`D`$-matrix models where $`B(\{R_{\mu \nu }\})`$ depends nontrivially only on the associated generalized Littlewood-Richardson coefficients. (On the contrary, the pattern (3.24) encodes the general Klebsch-Gordan coefficiens.)
#### 3.3.3 The ordering inside the inner $`\rho `$-products .
Finally, let us discuss the ordering inside the two inner $`\rho `$-products of the r.h. side of eq. (3.31). As it is shown in Appendix A, this ordering is entirely predetermined by the following characteristics $`G(\rho )`$ of a particular $`\rho `$-label called its cardinality. Consider the $`D(D1)/2`$ dimensional vector $`𝐌(\{\rho \})`$ defined so that its components $`M_k(\{\rho \})`$ are in one-to-one correspondence with the $`D(D1)/2`$ labels $`\mu \nu `$ (where $`1\mu <\nu D`$). The latter parametrize the $`\mu \nu `$th plaquette-holonomies $`U_{\mu \nu }`$ entering the characters in eq. (2.5). Let the $`\mu \nu `$th component of $`𝐌(\{\rho \})`$ is equal to the first link-label $`\mu `$ of the plaquette-label: $`M_{\mu \nu }(\{\rho \})=\mu `$. Then, the cardinality (ranging from $`0`$ to $`D1`$) is postulated to be $`G(\mu )=_{k=1}^{D(D1)/2}\delta [\mu ,M_k(\{\rho \})]`$, i.e. the number of times the particular $`\mu `$-label enters the entries of the vector $`𝐌(\{\rho \})`$. In eq. (2.5), it is always possible to arrange for a nondegenerate cardinality assignement $`\{G(\rho )\}`$ when $`G(\rho _1)G(\rho _2)`$ if $`\rho _1\rho _2`$. Then, given a nondegenerate set $`\{G(\rho )\}`$, the $`(\sigma _\rho ^{(\pm )}\widehat{1}_{[n_+/n_\rho ]})`$ factors in the eq. (3.31) are ordered (from the left to the right) according to the successively decreasing $`G(\rho )`$-assignements of their $`\rho `$-labels.
## 4 Schur-Weyl transformation of $`B(\{n_{\mu \nu }\})`$.
To transform the dual representation (2.6)/(3.34) of the TEK PF $`\stackrel{~}{X}_D`$ into the $`D2`$ stringy representation like (1.9), one is rewrite $`B(\{n_{\mu \nu }\})`$ entirely in terms of the symmetric groups’ variables which are suitable for the algebraic definition of the topological data associated to the mappings (1.11). For this purpose, we employ the following two useful identities (derived in Appendices B and D) which reflect the Schur-Weyl complementarity of the Lie and the symmetric groups. The first one trades the ubiquitous operator (3.9) for the product of the two elements of the associated $`S(n)`$ algebra
$$\begin{array}{c}\mathrm{\Lambda }_n^{(m)}=P_n^{(N)}(N^n\mathrm{\Omega }_n)^m,P_n^{(N)}=_{RY_n^{(N)}}P_R,\end{array}$$
(4.1)
where the projector $`(P_n^{(N)})^2=P_n^{(N)},P_n^{(N)}=1`$ if $`n<N`$, (that independently appeared within the method of ) reduces the $`S(n)`$ $`Y_n`$-variety of irreps to the $`Y_n^{(N)}`$-one of either U(N) or $`SU(N)`$. As for the second element
$$\begin{array}{c}\mathrm{\Omega }_n=_{\sigma S(n)}(1/N)^{nK_{[\sigma ]}}\sigma ;[\mathrm{\Omega }_n,\rho ]=0,\rho S(n),\end{array}$$
(4.2)
(belonging the center of the $`S(n)`$-algebra), it is defined by the equation
$$\begin{array}{c}(dimR)^m=\frac{\chi _R((N^n\mathrm{\Omega }_n)^m)}{d_R}\left(\frac{d_R}{n!}\right)^m;nK_{[\sigma ]}=_{k=1}^n(k1)p_k,\end{array}$$
(4.3)
where $`m𝐙`$, and the factor $`K_{[\sigma ]}=_{k=1}^np_k`$ in eq. (4.2) denotes the total number of various $`k`$-cycles in the cyclic decomposition of the conjugacy class $`[\sigma ]=[1^{p_1},2^{p_2},\mathrm{},n^{p_n}]`$, $`_{k=1}^nkp_k=n`$.
Generalizing (4.1), the second key-identity deals with the similar sum weighted this time by the factor $`e^\mathrm{\Gamma }`$ (which defines a generic model (1.1))
$$\begin{array}{c}_{RY_n^{(N)}}e^{\mathrm{\Gamma }(..,\{C_p(R)\})}d_R\left(\frac{n!dimR}{d_R}\right)^mC_R=P_n^{(N)}(N^n\mathrm{\Omega }_n)^mQ_n(\mathrm{\Gamma }).\end{array}$$
(4.4)
In the Heat-Kernal case (1.7) where $`\mathrm{\Gamma }(..)=\lambda C_2(R)/2N`$, the $`S(n)`$-algebra valued operator $`Q_n(\mathrm{\Gamma })`$ reads
$$\begin{array}{c}Q_n(\mathrm{\Gamma })=exp[\frac{\lambda }{2}(n\frac{n^2}{N^2})]exp[\frac{\lambda }{N}\widehat{T}_2^{(n)}],\widehat{T}_2^{(n)}=_{\tau T_2^{(n)}}\tau ,\end{array}$$
(4.5)
where $`T_2^{(n)}T_2`$ denotes the $`S(n)`$ conjugacy class $`[1^{n2}2^1]`$ of the simple transposition, and $`\lambda =\stackrel{~}{g}^2N`$. In a generic admissible model (1.1), as it is demonstrated in Appendix E, the operator $`Q_n(\mathrm{\Gamma })`$ generalizes to
$$\begin{array}{c}Q_n(\mathrm{\Gamma })=exp[_{\{p\}}\upsilon _{\{p\}}(\{\stackrel{~}{b}_k\},n,N)\widehat{T}_{\{p\}}^{(n)}],\widehat{T}_{\{p\}}^{(n)}=_{\xi ^{\{p\}}T_{\{p\}}^{(n)}}\xi ^{\{p\}},\end{array}$$
(4.6)
where $`\widehat{T}_{\{p\}}^{(n)}`$ denotes the sum of the $`S(n)`$ permutations belonging to a particular conjugacy class $`T_{\{p\}}^{(n)}`$ labelled by the partition $`\{p\}`$ of $`n`$: $`_{i=1}^nkp_k=n,`$. As for the weight $`\upsilon _{\{p\}}(..)`$, it assumes the form (which, in particular, results in the required asymptotics (1.6))
$$\begin{array}{c}\frac{\upsilon _{\{p\}}(\{\stackrel{~}{b}_k\},n,N)}{N^{_{k=1}^n(k1)p_k}}=_{m=0}^{M_{\{p\}}}_{l[m/2]}s_{\{p\}}(\{\stackrel{~}{b}_k\},m,l)N^{2l}n^m,\end{array}$$
(4.7)
where $`m,l𝐙_\mathrm{𝟎}`$, and $`[m/2]=m/2`$ or $`(m+1)/2`$ depending on whether $`m`$ is even or odd (while $`s_{\{p\}}(..)O(N^0)`$). The specific pattern of $`\upsilon _{\{p\}}(..)`$ implies that in eq. (1.1) the function $`\mathrm{\Gamma }(..)`$ satisfies certain conditions (see Appendix E) that ensure the consistent stringy interpretation of (4.7) to be discussed in Section 6 (for eq. (4.5)) and in Appendix E (for eqs. (4.6)/(4.7)).
Combining (4.1),(4.4) with the delta-function reformulation (2.8) of the $`Tr_{n_+}`$ character (derived in Appendix B), we arrive at the explicit $`_\varphi S(n_\varphi )`$ representation for the building block $`B(\{n_{\mu \nu }\})`$ of the TEK partition function (2.4). This central, for the present discussion of the lattice Gauge String, expression reads
$$\begin{array}{c}\delta _{n_+}(\mathrm{\Lambda }_{n_+}\left(_{\{\mu \nu \}}\frac{Q_{n_{\mu \nu }}\mathrm{\Lambda }_{n_{\mu \nu }}}{n_{\mu \nu }!}\right)_{\{\sigma _{\stackrel{~}{\rho }}\}}_{\{\rho \}}(\sigma _\rho \widehat{1}_{[\frac{n_+}{n_\rho }]})_{\{\mu \}}(\sigma _\mu ^1\mathrm{\Lambda }_{n_\mu }^{(1)}\widehat{1}_{[\frac{n_+}{n_\mu }]})),\end{array}$$
(4.8)
where, to all orders in 1/N, one can safely use the SU(N) variant of the representation (4.1) of $`\mathrm{\Lambda }_{n_\varphi }^{(m)}`$. It is noteworthy that the $`[𝐙_\mathrm{𝟐}]^{D(D1)/2}`$ invariance (with respect to $`R_{\mu \nu }\overline{R}_{\mu \nu }`$) of the sums in (2.1) defining the plaquette-factor (1.1) results in the invariance of eq. (4.8) under the simultaneous permutations $`\rho ,\mu \sigma (\rho ),\sigma (\mu ),\sigma S(D)`$ of the link-labels $`\rho ,\mu =1,\mathrm{},D,`$ in the two ordered inner products.
## 5 The stringy form of $`B(n)`$ in $`D=2`$.
To begin with, in the $`D=2`$ case (where $`n_+=n_{12}=n_1=n_2`$) all the involved into (4.8) $`S(n_\varphi )`$ operators act in one and the same $`S(n)`$-space that matches with the reduced formula (3.29). The resulting amplitude in the Heat-Kernal case (1.7) reads (with $`B(0)1`$)
$$\begin{array}{c}B(n)=\frac{e^{\frac{\lambda }{2}(n\frac{n^2}{N^2})}}{n!}_{\{\sigma _\rho \},T_2^{(n)}S(n)}\delta _n(P_n^{(N)}e^{\frac{\lambda }{N}\widehat{T}_2^{(n)}}(N^n\mathrm{\Omega }_n)^{12+1}[\sigma _1,\sigma _2])\end{array}$$
(5.1)
that is in complete agreement with the genus-one result of Baez and Taylor derived by a different method. It provides with the more compact reformulation of the Gross-Taylor stringy pattern (1.9), although the transformation (to be summarized by eqs. (5.7),(5.8) below) relating the two representations is not entirely manifest. Let us proceed recasting (5.1) into the form ’almost’ equivalent to the one of (1.9).
To make contact with the pattern of the $`f`$-mapping of eq. (1.9), first it is convenient to decompose
$$\begin{array}{c}P_n^{(N)}=_{T_{\{p\}}^{(n)}S(n)}P_n^{(N)}(T_{\{p\}}^{(n)})\widehat{T}_{\{p\}}^{(n)},\end{array}$$
(5.2)
where $`\widehat{T}_{\{p\}}^{(n)}`$ is defined in eq. (4.6). Then, expanding all the exponents except $`e^{\frac{\lambda }{2}n}=e^{\stackrel{~}{\sigma }_0n}`$, one is to rewrite (5.1) as
$$\begin{array}{c}B(n)=_{i,s,t0}_{T_{\{p\}}S(n)}_{f\stackrel{~}{M}}\frac{N^{2(t+s)i}}{|C_f(\{p\})|}P_n^{(N)}(T_{\{p\}})K_n(i,s,t),\end{array}$$
(5.3)
$$\begin{array}{c}K_n(i,s,t)=e^{\frac{\lambda }{2}n}\frac{(\lambda )^{i+s+t}}{i!s!t!}\frac{(1)^in^s(n^2n)^t}{2^{s+t}},\end{array}$$
(5.4)
where in the first exponent of (5.1) one is to decompose $`n^2=n/2+n(n1)/2`$, and for simplicity we denote $`T_{\{p\}}T_{\{p\}}^{(n)}`$ in the rest of the section. As for the symmetry factor,
$$\begin{array}{c}_{f\stackrel{~}{M}(\{p\},n,i)}\frac{1}{|C_f(\{p\})|}=_{\{\sigma _\rho S(n)\}}\frac{1}{n!}\delta _n(\widehat{T}_{\{p\}}(\widehat{T}_2)^i[\sigma _1,\sigma _2]),\end{array}$$
(5.5)
it emerges when one reformulates the r.h.s. of (5.5) as the sum over certain maps (1.9) (to be explicitly constructed below). The latter can be viewed as the topological mappings (i.e. immersions without singularities)
$$\begin{array}{c}f:f(\stackrel{~}{M}\{f^1(q_s)\})=M\{q_s\}.\end{array}$$
(5.6)
of the space $`\stackrel{~}{M}\{f^1(q_s)\}`$ onto the base-space torus M with $`i+1`$ deleted points $`\{q_s\}`$. These maps define the admissible (by the data in the r.h.s. of (5.5)) branched covering spaces $`\stackrel{~}{M}`$ of M which can be visualized as the Riemann surfaces $`\stackrel{~}{M}\stackrel{~}{M}(\{p\},n,i)`$ (to be identified with the worldsheets of the YM-flux) with $`n`$-sheets and $`i+1`$ branch points located at $`\{q_s\}`$.
As for the sums (5.3) over the nonnegative integers $`i,t,s`$, in addition to the number of the ’movable’ simple branch-points (where two sheets are identified), they refer to the extra ’movable’ singularities of the map (1.9). Namely, one is to attach to $`\stackrel{~}{M}(\{p\},n,i)`$ the $`t`$ collapsed to a point microscopic tubes (connecting two sheets) and the $`s`$ collapsed to a point handles (glued to a single sheet). Altogether, it results in the worldsheet $`\stackrel{~}{M}(\{p\},n,i|t,s)`$.
The important observation is that the factor $`|C_f(\{p\})|`$ is equal to the number of distinct automorphisms of the branched covering space $`\stackrel{~}{M}(\{p\},n,i|t,s)`$ in question. Complementary, in the absence of the $`P_n^{(N)}`$-twist (i.e. when $`T_{\{p\}}\widehat{1}`$), the 1/N factor enters (5.3) in the power equal to the $`G=1`$ option of the Riemann-Hurwitz formula $`h=n(2G2)+2(t+s)+i`$ calculating the overall genus of the corresponding (modified) Riemann surface $`\stackrel{~}{M}(\widehat{1},n,i|t,s)`$. Also, in what follows, we assume that that the contribution of the ’movable’ collapsed handles is reabsorbed into the redefinition $`\stackrel{~}{\sigma }_0=\lambda /2\lambda (11/N^2)/2`$ of the $`SU(N)`$ bare string tension which eliminates the corresponding singularities of the map (1.11). (In the case (4.6) of the generic model (1.1), the modified tension is given by the $`n=1`$ restriction of $`\upsilon _{\{p\}}(\{\stackrel{~}{b}_k\},n,N)`$ associated to $`\widehat{T}_{\{p\}}^{(n)}=\widehat{1}_{[n]}`$.)
¿From the general expression (4.6), it is clear that the pattern (1.9) emerges in a generic model (1.1) as well. In particular, owing to the pattern of eq. (4.7), the large $`N`$ asymptotics $`ln[Q_n(\mathrm{\Gamma })]=n\stackrel{~}{\sigma }_0(\{\stackrel{~}{b}_k\})(1+O(1/N))`$ is consistently provided by the $`\widehat{T}_{\{p\}}=\widehat{1}`$ term of (4.6). For $`T_{\{p\}}\widehat{1}`$, the leading $`l=0`$ term in eq. (4.7) describes the branch-point canonically parametrized (see e.g. ) by $`T_{\{p\}}`$. The latter point decreases the associated Euler character by $`_k(k1)p_k`$ which matches (akin to the pattern (4.3) of $`\mathrm{\Omega }_n`$) with power of the 1/N factor assigned in (4.7) to $`\widehat{T}_{\{p\}}`$ . As it is discussed in Appendix E (where the earlier results are summarized and reformulated), the $`l1`$ terms can be reinterpreted as the movable subsurfaces (of various topologies) collapsed to a point.
As for the factor $`P_n^{(N)}(T_{\{p\}})`$, inherited from the decomposition (5.2) of the projector $`P_n^{(N)}`$, its dependence on N is not particularly suitable for a manifest $`1/N`$ expansion like (1.9). It calls for a nontrivial resummation which would reproduce the well-defined large $`N`$ SC series (1.9) (obtained by Gross and Taylor without resort to (5.1)). The latter pattern effectively eliminates $`P_n^{(N)}(T_{\{p\}})`$ at the expense of working with the $`S(n^+)S(n^{})`$ double-representation, $`B(n)B(\{n^\pm \})`$, that refers to the two coupled sectors of the opposite worldsheet orientation. The prescription to reconstruct $`B(\{n^\pm \})`$ from the amplitude like (4.1) is quite simple: all the involved $`S(n)`$-structures are traded for their $`S(n^+)S(n^{})`$ counterparts
$$\begin{array}{c}\delta _n(..)\delta _{n^+\times n^{}}(..),N^n\mathrm{\Omega }_nN^{n^++n^{}}\mathrm{\Omega }_{n^+,n^{}},\sigma _\rho \sigma _\rho ^+\sigma _\rho ^{},\end{array}$$
(5.7)
$$\begin{array}{c}Q_nQ_{n^+,n^{}}=e^{\frac{\lambda }{2}(n^++n^{}((n^+)^2+(n^{})^22n^+n^{})/N^2)}e^{\frac{\lambda }{N}(\widehat{T}_2^{(n^+)}+\widehat{T}_2^{(n^{})})},\end{array}$$
(5.8)
where eq. (5.8) imples the Heat-Kernal case (1.7) and can be generalized to a generic model (1.1). (The definition and interpretation of $`\mathrm{\Omega }_{n^+,n^{}}`$ and other ingredients in eqs. (5.7),(5.8) can be found in .)
### 5.1 Construction of the branched covering spaces.
In the remaining subsections, we discuss the major issues related to the effective enumeration of the mappings (5.6) and their automorphisms. Let us proceed with an explicit algorithm which, given the symmetric group data in the r.h.s. of (5.5), reconstructs the topology of the Riemann surfaces $`\stackrel{~}{M}(\{p\},n,i)`$ in the l.h.s. of (5.5). As the elements $`\widehat{T}_{\{p\}},(\widehat{T}_2)^i`$ are associated to the branch points (BPs), it is convenient to start with the simpler case of the topological covering spaces (without BP’s singularities) removing the latter elements. The full branched covering spaces (BCSs) can be reproduced reintroducing the BPs onto the corresponding covering spaces (CSs).
As for a particular n-sheet CS $`\stackrel{~}{M}`$ of a given $`2d`$ surface $`M`$ (the 2-torus in what follows), it can be composed with the help of the cutting-gluing rules borrowed from the constructive topology. To begin with, cut a 2-torus M along the two uncontractible cycles $`\alpha (\rho )`$ trading the latter for the pairs of edges $`\alpha (\rho )\beta (\rho ),\rho =\mu ,\nu `$. It makes M into a rectangular $`H_{\mu \nu }`$ with the boundary edge-path represented as $`\alpha (\mu )\alpha (\nu )\beta ^1(\mu )\beta ^1(\nu )`$. Then, consider the trivial covering $`\stackrel{~}{H}_{\mu \nu }=H_{\mu \nu }\mathrm{{\rm Y}}_n`$ (where $`\mathrm{{\rm Y}}_n=\{1,\mathrm{},n\}`$) of $`H_{\mu \nu }`$ by $`n`$ copies of this rectagular with the boundary edge-paths given by $`\alpha (\mu _k)\alpha (\nu _k)\beta ^1(\mu _k)\beta ^1(\nu _k),k=1,\mathrm{},n`$. Perform the set of the pairwise reidentifications of the involved edges
$$\begin{array}{c}\alpha (\rho _k)=\beta (\rho _{\sigma _\rho (k)});\sigma _\rho :k\sigma _\rho (k),k\mathrm{{\rm Y}}_n,\end{array}$$
(5.9)
where $`\sigma _\mu \sigma _1,\sigma _\nu \sigma _2`$ are supposed to satisfy the $`\delta _n`$-constraint (5.5) (with the excluded contribution of $`\widehat{T}_{\{p\}}(\widehat{T}_2)^i`$).
Evidently, the two sets (5.9) of the reidentificatrions can be concisely represented as the two closed (i.e. without branch end-points) branch cuts $`\varpi _\rho `$ of the Riemann surface $`\stackrel{~}{M}(\{\sigma _\rho \})`$ with n sheets. According to the pattern (5.9), each connected component of $`\stackrel{~}{M}(\{\sigma _\rho \})`$ has the topology of 2-torus. In compliance with (1.9), it matches with the N-independence of the argument of the $`\delta _n`$-function (5.1) taking place after the exclusion of $`P_n^{(N)}exp[\lambda \widehat{T}_2^{(n)}/N]`$.
To reintroduce the branch points (encoded in the $`\widehat{T}_{\{p\}}(\widehat{T}_2)^i`$ factor of (5.5)), recall that each admissible BP is the end-point $`q_k`$ of the associated branch cut $`\varpi ^{(k)}`$ which should be included additionally to the closed cuts $`\varpi _\rho `$ of the CS $`\stackrel{~}{M}(\{\sigma _\rho \})`$. To implement $`\varpi ^{(k)}`$, we first cut $`\stackrel{~}{M}(\{\sigma _\rho \})`$ along the support of $`\varpi ^{(k)}`$. (Both $`\varpi ^{(k)}`$ and $`\varpi _\rho `$ are all supposed to terminate at a common base-point $`p=\alpha (\mu )\alpha (\mu )`$ of M.) Then, on the left and on the right sides of each cut $`\varpi ^{(k)}`$, the resulting two copies of the n new edges of the sheets are reidentified according to the prescription (5.9). The only modification is that, instead of $`\sigma _\rho `$, one is to substitute the appropriate permutations $`\xi ^{\{p\}}T_{\{p\}}`$ and $`\tau ^{(s)}T_2`$ entering respectively $`\widehat{T}_{\{p\}}`$ and the $`s`$th $`\widehat{T}_2`$-factor in the inner product $`(\widehat{T}_2)^i`$. It completes the construction of the admissible Riemann surfaces $`\stackrel{~}{M}(\{p\},n,i)`$ entering the l.h.s. of (5.5).
### 5.2 The homomorphism of $`\pi _1(M\{q_s\}|p)`$ into $`S(n)`$.
To enumerate the equivalence classes of BCSs and justify the asserted interpretation of $`|C_f(\{p\})|`$, one is to employ the relation to the following group homomorphism where the first homotopy group $`\pi _1(M\{q_s\}\}|p)`$
$$\begin{array}{c}\psi :\pi _1(M\{q_s\}|p)S(n),\end{array}$$
(5.10)
is mapped into $`S(n)`$. (In eq. (5.10), $`M\{q_s\}`$ denotes the base-surface M with the $`i+1`$ excluded points $`q_s`$, associated to the branch points, and with the base-point p.) Given an BCS encoded in the $`\delta _n`$-function (5.5), choose a set $`\mathrm{{\rm Y}}_n=\{1,2,\mathrm{},n\}`$ to label the $`n`$ sheets (at the base-point $`p`$). Consider the lift of the closed paths in $`M\{q_s\}`$ (defining $`\pi _1(M\{q_s\}|p)`$) into the covering space $`\stackrel{~}{M}\{f^1(q_s)\}`$. Then, the (equivalence classes of the) paths induce the permutations of the labels which determine the corresponding $`S(n)`$ operators acting on $`\mathrm{{\rm Y}}_n`$.
To make (5.10) explicit, let us first specify the pattern of the first homotopy group. Consider a topological space $`T\{q_1,\mathrm{},q_m\}`$ which, for our later purposes, is allowed to be a (CW) $`2d`$ cell-complex $`T`$ (not necessarily reduced to a $`2d`$ surface) with $`m`$ deleted points $`q_k`$. Recall that in this case the group $`\pi _1(T\{q_1,\mathrm{},q_m\})`$ (in what follows we will everywhere omit the specification of the base-point p) can be represented as the following abstract group . The generators of the latter group are associated to the homotopy equivalence classes (HEC) of the uncontractible closed paths based at a given point $`p`$ (supposed to be distinct from the set $`\{q_s\}`$). In the case at hand, additionally to the generators $`\alpha _r,r=1,\mathrm{},P`$ (corresponding to $`P`$ HECs of the uncontractible cycles of $`T`$), there are extra $`m`$ generators $`\gamma ^{(s)},s=1,\mathrm{},m,`$ which refer to the HECs of closed paths encircling a single deleted point $`q_s`$. Finally, there exists a constructive algorithm to find the complete set of $`K𝐙_\mathrm{𝟏}`$ relations $`\{F_l(\{\alpha _r,\gamma ^{(s)}\})=1\}_{l=1,\mathrm{},K}`$ that completes the intermediate mapping of the $`\pi _1(T\{q_1,\mathrm{},q_m\})`$ generators into the abstract group.
Returning to the case of a genus $`g`$ $`2d`$ surface $`T=M_g`$, there is a single relation (with $`P=2g,[\alpha _i,\alpha _k]\alpha _i\alpha _k\alpha _i^1\alpha _k^1`$) defining $`\pi _1(M_g\{q_s\})`$
$$\begin{array}{c}F(\{\alpha _r,\gamma ^{(s)}\})=\left(_{j=1}^g[\alpha _j,\alpha _{g+j}]_{s=1}^m\gamma ^{(s)}\right)=1.\end{array}$$
(5.11)
Comparing (5.11) with the pattern (5.1)/(5.5), one deduces the explicit form of the homomorphism (5.10) (with the identification $`g=1,m=i+1`$)
$$\begin{array}{c}\psi :\psi (\alpha _\rho )=\sigma _\rho ;\psi (\gamma ^{(1)})=\xi ^{\{p\}};\psi (\gamma ^{(s)})=\tau ^{(s1)},s2,\end{array}$$
(5.12)
where $`\xi ^{\{p\}}T_{\{p\}}`$ and $`\tau ^{(s1)}T_2`$ which enter respectively $`\widehat{T}_{\{p\}}`$ and the $`(s1)`$th $`\widehat{T}_2`$-factor (in the inner product $`(\widehat{T}_2)^i`$).
### 5.3 The symmetry factor.
Given the homomorphisms (5.12), one can readily enumerate the equivalence classes $`\stackrel{~}{M}(\{p\},n,i)`$ of the BCSs employing the following notion of equivalence of two homomorphisms $`\psi _1`$ and $`\psi _2`$. The latter are postulated to belong to the same equivalence class if there exists some $`\eta S(n)`$ so that
$$\begin{array}{c}\psi _1(\zeta )=\eta \psi _2(\zeta )\eta ^1,\zeta \pi _1(M\{q_1,\mathrm{},q_m\});\eta S(n).\end{array}$$
(5.13)
Then, the basic theorem of the topological coverings ensures that the inequivalent homomorphisms (5.12) are in one-to-one correspondence with the associated homeomorphically distinct branched covering spaces.
Finally, let $`\kappa :f\kappa =f,`$ denotes a particular automorphism of the branched covering space $`\stackrel{~}{M}`$. Being restricted to the n-set $`\mathrm{{\rm Y}}_n=f^1(\stackrel{~}{p})`$ ($`pq_s,s`$), the group of the automorphisms is isomorphic to (the conjugacy class of) the $`S(n)`$-subgroup $`C_f(\{p\})`$ that induces conjugations (5.13) leaving all the images $`\psi _1(\zeta )`$ invariant: $`\psi _1(\zeta )=\psi _2(\zeta ),\zeta ,\eta (\kappa )C_f(\{p\})`$. In turn, it justifies the required interpretation of $`|C_f(\{p\})|`$.
## 6 The stringy form of $`B(\{n_{\mu \nu }\})`$ in $`D3`$.
To rewrite the $`D3`$ amplitude (4.8) in the form generalizing the $`D=2`$ stringy pattern (5.3)-(5.5), we first expand each factor $`Q_{n_{\mu \nu }}`$ and select the $`i_{\mu \nu }`$th power $`(\widehat{T}_2^{(n_{\mu \nu })})^{i_{\mu \nu }}`$ (where the $`k_{\mu \nu }`$th $`\widehat{T}_2^{(n_{\mu \nu })}`$-factor in the latter product is supposed to be defined via eq. (4.5) in terms of $`\tau _{\mu \nu }^{(k_{\mu \nu })}T_2^{(n_{\mu \nu })},k_{\mu \nu }=1,\mathrm{},i_{\mu \nu }`$). Complementary, akin to eq. (4.6) one is to decompose
$$\begin{array}{c}\mathrm{\Lambda }_{n_\varphi }^{(m_\varphi )}=_{T_{\{p_\varphi \}}^{(n_\varphi )}S(n_\varphi )}\mathrm{\Lambda }_{n_\varphi }^{(m_\varphi )}(T_{\{p_\varphi \}}^{(n_\varphi )})_{\xi ^{\{p_\varphi \}}T_{\{p_\varphi \}}^{(n_\varphi )}}\xi ^{\{p_\varphi \}},\end{array}$$
(6.1)
and separate the contribution of a given $`\widehat{T}_{\{p_\varphi \}}^{(n_\varphi )}\widehat{T}_{\{p_\varphi \}}`$ (defined by eq. (4.6)).
Given the above expansions, one can prove (see below) that the associated building block of (4.8)
$$\underset{\{\sigma _\rho \}}{}\delta _{n_+}(\widehat{T}_{\{p_+\}}\underset{\{\mu \nu \}}{}\widehat{T}_{\{p_{\mu \nu }\}}\frac{(\widehat{T}_2^{(n_{\mu \nu })})^{i_{\mu \nu }}}{n_{\mu \nu }!}\underset{\{\rho \}}{}(\sigma _\rho \widehat{1}_{[\frac{n_+}{n_\rho }]})\underset{\{\mu \}}{}(\sigma _\mu ^1\widehat{T}_{\{p_\mu \}}\widehat{1}_{[\frac{n_+}{n_\mu }]}))=$$
$$\begin{array}{c}=_{\phi \stackrel{~}{M}(\{p_\varphi \},\{n_\varphi \},\{i_{\mu \nu }\})}\frac{1}{|C_\phi (\{p_\varphi \})|}\end{array}$$
(6.2)
can be rewritten as the sum over the relevant $`D2`$ mappings (1.11) to be reconstructed in the next subsection. Extending the $`D=2`$ theorem due to Gross and Taylor, each term of the latter sum is weighted by the inverse number $`|C_\phi (\{p_\varphi \})|`$ of distinct automorphisms associated to a given $`\stackrel{~}{M}_\phi \stackrel{~}{M}(\{p_\varphi \},\{n_\varphi \},\{i_{\mu \nu }\})`$.
As we will demonstrate, the involved into (6.2) spaces $`\stackrel{~}{M}_\phi `$ can be viewed as the generalized Riemann surfaces to be identified with the worldsheets of the YM-flux ’wrapped around’ the EK $`2d`$ cell-complex (2.2). Combining all the pieces together, the stringy reinterpretation of (4.8) essentially follows the steps discussed in the $`D=2`$ case. Namely, leaving aside the $`P_{n_\varphi }^{(N)}(T_{\{p_\varphi \}})`$ factors (4.1), the rest of the ingredients of (4.8) readily fit in the consistent $`D3`$ extension (1.11) of the Gross-Taylor stringy pattern (1.9). Indeed, by the same token as in the $`D=2`$ analysis, the involved (movable or nonmovable) branch points, microscopic tubes, and handles are weighted by the $`1/N`$\- and $`n_\varphi `$-dependent factors according to their contribution to the Euler character of $`\stackrel{~}{M}_\phi `$. The latter ’local’ matching is completed by eq. (6.2) together with the following ’global’ matching. To see it, let us remove temporarily the branch points and the collapsed subsurfaces via the substitution
$$\begin{array}{c}\mathrm{\Lambda }_{n_\varphi }^{(m_\varphi )}N^{n_\varphi m_\varphi };Q_{n_{\mu \nu }}exp[\lambda n_{\mu \nu }/2],\end{array}$$
(6.3)
into eq. (4.8). Then, the remaining $`\lambda n_{\mu \nu }`$-dependent factor and the overall power $`(1/N)^\epsilon `$ (inherited from eq. (4.8))
$$\begin{array}{c}\epsilon =\underset{\mu \nu =1}{\overset{D(D1)/2}{}}n_{\mu \nu }\underset{\rho =1}{\overset{D}{}}n_\rho +n_+=0;2n_+=_{\{\rho \}}n_\rho =2_{\{\mu \nu \}}n_{\mu \nu },\end{array}$$
(6.4)
match with the area and the 2-tora topology (revealed below) of each connected component of $`\stackrel{~}{M}_\phi `$ corresponding to the deformation (6.3).
By the same token as in the $`D=2`$ case, the manifest stringy representation of the full amplitude (4.8) (including the $`P_{n_\varphi }^{(N)}(T_{\{p_\varphi \}})`$-factors) calls for a resummation into the appropriate large $`N`$ SC series. Presumably, it eliminates these factors trading each remaining in (4.8) $`S(n_\varphi )`$-operator for its $`S(n_\varphi ^+)S(n_\varphi ^{})`$ descedant. It is suggestive that the short cut way for the latter reformulation is provided by the direct extention of the $`D=2`$ prescription (5.7),(5.8): one is to substitute the labels $`n,n^+,n^{}`$ by their properly associated $`\varphi `$-dependent counterparts $`n_\varphi ,n_\varphi ^+,n_\varphi ^{}`$ (where $`n_\varphi ,\varphi \{\mu \nu \},\{\rho \},+,`$ labels the relevant $`S(n_\varphi )`$-structures). Note also that, similarly to the $`D=2`$ case, a generic $`D3`$ model (1.1) complies with the pattern (1.9) as well.
In conclusion, we remark that the amplitude (4.8) exhibits purely topological assignement (see Section 6.2 for more details) of the $`S(n_\varphi )`$-structures after the formal deformation $`\{\mathrm{\Lambda }_{n_\varphi }^{(m_\varphi )}(N^{n_\varphi }\mathrm{\Omega }_{n_\varphi })^{m_\varphi }\},\{Q_{n_{\mu \nu }}1\}`$. In particular, consider the powers $`m_\varphi =\pm 1`$ in which the relevant $`S(n_\varphi )`$-twists $`\mathrm{\Omega }_{n_\varphi }`$ enter the $`\mathrm{\Lambda }_{n_\varphi }^{(m_\varphi )}`$ factors (4.1): $`\{m_{\mu \nu }=1\},\{m_\rho =1\},m_+=1`$. The latter are equal to the weights in the famous formula computing the Euler characteristic of a given cell-complex $`T`$: $`22G_T=n_pn_l+n_s=n_pm_p+n_lm_l+n_sm_s`$, where $`n_p,n_l,n_s`$ are the total numbers of respectively plaquettes, links, and sites of $`T`$. (The earlier heuristic arguments (in the case of the single chiral sector), consistent with the above pattern of the $`\mathrm{\Omega }_{n_\varphi }`$ assignement, can be found in ).
The topological nature of the considered defomation of the amplitudes like (6.2) can be made transparent by the conjecture building on the $`D=2`$ observation . One may expect that, after this deformation, the TEK partition function $`\stackrel{~}{X}_D`$ yields the generating functional for the orbifold Euler characters of the Hurwitz-like spaces associated to the following set. The latter includes all the generalized branched covering spaces $`\stackrel{~}{M}_\phi `$ (of the EK $`2d`$ cell-complex (2.2)) corresponding in eq. (6.2) to the considered deformation $`\{P_{n_\varphi }^{(N)}1\}`$. (For a close but somewhat distinct earlier conjecture, see .) Similar proposal can be made for the (weaker) ’topological’ deformation $`\{Q_{n_{\mu \nu }}1\}`$ of the PF $`\stackrel{~}{X}_D`$ reintroducing the $`P_{n_\varphi }^{(N)}`$ factors.
### 6.1 The generalized covering spaces of $`T_{EK}\{q_s\}`$.
Given the data in the argument of the $`\delta _{n_+}`$-function in the l.h.s of eq. (6.2), our aim is to reconstruct the associated topological spaces $`\stackrel{~}{M}_{EK}\stackrel{~}{M}_\phi `$ (i.e. the mappings (1.11)) which are summed over in the r.h.s. of (6.2). This problem, being inverse to what is usually considered in the framework of algebraic topology , will be resolved through the sequence of steps which appropriately generalizes the $`D=2`$ analysis of the previous section. It is noteworthy that the spaces $`\stackrel{~}{M}_{EK}`$ do not coincide with the canonical branched covering spaces (BCSs) of the TEK $`2d`$ cell-complex $`T_{EK}`$. To say the least, a generic canonical BCS of $`T_{EK}`$ is again a $`2d`$ cell-complex (with a number of branch points) but not a $`2d`$ surface like in the asserted mappings (1.11). As we will see, the basic amplitude (6.2) indeed complies with (1.11) and refers to the novel class of the associated to $`T_{EK}`$ spaces $`\stackrel{~}{M}_{EK}`$ to be called the generalized branched covering spaces (GBCSs) of $`T_{EK}`$.
Similarly to the $`D=2`$ analysis, we start with the simpler case of the generalized covering spaces (GCSs) of $`T_{EK}`$ corresponding to the deformation (6.3). Being specified by the immersions (1.11) without the (branch points’ and collapsed subsurfaces’) singularities, the GCSs can be reconstructed generalizing the $`D=2`$ surgery-construction of the covering spaces. Take $`D(D1)/2`$ $`\mu \nu `$-rectangulars $`H_{\mu \nu }`$ representing the ’decompactified’ plaquettes of the EK base-lattice (2.2). Then, one is to begin with trivial ($`n_{\mu \nu }`$-sheet) coverings $`\stackrel{~}{H}_{\mu \nu }=H_{\mu \nu }\mathrm{{\rm Y}}_{n_{\mu \nu }}`$ of $`H_{\mu \nu }`$ with the edge paths $`_q\alpha (\mu _q)\alpha (\nu _q)\beta ^1(\mu _q)\beta ^1(\nu _q),q=1,\mathrm{},n_{\mu \nu }`$. To reproduce the effect of the $`\sigma _\rho `$-permutations in (6.2), first let us denote by $`\{\alpha (\rho _k),k=1,\mathrm{},n_\rho =_{\nu \rho }^{D1}n_{\rho \nu }\}`$ the $`\rho `$-set of the $`\alpha (\rho _q)`$ edges (collected from the $`(D1)`$ different $`\rho \nu `$-plaquettes) ordered in accordance with the pattern (3.14) of the $`S(n_\rho )`$ basis $`|I_{n(\rho )}^{(\pm )}>`$. Similarly, we introduce the sets $`\{\beta (\rho _k),k=1,\mathrm{},n_\rho \}`$.
To reconstruct the associated to (4.8) GCS of $`T_{EK}`$, at each particular $`\rho `$-link one is to perform the pairwise $`\sigma _\rho `$-identifications of $`\alpha (\rho _k),\beta (\rho _k)`$ according to the prescription (5.9). In this way, we have constructed appropriate conglomerate of the generalized Riemann surfaces $`\stackrel{~}{M}(\{\sigma _\rho \})`$ (with the total number $`D`$ of the closed branch cuts) which are wrapped around $`T_{EK}`$ in compliance with the mapping (1.11). Evaluating the Euler characteristic $`\epsilon `$ of $`\stackrel{~}{M}(\{\sigma _\rho \})`$ in accordance with (5.9), each connected component of the latter surface reveals the topology of a 2-tora. In turn, it matches with the overall power (6.4) of the $`1/N`$ factor which is inherited via (6.3) from the product (4.8) of the $`\mathrm{\Lambda }_{n_\varphi }^{(m_\varphi )}`$ factors.
Given the GCSs $`\stackrel{~}{M}(\{\sigma _\rho \})`$ of $`T_{EK}`$, the branch points (encoded in (6.2) through the operators $`\widehat{T}_{\{p_{\mu \nu }\}}(\widehat{T}_2^{(n_{\mu \nu })})^{i_{\mu \nu }}`$) are reintroduced in essentially the same way as we did for the $`D=2`$ case. The only subtlety is that, to make the employed cutting-gluing rules well-defined, it is convenient to view the EK cell-complex $`T_{EK}`$ as the homotopy retract of a ’less singular’ $`2d`$ complex $`T_{EK}^{}`$ (possessing by construction the same fundamental group $`\pi _1(T_{EK})=\pi _1(T_{EK}^{})`$). We refer the reader to Appendix C for the details, and now proceed with the $`D3`$ generalization of the $`D=2`$ Gross-Taylor theorem concerning the interpretation of the symmetry factor entering the $`D2`$ stringy amplitudes like (1.9),(6.2).
### 6.2 The homomorphism of $`\pi _1(T_{EK}\{q_s\})`$ into $`S(n_+)`$.
To effectively enumerate the constructed GBCSs of $`T_{EK}`$ and their automorphisms, we first reconstruct what kind of homomorphism like (5.10) is encoded in the $`D3`$ pattern (6.2). It will provide with the precise mapping of the fundamental group $`\pi _1(T_{EK}\{q_s\})`$ (of $`T_{EK}`$ with a number of deleted points $`\{q_s\}`$ associated to the branch points) into the enveloping $`S(n_+)`$ group.
Observe first that the abstract group representation of $`\pi _1(T_{EK})`$ is defined, according to (2.2), in terms of the $`D`$ generators $`\alpha _\rho `$ corresponding to the uncontractible cycles (i.e. the compactified $`\rho `$-links) of $`T_{EK}`$. Having excluded from $`T_{EK}`$ a set of points $`\{q_s\}`$, it is convenient to recollect them into the three varieties of the $`\varphi `$-subsets $`\{q_s\}=_\varphi \{q_{k_\varphi }^{(\varphi )}\},\varphi \{\mu \nu \},\{\rho \},+,k_\varphi =1,\mathrm{},b_\varphi ,`$ belonging respectively to the interior of the $`\mu \nu `$-plaquette, to the interior of the $`\rho `$-link, and to the single site of $`T_{EK}`$. The set of the $`\pi _1(T_{EK}\{q_s\})`$ generators includes (additionally to the subset $`\{\alpha _\rho \}`$ inherited from $`\pi _1(T_{EK})`$) $`\gamma _{k_\varphi }^{(\varphi )}`$ associated to (the equivalence classes of) the closed paths encircling a single deleted point $`q_{k_\varphi }^{(\varphi )}`$. The representation of $`\pi _1(T_{EK}\{q_s\})`$ is then completed by the $`D(D1)/2`$ relations
$$\begin{array}{c}\left([\alpha _\mu ,\alpha _\nu ]_{\varphi =\mu \nu ,\mu ,\nu ,+}_{k_\varphi =1}^{b_\varphi }\gamma _{k_\varphi }^{(\varphi )}\right)=1;\mu \nu =1,\mathrm{},D(D1)/2,\end{array}$$
(6.5)
each of which can be viewed as the $`\mu \nu `$-’copy’ of (5.11).
Comparing the argument of the $`\delta _{n_+}`$-function in the l.h.s. of eq. (6.2) with the pattern of (6.5), it is straightforward to write down the precise form of the relevant $`D2`$ homomorphism generalizing (5.10):
$$\begin{array}{c}\psi :\psi (\alpha _\rho )=\sigma _\rho ;\psi (\gamma _1^{(\varphi )})=\xi ^{\{p_\varphi \}};\psi (\gamma _{k_{\mu \nu }}^{(\mu \nu )})=\tau _{\mu \nu }^{(k_{\mu \nu }1)},k_{\mu \nu }2,\end{array}$$
(6.6)
where $`\xi ^{\{p_\varphi \}}T_{n_\varphi },\tau _{\mu \nu }^{(k_{\mu \nu }1)}T_2^{(n_{\mu \nu })}`$. Note that the three $`\varphi `$-varieties of the generators $`\gamma _1^{(\varphi )}`$ match with the three species of $`\widehat{T}_{n_\varphi }`$ in eq. (6.2), while $`\gamma _{k_{\mu \nu }}^{(\mu \nu )}`$ represents the $`k_{\mu \nu }`$th $`\widehat{T}_2^{(n_{\mu \nu })}`$-factor in the product $`(\widehat{T}_2^{(n_{\mu \nu })})^{i_{\mu \nu }}`$. In particular, eq. (6.6) implies that in eq. (6.5) $`b_{\mu \nu }=i_{\mu \nu }+1`$ while $`b_+=b_\rho =1,\rho `$.
#### 6.2.1 The symmetry factor.
To deduce the announced interpretation of $`|C_\phi (\{p_\varphi \})|`$ (entering the r.h.s. of eq. (6.2)), let us start with the following observation. The summation in the l.h.s of (6.2) can be viewed as the sum over the associated homomorphisms (6.6) constrained by the condition: $`\xi ^{\{p_\varphi \}}T_{\{p_\varphi \}}`$ and $`\tau _{\mu \nu }^{(k_{\mu \nu }1)}T_2^{(n_{\mu \nu })}`$. Therefore, one is to identify the proper equivalence classes of the homomorphisms (6.6) to parametrize uniquely the topologically inequivalent spaces $`\stackrel{~}{M}_\phi `$ constructed in Section 6.1.
We assert that the two homomorphisms (6.6), $`\psi _1`$ and $`\psi _2`$, are equivalent (i.e. the corresponding generalized branched coverings $`\stackrel{~}{M}_\phi `$ are homeomorphic) if and only if there exists some $`\eta _{\mu \nu }S(n_{\mu \nu })`$ so that
$$\begin{array}{c}\psi _1(\zeta )=\eta \psi _2(\zeta )\eta ^1,\zeta \pi _1(T_{EK}^{}\{q_s\});\eta _{\mu \nu }S(n_{\mu \nu }),\end{array}$$
(6.7)
where the basis (3.32) for $`\eta S(n_+)`$ is implied. For a preliminary orientation, one observes that the conjugations (6.7) are induced by the $`_{\mu \nu }S(n_{\mu \nu })`$ permutations of the sheets (of $`\stackrel{~}{M}_\phi `$) separately within each of the $`D(D1)/2`$ distinct $`n_{\mu \nu }`$-subsets which covers (see Section 6.1) a given $`\mu \nu `$-tora $`E_{\mu \nu }`$ combined into $`T_{EK}^{}`$. (Complementary, $`_{\mu \nu }S(n_{\mu \nu })`$ is the largest subgroup of $`S(n_+)`$ providing with the conjugations leaving the argument of the $`\delta _{n_+}`$-function in eq. (6.2) invariant.) We refer to Appendix D for the justification of the notion (6.7) of the equivalence and now simply deduce its consequences. Consider a particular branched covering $`\stackrel{~}{M}_\phi `$, and let $`C_\phi (\{p_\varphi \})`$ is the group of the inequivalent automorphisms $`\kappa `$ of $`\stackrel{~}{M}_\phi `$: $`\phi \kappa =\phi `$. Take the restriction of a given $`\kappa `$ to the $`\mathrm{{\rm Y}}_{n_+}`$ space of $`\{\phi ^1(p)\})`$. Akin to the $`D=2`$ case, one shows that $`C_\phi (\{p_\varphi \})`$ is isomorphic to the conjugacy class (with respect to (6.7)) of the following subgroup of $`_{\mu \nu }S(n_{\mu \nu })`$. The latter is associated to such subset of conjugations (6.7) that leave all $`\psi _1(\zeta )`$ invariant: $`\psi _1(\zeta )=\psi _2(\zeta ),\zeta ,\eta (\kappa )C_\phi (\{p_\varphi \})`$. In sum, there are exactly $`(_{\mu \nu }n_{\mu \nu }!)/|C_\phi (\{p_\varphi \})|`$ distinct homomorphisms $`\{\psi _k\}`$ associated to the same topology-type of $`\stackrel{~}{M}_\phi `$ that justifies the basic identity (6.2).
## 7 The area-preserving homeomorphisms.
On a lattice, our $`D3`$ stringy proposal can be viewed as the confluence of the Wilson’s string-like reformulation of the (finite N) SC series with the power of the large N expansion. The feature, which sharply distinguishes the $`D3`$ Gauge String from the earlier $`D3`$ proposals in this direction, is the asserted invariance of the weights $`w[\stackrel{~}{M}]`$ under certain continuous group of the area-preserving worldsheet homeomorphisms. In $`D3`$, the latter symmetry is encoded in the following $`D3`$ ’descedant’ inherited from the $`D=2`$ renormgroup (RG) invariance of the $`YM_2`$ systems (1.1).
Recall first that a $`D3`$ $`YM_D`$ theory (1.1) (having some lattice $`𝐋^𝐃`$ as the base-space) can be equally viewed as the $`YM`$ theory being defined on the $`2d`$ skeleton $`𝐓^𝐃`$ of $`𝐋^𝐃`$ represented by the associated $`2d`$ cell-complex. Consider the partition function (PF) of (1.1) on $`𝐓^𝐃=_kE_k`$ composed from the associated $`2d`$ surfaces $`E_k`$ (of the areas $`A_k`$ and with certain boundaries) according to the corresponding incidence-matrix . Akin to the $`D=2`$ case, this PF is invariant under subdivisions of $`𝐓^𝐃`$ (preserving the total areas $`A_k`$) so that $`E_k`$ can be made into the associated smooth $`2d`$ manifolds $`M_k`$ (with boundaries) combined into a $`2d`$ cell-complex $`\stackrel{~}{𝐓}^𝐃=_kM_k`$ (homeomorphic to $`𝐓^𝐃`$). Therefore, refining the discretization, the $`YM_D`$ lattice theory (1.1) can be adjusted to merge with the following continuous system. To implement the latter, we first put the associated to (1.1) continuous $`YM_2`$ theory (1.4) (keeping free boundary conditions) on each $`M_k`$ . Then, one is to average over the gauge fields, ’living’ on the boundaries of $`\{M_k\}`$, in compliance with the associated to $`\stackrel{~}{𝐓}^𝐃`$ incidence-matrix.
Similarly to the $`D=2`$ case (see Section 1.1A), on the side of the proposed $`D3`$ Gauge String, the above relation to the continuous $`YM_2`$ system ensures the required invariance of the lattice string weights (entering the amplitudes defined by the $`D3`$ extension (1.11) of (1.9)). Indeed, on the one hand, both the parity $`P_\phi `$ and the symmetry factor $`|C_\phi |`$ depend only on the topology of the worldsheet $`\stackrel{~}{M}(T)`$ and on the one of the corresponding taget-space $`T=_kE_k`$. (Compare it with where the singular lattice geometry of T is ’built into’ the curvature-defects.) On the other hand, the sum over the mappings $`\phi `$ also supports the required invariance of the set of relevant $`\stackrel{~}{M}(T)`$. Indeed, the multiple integrals (rather than discrete sums as in ) over all admissible positions of the movable branch points (and certain collapsed subsurfaces) span the entire interior of the discretized $`2d`$ surfaces $`E_k`$ combined into T. The remaining ’discrete’ contributions into $`𝑑\phi `$ refer to the purely topological assignement of the nonmovable singularities anywhere in the interiors of the appropriate subspaces of T. Altogether, the relevant group of the worldsheet homeomorphisms continuously translates the positions of the admissible singularities of the map (1.11) within the interiors of the associated subspaces of T.
To be more specific, we compare the weights $`w[\stackrel{~}{M}]`$ associated to the $`YM`$ system (1.1) defined on a generic base-lattice $`F_{EK}`$ homeomorphic to the ’elementary’ EK $`2d`$ cell-complex (2.2). Let the total area (measured in the dimensionless units) of a given torus $`E_{\mu \nu }`$ is equal to $`A_{\mu \nu }`$. Combining the above general arguments with the analysis of Sections 3 and 4, the general weights on $`F_{EK}`$ can be deduced from the basic $`T_{EK}`$ amplitude (4.8) trading in each $`Q_{n_{\mu \nu }}`$-factor (4.5) the coupling constant $`\lambda `$ for $`\lambda A_{\mu \nu }`$. In particular, the $`(\mathrm{\Omega }_{n_\varphi })^{m_\varphi }`$-, $`\sigma _\rho `$-twist assignement is indeed purely topological: the corresponding conglomerates of the nonmovable branch points (and collapsed subsurfaces) can be placed anywhere (but without summation over positions) in the interior of the associated macroscopic $`\varphi `$-,$`\rho `$-cell of $`F_{EK}`$.
## 8 Smooth Gauge Strings vs. lattice ones.
It is clear that the method, we have developed in Sections 3-6 on the example of the TEK model (2.1), can be extended for the $`YM_D`$ theories (1.1) defined on a generic regular subspace T of the $`2d`$ skeleton of the D-dimensional base-lattice $`𝐋^𝐃`$. We defer the analysis of the generic weights (1.2) for a separate paper and now simply stress those implications of our present results which are T-independent and common for both the lattice and smooth realizations of the Gauge String. In this way, one actually decodes all the major peculiar features of the continuous theory of the smooth YM-fluxes which are novel compared to the conventional paradigm of the $`D3`$ ’fundamental’ strings.
To begin with, one observes that the relevant smooth mapping (see e.g. eq. (1.11) for the worldsheets $`\stackrel{~}{M}`$ without boundaries) are allowed to have singularities which usually are not included into the $`D3`$ string-pattern (where the maps (1.11) are restricted to be sheer immersions). Complementary, certain class of the selfintersecting worldsheets $`\stackrel{~}{M}`$ is endowed in eq. (1.2) with the extra factors $`J[\stackrel{~}{M}(T)|\{\stackrel{~}{b}_k\}]1`$ that does not have a direct counterpart in the known $`D3`$ string theories. In turn, it is the extra $`J[..]`$-weights which encode the data sufficient to reconstruct the associated continuous $`YM_D`$ model (1.3)/(1.4). As for a single $`D3`$ Wilson loop average, the major novel ingredient is due to the various movable branch points and the collapsed $`2d`$ subsurfaces (assigned with the $`\{\stackrel{~}{b}_k\}`$-dependent weights (1.10)). Their positions can be viewed as the zero modes associated to the minimal surface contribution (provided the latter properly selfintersects) resulting in the extra area-dependence in the preexponent of $`<W_C>`$. This pattern generalizes the one of the $`D=2`$ loop-averages .
Next, let us demonstrate that the total contribution of the worldsheets $`\stackrel{~}{M}`$ (and, consequently, of the taget-spaces T) with the backtrackings vanishes which matches with the absence of foldings in the $`D=2`$ Gross-Taylor representation (1.9). To make this feature manifest, we first recall that the relevant for (6.2) maps (1.11) (associated to the SU(N) TEK model) do not include any foldings of the worldsheets $`\stackrel{~}{M}`$ wrapped around $`T_{EK}`$. On the other hand, one could equally start with the large N U(N) TEK model where the proposed technique would result in all kinds of foldings ’covering’ various conglomerates of the plaquettes (with possibly different $`\mu \nu `$-orientation). The comparison of these two complementary large $`N`$ patterns justifies the above assertion.
On the side of the gauge theory defined on the standard cubic lattice $`𝐋^𝐃`$, the TEK foldings are associated to the more general backtrackings of $`\stackrel{~}{M}`$. To see it, let us call worldsheets $`\stackrel{~}{M}(T)`$ (or taget-spaces T) without any backtrackings regular. Then, to any regular taget $`T^{(r)}`$ one can associate the variety of taget-spaces creating all kinds of the backtracking (bounding a zero 3-volume) with the support not necessarily belonging the original regular taget $`T^{(r)}`$ (in contradistinction to the foldings discussed in ).
The irrelevance of the backtrackings is intimately related to the invariance (in gauge theories) of the Wilson loop averages $`<W_C>`$ under the zig-zag backtrackings of the boundary contour C that is in sharp contrast with the situation in the conventional Nambu-Goto string and most of its existing generalizations. (Recently Polyakov advocated the latter invariance as the crucial feature of the strings dual to gauge theories.) In the limit $`N\mathrm{}`$, the zig-zag symmetry of $`<W_C>`$ can be made manifest confronting as previously the two large N formulations of the open Gauge String: the $`U(N)`$ one with a given backtracking contour C should be compared with the $`SU(N)`$ one with $`\stackrel{~}{C}`$ obtained contracting zig-zags of $`C`$.
### 8.1 Correspondence with the WC Feynman diagrams.
Let us reveal the WC/SC correspondence between the continuous $`YM_D`$ models (1.4) in the WC phase and the associated smooth Gauge String in the SC regime. According to Section 7, once the backtrackings are absent, the relevant weights $`w[\stackrel{~}{M}(T)]`$ can be derived by putting the $`YM_2`$ system (1.4) onto a given $`2d`$ cell-complex (e.g. $`2d`$ surface) $`T^{(r)}=_kM_k`$ where the $`2d`$ surfaces $`M_k`$ are piecewise smooth. In the large N SC regime, the latter $`YM`$ system is represented by the conglomerates of the worldsheets properly wrapped around $`T^{(r)}`$. On the other hand, in the large N WC regime the WC perturbation theory represents our system through the set the Feynman diagrams. Employing the standard path integral representation of the propagator of the free particle (in a curved space), the diagrams are visualized as the fishnet (of the gluonic trajectories) appropriately ’wrapped’ around the same $`2d`$ complex $`T^{(r)}`$. The crucial observation is that, when $`T^{(r)}`$ is viewed as being embedded into $`𝐑^𝐃`$, the latter fishnet can be reinterpreted as the specific contribution of the WC perturbation theory in the D-dimensional continuous $`YM_D`$ model (1.4) (uniquely associated to the chosen $`YM_2`$ via eq. (1.3)). To select this contribution on the $`YM_D`$ side, not only the gluonic trajectories in $`𝐑^𝐃`$ should be constrained to have the space-time support on $`T^{(r)}`$ but also each colour $`a`$-component of the associated gluonic strength-tensor $`F_{\mu \nu }^a(𝐳),𝐳T^{(r)},`$ as the Lorentz (or $`O(D)`$) tensor should belong to the tangent space of $`T^{(r)}`$ at $`𝐳`$. (To circumvent gauge-fixing, tricky to explicitly match between the $`T^{(r)}`$\- and standard formulations, one is to introduce an infinitesimally small mass term for the gauge field and then perform the above comparison.)
### 8.2 Suppression of the selfintersections.
At this step, it is appropriate to discuss a number of crucial simplifications inherent in the dynamics of the considered $`D3`$ continuous flux-theory compared to its lattice counterpart. Consider first the subset of smooth closed $`2d`$ surfaces $`\stackrel{~}{M}`$ (resulting from the smooth immersions of a $`2d`$ manifold $`M`$ into $`𝐑^𝐃`$ which alternatively can be represented by the maps (1.11) with $`T𝐑^𝐃`$) with selfintersections on submanifolds of the dimensionality $`d>d_{cr}=(4D)`$. The latter subset is of measure zero in the set of all smooth immersions $`M𝐑^𝐃`$ (with arbitrary selfintersections). Complementary, the smooth immersions $`M𝐑^𝐃,D4,`$ are dense in the space of all piecewise smooth immersions $`M𝐑^𝐃`$. In particular, in $`D5`$ the subset of $`\phi `$-maps resulting in smooth embeddings (i.e. in closed nonselfintersecting $`2d`$ manifolds $`\stackrel{~}{M}`$) is the open, dense subspace of the space of all piecewise smooth immersions $`M𝐑^𝐃`$. (Another example, in $`D3`$ the backtrackings (being viewed as two-dimensional selfintersections) are of measure zero, i.e. unstable.)
Next, suppose that the redefinition (discussed in the very end of Section 1 and after eq. (5.6)) of the bare $`SU(N)`$ string tension $`\stackrel{~}{\sigma }_0`$ is performed that eliminates certain types of the collapsed $`2d`$ subsurfaces originally built into (1.11). Then, by virtue of the Whitney immersion theorem , the singularities of the generic (piecewise smooth) mapping (1.11) in $`D4`$ can be restricted to the ones listed after eq. (1.11) with the exclusion (in the Heat-Kernal case (1.7)) of the ’movable’ collapsed subsurfaces attached to a single sheet of the covering. As for the remaining admissible collapsed subsurfaces and the branch-points, being necessarily attached to nontrivial selfintersections of the worldsheet, they therefore correspond in $`D5`$ to ’measure-zero’ limiting points of the dense subspace of the $`2d`$ manifolds $`\stackrel{~}{M}`$ (induced by the embeddings (1.11)) without boundaries. As for $`D=4`$, the stable (i.e. of nonzero measure) selfintersections of a closed smooth surface $`\stackrel{~}{M}`$ can occur only at a set of isolated points. The advantage of the Gauge String is that, as it is clear from the previous sections, such zero-dimensional selfintersections are not weighted by any extra factors so that the corresponding worldsheets are still assigned with the $`J[\stackrel{~}{M}(T)|\{\stackrel{~}{b}_k\}]=1,b=0,`$ reduction of the weight (1.2). Altogether, it justifies the assertion made in the very end of Section 1.
To see how the $`J[\stackrel{~}{M}(T)|\{\stackrel{~}{b}_k\}]1`$ pattern is observable in $`D3`$, recall first the basic theorem on the stability of the smooth immersions: the stability is equivalent to the local stability. In particular, it implies that the unstable in the case of a closed $`2d`$ surface selfintersections can not be stabilized (in the bulk) introducing some selfintersecting boundary contour(s). To be more specific, consider the framework of the quasiclassical expansion for $`<W_C>`$ where the weight of the minimal surface enters as the isolated, discrete contribution disparate from the continuum of the string fluctuations. As a result of the above theorem, in $`D4`$ the only place where a nontrivial $`J[\stackrel{~}{M}(T)|\{\stackrel{~}{b}_k\}]1`$ factor may be observable (beyond the considered $`SU(N)`$ redefinition of $`\stackrel{~}{\sigma }_0`$) for a nonbacktracking contour $`C`$ seems to be the contribution of the selfintersecting minimal surface. The simplest option, where the weights (1.10) of the movable branch points can be ’measured’, is to take contours winding a number of times around a nonselfintersecting loop C. On the other hand, let C is any nonbacktracking loop associated to some nonselfintersecting ’minimal-area’ surface(s). In this case, we expect that the simplest $`J[\stackrel{~}{M}(T)|\{\stackrel{~}{b}_k\}]=1,b=1`$ pattern of (1.2) (with the above redefinition of $`\stackrel{~}{\sigma }_0`$ and with the worldsheets represented by the smooth strict immersions $`\stackrel{~}{M}:M𝐑^𝐃`$) in $`D4`$ is sufficient to reproduce the correct result for the contribution of the genus h smooth worldsheets to the average $`<W_C>`$.
## 9 Conclusions.
We propose the correspondence between the smooth Gauge String (1.2), induced from the lattice models (1.1), and the associated via (1.3) continuous $`YM_D`$ theory (1.4) (with a finite UV cut off) in the large N SC phase. This duality implies the concrete prediction (1.6) for the bare string tension $`\sigma _0=\mathrm{\Lambda }^2\stackrel{~}{\sigma }_0`$ as the function of the coupling constants entering the $`YM_D`$ lagrangian (1.4). In particular, it readily allows for a number of nontrivial large N predictions in the extreme $`SC`$ limit where $`\sigma _0`$ merges with the leading asymptotics of the physical string tension $`\sigma _{ph}`$. More generally, the correspondence asserts that the generic continuous $`YM_D`$ model (1.4) is confining in $`D3`$ at least in the large N SC regime accessible by our approach.
In turn, it suggests the mechanism of confinement in the standard weakly coupled (WC) $`D=4`$ continuous gauge theory (1.8) at large N. Consider the effect of the Wilsonian renormgroup flow on the original $`YM_D`$ theory in the WC phase at the UV scale (where the large N SC expansion, in terms of the proposed microscopic gauge strings, fails). The idea is that the latter $`YM_D`$ theory, at the sufficiently low-energy scale, may be superseded by such effective strongly coupled $`YM_D`$ system where the (effective) Gauge String representation is already valid. As the effective $`YM_D`$ system is quasilocal, we assert that in the effective Gauge String the Nambu-Goto term in (1.2) is traded for the whole operator expansion (OPE) running in terms of the extrinsic and intrinsic curvatures of the worldsheet. (Complementary, the weights of the movable branch points are modified in such a way that the integration, over the positions of the latter points, results in exactly the same pattern of the worldsheet OPE as for the descedant of the Nambu-Goto term.)
Finally, let us make contact with the two alternative stringy proposals ,. As for , Witten argues that in continuous $`YM_4`$ theory (with a finite UV cut off $`\mathrm{\Lambda }`$) the physical string tension $`\sigma _{ph}`$ in the extreme large N SC limit $`g^2N\mathrm{},g^2O(1/N),`$ scales as $`\sigma _{ph}g^2N\mathrm{\Lambda }^2`$. On the other hand, eq. (1.6) predicts similar large N $`SC`$ $`O(N\stackrel{~}{g}^2)`$-scaling of the bare string tension $`\sigma _0`$ (where $`\stackrel{~}{g}^2\stackrel{~}{g}^2(\{\stackrel{~}{b}_k\},N)O(N^1\mathrm{\Lambda }^0)`$) in the following subclass of the $`YM_D`$ systems. The latter are defined by the subset of the actions (1.4) with the coefficients $`g_r`$ constrained by
$$\begin{array}{c}\stackrel{~}{g}^2N\mathrm{}:[g_r]^1O(\mathrm{\Lambda }^{2n+4}[N^{\frac{1}{2}}\stackrel{~}{g}]^{n+\gamma (r)}N^{2_{k=1}^np_k}),\end{array}$$
(9.1)
where $`\gamma (r)0`$ (and there is at least one irrep $`r`$ for which $`\gamma (r)=0`$). According to the structure of the (abelianized) Born-Infeld action, the Witten’s prescription is supposed to induce the $`YM_D`$ action with the local part belonging to the variety (9.1) (with possible addition of the commutator-terms which does not alter the conclusions). In sum, our prediction (1.6) for the string tension in the extreme large N SC limit semiquantitatively matches with the $`D=4`$ pattern of motivated by the $`AdS/CFT`$ correspondence.
As for the Polyakov’s $`D=4`$ Ansatz , it is aimed at a stringy reformulation of the continuous $`YM_4`$ theory (1.8) in the large N WC regime. Assuming that the general pattern of the Ansatz is applicable to the SC phase as well, an indirect comparison might be possible. In particular, the advocated by Polyakov invariance of the worldsheet action under the extended group of the diffeomorphisms (with the singularities due to the zero Jacobian) matches with the two key features of the Gauge String: the vanishing contributions of the worldsheets $`\stackrel{~}{M}`$ with the backtrackings and the invariance of $`w[\stackrel{~}{M}]`$ under the group of the area-preserving diffeomorphisms specified in the Section 7.
Acknowledgements.
This project was started when the author was the NATO/NSERC Fellow at University of British Columbia, and I would like to thank all the staff and especially Gordon Semenoff for hospitality.
## Appendix A: Reconstruction of $`\mathrm{\Xi }_{4n_+}(\{R_{\mu \nu }\})`$.
To derive the $`S(4n_+)`$ operator $`\mathrm{\Xi }_{4n_+}(\{R_{\mu \nu }\})`$ of eq. (3.3), we first rewrite the characters of the original expression (2.5) in terms of certain $`S(4n_{\mu \nu })`$ operator $`𝐃(\mathrm{\Xi }_{4n_{\mu \nu }}(R_{\mu \nu }))`$
$$\begin{array}{c}\chi _{R_{\mu \nu }}(U_\mu U_\nu U_\mu ^+U_\nu ^+)=Tr_{n_{\mu \nu }}[𝐃(\mathrm{\Xi }_{4n_{\mu \nu }}(R_{\mu \nu }))𝐃_2(\{U_\rho U_\rho ^+\})],\end{array}$$
(A.1)
where $`𝐃_2(\{U_\rho U_\rho ^+\})`$ is the $`D=2`$ option of (3.4). In the $`|\stackrel{~}{I}_{4n(\mu \nu )}>`$-basis (3.22), $`𝐃(\mathrm{\Xi }_{4n_{\mu \nu }}(R_{\mu \nu }))`$ reads explicitly
$$\begin{array}{c}_{\sigma S(n_{\mu \nu })}\frac{\chi _{R_{\mu \nu }}(\sigma )}{n_{\mu \nu }!}\delta _{l_{\sigma (3n+1)}}^{p_1}..\delta _{l_{\sigma (4n)}}^{p_n}\delta _{l_1}^{p_{n+1}}..\delta _{l_n}^{p_{2n}}\delta _{l_{n+1}}^{p_{2n+1}}..\delta _{l_{2n}}^{p_{3n}}\delta _{l_{2n+1}}^{p_{3n+1}}..\delta _{l_{3n}}^{p_{4n}}\end{array}$$
(A.2)
which can be represented in the concise form of the inner-product
$$\begin{array}{c}\mathrm{\Xi }_{4n_{\mu \nu }}(R_{\mu \nu })=\mathrm{\Psi }_{4n_{\mu \nu }}\stackrel{~}{P}_{4n_{\mu \nu }}(R_{\mu \nu });\stackrel{~}{P}_{4n}(R)=\widehat{1}_{[n]}^3C_R,\end{array}$$
(A.3)
where $`\mathrm{\Psi }_{4n_{\mu \nu }}`$ is defined by eq. (3.20). Each of the four (ordered) $`S(n_{\mu \nu })`$-operators in the outer product composition of $`\stackrel{~}{P}_{4n_{\mu \nu }}(R_{\mu \nu })`$ are postulated to act on the corresponding four (ordered) $`S(n_{\mu \nu })`$-subspaces (3.22) of the space $`|\stackrel{~}{I}_{4n(\mu \nu )}>`$. Finally, in the $`S(4n_+)`$ basis (3.17), one evidently obtains $`\mathrm{\Xi }_{4n_+}(\{R_{\mu \nu }\})=_{\{\mu \nu \}}\mathrm{\Xi }_{4n_{\mu \nu }}(R_{\mu \nu })`$ that matches with eq. (3.19) modulo a slight deviation of $`\stackrel{~}{P}_{4n_{\mu \nu }}(R_{\mu \nu })`$ from $`P_{4n_{\mu \nu }}(R_{\mu \nu })`$ (defined by eq. (3.23)).
The possibility to substitute in eq. (3.3) the operator $`\stackrel{~}{P}_{4n_{\mu \nu }}`$ by $`P_{4n_{\mu \nu }}`$ is ensured by the following feature of the pattern (A.1). Owing to the basic commutativity $`[𝐃(\sigma ),U^n]=0`$$`\sigma S(n)`$, any of the unity operators $`\widehat{1}_{[n]}`$ (the operator $`\stackrel{~}{P}_{4n}`$ is composed of) can be substituted by the $`C_R`$-factor properly weighted according to the multiplication rule $`(C_R)^2=C_R/d_R`$. Remark also that in eq. (A.3) the relative order of the factors $`\mathrm{\Psi }_{4n_{\mu \nu }}`$ and $`\stackrel{~}{P}_{4n_{\mu \nu }}`$ is immaterial since $`[\mathrm{\Psi }_{4n},_{k=1}^4\sigma _k]=0,\sigma _kS(n)`$.
Next, let us prove that, in the dual representation (3.24) of (2.5), the substitution (3.27) (of $`_{\rho =1}^D\mathrm{\Phi }_{2n_\rho }`$ by its square) doesn’t change the character (3.5). The latter ’symmetry’ can be traced back to the fact that the multiple integral (2.5), represented by the $`S(4n_+)`$ element (3.24), is a real-valued function. To take advantage of this fact, observe that $`\stackrel{~}{\mathrm{\Phi }}_{2m}=(\mathrm{\Phi }_{2n}\mathrm{\Phi }_{2n}),m=2n,`$ is the operator which represents the complex conjugation of the characters $`\chi _R(U_{\mu \nu })=Tr_{4n}[𝐃(\mathrm{\Xi }_{4n})\stackrel{~}{U}_{\mu \nu }^n]`$ entering (2.5):
$$\begin{array}{c}Tr_{4n_{\mu \nu }}[𝐃(\mathrm{\Xi }_{4n_{\mu \nu }})\stackrel{~}{U}_{\mu \nu }^{n_{\mu \nu }}]^{}=Tr_{4n_{\mu \nu }}[𝐃(\mathrm{\Xi }_{4n_{\mu \nu }}\stackrel{~}{\mathrm{\Phi }}_{2m_{\mu \nu }})\stackrel{~}{U}_{\mu \nu }^{n_{\mu \nu }}],\end{array}$$
(A.4)
where $`\mathrm{\Xi }_{4n_{\mu \nu }}\mathrm{\Xi }_{4n_{\mu \nu }}(R_{\mu \nu })`$ is defined in eq. (A.3), and $`\stackrel{~}{U}_{\mu \nu }`$ is introduced in eq. (3.28). As the master-integral (2.5) is invariant under the simultaneous transformation (A.4) of all the involved characters while $`_{\mu \nu }\stackrel{~}{\mathrm{\Phi }}_{2m_{\mu \nu }}=_{\rho =1}^D\mathrm{\Phi }_{2n_\rho }`$, we arive at the required invariance under (3.27).
As for eq. (A.4), it can be deduced by linearity from the following more elementary identity. To formulate the latter, let us first introduce the representation of the basic traces (the characters (A.4) are composed of)
$$\begin{array}{c}tr((U_\mu U_\nu U_\mu ^+U_\nu ^+)^n)Tr_{4n}[𝐃(\mathrm{\Gamma }_{4n}\mathrm{\Psi }_{4n})\stackrel{~}{U}_{\mu \nu }^n]=Tr_{4n}[𝐃(\mathrm{\Gamma }_{4n}^1\mathrm{\Psi }_{4n})\stackrel{~}{U}_{\mu \nu }^n],\end{array}$$
(A.5)
where $`\mathrm{\Gamma }_{4n}=(c_n\widehat{1}_{[n]}^3)`$ in the $`S(4n)`$ basis (3.22), and $`c_n`$ is the n-cycle permutation. Then, the required identity reads (with $`m=2n`$)
$$\begin{array}{c}Tr_{4n}[𝐃(\mathrm{\Gamma }_{4n}\mathrm{\Psi }_{4n})\stackrel{~}{U}_{\mu \nu }^n]^{}=Tr_{4n}[𝐃(\mathrm{\Gamma }_{4n}\mathrm{\Psi }_{4n}\stackrel{~}{\mathrm{\Phi }}_{2m})\stackrel{~}{U}_{\mu \nu }^n].\end{array}$$
(A.6)
For its justification. we first introduce the tensor representation of the relevant complex conjugation: $`Tr_n[𝐃(\sigma )_{k=1}^nU_k]^{}=Tr_n[𝐃(\sigma ^1)_{k=1}^nU_k^+]`$ where the orderings of the $`U_k`$ factors in the l.h.s. is the same as that of $`U_k^+`$ in the r.h.side. The key-observation is that, when $`_{k=1}^nU_k^+=\stackrel{~}{U}_{\mu \nu }^{n_{\mu \nu }}`$, the substitution $`U_kU_k^+`$ can be performed as the conjugation with respect to $`(\mathrm{\Phi }_{2n_{\mu \nu }}\mathrm{\Phi }_{2n_{\mu \nu }})=\stackrel{~}{\mathrm{\Phi }}_{2m_{\mu \nu }}=\stackrel{~}{\mathrm{\Phi }}_{2m_{\mu \nu }}^1`$: $`\stackrel{~}{\mathrm{\Phi }}_{2m_{\mu \nu }}^1\stackrel{~}{U}_{\mu \nu }^{n_{\mu \nu }}\stackrel{~}{\mathrm{\Phi }}_{2m_{\mu \nu }}=\stackrel{~}{V}_{\mu \nu }^{n_{\mu \nu }}`$ where $`\stackrel{~}{V}_{\mu \nu }=U_\mu ^+U_\mu U_\nu ^+U_\nu `$. This is because, owing to eq. (3.16), $`\mathrm{\Phi }_{2n}`$ interchanges (either upper or lower) indices between the $`U^n`$ and $`(U^+)^n`$ blocks of $`U^n(U^+)^n`$. As in eq. (A.5) $`\mathrm{\Gamma }_{4n}`$ can be substituted by $`\mathrm{\Gamma }_{4n}^1`$ (while $`\mathrm{\Gamma }_{4n}`$ commutes with both $`\stackrel{~}{\mathrm{\Phi }}_{2m}`$ and $`\mathrm{\Psi }_{4n}`$), all what remains to be proved is that $`\mathrm{\Psi }_{4n}\stackrel{~}{\mathrm{\Phi }}_{2m}=\stackrel{~}{\mathrm{\Phi }}_{2m}\mathrm{\Psi }_{4n}=\mathrm{\Psi }_{4n}^1,n=2m`$. The last identity, owing to the specific patterns (3.16) and (3.20) of $`\mathrm{\Phi }_{2n}`$ and $`\mathrm{\Psi }_{4n}`$, directly follows from its reduced $`m=2`$ variant. This completes the justification of (A.6).
Finally, let us sketch the proof of the basis formula (3.31). To make sure that the contraction of the $`\sigma _\rho ^{(\pm )}`$ indices is the same in the both sides of (3.31), the following decomposition of the index-structure is helpful. Building on the prescription (3.6), we rewrite the $`S(n_\rho )`$ operators $`D(\sigma _\rho ^{(\pm )})`$ in the form
$$\begin{array}{c}D(\sigma _\rho ^{(\pm )})_{\{i^{n_\rho }\}}^{\{j^{n_\rho }\}}D(\sigma _\rho ^{(\pm )})_{\{i^{n_{\rho \mu }}\}\mathrm{}\{i^{n_{\rho \nu }}\}}^{\{j^{n_{\rho \mu }}\}\mathrm{}\{j^{n_{\rho \nu }}\}},\mu ,\mathrm{},\nu \rho ,\end{array}$$
(A.7)
where each of the $`D1`$ $`n_{\rho \lambda }`$-subsets of the indices ($`\lambda \rho `$) acts on the associated $`S(n_{\rho \lambda })`$ subspace (3.32) of the enveloping space $`S(n_+)`$. Let consider any $`n_{\rho \nu }`$-subset, associated to $`\sigma _\rho ^{(\pm )}`$, as the composite block-index $`j_{\rho \nu }^{(\pm )}`$. A direct inspection then reveals that, in both sides of eq. (3.31), for a given $`\mu ,\nu `$ the four block-indices $`j_{\mu \nu }^{(\pm )},j_{\nu \mu }^{(\pm )}`$ (corresponding to either $`\rho =\mu `$ or $`\rho =\nu `$) are $`c_4`$-cyclically contracted according to the pattern (3.19)/(3.20) of $`_{\mu \nu }\mathrm{\Psi }_{\mu \nu }`$.
As for the ordering of the inner $`\rho `$-products in the r.h.s. of (3.31), it should be deduced from the particular ordering of the $`|i_\pm (\rho )>`$-blocks composed into the elementary subspace (3.21) used to define the $`D(D1)/2`$ operators $`\mathrm{\Psi }_{4n_{\mu \nu }}`$. To justify the prescription stated in Section 3, one is to combine the above analysis with the specified pattern of the ordering in the elementary $`D=2`$ case (3.29).
## Appendix B: Tensor representation vs. Regular one.
Let us first derive the identity (2.8). Actually, this equation is nothing but the transformation of the trace of the tensor $`𝐃(\sigma )`$-representation (2.7) into that of the canonical regular representation $`X_{REG}`$ (both associated to a given $`S(n)`$-algebra). Recall that $`X_{REG}`$ is defined on the vector space $`\mathrm{\Theta }=\{\sigma _i;\sigma _iS(n)\}`$ by the homomorphism $`S(n)M_iEnd(\mathrm{\Theta })`$ of the $`S(n)`$-group into the group of the endomorphisms of $`\mathrm{\Theta }`$: $`\sigma _i\sigma _j=\sigma _k\mathrm{\Delta }_{ij}^k(M(\sigma _i))_j^k(M_i)_j^k=\mathrm{\Delta }_{ij}^k`$. Here $`\sigma _i\sigma _j=\sigma _q`$, and $`\mathrm{\Delta }_{ij}^k=1`$ or $`0`$ depending on whether $`q=k`$ or $`qk`$. Defined in this way, the matrices $`M_i`$ satisfy the same relations $`M_iM_j=M_q`$ as the origibal $`S(n)`$ group-elements $`\sigma _i`$.
To make contact with the tensor representation (2.7), recall first that
$$\begin{array}{c}\chi _{REG}(\sigma )=_{RY_n}d_R\chi _R(\sigma )=n!\delta _n(\sigma ),\chi _{REG}(P_R\sigma )=d_R\chi _R(\sigma ),\end{array}$$
(B.1)
where $`tr[M(\sigma )]\chi _{REG}(\sigma )`$ and $`P_R=d_RC_R`$. As for $`\delta _n(..)`$, on the $`S(n)`$-group it reduces to the standard Kronecker $`\delta `$-function: $`\delta _n(\sigma )=\delta [\sigma ,\widehat{1}_{[n]}]`$ (with $`\widehat{1}_{[n]}`$ being the ’trivial’ unity-permutation of $`S(n)`$). By linearity, it is then extended to the $`S(n)`$-algebra. Finally, rewriting the $`V=\widehat{1}`$ reduction of the second Frobenius formula (see e.g. ) in terms of $`\chi _{REG}(C_R\sigma )`$
$$\begin{array}{c}Tr_n[𝐃(\sigma )]=_{RY_n^{(N)}}\chi _R(\sigma )\chi _R(V)|_{V=\widehat{1}}=_{RY_n^{(N)}}dimR\chi _{REG}(C_R\sigma ),\end{array}$$
(B.2)
and employing the definition (3.9) of $`\mathrm{\Lambda }_n^{(1)}`$ (where the sum runs over the chiral U(N) irreps R), we arrive at the basic identity (2.8). (Omitting the $`P_n^{(N)}`$ projector, the latter identity was discussed in .)
In conclusion, let us briefly sketch the derivation of eqs. (4.1), (4.4). As for (4.1), in its l.h.s. we use first (4.3) to substitute $`d_R(n!dimR/d_R)^m`$ by $`\chi _R((N^n\mathrm{\Omega }_n)^m)`$. Replacing $`C_R=_{\sigma S(n)}\chi _R(\sigma ^1)\sigma /n!`$, we then combine the two resulting characters into one
$$\begin{array}{c}\frac{1}{n!}_{\sigma S(n)}_{RY_n}d_R\chi _R((N^n\mathrm{\Omega }_n)^mP_n^{(N)}\sigma ^1)𝐃(\sigma ).\end{array}$$
(B.3)
As $`[\mathrm{\Omega }_n,\sigma ]=0,\sigma S(n)`$, for the derivation of (B.3) we used the identity $`\chi _R(\mathrm{\Psi }\sigma )=\chi _R(\sigma )\chi _R(\mathrm{\Psi })/d_R,\rho S(n)`$, if $`[\mathrm{\Psi },\rho ]=0`$ (which directly follows from the possibility to expand any such $`\mathrm{\Psi }`$, in the center of $`S(n)`$, in terms of the Young projectors $`\mathrm{\Psi }=_{RY_n}\psi _RP_R`$). We have used also that $`\chi _R(P_n^{(N)}\sigma )=\chi _R(\sigma )`$ or $`0`$ depending on whether or not $`RY_n^{(N)}`$. Applying the completeness condition $`_{\sigma S(n)}\delta _n(\sigma ^1\mathrm{\Phi })𝐃(\sigma )=𝐃(\mathrm{\Phi })`$ (where $`\mathrm{\Phi }`$ is any element of the $`S(n)`$-algebra) we arrive at the announced result (4.1).
Concerning eq. (4.4), first one is to represent (see e.g. ): $`C_2(R)=(nN+2\chi _R(\widehat{T}_2^{(n)})/d_rn^2/N),RY_n,`$ which is a particular case of the general relation (E.1). Expanding $`\chi _R(\widehat{T}_2^{(n)})/d_r`$ (where $`\widehat{T}_2^{(n)}`$ is defined by eq. (4.5)) into the preexponent, the derivation of (4.4) is given by a minor modification of the steps developed in the context of (4.1).
## Appendix C: Reintroducing the branch points.
Given the generalized covering spaces (GCSs) $`\stackrel{~}{M}(\{\sigma _\rho \})`$ of $`T_{EK}`$ (resulting after the deformation (6.3) of (4.8)), the generalized branched covering spaces (GBCSs) of $`T_{EK}`$ (encoded in the full expression (6.2)) can be reconstructed reintroducing onto $`\stackrel{~}{M}(\{\sigma _\rho \})`$ the relevant branch points (BPs). Alternatively, the GBCSs of $`T_{EK}`$ are reinterpreted as the GCSs of $`T_{EK}\{q_s\}`$ where the deleted (from $`T_{EK}`$) set of points $`\{q_s\}`$ is associated to the corresponding BP permutations $`\xi ^{\{p_\varphi \}},\tau _{\mu \nu }^{(k_{\mu \nu })}`$ in compliance with the homomorphism (6.6). Then, a particular GBCS of $`T_{EK}`$ is reproduced from the associated GBC (of $`T_{EK}\{q_s\}`$) ’closing’ the temporarily deleted points $`\{\phi ^1(q_s)\}`$. At a given $`q_s`$, the closure is performed as the mapping of the $`n_{\varphi (s)}`$-set $`\phi ^1(q_s)\mathrm{{\rm Y}}_{\varphi (s)}`$ onto set of points matching with the number of cycles in the cyclic decomposition of the associated to $`q_s`$ permutation (6.6).
To make the above procedure well-defined, I propose to view $`T_{EK}`$ as the homotopy retract of a ’less singular’ $`2d`$ cell-complex $`T_{EK}^{}`$ with the same fundamental group: $`\pi _1(T_{EK}^{})=\pi _1(T_{EK})`$. As the proposed retraction preserves the basic homomorphism (6.6), the GBCSs of $`T_{EK}`$ can be consistently treated as the limiting case of the corresponding GBCSs of $`T_{EK}^{}`$.
First, let us thicken each $`\rho `$-link into an infinitesimally thing cylinder $`\overline{Z}_\rho =_{\nu \rho }E_{\rho \nu }`$ shared by the $`D1`$ corresponding 2-tora $`E_{\rho \nu }`$. (To match with the canonical construction of the branch points, (for a given $`\rho `$) all $`E_{\rho \nu },\nu \rho ,`$ can be always adjusted to have the same orientation when restricted to $`\overline{Z}_\rho `$.) Complementary, the intersection $`\overline{Z}_+=_\rho \overline{Z}_\rho =_{\rho \nu }E_{\rho \nu }`$ of the $`D`$ different cylinders $`\overline{Z}_\rho `$ (or, equivalently, of all the $`D(D1)/2`$ $`\mu \nu `$-tora $`E_{\mu \nu }`$) is thicken into an infinitesimal disc $`\overline{Z}_+`$. (Again, all the cylinders $`\overline{Z}_\rho `$ can be adjusted to have the same orientation when restricted to $`\overline{Z}_+`$.)
Choose a base-point $`pq_s,s,`$ (common for all the equivalence classes of the pathes represented in the constraint (6.5)) in the interior $`Z_+`$ of $`\overline{Z}_+`$. We require that the location of the set $`\{q_s\}=_\varphi \{q_{k_\varphi }^{(\varphi )}\}`$ complies with
$$\begin{array}{c}q_{k_{\mu \nu }}^{(\mu \nu )}(E_{\mu \nu }(\overline{Z}_\mu \overline{Z}_\nu ));q_{k_\rho }^{(\rho )}(Z_\rho \overline{Z}_+);q_{k_+}^{(+)}Z_+,\end{array}$$
(C.1)
where $`Z_\varphi `$ is the interior of the closed space $`\overline{Z}_\varphi ,\varphi \{\mu \nu \},\{\rho \},+`$. Also, the introduced in Section 6 closed branch cuts $`\varpi _\rho `$ (inherited from the GCS of $`T_{EK}^{}`$) are supposed to satisfy: $`\varpi _\rho Z_\rho `$, while $`\{q_{k_\varphi }^{(\varphi )}\}\varpi _\rho =0,\rho ,k_\varphi `$.
To reintroduce the branch points onto the GCSs $`\stackrel{~}{M}(\{\sigma _\rho \})`$ of $`T_{EK}^{}`$, one is to consider the additional (to $`\varpi _\rho `$) branch cuts $`\varpi _{\varphi (s)}`$ outgoing from the corresponding points $`q_s`$ of $`T_{EK}^{}`$. (Combining all the cuts together, one is supposed to arrive at a network $`\mathrm{\Omega }=\{\varpi _{\varphi (s)},\varpi _\rho \}`$ whose global consistency is ensured by the $`\delta _{n_+}`$-constraint (6.2).) To implement $`\varpi _{\varphi (s)}`$, we cut the GCS $`\stackrel{~}{M}(\{\sigma _\rho \})`$ along the $`\phi ^1`$-image of $`\varpi _{\varphi (s)}`$ (starting at $`q_s`$ and terminating either at some other point $`q_k`$ or, possibly, at an auxiliary vertex of the network $`\mathrm{\Omega }`$). Denote $`\alpha (s_k)`$ and $`\beta (s_k)`$ (with $`k=1,\mathrm{},n_{\varphi (s)}`$) the resulting $`n_{\varphi (s)}`$ boundaries of the sheets (of $`\stackrel{~}{M}(\{\sigma _\rho \})`$) respectively on the left and on the right sides of $`\varpi _{\varphi (s)}`$. To obtain the GCS of $`T_{EK}^{}\{q_s\}`$ corresponding to the homomorphism (6.6), one is to perform the pairwise identifications of these boundaries according to the developed prescription (5.9). Matching with (6.2), one simply substitutes in (5.9): $`\rho s`$. so that $`\sigma _s`$ is either $`\xi ^{\{p_\varphi \}}`$ or $`\tau _{\mu \nu }^{(k_{\mu \nu })}`$ depending on the image of the homomorphism (6.6) assigned to a given branch point $`q_s`$ constrained by (C.1).
## Appendix D: Counting the generalized coverings.
The proof, that inequivalent spaces $`\stackrel{~}{M}_\phi `$ in (6.2) are parametrized by the equivalence classes (6.7) of the homomorphisms (6.6), appropriately generalizes the analogous proof for the canonical coverings employed in . First, one is to demonstrate that the relevant inequivalent spaces $`\stackrel{~}{M}_\phi `$ entering (6.2) are uniquely parametrized by certain $`_{\mu \nu }S(n_{\mu \nu })`$ conjugacy classes of the $`\pi _1(T_{EK}^{}\{q_s\}|p)`$ subgroups. Then, one proves the one-to-one correspondence between the latter classes of the subgroups and the equivalence classes (6.7) of the homomorphisms (6.6). Let us outline the above proofs with the emphasis on the subtleties novel compared to .
To begin with, we observe that (by the same token as in the canonical case) any subgroup of $`\pi _1(T_{EK}^{}\{q_s\}|p)`$ can be viewed as the image
$$\begin{array}{c}\phi _{}(\pi _1(\stackrel{~}{M}_\phi \{\phi ^1(q_s)\}|\stackrel{~}{p})),\stackrel{~}{p}\phi ^1(p)\mathrm{{\rm Y}}_{n_+},\end{array}$$
(D.1)
induced by the $`\phi `$-mapping (1.11): $`\stackrel{~}{M}_\phi T_{EK}^{}`$, reconstructed in Section 6.1 and Appendix C. The consistency of (D.1) implies that a given $`\stackrel{~}{M}_\phi `$ yields some branched covering (with $`n_+`$-sheets) ’above’ the single site s of $`T_{EK}`$ which is ’regularized’, $`sZ_+`$ (see eq. (C.1) of Appendix C), into an infinitesimal disc $`Z_+`$ of $`T_{EK}^{}`$ so that $`pZ_+`$. (In contradistinction to the canonical case, the point p is not allowed to leave $`Z_+`$.) Perform a shift of the base-point $`\stackrel{~}{p}\phi ^1(p)`$, located on a $`j`$th sheet of the branched covering of $`Z_+`$, to $`\stackrel{~}{p}^{}\phi ^1(p)`$ on some other $`k`$th sheet along the path $`\epsilon \stackrel{~}{M}_\phi \{\phi ^1(q_s)\}`$. It induces the associated conjugation of the $`\pi _1(T_{EK}^{}\{q_s\}|p)`$ elements (and, hence, its subgroups) with respect to the element given by the image $`\phi (\epsilon )\pi _1(T_{EK}^{}\{q_s\}|p)`$ of the path $`\epsilon `$.
Next, one notes that (similarly to the canonical construction ) the shift $`\epsilon `$ acts on the $`n_+`$-set $`\phi ^1(p)\mathrm{{\rm Y}}_{n_+}`$ as the simple $`S(n_+)`$-transposition $`\psi (\phi (\epsilon ))T_2^{(n_+)}`$ of the two corresponding entries: $`jk,j,k,=1,\mathrm{},n_+`$. As a result, $`\epsilon `$ induces the $`S(n_+)`$ conjugation (6.7) of (6.6) with $`\eta =\psi (\phi (\epsilon ))`$ that can be visualized as the permutation of the two associated sheets of $`\stackrel{~}{M}_\phi `$. The latter does not necessarily leaves intact the topology of $`\stackrel{~}{M}_\phi `$. The cutting-gluing technique of Section 6.1 implies that interchanges of the sheets encoded in the conjugations with $`\eta S(n_+)/_{\mu \nu }S(n_{\mu \nu })`$ (i.e. between the $`n_{\mu \nu }`$-sheet coverings of different $`\mu \nu `$-tora $`E_{\mu \nu }`$ of $`T_{EK}^{}`$) should be excluded when collecting an equivalence class of maps (6.7) corresponding to a given topology of $`\stackrel{~}{M}_\phi `$.
Consider the $`_{\mu \nu }S(n_{\mu \nu })`$ conjugacy classes of the $`\pi _1(T_{EK}^{}\{q_s\}|p)`$ subgroups induced by the combination of the elements $`\phi _{}(\epsilon )`$ which are the images of those of the shifts $`\epsilon `$ that connect the sheets within the $`n_{\mu \nu }`$-sheet covering of a given torus $`E_{\mu \nu }`$ of $`T_{EK}^{}`$. Taking into account the above discussion, one readily modifies the canonical construction to prove that the considered classes of the subgroups are indeed in the 1-to-1 correspondence with the inequivalent generalized covering spaces $`\stackrel{~}{M}_\phi \{\phi ^1(q_s)\}`$ of $`T_{EK}^{}\{q_s\}`$ (and, therefore, with inequivalent GBCSs $`\stackrel{~}{M}_\phi `$ of $`T_{EK}^{}`$). In particular, the cutting-gluing rules of Section 6.1 ensure that the mapping from $`\stackrel{~}{M}_\phi `$ to the corresponding classes of the $`\pi _1(T_{EK}^{}\{q_s\}|p)`$ subgroups is onto.
Finally, let us reveal the second 1-to-1 correspondence of the former classes of the subgroups with the equivalence classes (6.7) of the homomorphisms (6.6). From above, it is clear that the $`\psi `$-images (6.6) of these $`\pi _1(T_{EK}^{}\{q_s\}|p)`$ subgroups can be related by the $`_{\mu \nu }S(n_{\mu \nu })`$ permutations corresponding to the interchanges of the sheets separately within each of the $`D(D1)/2`$ $`n_{\mu \nu }`$-sheet coverings associated to particular $`E_{\mu \nu }`$-tora of $`T_{EK}^{}`$. As a result, akin to the $`D=2`$ case , the latter permutations are represented by the required conjugations (6.7) with $`\eta _{\mu \nu }S(n_{\mu \nu })`$. Conversly, equivalent homomorphisms determine equivalent $`_{\mu \nu }S(n_{\mu \nu })`$ classes of the $`\pi _1(T_{EK}^{}\{q_s\}|p)`$ subgroups (according to the discussed above construction of $`\eta =\psi (\phi (\epsilon ))`$). Summarizing, it completes the proof of the two asserted correspondences.
## Appendix E: The higher Casimirs’ actions.
To derive the generalized form (4.6),(4.7) of $`Q_n(\mathrm{\Gamma })`$ in (4.4), first one is to begin with the (Schur-Weyl) duality between the Casimir eigenvalues
$$\begin{array}{c}\frac{C_q(R)}{N^{q1}}=_{T_{\{p\}}S(n)}a_q(N,n,T_{\{p\}})\frac{\chi _R(\widehat{T}_{\{p\}})}{d_R},RY_n^{(N)},\end{array}$$
(E.1)
and the symmetric group characters $`\{\chi _R(\widehat{T}_{\{p\}})\}`$ (where $`\widehat{T}_{\{p\}}^{(n)}\widehat{T}_{\{p\}}`$ is defined by eq. (4.6), while $`1qN`$). The coefficients $`a_q(\mathrm{})`$ are defined by the formular which can be deduced from eq. (4.7) via the substitution: $`\upsilon _{\{p\}}(..,n,N)a_q(N,n,T_{\{p\}}),s_{\{p\}}(\{\stackrel{~}{b}_r\},m,l)s_q(T_{\{p\}},m,l),M_{\{p\}}M_0(T_{\{p\}})`$. In particular, (keeping $`nO(N^0)`$) the leading term of the formal $`1/N`$ expansion of the $`p`$th order Casimir operator reads: $`C_q(R)=N^{q1}(n+O(N^1)),RY_n^{(N)}`$, which corresponds in (E.1) to the trivial unity permutation $`T_{\{p\}}=\widehat{1}_{[n]}`$. As for the remaining $`T_{\{p\}}\widehat{1}_{[n]}`$ contributions, the branch point interpretation of the $`l=0`$ term in expression like (4.7) is discussed in the end of Section 5 (in fact $`s_q(T_{\{p\}},m,0)=0`$ for $`m2`$). To interprete a given $`l1`$ term, first one is to resum the powers $`n^m`$ in terms of the numbers of inequivalent subdivisions of n objects into two subsets containing $`t`$ and $`(nt)`$ objects: $`n^m=_{t=1}^mf_{m,n}n!/t!(nt)!`$. As a result, similarly to the latter contributions include (in addition to the $`T_{\{p\}}`$ branch point) the extra collapsed to a point $`2d`$ subsurfaces. Being cut out, these surfaces can be viewed as having genus $`g=l[t/2]`$ and $`2[t/2]`$ holes.
Finally, to select the admissible class of the functions $`\mathrm{\Gamma }`$ (defining a given model (1.1)), one is to require that in the associated large $`N`$ SC series (1.9) the factor $`N^{22h}`$ matches with the genus $`h`$ of the associated worldsheet $`\stackrel{~}{M}`$. Then, to ensure eq. (4.7), the admissible polynomial pattern of the function $`\mathrm{\Gamma }`$ should belong to the variety (with $`M𝐙_\mathrm{𝟏}`$)
$$\begin{array}{c}\mathrm{\Gamma }(\{\stackrel{~}{b}_r\},N,\{C_q(R)\})=_{M,\{q_k\}}_{\overline{l}[\frac{M}{2}]}s(\{\stackrel{~}{b}_r\},\{q_k\},\overline{l})N^{2\overline{l}}_{k=1}^M\frac{C_{q_k}(R)}{N^{q_k1}},\end{array}$$
(E.2)
where the (properly weighted by the 1/N factors) $`\overline{l}1`$ terms encode the additional, compared to those encoded in (E.1), collapsed subsurfaces connecting a few sheets according to the same rules as for (E.1). In particular, the pattern (E.2)/(E.1) ensures the existence of the proper ’asymptotics’ (1.6) of $`\mathrm{\Gamma }(\mathrm{})`$ defining the bare string tension $`\sigma _0(\{\stackrel{~}{b}_r\})`$ (to which only linear in $`C_q(R)`$ terms in (E.2) contribute).
Finally, with the help of some elementary identities from the theory of $`\chi _R(\widehat{T}_{\{p\}})`$ characters, the generic function (E.2) results in the generic pattern (4.6) of the generalized operator $`Q_n(\mathrm{\Gamma })`$. One can argue that the required pattern (4.7) of the N-dependence takes place provided, in the standard $`F_{\mu \nu }`$-representation (1.4), the $`1/N`$ scaling of the coefficients $`g_r`$ is constrained to yield the conventional ’t Hooft pattern of the $`1/N`$ topological expansion. |
warning/0002/math0002141.html | ar5iv | text | # Ergodicity of the action of the positive rationals on the group of finite adeles and the Bost-Connes phase transition theorem
## Abstract
For each $`\beta (0,+\mathrm{})`$ there exists a canonical measure $`\mu _\beta `$ on the ring $`𝒜_f`$ of finite adeles. We show that $`_+^{}`$ acts ergodically on $`(𝒜_f,\mu _\beta )`$ for $`\beta (0,1]`$, and then deduce from this the uniqueness of KMS<sub>β</sub>-states for the Bost-Connes system.
Bost and Connes \[BC\] constructed a remarkable C-dynamical system which has a phase transition with spontaneous symmetry breaking involving an action of the Galois group Gal($`^{ab}/`$), and whose partition function is the Riemann $`\zeta `$ function. In their original definition the underlying algebra arises as the Hecke algebra associated with an inclusion of certain $`ax+b`$ groups. Recently Laca and Raeburn \[LR, L2\] have realized the Bost-Connes algebra as a full corner of the crossed product algebra $`C_0(𝒜_f)_+^{}`$. This new look at the system has allowed to simplify significantly the proof of the existence of KMS-states for all temperatures, and the proof of the phase transition theorem for $`\beta >1`$ \[L1\]. On the other hand, for $`\beta 1`$ the uniqueness of KMS<sub>β</sub>-states implies the ergodicity of the action of $`_+^{}`$ on $`𝒜_f`$ for certain measures (in particular, for the Haar measure). The aim of this note is to give a direct proof of the ergodicity, and then to show that the uniqueness of KMS<sub>β</sub>-states easily follows from it. Though the proof of the Bost-Connes phase transition theorem (for $`\beta 1`$) thus obtained differs from the proofs given in \[BC\] and \[L1\], it is entirely based on these papers. In particular, the key point is an application of Dirichlet’s theorem.
So let $`𝒫`$ be the set of prime numbers, $`𝒜_f`$ the restricted product of the fields $`_p`$, $`p𝒫`$, of $`p`$-adic numbers, $`=_p_p`$ its maximal compact subring, $`W=^{}=_p_p^{}`$. The group $`_+^{}`$ of positive rationals is embedded diagonally into $`𝒜_f`$, and so acts by multiplications on the additive group of finite adeles. Then the Bost-Connes algebra $`𝒞_{}`$ is the full corner of $`C_0(𝒜_f)_+^{}`$ determined by the characteristic function of $``$ \[L2\]. The dynamics $`\sigma _t`$ is defined as follows \[L1\]: it is trivial on $`C_0(𝒜_f)`$, and $`\sigma _t(\lambda (q))=q^{it}\lambda (q)`$, where $`\lambda (q)`$ is the multiplier of $`C_0(𝒜_f)_+^{}`$ corresponding to $`q_+^{}`$. Then (\[L1\]) there is a one-to-one correspondence between $`(\beta ,\sigma _t)`$-KMS-states on $`𝒞_{}`$ and measures $`\mu `$ on $`𝒜_f`$ such that
$$\mu ()=1\text{and}q_{}\mu =q^\beta \mu \text{for all}q_+^{}(\text{i.e.,}\mu (q^1X)=q^\beta \mu (X)).$$
$`(1\beta )`$
Namely, the KMS-state corresponding to $`\mu `$ is the restriction of the dual weight on $`C_0(𝒜_f)_+^{}`$ to $`𝒞_{}`$.
Note that if $`\beta >1`$ and $`\mu `$ is a measure with the property (1$`\beta `$) then $`\mu (W)=\frac{1}{\zeta (\beta )}>0`$, since $`W=\backslash _pp`$. Moreover, the sets $`qW`$, $`q_+^{}`$, are disjoint, and their union is a set of full measure (since $`_n\mu (nW)=1`$). Thus there exists a one-to-one correspondence between probability measures on $`W`$ and measures on $`𝒜_f`$ satisfying (1$`\beta `$\[L1\]. On the other hand, if $`\beta 1`$ then $`\mu (W)=0`$.
For each $`\beta (0,+\mathrm{})`$ there is a unique $`W`$-invariant measure $`\mu _\beta `$ satisfying (1$`\beta `$\[BC, L1\]. Explicitly, $`\mu _\beta =_p\mu _{\beta ,p}`$, where $`\mu _{\beta ,p}`$ is the measure on $`_p`$ such that $`\mu _{1,p}`$ is the Haar measure ($`\mu _{1,p}(_p)=1`$), and
$$\frac{d\mu _{\beta ,p}}{d\mu _{1,p}}(a)=\frac{1p^\beta }{1p^1}|a|_p^{\beta 1}\text{for}a_p.$$
In fact, for the proof of Proposition below we will only need to know that the restriction of $`\mu _{\beta ,p}`$ to $`_p^{}`$ is a (non-normalized) Haar measure.
Proposition. The action of $`_+^{}`$ on $`(𝒜_f,\mu _\beta )`$ is ergodic for $`\beta (0,1]`$.
Proof. Consider the space $`L^2(,d\mu _\beta )`$ and the subspace $`H`$ of it consisting of the functions that are constant on $``$-orbits. In other words, $`H=\{fL^2(,d\mu _\beta )|V_nf=f,n\}`$, where $`(V_nf)(x)=f(nx)`$. Since any $`_+^{}`$-invariant subset of $`𝒜_f`$ is completely determined by its intersection with $``$, it suffices to prove that $`H`$ consists of constant functions. For this we will compute the action of the projection $`P:L^2(,d\mu _\beta )H`$ on a basis of $`L^2(_{pB}_p,_{pB}\mu _{\beta ,p})`$ (considered as a subspace of $`L^2(,d\mu _\beta )`$) for each finite subset $`B`$ of $`𝒫`$.
Let $`\chi `$ be a character of $`_{pB}_p^{}`$. Consider $`\chi `$ first as a function on $`_{pB}_p`$ by letting $`\chi =0`$ outside of $`_{pB}_p^{}`$. Then using the projection $`_{pB}_p`$, consider $`\chi `$ as a function on $``$. Let $`_B`$ be the unital multiplicative subsemigroup of $``$ generated by $`pB`$. Note that the sets $`n_{pB}_p^{}`$, $`n_B`$, are disjoint, their union is a subset of $`_{pB}_p`$ of full measure, and the operator $`n^{\beta /2}V_n^{}`$ maps isometrically $`L^2(_{pB}_p^{},_{pB}\mu _{\beta ,p})`$ onto $`L^2(n_{pB}_p^{},_{pB}\mu _{\beta ,p})`$ for any $`n_B`$. Hence the functions $`V_n^{}\chi `$, $`n_B`$, $`\chi (_{pB}_p^{})\widehat{}`$, form an orthogonal basis for $`L^2(_{pB}_p,_{pB}\mu _{\beta ,p})`$. So we have to compute $`PV_n^{}\chi `$. But if $`gH`$ then $`(V_n^{}\chi ,g)=(\chi ,g)`$, whence $`PV_n^{}\chi =P\chi `$. Thus we have only to compute $`P\chi `$.
For a finite subset $`A`$ of $`𝒫`$, let $`H_A`$ be the subspace consisting of the functions that are constant on $`_A`$-orbits, $`P_A`$ the projection onto $`H_A`$. Then $`P_AP`$ as $`A𝒫`$. Set
$$W_A=\underset{pA}{}_p^{}\times \underset{q𝒫\backslash A}{}_q.$$
Note, as above, that $`_{n_A}nW_A`$ is a subset of $``$ of full measure. We state that
$$P_Af|_{_Ax}\frac{1}{\zeta _A(\beta )}\underset{n_A}{}n^\beta f(nx)\text{for}xW_A,$$
$`(2)`$
where $`\zeta _A(\beta )=_{n_A}n^\beta =_{pA}(1p^\beta )^1`$. Indeed, denoting the right hand part of (2) by $`f_A`$, for $`gH_A`$ we obtain
$`(f_A,g)`$ $`=`$ $`{\displaystyle \underset{n_A}{}}{\displaystyle _{nW_A}}f_A(x)\overline{g(x)}𝑑\mu _\beta (x)={\displaystyle \underset{n_A}{}}n^\beta {\displaystyle _{W_A}}f_A(x)\overline{g(x)}𝑑\mu _\beta (x)`$
$`=`$ $`\zeta _A(\beta ){\displaystyle _{W_A}}f_A(x)\overline{g(x)}𝑑\mu _\beta (x)={\displaystyle \underset{n_A}{}}n^\beta {\displaystyle _{W_A}}f(nx)\overline{g(x)}𝑑\mu _\beta (x)`$
$`=`$ $`{\displaystyle \underset{n_A}{}}{\displaystyle _{nW_A}}f(x)\overline{g(x)}𝑑\mu _\beta (x)=(f,g).`$
Returning to the computation of $`P\chi `$, we see that
$$P_A\chi |_{_Ax}\frac{\chi (x)}{\zeta _A(\beta )}\underset{n_A}{}n^\beta \chi (n)=\chi (x)\underset{pA}{}\frac{1p^\beta }{1\chi (p)p^\beta }\text{for}xW_A.$$
Thus if $`\chi `$ is trivial then $`P_A\chi _{pB}(1p^\beta )`$ for all $`AB`$, hence $`P\chi `$ is a constant. If $`\chi `$ is non-trivial then since $`P_A\chi _{\mathrm{}}1`$ and the product $`_{p:Re\chi (p)<0}(1p^\beta )`$ diverges by Dirichlet’s theorem \[S\], we have $`P\chi =0`$.
Corollary.\[BC\] For $`\beta (0,1]`$ there exists a unique $`(\beta ,\sigma _t)`$-KMS state on $`𝒞_{}`$.
Proof. Let $`\varphi _\beta `$ be the KMS<sub>β</sub>-state corresponding to $`\mu _\beta `$. Since $`L^{\mathrm{}}(𝒜_f,d\mu _\beta )_+^{}`$ is a factor by Proposition, and $`\pi _{\varphi _\beta }(𝒞_{})^{\prime \prime }`$ is its reduction, $`\varphi _\beta `$ is a factor state. This and the discussion before Proposition show that
* (i) $`\varphi _\beta `$ is an extremal KMS<sub>β</sub>-state;
* (ii) $`\varphi _\beta `$ is a unique $`W`$-invariant KMS<sub>β</sub>-state.
Now the proof is finished as in \[BC, Theorem 25\]:
If $`\psi `$ is an extremal KMS<sub>β</sub>-state then $`_Ww_{}\psi 𝑑w=\varphi _\beta `$. Since KMS<sub>β</sub>-states form a simplex, we conclude that $`\psi =\varphi _\beta `$.
Remarks.
(i) The expression for $`P\chi `$ in the proof of Proposition shows that the divergence of the product
$$\underset{p𝒫}{}\left|\frac{1p^\beta }{1\chi (p)p^\beta }\right|$$
for non-trivial $`\chi `$ is a necessary condition for the ergodicity (otherwise $`P\chi `$ would be a non-zero function, which can not be constant since $`_{}P\chi 𝑑\mu _\beta =_{}\chi 𝑑\mu _\beta =0`$), hence for the uniqueness of KMS<sub>β</sub>-states. So the appearance of (some form of) Dirichlet’s theorem in the proofs is not an accident.
(ii) By \[BC, Theorem 5\] $`\pi _{\varphi _\beta }(𝒞_{})^{\prime \prime }`$ is a factor of type III<sub>1</sub> for $`\beta (0,1]`$. Then the factor $`L^{\mathrm{}}(𝒜_f,d\mu _\beta )_+^{}`$ is also of type III<sub>1</sub>. Hence its smooth flow of weights is trivial, that means that the action of $`_+^{}`$ on $`(_+\times 𝒜_f,dtd\mu _\beta )`$ is ergodic \[CT\]. In particular, the spectral subspaces of $`L^{\mathrm{}}(𝒜_f,d\mu _\beta )`$ corresponding to the characters $`qq^{it}`$ of $`_+^{}`$ have to be trivial for all $`t0`$. But the projection $`P_t`$ onto the subspace $`\{f|V_nf=n^{it}f\}`$ of $`L^2(,d\mu _\beta )`$ is computed with the same ease as in the proof of Proposition:
$$P_t=s\underset{A𝒫}{lim}P_{t,A},(P_{t,A}f)(mx)=\frac{m^{it}}{\zeta _A(\beta )}\underset{n_A}{}n^{\beta it}f(nx)\text{for}xW_A,m_A.$$
Thus the product
$$\underset{p𝒫}{}\left|\frac{1p^\beta }{1\chi (p)p^{\beta it}}\right|$$
has to be divergent for all $`t0`$ and all number characters $`\chi `$ modulo $`m`$.
Institute for Low Temperature Physics & Engineering
Lenin Ave 47
Kharkov 310164, Ukraine
neshveyev@ilt.kharkov.ua |
warning/0002/astro-ph0002024.html | ar5iv | text | # ISO-LWS spectroscopy of Centaurus A: extended star formation
## 1 Introduction
Centaurus A (NGC 5128) is the nearest (d = 3.5 Mpc; 1 <sup>′′</sup>$``$17pc, Hui et al. 1993) example of a giant elliptical galaxy associated with a powerful radio source. The large-scale radio morphology consists of twin radio lobes separated by $``$ 5 degrees on the sky. The compact ($``$ milliarcsecond) radio nucleus is variable and has a strong jet extending $``$ 4 arcminutes towards the northeast lobe. The spectacular optical appearance is that of a giant elliptical galaxy that appears enveloped in a nearly edge on, warped dust lane. There is also a series of faint optical shells. The stellar population in the dominant elliptical structure is old, whilst that of the twisted dust lane is young, sporadically punctuated by HII regions, dust and gas (Graham 1979). The overall structure of Cen A resembles that of a recent ($`<4\times 10^8`$ years, Tubbs 1980) merger, between a spiral and a large elliptical galaxy. The dust lane is the source of most (90 %) of the far-infrared luminosity (L$`{}_{\mathrm{FIR}}{}^{}3\times 10^9`$ L) and is thought to be re-radiated starlight from young stars in the dusty disk (Joy et al. 1988).
In Sect. 2 we describe the observations and data analysis. Sect. 3 looks at the general FIR properties and proceeds to model the HII regions and the PDRs in the dust lane. Sect. 4 summarises the results and presents our conclusions.
## 2 Observations
Cen A was observed with the LWS grating ($`R=\lambda /\mathrm{\Delta }\lambda 200`$) as part of the LWS consortium’s guaranteed time extragalactic programme. A full grating observation (43 - 196.7 $`\mu `$m) was taken of the nucleus at the centre of the dust lane and a series of line observations were taken at two positions in the SE and NW regions of the dust lane. A short \[Cii\] 157 $`\mu `$m line observation was taken off-source at position #4 (see Table 1) to estimate the Galactic emission near the source. Position #1 was intended to provide a deeper integration coincident with position #2, but was accidently offset.
A series of half-second integration ramps were taken at each grating position with four samples per resolution element ($`\mathrm{\Delta }\lambda =0.29\mu `$m $`\lambda \lambda 4393\mu `$m and $`\mathrm{\Delta }\lambda =0.6\mu `$m $`\lambda \lambda 84196\mu `$m). The total integration time per resolution element and per pointing were: position #1 88s for the \[Oiii\] 52$`\mu `$m and 34s for the \[Niii\] 57$`\mu `$m; position #2 (the centre), 30s for the range 43–196 $`\mu `$m; positions NW and SE (2 point raster map) 22s for the the \[Oi\] 63$`\mu `$m, 14s for the \[Oiii\] 88$`\mu `$m, 12s for the \[Nii\] 122$`\mu `$m, 28s for the \[Oi\] 145$`\mu `$m and 12s for the \[Cii\] 158$`\mu `$m; position #4 12s for the \[Cii\] 158$`\mu `$m.
The data were processed with RAL pipeline 7 and analysed using the LIA and ISAP packages. The LWS flux calibration and relative spectral response function (RSRF) were derived from observations of Uranus (Swinyard et al. 1998). The full grating spectrum at the centre enabled us to estimate the relative flux uncertainty between individual detectors arising from uncertainties in the relative responsivity and the dark-current subtraction. The offsets between the detectors (excluding detector SW1) was $`10`$ %. The \[Oiii\] 88 $`\mu `$m line on detectors SW5 and LW1 had a 15 % systematic uncertainty and the \[Cii\] line on detectors LW3 and LW4 had a 10 % systematic uncertainty. We therefore adopt a relative flux uncertainty of $``$ 15%. Because we only took spectra of individual lines at the NW and SE positions there is no corresponding overlap in wavelength coverage at these positions. One indicator of relative flux uncertainty is a discrete step down in flux, of $``$ 25 %, at $``$ 125$`\mu `$m at the SE position. The relative flux uncertainty is assumed to be $``$ 25 % at these positions.
The absolute flux calibration w.r.t. Uranus for point like objects observed on axis is better than 15 % (Swinyard et al. 1998). However, extended sources give rise either to channel fringes or to a spectrum that is not a smooth function of wavelength. This is still a calibration issue. For example, in Fig. 2, detectors SW5, LW1, LW2 have slopes that differ from those of their neighbours in the overlap region. This may account for the continuum shape, which is discussed in Sect. 3.1. The LWS beam profile is known to be asymmetric and is still under investigation. We therefore adopt a value for the FWHM of 70<sup>′′</sup>at all wavelengths, believing that a more sophisticated treatment would not significantly affect our conclusions. We also note that there is good cross calibration between the ISO-LWS results and the Far-infrared Imaging Fabry-Perot Interferometer (FIFI) (Madden et al. 1995); the \[Cii\] peak fluxes agree to within $``$ 10 %.
## 3 Results and discussion
### 3.1 General FIR properties
The far-infrared continuum at each position is modelled with a single-temperature blackbody spectrum of the form F<sub>λ</sub> $`\alpha `$ $`\mathrm{\Omega }`$ B<sub>λ</sub>(T)(1-e$`^{\tau _{\mathrm{dust}}}`$), where the solid angle, $`\mathrm{\Omega }`$, is constrained to equal the LWS beam, B<sub>λ</sub>(T) is the Planck function at temperature T and $`\tau _{\mathrm{dust}}`$ $`\alpha `$ $`\lambda ^{1.5}`$. The result for the central position is shown as the dashed curve in Fig. 2. Although the observed continuum is not a simple function of wavelength and the single temperature blackbody is not an especially good fit, particularly at wavelengths $`>100\mu `$m, a better calibration of straylight and the beam profile is required for anything more sophisticated. The best FIR temperature at each position is $``$ 30 K.
The luminosities quoted here are derived from the line fluxes listed in Table 2. At the central position, the total luminosity in all of the far-infrared lines is 2.6 $`\times 10^7`$ L, which is $``$ 1 % of the total FIR luminosity (L$`{}_{43197\mu m}{}^{}=3.2\times 10^9`$ L). The \[Cii\] luminosity is 1.1 $`\times 10^7`$ L (0.4 % FIR) and the \[Oi\] luminosity is 7.5 $`\times 10^6`$ L (0.2 % FIR). Because full spectra are not available at the NW and SE positions we estimate the FIR continuum luminosity by integrating under the single-temperature blackbody fit to the data. At the NW position the total FIR luminosity, L$`{}_{43197\mu m}{}^{}=2.2\times 10^9`$ L. The \[Cii\] luminosity is 9.3 $`\times 10^6`$ L (0.4 % FIR) and the \[Oi\] luminosity is 3.4 $`\times 10^6`$ L (0.2 % FIR). At the SE position the total FIR luminosity, L$`{}_{43197\mu m}{}^{}=8.7\times 10^8`$ L. The \[Cii\] luminosity is 3.5 $`\times 10^6`$ L (0.4 % FIR) and the \[Oi\] luminosity is 2.0 $`\times 10^6`$ L (0.2 % FIR). These are typical values for starburst galaxies (c.f. Lord et al. 1996).
### 3.2 Ionized gas lines
Photons of energy 35.12, 29.60 and 14.53 eV are required to form O++, N++ and N+, respectively, so that the observed \[Oiii\] , \[Niii\] and \[Nii\] emission must originate in or around HII regions. The \[Oiii\] line ratio is a sensitive function of density in the range $`3010^4`$ cm<sup>-3</sup>. For the central position this ratio is $``$ 0.9, corresponding to an electron density, n$`{}_{\mathrm{e}}{}^{}`$ 100 cm<sup>-3</sup> (Rubin et al. 1994). The \[Oiii\] lines indicate a higher electron density, n$`{}_{\mathrm{e}}{}^{}`$ 250 cm<sup>-3</sup>, for the starburst nuclei of M82 (Colbert et al. 1999) and M83 (Stacey et al. 1999). In contrast, the \[Nii\] 205 $`\mu `$m / 122 $`\mu `$m line intensity ratio for the Galaxy gives an average electron density, of only $``$ 3 cm<sup>-3</sup> (Petuchowski & Bennett 1993). The thermal pressure of the ionized material in the Cen A dust lane is therefore closer to that of starburst galaxies than to that of the Milky Way.
Since the \[Niii\] 57 $`\mu `$m and the \[Nii\] 122 $`\mu `$m lines arise from different ionization states of the same element, the line intensity ratio is sensitive to the hardness of the interstellar UV field and therefore to the spectral type of the hottest main sequence star. For the central position \[Niii\] / \[Nii\] $``$ 1.6. This is larger than the value of $``$ 0.9 for M83 (Stacey et al. 1999) but smaller than the value of $``$ 2.1 for M82 (Colbert et al. 1999). Assuming that the region is ionization bounded, with an electron density, n<sub>e</sub> $``$ 100 cm<sup>-3</sup> the \[Niii\] / \[Nii\] line intensity ratio for Cen A corresponds to an abundance ratio N++/N+ of $``$ 0.3; this corresponds to an effective temperature, T$`{}_{\mathrm{eff}}{}^{}\mathrm{35\hspace{0.17em}500}`$ K (Rubin et al. 1994). Applying the same corrections to M82 and M83 with n<sub>e</sub> $``$ 250 cm<sup>-3</sup> implies an effective temperature, T$`{}_{\mathrm{eff}}{}^{}\mathrm{34\hspace{0.17em}500}`$ K for M83 and T$`{}_{\mathrm{eff}}{}^{}`$ 35 500 K for M82. If the effective temperature in Cen A corresponds to the tip of the main sequence formed in a single starburst, we are observing O8.5 stars and the burst is $``$ 6 $`\times 10^6`$ years old. If the burst was triggered by the spiral-elliptical galaxy merger then its occurance was very recent. Alternatively, the merger triggered a series of bursts of star formation, of which we are witnessing the most recent.
The N++ and O++ coexist in roughly the same ionization zones, and the \[Oiii\] 52 $`\mu `$m and \[Niii\] 57 $`\mu `$m lines have roughly the same critical density. As a result the ratio of these lines is an indicator, to within $`50`$%, of the N++/O++ abundance ratio, which itself, is nearly equal to the N/O ratio in the hard UV field environments we are seeing here (Rubin et al. 1988). The line ratio we observe at the centre of the dust lane is $``$ 0.3 - the same as found in the nucleus of M82 (Colbert et al. 1999), but much smaller than that found for the nucleus of M83 ($``$ 0.67 Stacey et al. 1999).
A more precise determination of the abundance ratio requires the observed line ratio to be divided by the volume emissivity ratio. The latter ratio is dependent on the electron density because the two lines have slightly different critical densities. Using our value for the electron density $``$ 100 cm<sup>-3</sup> and Fig. 3 of Lester et al. (1987) we estimate that the N/O abundance ratio to be $``$ 0.2 in Cen A. This value is consistent with the range of $``$ 0.2 - 0.3 found for Galactic HII regions (Rubin et al. 1988). The nitrogen to oxygen abundance ratio is a measure of the chemical evolution and we expect it to increase with time (cf. the solar value of $``$ 0.12).
### 3.3 Neutral gas lines
Carbon has a low ionization potential (11.4 eV), which is less than that of hydrogen. \[Cii\] 157$`\mu `$m line emission is therefore observed from both neutral and ionized hydrogen clouds. We model the \[Cii\] line emission with three components: Photodissociation regions (PDRs) on the surfaces of UV exposed molecular clouds; cold (T $``$ 100 K) HI clouds (i.e. the cold neutral medium (CNM) Kulkarni & Heiles 1987, Wolfire et al. 1995); and diffuse HII regions (i.e. the warm ionized medium (WIM) Heiles 1994).
#### 3.3.1 HI clouds
It can be shown that the intensity in the \[Cii\] line emitted from gas clouds with density, n(H) and temperature (T) is given by (c.f. Madden et al. 1993)
$$\mathrm{I}_{c^+}=2.35\times 10^{21}\left[\frac{2exp(\frac{91.3}{T})}{1+2exp(\frac{91.3}{T})+\frac{n_{crit}}{n_H}}\right]X_{c^+}N(HI)$$
(1)
where the critical density for collisions with H, n$`{}_{\mathrm{crit}}{}^{}3000`$ cm<sup>-3</sup> (Launay & Roueff 1977) and the fractional C<sup>+</sup> relative to hydrogen is $`\mathrm{X}_{\mathrm{c}^+}\mathrm{X}_\mathrm{c}=1.4\times 10^4`$ (Sofia et al. 1997).
N(HI) is estimated from the HI 21cm map of Van Gorkom et al. (1990) to be 18.8 $`\times 10^{20}`$ atom cm<sup>-2</sup> at the SE position. The central and NW positions are difficult to estimate due to HI absorption against the nuclear continuum. There may be a central hole in the HI and the column density is certainly not higher than the peak observed in the SE region of the dust lane (Van Gorkom et al. 1990)
Assuming typical Galactic values for the temperature, T $``$ 80 K, and hydrogen density, n $``$ 30 cm<sup>-3</sup>, results in an estimated \[Cii\] flux of $`3.5\times 10^{20}`$ W cm<sup>-2</sup> in a 70<sup>′′</sup>LWS beam at the SE position. This corresponds to 4 %, 1 % and 1 % of the observed \[Cii\] flux at positions SE, Centre and NW respectively. The peak HI emission line flux corresponds to $`1.9\times 10^{19}`$ W cm<sup>-2</sup> which is only 6% and 8% of the \[Cii\] flux at the centre and NW positions respectively. We conclude that there is very little \[Cii\] emission in our beams from HI clouds.
#### 3.3.2 Diffuse HII regions
Ionized carbon can be found in both neutral gas and ionized gas clouds, and is an important coolant for each. We detected \[Oiii\] 88 $`\mu `$m in all 3 beam positions so there is an ionized gas component in each beam. Using the constant density HII region model of Rubin (1985) with the Kurucz abundances, 10<sup>49</sup> ionizing photons per second and our derived density, n$`{}_{e}{}^{}`$ 100 cm<sup>-3</sup> and effective temperature, T$`{}_{\mathrm{eff}}{}^{}`$ 35 500 K we can estimate the \[Cii\] emission from the HII regions. Applying the model \[Oiii\] 88 $`\mu `$m / \[Cii\] 158 $`\mu `$m line ratio of 0.35 to the observed \[Oiii\] 88 $`\mu `$m line flux at each position results in $``$ 10 % contribution to the observed \[Cii\] line flux in each beam. Scaling the model fluxes to the distance of Cen A gives $``$ 3000 HII regions in the central and NW regions and $``$ 1000 HII regions in the SE region.
The estimate above assumes that the observed lines have the same filling factor in the large LWS beam. If, alternatively, we were to assume that the ionized component was instead dominated by a contribution from an extended low density warm ionized medium (ELDWIM) with n$`{}_{e}{}^{}`$ 3 cm<sup>-3</sup>, then the \[Cii\] flux can be estimated from the ratio of the \[Cii\] /\[Nii\] lines to be $``$ 18 % at the central position. The observations of the \[Nii\] 121.9$`\mu `$m line at the NW and SE positions (with lower signal to noise) indicate a similar fractional component (21 % and $`56`$ %, respectively).
We have estimated the density in the HII regions in the centre of Cen A to be $``$ 100 cm<sup>-3</sup> with an effective temperature, T$`{}_{\mathrm{eff}}{}^{}\mathrm{35\hspace{0.17em}500}`$ K. Based on the HII region models of Rubin (1985) we estimate that $``$ 10 % of the observed \[Cii\] arises in the WIM.
#### 3.3.3 PDRs
Far-UV photons (6 eV $`<`$ h$`\nu `$ 13.6 eV) from either O/B stars or an AGN will photo-dissociate H<sub>2</sub> and CO molecules and photo-ionize elements with ionization potentials less than the Lyman limit (e.g. C<sup>+</sup> ionization potential = 11.26 eV). The gas heating in these photodissociation regions (PDRs) is dominated by electrons ejected from grains due to the photoelectric effect. Gas cooling is dominated by the emission of \[Oi\] 63$`\mu `$m and \[Cii\] 158$`\mu `$m emission. Observations of these lines, the \[Oi\] 146$`\mu `$m and CO(J=1-0) 2.6 mm lines and the FIR continuum can be used to model the average physical properties of the neutral interstellar medium (Wolfire et al. 1990). Kaufman et al. (1999) have computed PDR models over a wide range of physical conditions. The new code accounts for gas heating by small grains/PAHs and large molecules, and uses a lower, gas phase carbon abundance (X<sub>C</sub> = 1.4x10<sup>-4</sup>, Sofia et al. 1997) and oxygen abundance (X<sub>O</sub> = 3.0x10<sup>-4</sup>, Meyer et al. 1998). The \[Oi\] 63 $`\mu `$m / \[Cii\] 158 $`\mu `$m line ratio and either the \[Oi\] 146 $`\mu `$m / \[Oi\] 63 $`\mu `$m line ratio or the (\[Oi\] 63 $`\mu `$m + \[Cii\] 158 $`\mu `$m) / FIR continuum can be used as PDR diagnostics to determine the average gas density (i.e the proton density, n cm<sup>-3</sup>), the average incident far-UV flux (in units of the Milky Way flux, G<sub>o</sub> = 1.6 $`\times 10^3`$erg cm<sup>-2</sup> s<sup>-1</sup>) and the gas temperature.
We assume that the measured \[Cii\] flux at each position should have $``$ 10 % subtracted, due to the HI and WIM components, before it is used to model the PDRs (if, alternatively, a 20 % ELDWIM contribution is subtracted it would not significantly affect the PDR parameters derived below). The PDR lines are plotted in Fig. 3 and the line intensity ratios are given in Table 3.
The results for the three regions are consistent with each other, having a gas density, n $``$ 10<sup>3</sup> cm<sup>-3</sup>, and an incident far-UV field, G $``$ 10<sup>2</sup>.
At the NW position, only the combination of the \[Oi\] 63 $`\mu `$m / \[Cii\] 158 $`\mu `$m ratio and the (\[Oi\] 63 $`\mu `$m + \[Cii\] 158 $`\mu `$m) /FIR continuum ratio gives a meaningful solution for G and n. The \[Oi\] 146 $`\mu `$m line is clearly detected but with a very rippled baseline due to channel fringes. The observed \[Oi\] 146 $`\mu `$m line flux would need to be reduced by $``$ 60 % in order to obtain a consistent result with the \[Oi\] 146 $`\mu `$m / \[Oi\] 63 $`\mu `$m line ratio predicted by the PDR model.
The LWS results for the nucleus confirm those previously derived from IR, submm and CO observations. The consistent set of derived PDR conditions for all three positions suggest that the observed FIR emission in a 70<sup>′′</sup>beam centred on the nucleus is dominated by star formation and not AGN activity. Joy et al. (1988) mapped Cen A at 50 and 100 $`\mu `$m on the KAO. They concluded that the extended FIR emission was from dust grains heated by massive young stars distributed throughout the dust lane, not the compact nucleus. Hawarden et al. (1993) mapped Cen A at 800 $`\mu `$m and 450 $`\mu `$m with a resolution of $``$10 <sup>′′</sup>. They attribute the large scale 800 $`\mu `$m emission to thermal emission from regions of star formation embedded in the dust lane. They note that the H<sub>2</sub> emission within a few arcseconds of the nucleus, observed by Israel et al. (1990), indicates that significant UV radiation from the nucleus does not reach large radii in the plane of the dust lane i.e. the nuclear contribution to exciting the extended gas and dust disk is small.
Eckart et al. (1990) and Wild et al. (1997) mapped Cen A in <sup>12</sup>CO J=1-0, <sup>12</sup>CO J=2-1 and <sup>13</sup>CO J=1-0. All three maps have two peaks separated by $``$ 90 <sup>′′</sup>centred on the nucleus. It is interesting to note that our SE position only clips the lowest contours of the CO (1-0) and CO (2-1) maps of Wild et al. (1997). In spite of this the derived PDR parameters are consistent with those encompassing the bulk of the molecular emission. There must be extended low level CO (1-0) emission beyond the sensitivity limits of the Wild et al. (1997) maps. The lowest contour is 17.5 K kms<sup>-1</sup>, corresponding to M$`_{\mathrm{H}_2}`$ $``$ 10<sup>8</sup> M if the material filled the LWS beam.
## 4 Summary and conclusions
We present the first full FIR spectrum from 43 - 196.7 $`\mu `$m of Cen A. We detect seven fine structure lines (see Table 2), the strongest being those generated in PDRs. At the central position, the total flux in the far-infrared lines is $``$ 1 % of the total FIR luminosity (L$`{}_{43197\mu m}{}^{}=3.2\times 10^9`$ L for a distance of 3.5 Mpc). The \[Cii\] line flux is $``$0.4 % FIR and the \[Oi\] line flux is $``$ 0.2 % FIR. These are typical values for starburst galaxies (Lord et al. 1996). The \[Oiii\] 52 $`\mu `$m / \[Oiii\] 88 $`\mu `$m line intensity ratio is $``$ 0.9, which corresponds to an electron density, n$`{}_{\mathrm{e}}{}^{}`$ 100 cm<sup>-3</sup> (Rubin et al. 1994). The thermal pressure of the ionized medium in the Cen A dust lane is closer to that of starburst galaxies (n$`{}_{e}{}^{}`$ 250 cm<sup>-3</sup> in M82 (Colbert et al. 1999) and M83 (Stacey et al. 1999)) than that of the Milky Way (n$`{}_{e}{}^{}`$ 3 cm<sup>-3</sup> (Pettuchowski & Bennett 1993)).
The \[Niii\] / \[Nii\] line intensity ratio is $``$ 1.6, giving an abundance ratio N++/N+ $``$ 0.3, which corresponds to an effective temperature, T$`{}_{\mathrm{eff}}{}^{}`$ 35 500 K (Rubin et al. 1994). Assuming a coeval starburst, then the tip of the main sequence is headed by O8.5 stars, and the starburst is $``$ 6 $`\times 10^6`$ years old. If the burst in Cen A was triggered by the spiral-elliptical galaxy merger then its occurance was very recent. Alternatively, the merger triggered a series of bursts of star formation and we are witnessing the most recent activity.
We estimate that the N/O abundance ratio is $``$ 0.2 in the HII regions in Cen A. This value is consistent with the range of $``$ 0.2 - 0.3 found for Galactic HII regions (Rubin et al. 1988). N/O is a measure of the chemical evolution and we expect it to increase with time (c.f. the solar value of $``$ 0.12).
We estimate that $``$ 10 % of the observed \[Cii\] arises in the WIM. The CNM contributes very little ($`<5`$ %) \[Cii\] emission at our beam positions. The bulk of the emission is from the PDRs.
We derive the average physical conditions for the PDRs in Cen A for the first time. There is active star formation throughout the dust lane and in regions beyond the bulk of the molecular material. The FIR emission in the 70<sup>′′</sup>LWS beam at the nucleus is dominated by emission from star formation rather than AGN activity. On scales of $``$ 1 kpc the average physical properties of the PDRs are modelled with a gas density, n $``$ 10<sup>3</sup> cm<sup>-3</sup>, an incident far-UV field, G $``$ 10<sup>2</sup> and a gas temperature of $``$ 250 K.
## Acknowledgements
Many thanks to the dedicated efforts of the LWS instrument team. The ISO Spectral Analysis Package (ISAP) is a joint development by the LWS and SWS Instrument Teams and Data Centers. Contributing institutes are CESR, IAS, IPAC, MPE, QMW, RAL and SRON.
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warning/0002/astro-ph0002491.html | ar5iv | text | # Hot Stars and Cool Clouds: The Photodissociation Region M161
## 1. Introduction
Star formation occurs in dense regions of the interstellar medium, and as a result, the environments of early-type stars often contain molecular clouds. In addition to creating an H II region, where hydrogen is predominantly ionized, stellar ultraviolet radiation creates a photodissociation region (PDR) when it interacts with molecular material. The feedback between extant stars and the nearby clouds in PDRs is important in regulating subsequent episodes of star formation. The UV flux determines the ionization fraction in molecular clouds, which in turn sets the ambipolar diffusion rate and thereby the rate of low-mass star formation (McKee (1989)).
Hollenbach & Tielens (1997, 1999) comprehensively review the physical and chemical processes of PDRs. Far-ultraviolet (FUV) photons ($`6<h\nu <13.6\mathrm{eV}`$) dissociate molecules and excite the gas on the surface ($`A_V<3`$) of the molecular cloud. The result is spatial stratification of molecular, atomic, and ionic components of a given element. At higher densities ($`n\genfrac{}{}{0pt}{}{_>}{^{}}10^4\mathrm{cm}^3`$), collisional excitation also populates the excited states (Sternberg & Dalgarno (1989)).
The columns of M16 (the Eagle Nebula) are PDRs. Most of the mass in the nebula is in molecular clouds, and a variety of molecular transitions are observed directly (Pound (1998); White et al. (1999)). H<sub>2</sub> is observed on the cloud surfaces (Allen et al. (1999)), and at the interface with the H II region, the columns are prominent in optical emission lines (Hester et al. (1996)). The young cluster NGC 6611 contains many massive stars, which illuminate the columns with total FUV flux $`16\mathrm{erg}\mathrm{s}^1\mathrm{cm}^2`$ (Allen et al. (1999)) from a distance of about 2 pc (Hester et al. (1996)) .
We simultaneously observe the spatial and spectral relationship of the excited molecular and ionized phases of the PDR. Based on these observations (described in §2), we determine physical conditions in §3, and summarize the conclusions in §4.
## 2. Observations and Data Reduction
We obtained images and high resolution longslit spectra of the northernmost column of M16 using NIRSPEC, the facility near-infrared spectrometer on the Keck II 10-m telescope (McLean et al. (1998)). Figure 1a identifies the observed region on the Hubble Space Telescope H$`\alpha `$ image of M16 (Hester et al. (1996)) with contours of integrated CO intensity (Pound (1998)). We used the slit-viewing camera to image M16 on 1999 July 6 (UT), obtaining 9 exposures of 60 s each. In addition to the H<sub>2</sub> 1-0S(1) filter, (bandpass 2.110–2.129 $`\mu `$m) the NIRSPEC-6 filter (1.558–2.315 $`\mu `$m) was in the optical path. We combined several observations of a blank field in the NIRSPEC-6 filter alone to use as a flat field and scaled this for sky subtraction from this crowded field. Figure 1b shows the resultant 2.12 $`\mu `$m image and the spectroscopy slit.
We employed the NIRSPEC-6 blocking filter and selected echelle and cross disperser positions to detect several H<sub>2</sub> rovibrational transitions, Br$`\gamma `$, and He I $`\lambda `$2.06 $`\mu `$m in our spectroscopic observations. We used the $`0\stackrel{}{\mathrm{.}}432\times 24^{\prime \prime }`$ slit at position angle 135 for resolution $`\lambda /\mathrm{\Delta }\lambda =26,000`$, where $`\mathrm{\Delta }\lambda `$ is the observed $`FWHM`$ of an unresolved emission line. We obtained six exposures of 300 s each on the night of 1999 April 30 (UT), nodding along the slit between each exposure.
We subtracted a median dark frame from each image, divided by a flat field, and interpolated over deviant pixels, removing both cosmic rays and bad detector pixels. The spectra required rectification in both the spectral and spatial dimensions. We used the OH night sky lines and wavenumbers tabulated by Abrams et al. (1994) for wavelength calibration and the continuum emission of the telluric standard for spatial rectification. We observed HD 161056 and reduced the spectra similarly, then modelled this B1.5V star as a $`T=24,000`$ K blackbody and interpolated over the stellar Br$`\gamma `$ feature. We divided the M16 data by this spectrum to correct for atmospheric absorption and normalized the stellar spectrum from this $`K=5.4`$ standard at 2.2 $`\mu `$m to flux calibrate the data.
Figure 3.1 contains the resultant two-dimensional spectra of several bright lines. We have subtracted off-source emission from each of these spectra. For the molecular emission, the background is near the edge of the slit toward the ionizing source, while for Br$`\gamma `$, the background lies away from the ionizing source. We list the complete set of detected lines, their vacuum wavelengths, their total fluxes in the slit, and their peak fluxes in Table 1.
## 3. Physical Conditions
### 3.1. Geometry
The gross structure of the photodissociation region is evident in the images (Fig. 1) and two-dimensional spectra (Fig. 3.1). Fewer stars are observed directly through the cloud than in the surrounding region (Fig. 1b), implying $`A_K>1`$. From the center of the molecular cloud out to the photoionizing source, distinct regions of quiescent molecular, excited molecular, atomic, and ionized material appear. The FUV photons do not fully penetrate the molecular cloud, so the excited H<sub>2</sub> resides only on the cloud’s surface. The optically-thin ionized emission is brightest in the regions of highest density, near the cloud, and it is spatially more extended than the H<sub>2</sub> emission. The slit is aligned toward the ionizing source and covers several edges on the irregular surface of the cloud. Thus, we observe several distinct regions of photoevaporative flow. The greatest intensities are detected in the regions that are viewed edge-on.
In these observations, each of the bright regions within the PDR consists of spatially segregated molecular and recombination line emission. The two brightest regions at slit positions $`1\stackrel{}{\mathrm{.}}2`$ and $`3\stackrel{}{\mathrm{.}}2`$ in the 1-0S(1) emission, for example, are physically associated with the ionized emission at $`7\stackrel{}{\mathrm{.}}1`$ and $`3\stackrel{}{\mathrm{.}}6`$, respectively (Fig. 3.1d). (Note, however, that the more spatially extended Br$`\gamma `$ peak at the top of slit includes some contribution from a nearby surface north of the slit, which we do not observe in H<sub>2</sub>.) Moderate emission peaks appear at $`2\stackrel{}{\mathrm{.}}8`$ and $`1\stackrel{}{\mathrm{.}}8`$, and fainter surfaces at $`9\stackrel{}{\mathrm{.}}3`$, $`7\stackrel{}{\mathrm{.}}9`$, and $`5\stackrel{}{\mathrm{.}}3`$ in the 1-0S(1) line. The first two of these are blended in the spatially extended Br$`\gamma `$ emission, while the latter correspond to local emission maxima at $`2\stackrel{}{\mathrm{.}}8`$, $`0\stackrel{}{\mathrm{.}}2`$, and $`11\stackrel{}{\mathrm{.}}3`$, so the projected separations of the H<sub>2</sub> and Br$`\gamma `$ emission regions in the plane of the sky range from 1.8 to $`2.5\times 10^{17}`$ cm at the 2000 pc distance of M16 (Humphreys (1978); Hillenbrand et al. (1993)). The transition from the atomic to molecular zone typically occurs at $`A_V2`$ (Tielens & Hollenbach (1985)), or in terms of column density, $`N_H4\times 10^{21}\mathrm{cm}^2`$. Thus, the average gas density in the bright regions of the PDR is 1.6–$`2.2\times 10^4\mathrm{cm}^3`$.
We determine the density in the ionized region and use the molecular line width (discussed below) to calculate the average density in the molecular region of the PDR. Bertoldi & Draine (1996) show $`n_i=S_{Ly}/(4\pi R^2qc_i)`$, where $`S_{Ly}`$ is the Lyman continuum flux, $`R`$ is the distance to the ionizing source, $`q`$ accounts for attenuation of the ionization in the evaporative flow, and $`c_i`$ is the isothermal sound speed in the ionized medium. We determine the value of $`q4.0\times 10^{14}\sqrt{S_{Ly}r}/R`$ (Bertoldi (1989)), where $`r`$ is the cloud radius. We find the density of the ionized region $`n_i=5.2\times 10^3\mathrm{cm}^3`$, for $`c_i=12\mathrm{km}\mathrm{s}^1`$ (calculated for $`T_{H\mathrm{ii}}=9500`$ K, below) and $`r=25^{\prime \prime }`$. For pressure balance across the ionization front, $`n2n_ic_i^2/(\mathrm{\Delta }v/2.35)^2`$, where $`\mathrm{\Delta }v`$ is the $`FWHM`$ of the molecular line width, so in the molecular region, $`n3.2\times 10^5\mathrm{cm}^3`$. This is significantly greater than the density calculated above, which integrated over contributions from the ionized to the molecular region. Here, $`n`$ is determined in the molecular zone alone, which should be denser.
### 3.2. Kinematics
The various gas phases exhibit distinct kinematic structure. We measure the line profiles at slit position $`3\stackrel{}{\mathrm{.}}2`$ in 1-0S(1) emission, which is clearly distinguished in the Br$`\gamma `$ and He emission at $`3\stackrel{}{\mathrm{.}}6`$. Because we observe this PDR nearly edge-on, the central velocities of the H$`{}_{}{}^{}{}_{2}{}^{}`$ and ionized material are not significantly different; the average motion occurs in the plane of the sky. We measure $`v_{LSR}=29\mathrm{km}\mathrm{s}^1`$ for this 1-0S(1) feature and $`v_{LSR}=28\mathrm{km}\mathrm{s}^1`$ in Br$`\gamma `$. In H<sub>2</sub>, $`\mathrm{\Delta }v=4.6\pm 0.1\mathrm{km}\mathrm{s}^1`$, $`\mathrm{\Delta }v=29.0\pm 0.2\mathrm{km}\mathrm{s}^1`$ in Br$`\gamma `$, and $`\mathrm{\Delta }v=18.7\pm 0.2\mathrm{km}\mathrm{s}^1`$ in He I $`\lambda 2.06\mu \mathrm{m}`$, after correcting for the instrumental profile. In the molecular gas, the velocity width corresponds to a kinetic temperature $`T_{H_2}=930\pm 50`$ K. This value is reasonable near the edge of the cloud in the H$`{}_{}{}^{}{}_{2}{}^{}`$ zone where maximum temperature of the PDR occurs (Tielens & Hollenbach (1985)).
Accounting for the mass difference between H and He demonstrates that these line widths are not purely thermal, however. Assuming that they each consist of a macroscopic component that is independent of molecular weight in addition to the thermal component, we find $`\mathrm{\Delta }v_{macro}=12.2\mathrm{km}\mathrm{s}^1`$. For Br$`\gamma `$, $`\mathrm{\Delta }v_{therm}=25.6\mathrm{km}\mathrm{s}^1`$, and $`T_{H\mathrm{ii}}=14000\pm 230`$ K. This temperature is hotter than typically observed in H II regions, but several factors tend to increase $`T_{H\mathrm{ii}}`$ at the PDR-H II interface: increased photoelectric heating near a greater population of grains; the harder radiation field, which reduces the contribution of significant low-ionization coolants, such as O<sup>+</sup>; and increased density, which reduces the radiative cooling rate by increasing the collisional de-excitation rate (e.g., Spitzer (1978)).
The non-thermal recombination line widths may also be attributable to bulk motion of material flowing off the cloud surface. In the bright region discussed above, we view the PDR nearly edge-on, along a line of sight that includes both redshifted and blueshifted velocity components of the material that evaporates off the curved surface of the cloud. When the view of the PDR is closer to face-on, the intensity is less enhanced, and only the blueshifted velocity component arises. This is the case for the minor emission peaks at the bottom of the slit, which occur on the near surface of the cloud. These central velocities are slightly lower (by approximately 7$`\mathrm{km}\mathrm{s}^1`$) than those of the brightest peaks, and the observed widths are narrower. At slit position $`2\stackrel{}{\mathrm{.}}8`$, for example, we find $`\mathrm{\Delta }v=20.9\pm 0.2`$ and $`10.1\pm 0.2\mathrm{km}\mathrm{s}^1`$ in Br$`\gamma `$ and He I, respectively. These both correspond to kinetic temperature $`T_{H\mathrm{ii}}=9500\pm 180`$ K. If this is a more accurate measurement of the temperature at the edge of the H II region, then we have underestimated the contribution of the non-thermal velocity component above. While $`T_{H\mathrm{ii}}`$ at the interface is likely higher than the typical temperatures in the diffuse central portions of H II regions, as argued above, it may be less than 14,000 K.
### 3.3. H<sub>2</sub> Excitation
In a PDR, UV pumping is the primary source of H<sub>2</sub> excitation. In M16, such fluorescent emission obviously occurs. A variety of lines that originate in high vibrational levels $`v`$ are observed, including 7-5O(4) and 3-2S(5) (Table 1), whereas thermal population of these high-$`v`$ levels would dissociate the molecules. We detect emission from all of the strong transitions in our bandpass that Black and van Dishoeck (1987) predict, except those originating in levels $`v9`$ or lie in regions of poor atmospheric transparency.
Collisional excitation is also important, however. The ratio of line intensities $`I`$(1-0S(1))/$`I`$(2-1S(1))$`=2`$–5, and the average ratio over the slit is 4. As we found above in §3.1, this requires $`n>10^4\mathrm{cm}^3`$ (Burton, Hollenbach, & Tielens (1990)), so collisions aid in populating the $`v=1`$ level. We do not, however, observe line ratios as high as Allen et al. (1999), who find $`I`$(1-0S(1))/$`I`$(2-1S(1))$`>10`$ in similar regions with Fabry-Perot imaging. We expect our results to be more accurate, having better flux calibration and simultaneous sky measurements.
Line ratios also directly measure excitation temperatures. Across the slit, we find from the 1-0S(1) and 2-1S(1) lines vibrational temperature $`T_{vib}=3000`$–3800 K, with a median $`T_{vib}=3400`$ K. Using $`I`$(1-0S(1))/$`I`$(1-0S(3)), we find rotational temperature $`T_{rot}=1600`$–2100 K and median $`T_{rot}=1800`$ K. Comparing $`I`$(1-0S(1))/$`I`$(1-0S(0)), $`T_{rot}=450`$–870 K, with a median $`T_{rot}=580`$ K. The relatively low value of $`T_{vib}`$ and relatively high value of $`T_{rot}`$ (particularly in the brightest regions) distinguish the M16 PDR from purely fluorescent emission in a low density medium, where $`T_{vib}5500`$ and $`T_{rot}1400`$ from $`I`$(1-0S(1))/$`I`$(1-0S(3)) are typical.
We compare the line intensities of the three brightest regions with the models of Draine & Bertoldi (1996). Considering all observed lines, all three regions are best fit with $`\chi /n=0.1\mathrm{cm}^3`$, and $`T_0=1000`$ K, where $`\chi `$ measures the incident UV flux ($`\chi =1`$ is the Habing flux), and $`T_0`$ is the temperature at the edge of the PDR. The observed line ratios clearly rule out models of a weak field ($`\chi 10^3`$) and low density ($`n10^4\mathrm{cm}^3`$). Both $`\chi =10^4`$ and $`\chi =10^5`$ are acceptable. For consistency with other density measurements and the UV field expected from members of NGC 6611, $`\chi =10^4`$ and $`n=10^5\mathrm{cm}^3`$ is the preferred model. Although the grid of models is coarse, we identify trends by contrasting these regions. For example, the variation of density-sensitive line ratios along the spatial extent of the slit demonstrates that the cloud density is not constant. The decreased ratios $`I`$(1-0S(1))/$`I`$(1-0S(0)) and $`I`$(1-0S(1))/$`I`$(2-1S(1)) at slit position $`3\stackrel{}{\mathrm{.}}2`$ relative to slit position $`1\stackrel{}{\mathrm{.}}2`$ show that the former region has lower density.
In thermal equilibrium, the ratio of ortho-hydrogen (having odd rotational level, $`J`$) to para-H<sub>2</sub> (having even $`J`$) is 3, and only collisions can change the ortho:para ratio. This abundance ratio is therefore a significant parameter in PDRs and affects the emission line ratios. We calculate the ratio of ortho-H<sub>2</sub> and para-H<sub>2</sub> from the 1-0S(1) and 1-0S(0) lines, finding ratios of 1.3–2.2, with an median value of 1.7. Because UV pumping of the damping wings of optically thick lines populate the excited levels, the expected ratio of column densities in the vibrationally excited levels is 1.7, while the true ortho:para abundance is 3 (Sternberg & Neufeld (1999)).
The photoelectric effect on interstellar grains is likely to be the primary heating source in PDRs, but its exact process is uncertain. We determine the efficiency of converting FUV flux to photoelectric heating, $`ϵ`$, assuming thermal equilibrium. Emission in the \[O I\] 63 $`\mu `$m and \[C II\] 158 $`\mu `$m fine structure lines cools the gas. For the measured temperature and density, we calculate a net cooling rate $`\mathrm{\Lambda }=4.2\times 10^{17}\mathrm{erg}\mathrm{s}^1\mathrm{cm}^3`$ in M16, with 96% due to \[O I\]. In the parameterization of Bakes & Tielens (1994), the heating function $`\mathrm{\Gamma }=10^{24}nϵ\chi \mathrm{erg}\mathrm{s}^1\mathrm{cm}^3`$. Thus, we find $`ϵ=0.042`$, for $`\chi =9700`$ (Allen et al. (1999)). The theoretical prediction, $`ϵ=0.016`$ (Bakes & Tielens (1994)), is a function of $`\chi \sqrt{T}/n_e`$, where $`n_e`$ is the electron density. We have assumed that the ionization of carbon provides the electrons; increasing $`n_e`$ would increase $`ϵ`$, without significantly altering the cooling rate. Alternatively, decreasing the gas to dust ratio would allow a lower efficiency in equilibrium with the cooling rate we predict.
## 4. Conclusions
Spectroscopy of M16 with NIRSPEC offers the advantage of simultaneous observation of both the molecular and ionized phases of its photodissociation regions. At high resolution, we detect the spatial and kinematic signatures of photoexcitation of molecular material. Most of the emission is due to distinct regions of photodissociation on the molecular cloud’s surface that we observe edge-on. The spectrum includes several lines that originate in high vibrational levels, which demonstrates that fluorescent excitation predominates. The density $`n10^5\mathrm{cm}^3`$ and varies among the emitting regions, so collisions also preferentially populate low-lying vibrational levels. We measure kinetic temperature $`T_{H_2}=930`$ K in the molecular gas, which requires efficient photoelectric heating near the cloud’s surface. In the spatially-segregated ionized region, $`T_{H\mathrm{ii}}=9500`$ K, and the velocity widths of the brightest regions include a macroscopic contribution.
It is a pleasure to acknowledge the hard work of past and present members of the NIRSPEC instrument team at UCLA: Maryanne Angliongto, Oddvar Bendiksen, George Brims, Leah Buchholz, John Canfield, Kim Chin, Jonah Hare, Fred Lacayanga, Samuel B. Larson, Tim Liu, Nick Magnone, Gunnar Skulason, Michael Spencer, Jason Weiss, and Woon Wong. In addition, we thank the Keck Director Fred Chaffee, CARA instrument specialist Thomas A. Bida, and all the CARA staff involved in the commissioning and integration of NIRSPEC. We especially thank our Observing Assistants Joel Aycock, Gary Puniwai, Charles Sorenson, Ron Quick, and Wayne Wack for their support. We thank Amiel Sternberg for helpful discussions and Marc Pound for providing his CO data. |
warning/0002/hep-ph0002104.html | ar5iv | text | # Constraints on a Model with Pure Right-Handed Third Generation Couplings
## I Introduction
The standard $`SU(2)_L\times U(1)`$ model of the weak interactions has achieved great success. Nevertheless, viable competing models with the $`SU(2)_L\times SU(2)_R\times U(1)`$ gauge group have been proposed . In these models the left-handed Cabbibo-Kobayashi-Maskawa (CKM) mixing matrix is that of the standard model and the parameters of the right-handed mixing matrix as well as the right handed coupling and the mass of the right-handed gauge bosons are constrained by experiment. In 1992, Gronau and Wakaizumi (GW) presented a model with this gauge group in which the flavor changing third generation decays occur only through right-handed currents . This model owes its feasibility to the difficulty in differentiating between $`(VA)(VA)`$ quark-lepton couplings and $`(V+A)(V+A)`$ couplings. Experimental evidence has since ruled out the GW model as a possible alternative to the standard model. However, a more general choice than that chosen by GW for the right-handed mixing matrix, although tightly constrained by experiment, can not be entirely excluded on phenomenological grounds.
CP violation in the GW model and its more general extensions has been studied previously . Previous authors have shown that the GW model parameters are constrained by the CP violating observables $`ϵ`$ and $`ϵ^{}`$ and the neutron dipole moment. They have also shown asymmetry values in nonleptonic neutral $`B`$ decays differing from standard model predictions. It is, however, necessary to reexamine the constraints imposed by CP violation on these models in light of recent experiments .
In this paper we take the following approach. First, we briefly review non-symmetric left-right models and the GW model. Then, in section 3 we constrain the angles of the most general right-handed mixing matrix from observables not related to CP violation. We find that there is a tightly constrained region in which this model is viable, and we make a particular choice of angles. We then place constraints on the phases from CP violating observables in section 4. With constraints so imposed we examine various predictions in $`B`$ decays in section 5. We summarize our results in section 6.
## II Review of $`SU(2)_L\times SU(2)_R\times U(1)`$ models
Langacker and Sankar have reviewed $`SU(2)_L\times SU(2)_R\times U(1)`$ models . In discussing these models below we follow much of their notation.
In $`SU(2)_L\times SU(2)_R\times U(1)`$ models, the left and right-handed quarks and leptons transform under doublets of separate $`SU(2)`$ gauge groups. This gives rise to a covariant derivative of the form
$$D^\mu =^\mu +\frac{i}{2}(g_L\tau ^aW_L^{\mu a}+g_R\tau ^aW_R^{\mu a}+g^{}YB^\mu ),$$
(1)
where $`g^{}`$ is the $`U(1)`$ gauge coupling, $`\tau ^a`$ are the Pauli spin matrices, $`W_{L,R}^a`$ and $`B`$ are the gauge boson fields and $`g_{L,R}`$ are the $`SU(2)_L`$ and $`SU(2)_R`$ gauge coupling constants. The gauge symmetry is spontaneously broken by introducing a Yukawa interaction with some Higgs sector and giving the Higgs a vacuum expectation value. This gives masses to the quarks, leptons and gauge bosons. We take a Higgs, $`\mathrm{\Phi }`$, that transforms as $`\mathrm{\Phi }L\mathrm{\Phi }R^{}`$ under $`SU(2)_L`$ and $`SU(2)_R`$ and is neutral under hypercharge. A general choice for the vacuum expectation value gives
$$\mathrm{\Phi }=\left(\begin{array}{cc}\varphi _1^0\hfill & \varphi _1^+\hfill \\ \varphi _2^{}\hfill & \varphi _2^0\hfill \end{array}\right)\left(\begin{array}{cc}k\hfill & 0\hfill \\ 0\hfill & k^{}\hfill \end{array}\right).$$
(2)
With this Higgs we have the relation $`M_R=g_R/g_LM_L`$, where $`M_L`$ and $`M_R`$ are the masses of the left and right-handed charged gauge bosons respectively. Taking the ratio $`g_R/g_L`$ to be $`𝒪(1)`$, it is necessary to introduce additional Higgs to arrange for these masses to be much different. Minimally, one introduces two doublets or triplets under $`SU(2)_L`$ and $`SU(2)_R`$ which carry a hypercharge of 1. These obtain vacuum expectation values $`v_L`$ and $`v_R`$.
The Yukawa couplings to the quarks are given as
$$_Y=\underset{i,j}{}(\overline{f}_{iL}^{}(r_{ij}\mathrm{\Phi }+s_{ij}\stackrel{~}{\mathrm{\Phi }})f_{jR}^{}+\mathrm{h}.\mathrm{c}.),$$
(3)
where $`f^{}`$ are the gauge eigenstate quark fields, $`r`$ and $`s`$ are general complex matrices and $`\stackrel{~}{\mathrm{\Phi }}=\tau ^2\mathrm{\Phi }^{}\tau ^2`$. This term gives rise to the mass matrices, $`M^u=rk+sk^{}`$ and $`M^d=rk^{}+sk^{}`$. In the mass basis of quarks and leptons the charged current interaction is given by
$$_{CC}=\frac{g_L}{\sqrt{2}}\overline{u}_{iL}\gamma _\mu V_{ij}^Ld_{jL}W_L^{\mu +}+\frac{g_R}{\sqrt{2}}\overline{u}_{iL}\gamma _\mu V_{ij}^Rd_{jL}W_R^{\mu +}+\mathrm{h}.\mathrm{c}.,$$
(4)
where $`V^L`$ and $`V^R`$ are the unitary mixing matrices for the quarks, the elements of which are
$$V^{L,R}=\left(\begin{array}{ccc}V_{ud}^{L,R}& V_{us}^{L,R}& V_{ub}^{L,R}\\ V_{cd}^{L,R}& V_{cs}^{L,R}& V_{cb}^{L,R}\\ V_{td}^{L,R}& V_{ts}^{L,R}& V_{tb}^{L,R}\end{array}\right).$$
(5)
The kinetic term for the Higgs gives a mass structure to the gauge bosons. There are two heavy neutral gauge bosons, the massless photon and charged gauged bosons from the left and right-handed sectors. The non-diagonal mass matrix of the charged left and right gauge bosons is
$$M_W^2=\left(\begin{array}{cc}\frac{1}{2}g_L^2(|v_L|^2+|k|^2+|k^{}|^2)& g_Lg_Rk^{}k^{}\\ g_Lg_Rk^{}k& \frac{1}{2}g_R^2(|v_R|^2+|k|^2+|k^{}|^2)\end{array}\right).$$
(6)
where $`M_L`$ and $`M_R`$ are the upper and lower diagonal elements respectively. This matrix gives the mixing between the mass and the gauge eigenstates. In terms of the mixing angle this is
$$\left(\begin{array}{c}W_L^+\\ W_R^+\end{array}\right)=\left(\begin{array}{cc}\mathrm{cos}\zeta & \mathrm{sin}\zeta \\ e^{i\omega }\mathrm{sin}\zeta & e^{i\omega }\mathrm{cos}\zeta \end{array}\right)\left(\begin{array}{c}W_1^+\\ W_2^+\end{array}\right),$$
(7)
with
$$\mathrm{tan}2\zeta =\frac{2g_Lg_R|k^{}k|}{M_R^2M_L^2}.$$
(8)
In this paper we will often use the following quantities:
$$\zeta _g\frac{g_R}{g_L}\zeta ,\beta _g\frac{g_R^2}{g_L^2}\beta =\frac{g_R^2M_1^2}{g_L^2M_2^2}.$$
(9)
where $`M_1`$ and $`M_2`$ are the eigenvalues of the mass matrix (6). In the case where $`M_RM_L`$ one has $`M_1M_L`$ and $`M_2M_R`$. $`\zeta _g`$ is the mixing parameter which determines the strength of the interactions due to mixing between left-handed and right-handed currents relative to pure left-handed current interactions. $`\beta _g`$ determines the relative strength of right-handed to left-handed interactions.
Gronau and Wakaizumi proposed to modify the mixing matrices such that the third generation of quarks couples to the other generations only through the right-handed $`W`$ bosons . The specific parametrization is
$$V^L=\left(\begin{array}{ccc}\mathrm{cos}\theta _c& \mathrm{sin}\theta _c& 0\\ \mathrm{sin}\theta _c& \mathrm{cos}\theta _c& 0\\ 0& 0& 1\end{array}\right),$$
(10)
for the left-handed mixing matrix, where $`\theta _c`$ is the Cabibbo angle. In the original GW model the right-handed mixing matrix was parametrized by a single angle and CP violation accommodated by a single phase. We choose to study the most general form of the right-handed mixing matrix, of which the GW model is a particular choice of angles and phases. This matrix involves three angles and four phases:
$$V^R=\left(\begin{array}{ccc}c_{12}c_{13}e^{i(\alpha +\beta )}& c_{13}s_{12}e^{i(\gamma +\alpha )}& s_{13}e^{i(\beta +\gamma \alpha )}\\ (c_{12}s_{23}s_{13}+s_{12}c_{23}e^{i\delta })e^{i(\beta \alpha )}& (s_{12}s_{23}s_{13}+c_{12}c_{23}e^{i\delta })e^{i(\gamma \alpha )}& s_{23}c_{13}e^{i(\beta +\gamma +\alpha )}\\ (c_{12}c_{23}s_{13}s_{12}s_{23}e^{i\delta })e^{i\beta }& (s_{12}c_{23}s_{13}c_{12}s_{23}e^{i\delta })e^{i\gamma }& c_{23}c_{13}e^{i(\beta +\gamma )}\end{array}\right),$$
(11)
where $`c_{ij}`$ denotes $`\mathrm{cos}\theta _{ij}`$ and $`s_{ij}`$ denotes $`\mathrm{sin}\theta _{ij}`$. We will denote the model employing (10) and (11) as the generalized GW model. In the following section we will choose the angles in this matrix to satisfy experimental constraints.
## III General Constraints on the Model
The $`b`$ quark decays through pure right-handed couplings in this model. This allows us to use constraints from semi-leptonic $`B`$ decays. In particular the decay $`bcl\nu `$ gives the important relation
$$|V_{cb}^R|(\beta _g^2+\zeta _g^2)^{1/2}=|V_{cb}^{SM}|=0.0360.042,$$
(12)
where SM denotes the standard model value. Results from CLEO measure the asymmetry in the decay $`BD^{}l\nu `$ assuming a pure left-handed lepton current . This puts an upper bound on the ratio
$$\left(\frac{\zeta _g}{\beta _g}\right)^2<0.30.$$
(13)
D0 performed direct searches for $`W^R`$ . For $`g_R=g_L`$ they obtain $`M_{W_R}>720`$ GeV, which in turn implies $`\beta _g<0.012`$. Because $`|V_{cb}^R|<1`$ by unitarity, condition (12) gives $`\beta _g>0.03`$. It is obvious that these two conditions can not simultaneously apply. It has been noted that the form of the right-handed mixing matrix and the size of the ratio $`g_R/g_L`$ affect the lower bound on the mass of $`W_R`$ . In particular for $`|V_{ud}^R|1`$ and $`g_R>g_L`$ one lowers the bound on the $`W_R`$ mass. In the region of parameter space of $`V^R`$ in which $`|V_{ud}^R|1`$ we apply the constraints to be discussed below and find that the neutral $`B_s`$ mass difference is too small to satisfy the experimental lower bound. We are left with the region where $`g_R>g_L`$, which may be unnatural in grand unified models . D0 presents constraints on $`M_R`$ for $`g_R/g_L=\sqrt{2}`$ but for no higher values of the ratio. It is not unreasonable to accept of $`\beta _g`$ in the range $`0.030.04`$ for $`g_R/g_L2`$. In this paper we choose $`\beta _g=0.035`$. Additionally, there is the ratio from semi-leptonic $`B`$ decays
$$\left|\frac{V_{ub}^R}{V_{cb}^R}\right|=\left|\frac{V_{ub}^{SM}}{V_{cb}^{SM}}\right|=0.09\pm 0.03.$$
(14)
This gives the constraint $`s_{13}/s_{23}c_{13}=0.09\pm 0.03`$.
We now turn to the constraints imposed by the mass difference of the neutral mesons $`K`$, $`B_d`$ and $`B_s`$. The neutral $`K`$ mass difference is given as
$$\mathrm{\Delta }m=2eK^0|H_{\mathrm{\Delta }S=2}|\overline{K^0}.$$
(15)
The effective $`\mathrm{\Delta }S=2`$ Hamiltonian arises through the box diagram. In this model there are contributions due to the exchange of two $`W_R`$’s and the exchange of a $`W_R`$ and a $`W_L`$ in addition to the two $`W_L`$ exchange familiar from the standard model. We use the result of Mohapatra et al. for the box diagram without QCD corrections . To calculate the matrix element of quark field operators we use the vacuum insertion approximation with bag factors equal to unity. The neutral $`B_d`$ and $`B_s`$ mass differences are determined by relations similar to (15) and we again use the results . We use estimates of the $`B_d`$ and $`B_s`$ decay constants from the lattice . For the experimental values
$`\mathrm{\Delta }m_K`$ $`=`$ $`3.49\times 10^{12}\mathrm{MeV},`$ (16)
$`\mathrm{\Delta }m_{B_d}`$ $`=`$ $`3.05\times 10^{10}\mathrm{MeV},`$ (17)
$`\mathrm{\Delta }m_{B_s}`$ $`>`$ $`8.16\times 10^9\mathrm{MeV},`$ (18)
we find the following satisfactory choice of angles:
$$\theta _{13}=0.08,\theta _{12}=0.04,\theta _{23}=1.8,$$
(19)
where tuning of $`𝒪(10^2)`$ is necessary for $`\theta _{13}`$ and $`\theta _{12}`$. Tuning of $`𝒪(10^1)`$ is necessary for $`\theta _{23}`$. With these angles we satisfy the experimental conditions within theoretical uncertainties. We will use these values in the remainder of this paper. At this point there are no constraints on the four phases in this model. These will be adjusted by the CP violating observables discussed in the next section.
## IV Constraints from CP Violation
CP violation has been measured in the neutral $`K`$ sector in the form of the observables $`ϵ`$ and $`ϵ^{}`$. In addition, CP violation should give rise to a nonzero electric dipole moment of the neutron. We will now use the measurements of $`ϵ`$, $`ϵ^{}/ϵ`$ and the upper bound on the neutron dipole moment to constrain the phases in our model.
The parameter $`ϵ`$ is related to the $`\mathrm{\Delta }S=2`$ Hamiltonian.
$$ϵ=\frac{e^{i\pi /4}}{\sqrt{2}}\frac{mK^0|H_{\mathrm{\Delta }S=2}|\overline{K^0}}{\mathrm{\Delta }m_K}.$$
(20)
The effective Hamiltonian is that used to calculate the neutral $`K`$ mass splitting. Using this and the choice of angles from the previous section we find several terms of $`𝒪(10^1)`$. The dominant contribution is from the box diagram involving the exchange of two $`W_R`$ and two top quarks. This gives a strong contribution in this model because $`V_{td}^R0.23`$ and $`V_{ts}^R1`$, whereas in the standard model the two top exchange diagram is CKM suppressed. The leading terms in $`ϵ`$ are
$$|ϵ|=\left|0.14\mathrm{cos}2(\beta \gamma )\mathrm{sin}\delta 0.14\mathrm{sin}2(\beta \gamma )\mathrm{cos}\delta +0.23\mathrm{sin}(\beta \gamma \delta )\right|.$$
(21)
With the experimental value of $`(ϵ=2.28\pm 0.02)\times 10^3`$ and assuming no cancellation between terms, this suggests
$$\mathrm{sin}(\beta \gamma )𝒪(10^2),\mathrm{sin}\delta 𝒪(10^2).$$
(22)
There are other terms in the calculation of $`ϵ`$ of order less than or equal to $`10^2`$ which impose no further constraints on the phases.
The parameter $`ϵ^{}`$ is given in terms of the $`K`$ decay amplitude to two pions as
$$ϵ^{}=\frac{e^{i(\frac{\pi }{2}+\delta _2\delta _0)}}{\sqrt{2}}\frac{eA_2}{eA_0}\left(\frac{mA_2}{eA_2}\frac{mA_0}{eA_0}\right),$$
(23)
with
$$A_i=(\pi \pi )_{I=i}|H_{\mathrm{\Delta }S=1}|K^0,$$
(24)
where $`i`$ denotes the isospin channel and $`\delta _i`$ is the hadronic phase shift. The problem of calculating $`ϵ^{}`$ is then to calculate the $`\mathrm{\Delta }S=1`$ Hamiltonian and with this, to estimate the decay amplitudes. Of course, this problem is plagued with hadronic uncertainties. The calculation of $`ϵ^{}`$ in this model is interesting but lengthy. We relegate it to the appendix.
We find in the resulting expression for $`ϵ^{}/ϵ`$ terms proportional to $`\beta _g`$ and $`\zeta _g`$. The terms proportional to $`\beta _g`$ are too small to accommodate the measured value of $`ϵ^{}/ϵ`$ of $`(21.2\pm 4.6)\times 10^4`$. Terms proportional to $`\zeta _g`$ must provide the dominant contribution. We will see the effects of a non-zero $`\zeta _g`$ in the following section. The dominant terms are
$$\left|ϵ^{}/ϵ\right|=\left|\zeta _g\left(1.8R_c^{LR}\mathrm{sin}(\alpha \beta )4.0R_u^{LR}\mathrm{sin}(\alpha +\beta )\right)\right|,$$
(25)
where $`R_u^{LR}`$ and $`R_c^{LR}`$ are ratios of operators defined in the appendix. They are estimated to be $`𝒪(1)`$ and $`𝒪(10^1)`$ respectively. The constraint from $`ϵ^{}/ϵ`$ then requires either small phases or a small $`\zeta _g`$.
The electric dipole moment of the neutron arises in this model at the one loop level due to mixing of the $`W_L`$ and $`W_R`$. In the standard model one loop diagrams do not contribute because they are proportional to the magnitude of CKM elements and so are real. $`W_L`$-$`W_R`$ mixing permits imaginary coefficients in the loop diagram, allowing for a non-zero edm. Electric dipole moments of the $`u`$ and $`d`$ quark arise from diagrams involving the creation of a virtual $`W`$ and the emission of a photon. These contributions have been calculated and are given as
$`d_u`$ $`=`$ $`{\displaystyle \frac{eG_F}{4\pi ^2}}\zeta _g{\displaystyle \underset{j=d,s,b}{}}m_jm(V_{uj}^LV_{uj}^R)f_1\left({\displaystyle \frac{m_j^2}{M_L^2}}\right),`$ (26)
$`d_d`$ $`=`$ $`{\displaystyle \frac{eG_F}{4\pi ^2}}\zeta _g{\displaystyle \underset{j=u,c,t}{}}m_jm(V_{jd}^LV_{jd}^R)f_2\left({\displaystyle \frac{m_j^2}{M_L^2}}\right),`$ (27)
where $`f_1`$ and $`f_2`$ are dimensionless functions of the quark masses. In addition to the loop diagram there is also a contribution from the exchange of a mixed $`W_L`$-$`W_R`$ from the $`u`$ to the $`d`$ with the emission of a photon. This contribution has large hadronic uncertainties and is estimated in the harmonic oscillator parton model of the neutron . It is given by
$$d_{ex}=\frac{eG_F}{3\pi ^{3/2}}\zeta _g\sqrt{2m_q\omega }(1\beta _g)m(V_{ud}^LV_{ud}^R),$$
(28)
where we use $`\sqrt{m_q\omega }=0.3`$ GeV. The edm of the neutron is related to these contributions by
$$d_n=\frac{4}{3}d_d\frac{1}{3}d_u+d_{ex},$$
(29)
Evaluated with our choice of angles we find
$$d_n=\zeta _g\left(3.9\times 10^{21}\mathrm{sin}(\alpha +\beta )2.3\times 10^{22}\mathrm{sin}(\alpha +\gamma )\right)\text{e-cm}$$
(30)
The experimental upper bound is $`d_n<2.6\times 10^{25}`$ e-cm. This constraint could be accommodated by a small $`\zeta _g𝒪(10^5)`$. However, this would make $`ϵ^{}/ϵ`$ too small. Together the edm and $`ϵ^{}/ϵ`$ require
$$\zeta _g𝒪(10^2),\mathrm{sin}(\alpha +\beta )𝒪(10^3).$$
(31)
We see that $`\zeta _g`$ must be close to the upper bound of (13).
## V Discussion on $`B`$ decays
We have now found a satisfactory albeit tightly constrained region in which the generalized GW model is valid. If we make the simplifying assumption
$$\beta =\gamma =\alpha ,\delta =\pi ,$$
(32)
which is consistent with the constraints and necessarily correct to at least $`𝒪(10^2)`$, the right-handed mixing matrix becomes
$$V^R=\left(\begin{array}{ccc}0.996& 0.0399& 0.0799e^{i3\beta }\\ 0.0862e^{i2\beta }& 0.195e^{i2\beta }& 0.977e^{i\beta }\\ 0.233e^{i\beta }& 0.979e^{i\beta }& 0.198e^{i2\beta }\end{array}\right).$$
(33)
Although we found no others in our search, we do not suggest that this is the only possible choice of the seven free parameters in this matrix that satisfy the experimental constraints. We merely point out that this particular choice is phenomenologically acceptable and as such the GW ansatz of pure right-handed $`b`$ decays is not completely ruled out. We now examine some consequences of this choice of parameters.
The ratio of branching ratios
$$R=\frac{\mathrm{Br}(B^{}\psi \pi ^{})}{\mathrm{Br}(B^{}\psi K^{})}=0.052\pm 0.024.$$
(34)
has been measured. In the limit of dominant right-handed tree contributions to the decay we have
$$R\left|\frac{V_{cd}^R}{V_{cs}^R}\right|^2=0.2,$$
(35)
where the ratio has been evaluated according to (33). To what degree should we trust the discrepancy here between experiment and our model? Certainly the ratios are the same within an order of magnitude. Penguin contributions will affect the theoretical prediction. There will also be a strong contribution due to $`W_L`$-$`W_R`$ mixing because of the relatively large value of $`\zeta _g`$. To estimate the mixing contribution in the decay $`B^{}\psi K^{}`$ we look at the ratio of mixed to unmixed tree level contributions
$$\frac{\zeta _g}{\beta _g}\times \frac{V_{cs}^L}{V_{cs}^R}\times \frac{O^{LR}}{O^{RR}}3.5,$$
(36)
where we have used the upper bound (13) and set the ratio of matrix elements
$$\frac{O^{LR}}{O^{RR}}=\frac{K^{}\psi |\overline{b}\gamma ^\mu (1+\gamma ^5)c\overline{c}\gamma _\mu (1\gamma ^5)s|B^{}}{K^{}\psi |\overline{b}\gamma ^\mu (1+\gamma ^5)c\overline{c}\gamma _\mu (1+\gamma ^5)s|B^{}}=1.4,$$
(37)
found in the vacuum insertion approximation. There is a substantial and possibly dominant contribution to the decay due to mixing. Our choice of phases and angles is then consistent with the ratio (34).
In the neutral $`B`$ meson system one can write the physical mass eigenstates in terms of the gauge eigenstates as
$$|B_{1,2}^0=p|B^0\pm q|\overline{B^0}.$$
(38)
Their decay amplitudes are
$`A`$ $`=`$ $`f|H|B^0\mathrm{and}`$ (39)
$`\overline{A}`$ $`=`$ $`f|H|\overline{B^0},`$ (40)
where $`f`$ is a CP eigenstate. If there is a single dominant decay process (e.g. no strong penguin processes), then the decay asymmetry becomes
$$a_f=\frac{\mathrm{\Gamma }(B_{phys}^0f)\mathrm{\Gamma }(\overline{B}_{phys}^0f)}{\mathrm{\Gamma }(B_{phys}^0f)+\mathrm{\Gamma }(\overline{B}_{phys}^0f)}m\left(\frac{q\overline{A}}{pA}\right),$$
(41)
where $`B_{phys}^0(\overline{B}_{phys}^0)`$ denotes the time evolved $`B^0(\overline{B^0})`$ meson. If the final state, $`f`$, is not a CP eigenstate but a neutral meson such as $`K_S`$, then it also contributes a mixing term to the asymmetry.
In the absence of mixing $`(\zeta _g=0)`$ all $`B`$ decays occur through pure right-handed interactions. There are both tree and penguin diagrams contributing to the decay. However, due to the phase constraints imposed on this model, tree diagrams and penguin diagrams contribute with the same phase up to corrections to which (32) holds. This allows for clean asymmetry predictions which have been previously discussed .
However, in a situation in which mixing is large it is necessary to consider pollution from mixing. Previous work is no longer applicable to this situation. In the case of $`B_d^0\psi K_S`$, we see that the ratio of left-right to right-right tree amplitudes is given by
$$\left|\frac{T^{LR}}{T^{RR}}\right|=\left|\frac{V_{cb}^RV_{cs}^L}{V_{cb}^RV_{cs}^R}\right|\times \left|\frac{\zeta _g}{\beta _g}\right|\times \left|\frac{O^{LR}}{O^{RR}}\right|,$$
(42)
where
$`O^{RR}`$ $`=`$ $`K_S\psi |\overline{b}\gamma ^\mu (1+\gamma ^5)c\overline{c}\gamma _\mu (1+\gamma ^5)s|B^0,`$ (43)
$`O^{LR}`$ $`=`$ $`K_S\psi |\overline{b}\gamma ^\mu (1+\gamma ^5)c\overline{c}\gamma _\mu (1\gamma ^5)s|B^0.`$ (44)
In the vacuum insertion approximation we find $`O^{LR}/O^{RR}`$ to be $`𝒪(1)`$. With the upper bound (13) this gives
$$\frac{T^{LR}}{T^{RR}}3.55.$$
(45)
Pollution is possibly over 100% and this decay ceases to be predictive. (It is clean in the standard model due to CKM suppression of the penguins.)
In the same way we examine the decays $`B_dD_1^0\pi ^0,D^+D^{},K_S\pi ^0,\varphi K_S,K_SK_S`$ and $`B_sD_S^+D_S^{},D_1^0K_S,\psi K_S,\rho ^0K_S,K^+K^{},\eta ^{}\eta ^{},\varphi K_S`$. For pure penguin decay processes we determine the mixing contribution by assuming that the ratio of left-right to right-right matrix elements is $`𝒪(1)`$. For all of these decays we find pollution due to mixing on the order of 100%. There is little predictive power left from CP asymmetries in the neutral $`B`$ sector. We stress that although disappointing this is a new result for models employing the GW ansatz. This model is not inconsistent with any values of the various CP asymmetries in the neutral $`B`$ sector. In fact, this model is consistent with no correlations of any kind among the phases in these decays.
## VI Conclusions
In this $`SU(2)_L\times SU(2)_R\times U(1)`$ model where the third generation interacts weakly through pure right-handed couplings the parameters are highly constrained. Nevertheless, we have found a region in parameter space in which this model is consistent with measurements of the neutral meson mass differences $`\mathrm{\Delta }m_K`$, $`\mathrm{\Delta }m_{B_d}`$ and $`\mathrm{\Delta }m_{B_s}`$ and semi-leptonic $`B`$ decays. We find that it is necessary that the ratio of coupling constants, $`g_R/g_L`$, be on the order of two. Constraining the phases with CP violating observables leads to a second undesirable result. There are three fine tuning conditions on the four phases in the right-handed mixing matrix.
Constraints from CP violating observables also require that mixing between the left and right-handed $`W`$’s is not small, but $`\zeta _g𝒪(10^2)`$. This leads to pollution in CP asymmetries in $`B`$ decays to CP eigenstates on the order of 100%. There are no definite predictions or clean phase measurements in these decays. If discrepancies between experiment and the standard model are found in the $`B`$ decay asymmetries this model can not be ruled out. However, increasing the lower bound on the right-handed $`W`$ from direct searches and a more stringent limit on $`\zeta _g`$ could decisively determine the fate of this model.
###### Acknowledgements.
I thank Adam Falk for introducing the topic, helpful direction and many useful discussions. This work was supported by the United States National Science Foundation under Grant No. PHY-9404057 and by an Owen Fellowship of the Johns Hopkins University, Zanvyl Krieger School of Arts and Sciences.
## A Calculation of $`ϵ^{}/ϵ`$
The expression (23) relates $`ϵ^{}`$ to $`A_I`$, the neutral $`K`$ decay amplitude to pions with isospin $`I`$, given in (24). To calculate $`ϵ^{}/ϵ`$ we will use this expression employing an isoconjugate simplification due to Mohapatra and Pati . In this procedure the $`\mathrm{\Delta }S=1`$ weak decay Hamiltonian is decomposed into scalar ($`S`$) and pseudoscalar terms ($`P`$) which are further decomposed into CP even and odd components (denoted by superscript $`+`$ or $``$).
$$H^{\mathrm{\Delta }S=1}=S^++S^{}+P^++P^{}.$$
(A1)
If a relationship can be found such that
$$[I_3,P^{}]=i\alpha P^+,$$
(A2)
where $`\alpha `$ is a real constant, then it can be shown that the ratio $`mA_I/eA_I`$ is independent of $`I`$ and $`ϵ^{}=0`$.
To see this notice that with $`|K_{1,2}`$ as CP eigenstates, $`I_3|K_1=1/2|K_2`$ and $`I_3|\pi ^i\pi ^j=0`$ where $`(i,j)`$ denote $`(+,)`$ or $`(0,0)`$. Now
$$\pi ^i\pi ^j|P^{}|K_2=i\alpha \pi ^i\pi ^j|P^+|K_1$$
(A3)
holds independent of $`i`$ and $`j`$. The amplitude, $`A_I`$, can then be written as a real factor containing matrix elements multiplied by a complex factor independent of $`I`$. The matrix elements cancel in $`mA_I/eA_I`$ and we have the desired result.
In the GW model the $`\mathrm{\Delta }S=1`$ Hamiltonian can be split into three terms pertaining to the tree and penguin amplitudes for pure left-handed couplings, pure right-handed couplings and mixed couplings. This can be written as
$$H=(T^{LL}+P^{LL})+(T^{RR}+P^{RR})+(T^{LR}+P^{LR}),$$
(A4)
where $`T`$ and $`P`$ denote tree and penguin contributions respectively. Because of the pure right-handed nature of the third generation couplings, of the penguin diagrams only $`P^{RR}`$ has a contribution from the top quark. The elements have been calculated . They are
$`T^{LL}`$ $`=`$ $`{\displaystyle \frac{4G_F}{\sqrt{2}}}(V_{ud}^LV_{us}^L)\overline{s}_L\gamma ^\mu u_L\overline{u}_L\gamma _\mu d_L+\mathrm{h}.\mathrm{c}.,`$ (A5)
$`T^{RR}`$ $`=`$ $`{\displaystyle \frac{4G_F}{\sqrt{2}}}\beta _g(V_{ud}^RV_{us}^R)\overline{s}_R\gamma ^\mu u_R\overline{u}_R\gamma _\mu d_R+\mathrm{h}.\mathrm{c}.,`$ (A6)
$`T^{LR}`$ $`=`$ $`{\displaystyle \frac{4G_F}{\sqrt{2}}}\zeta _g\left(V_{ud}^LV_{us}^R\overline{s}_R\gamma ^\mu u_R\overline{u}_L\gamma _\mu d_L+V_{ud}^RV_{us}^L\overline{s}_L\gamma ^\mu u_L\overline{u}_R\gamma _\mu d_R\right)+\mathrm{h}.\mathrm{c}.,`$ (A7)
$`P^{LL}`$ $`=`$ $`{\displaystyle \frac{4G_F}{\sqrt{2}}}{\displaystyle \frac{\alpha _s(\mu )}{24\pi }}\left({\displaystyle \underset{q=u,c}{}}V_{qd}^LV_{qs}^Lf\left({\displaystyle \frac{m_q^2}{M_L^2}}\right)\right)(\overline{u}\gamma _\mu \tau ^au+\overline{d}\gamma _\mu \tau ^ad)\overline{s}_L\gamma ^\mu \tau ^ad_L+\mathrm{h}.\mathrm{c}.,`$ (A8)
$`P^{RR}`$ $`=`$ $`{\displaystyle \frac{4G_F}{\sqrt{2}}}\beta _g{\displaystyle \frac{\alpha _s(\mu )}{24\pi }}\left({\displaystyle \underset{q=u,c,t}{}}V_{qd}^RV_{qs}^Rf\left({\displaystyle \frac{m_q^2}{M_R^2}}\right)\right)(\overline{u}\gamma _\mu \tau ^au+\overline{d}\gamma _\mu \tau ^ad)\overline{s}_R\gamma ^\mu \tau ^ad_R+\mathrm{h}.\mathrm{c}.`$ (A9)
$`P^{LR}`$ $`=`$ $`{\displaystyle \frac{4G_F}{\sqrt{2}}}\zeta _g{\displaystyle \frac{\alpha _s(\mu )}{8\pi }}(\overline{u}\gamma ^\mu \tau ^au+\overline{d}\gamma ^\mu \tau ^ad){\displaystyle \frac{k^\nu }{k^2}}`$ (A11)
$`\times \overline{s}i\sigma _{\mu \nu }\left({\displaystyle \underset{q=u,c}{}}(V_{qd}^LV_{qs}^R\gamma _L+V_{qd}^RV_{qs}^L\gamma _R)g\left({\displaystyle \frac{m_q^2}{M_L^2}}\right)m_q\right)\tau ^ad+\mathrm{h}.\mathrm{c}.,`$
where $`f`$ and $`g`$ are dimensionless functions of quark masses and $`\gamma _{R/L}=(1\pm \gamma ^5)/2`$.
These terms are now separated into scalar and pseudoscalar components. In dealing with kaon decays to two pions only the pseudoscalar terms are relevant. We want to decompose the Hamiltonian into a part which satisfies an isoconjugate relation, $`H_0`$, and a part which accounts for the nonzero $`ϵ^{}`$. To do this we use the unitarity relations,
$$V_{ud}^LV_{us}^L=V_{cd}^LV_{cs}^L,$$
(A12)
$$V_{cd}^RV_{cs}^R=V_{ud}^RV_{us}^RV_{td}^RV_{ts}^R,$$
(A13)
to remove these two factors from the Hamiltonian. The pseudoscalar part of the Hamiltonian which satisfies an isoconjugate relation can be written as
$$H_0=(V_{ud}^LV_{us}^L+\beta _gV_{ud}^RV_{us}^R)P+\mathrm{h}.\mathrm{c}.,$$
(A14)
where
$`P={\displaystyle \frac{G_F}{\sqrt{2}}}[\overline{s}\gamma ^\mu \gamma ^5u\overline{u}\gamma _\mu d+\overline{s}\gamma ^\mu u\overline{u}\gamma _\mu \gamma ^5d`$ (A16)
$`+{\displaystyle \frac{\alpha _s(\mu )}{12\pi }}(\overline{u}\gamma _\mu \tau ^au+\overline{d}\gamma ^\mu \tau ^ad)\overline{s}\gamma _\mu \gamma ^5\tau ^ad(f\left({\displaystyle \frac{m_u^2}{M_L^2}}\right)f\left({\displaystyle \frac{m_c^2}{M_L^2}}\right))].`$
We have used the fact that the relation,
$$f\left(\frac{m_u^2}{M_L^2}\right)f\left(\frac{m_c^2}{M_L^2}\right)=f\left(\frac{m_u^2}{M_R^2}\right)f\left(\frac{m_c^2}{M_R^2}\right),$$
(A17)
holds up to negligible corrections of $`𝒪(10^4)`$. Splitting this term and its Hermitian conjugate into CP even and odd states, $`P^+`$ and $`P^{}`$ respectively, we arrive at the relationship
$$[I_3,P^{}]=\frac{im(V_{ud}^LV_{us}^L+\beta _gV_{ud}^RV_{us}^R)}{2e(V_{ud}^LV_{us}^L+\beta _gV_{ud}^RV_{us}^R)}P^+.$$
(A18)
We now examine the mixing terms and the term proportional to $`V_{td}^RV_{ts}^R`$. We define the operators
$`O^\pm `$ $`=`$ $`\overline{s}\gamma ^\mu \gamma ^5u\overline{u}\gamma _\mu d\pm \overline{s}\gamma ^\mu u\overline{u}\gamma _\mu \gamma ^5d,`$ (A19)
$`O^5`$ $`=`$ $`\left(\overline{u}\gamma _\mu \tau ^au+\overline{d}\gamma _\mu \tau ^ad\right)\overline{s}\gamma ^\mu \gamma ^5\tau ^ad,`$ (A20)
$`O_{LR}^P`$ $`=`$ $`{\displaystyle \frac{k^\nu }{k^2}}\overline{s}i\sigma _{\mu \nu }\gamma ^5\tau ^ad(\overline{u}\gamma ^\mu \tau ^au+\overline{d}\gamma ^\mu \tau ^ad).`$ (A21)
The dominant contributions to the ratios $`mA_I/eA_I`$ cancel in the difference of isospin channels. In terms of the above operators the following ratios are needed.
$`R^P`$ $`=`$ $`{\displaystyle \frac{\frac{\alpha _s(\mu )}{12\pi }\left(f\left(\frac{m_t^2}{M_R^2}\right)f\left(\frac{m_c^2}{M_R^2}\right)\right)O_{1/2}^5}{O_{1/2}^++\frac{\alpha _s(\mu )}{12\pi }\left(f(\frac{m_u^2}{M_L^2})f\left(\frac{m_c^2}{M_L^2}\right)\right)O_{1/2}^5}},`$ (A22)
$`R_u^{LR}`$ $`=`$ $`{\displaystyle \frac{O_{1/2}^{}\frac{\alpha _s(\mu )}{4\pi }m_qg(\frac{m_q^2}{M_L^2})O_{LR1/2}^P}{O_{1/2}^++\frac{\alpha _s(\mu )}{12\pi }\left(f\left(\frac{m_u^2}{M_L^2}\right)f\left(\frac{m_c^2}{M_L^2}\right)\right)O_{1/2}^5}}{\displaystyle \frac{O_{3/2}^{}}{O_{3/2}^+}},`$ (A23)
$`R_c^{LR}`$ $`=`$ $`{\displaystyle \frac{\frac{\alpha _s(\mu )}{4\pi }m_qg(\frac{m_q^2}{M_L^2})O_{LR1/2}^P}{O_{1/2}^++\frac{\alpha _s(\mu )}{12\pi }\left(f\left(\frac{m_u^2}{M_L^2}\right)f\left(\frac{m_c^2}{M_L^2}\right)\right)O_{1/2}^5}},`$ (A24)
where we use $`O_{\mathrm{\Delta }I}=\pi \pi _I|O|K`$. To first order in $`\zeta _g`$ and $`\beta _g`$, $`ϵ^{}`$ is given by
$`ϵ^{}`$ $`=`$ $`{\displaystyle \frac{\omega }{\sqrt{2}}}{\displaystyle \frac{1}{V_{ud}^LV_{us}^L}}(\beta _gm(V_{td}^RV_{ts}^R)R^P`$ (A26)
$`+\zeta _gm(V_{ud}^LV_{us}^RV_{ud}^RV_{us}^L)R_u^{LR}+\zeta _gm(V_{cd}^LV_{cs}^RV_{ud}^RV_{us}^L)R_c^{LR}),`$
where $`\omega =eA_2/eA_01/20`$.
There are hadronic uncertainties in the ratios of matrix elements. To restrict the parameters in this model we need to obtain order of magnitude estimates of these ratios at the least. Assuming an $`𝒪(1)`$ estimate for $`O^5/O^+`$ and $`O^{}/O^+`$ we obtain $`R^P𝒪(10^1)`$. Assuming $`k^\nu /k^21\mathrm{GeV}^1`$ we estimate $`R_c^{LR}𝒪(10^1)`$ and $`R_u^{LR}𝒪(1)`$.
Evaluating the coefficients of the operators after imposing (22), the constraint from $`ϵ`$, we find the following expression.
$$ϵ^{}/ϵ=\zeta _g\left(1.8R_c^{LR}\mathrm{sin}(\alpha \beta )4.0R_u^{LR}\mathrm{sin}(\alpha +\beta )\right)$$
(A27)
We use this result in section IV. |
warning/0002/hep-th0002142.html | ar5iv | text | # Logarithmic Conformal Field Theories Near a Boundary
## 1 Introduction
Gurarie has pointed out that it is possible to construct consistent conformal field theories which have logarithmic terms in their correlation functions. In these theories representation of the Virasoro algebra is not diagonalizable. Logarithmic conformal field theories (LCFTs) contain logarithmic operators which have logarithms as well as powers in their operator product expansion. Such operators do not appear in unitary CFTs. However this does not mean that they are not relevant. By now a large number of examples of the applications of such theories have been found. For some examples see citations in . Extensive work has been done on the structure of LCFTs , although a lot remains to be worked out.
Conformal invariance in physical systems arises at a fixed point. It is also known that critical behavior near boundaries is affected by the geometry of the boundary. Indeed correlators of conformal field theories near boundaries have extra structure which gives rise to the surface critical behavior . Conformal theories of turbulence near a boundary were discussed by Chung et al .
In this paper we investigate such theories near a boundary. For definiteness, we consider a semi-infinite $`d`$-dimensional system bounded by a ($`d1`$)-dimensional plane surface. The results of this paper fall into two classes. For arbitrary dimensions $`d`$, we show that conformal invariance determines the one-point functions up to some constants and restricts the form of two-point functions up to some unknown functions of a single scaling variable. In two-dimensions our result is much stronger. We show that in this case the n-point functions in the semi-infinite geometry satisfy the same equations as the 2n-point functions in the bulk. In the end, we apply the results derived in sections II and III, to the 2D MHD problem and find the mean values of velocity and magnetic field.
## 2 Arbitrary Dimensions
In logarithmic conformal field theories, there exist blocks of fields which constitute a non-diagonalizable representation of the conformal group. Under a conformal transformation $`𝐫𝐫^{}=f(𝐫)`$ fields of a block trasmform as :
$$\mathrm{\Phi }^{}(𝐫^{})=\left|\frac{𝐫^{}}{𝐫}\right|^T\mathrm{\Phi }(𝐫)$$
(1)
where
$$\mathrm{\Phi }=\left(\begin{array}{c}\varphi _1\\ \varphi _2\\ \mathrm{}\\ \varphi _n\end{array}\right)$$
(2)
and T is Jordanian matrix
$$T=\left(\begin{array}{cccc}\frac{\mathrm{\Delta }_\varphi }{d}& 0& \mathrm{}& 0\\ 1& \frac{\mathrm{\Delta }_\varphi }{d}& \mathrm{}& \mathrm{}\\ 0& 1& \mathrm{}& 0\\ \mathrm{}& \mathrm{}& \mathrm{}& 0\\ 0& \mathrm{}& 1& \frac{\mathrm{\Delta }_\varphi }{d}\end{array}\right).$$
(3)
Here $`d`$ is dimension of space-time and $`\mathrm{\Delta }_\varphi `$ is conformal dimension of $`\varphi _i`$’s. As an example, for a two-dimensional Jordan cell $`\mathrm{\Phi }=\varphi _1`$ and $`\mathrm{\Psi }=\varphi _2`$, equation (1) reduces to:
$`\mathrm{\Phi }^{}(𝐫^{})`$ $`=`$ $`\left|{\displaystyle \frac{𝐫^{}}{𝐫}}\right|^{\frac{\mathrm{\Delta }_\varphi }{d}}\mathrm{\Phi }(𝐫)`$
$`\mathrm{\Psi }^{}(𝐫^{})`$ $`=`$ $`\left|{\displaystyle \frac{𝐫^{}}{𝐫}}\right|^{\frac{\mathrm{\Delta }_\varphi }{d}}\left(\mathrm{\Psi }(𝐫)+\mathrm{log}|{\displaystyle \frac{𝐫^{}}{𝐫}}|\mathrm{\Phi }(𝐫)\right).`$ (4)
In $`d`$ dimensions, invariance under conformal transformations determines two- and three-point functions up to some constants and four-point functions are determined up to functions of crossing ratios. These invariances exist only if the geometry of the problem is not changed by the transformations. So if there is a boundary in the problem, correlation functions are invariant only under transformations which preserve the boundary. In this case we have less symmetries and correlation functions are not easily determined . To be more specific, we investigate correlation functions of a two-dimensional Jordanian cell with a simple boundary – a $`(d1)`$-dimensional hyper-plane. The symmetry group that preserves the geometry is made up of translations along the boundary, rotations in the boundary, dilatation and special conformal transformation along the boundary. We take as the boundary, the hyper-plane $`y=0`$.
Consider the one-point functions $`\mathrm{\Phi }(𝐫)=f_1(𝐫)`$ , $`\mathrm{\Psi }(𝐫)=f_2(𝐫)`$ near the boundary. Invariance under translations along the surface implies that $`f_1`$ and $`f_2`$ depend only on y. Also under an infinitesimal dilatation transformation:
$$𝐫^{}=𝐫+ϵ𝐫$$
(5)
we must have:
$`f_1(y)`$ $`=`$ $`(1+ϵ)^{\mathrm{\Delta }_\varphi }f_1(y^{}),`$
$`f_2(y)`$ $`=`$ $`(1+ϵ)^{\mathrm{\Delta }_\varphi }(f_2(y^{})+\mathrm{log}(1+ϵ)^df(y^{})).`$ (6)
Expanding these equations to the first order in $`ϵ`$, $`f_1`$ and $`f_2`$ satisfy the differential equations:
$`y{\displaystyle \frac{f_1}{y}}+\mathrm{\Delta }_\varphi f_1`$ $`=`$ $`0`$
$`y{\displaystyle \frac{f_2}{y}}+\mathrm{\Delta }_\varphi f_2+df`$ $`=`$ $`0`$ (7)
which yield:
$`\mathrm{\Phi }(𝐫)`$ $`=`$ $`{\displaystyle \frac{C_1}{y^{\mathrm{\Delta }_\varphi }}}`$
$`\mathrm{\Psi }(𝐫)`$ $`=`$ $`{\displaystyle \frac{1}{y^{\mathrm{\Delta }_\varphi }}}(C_2dC_1\mathrm{log}y)`$ (8)
Invariance under rotation and special conformal transformations reveals no other constraint on these one-point functions.
Now consider the two-point function $`G_1(𝐫_1,𝐫_2)=\mathrm{\Phi }(𝐫_1)\mathrm{\Phi }(𝐫_2)`$ . The points $`𝐫_1`$ and $`𝐫_2`$ lie in a unique plane perpendicular to the surface. Within this plane we can specify the points by coordinates $`(x_1,y_1)`$, $`(x_2,y_2)`$. By translational invariance parallel to the surface, $`G_1`$ depends on $`x_1x_2`$, $`y_1`$ and $`y_2`$. Invariance under dilatation implies:
$$G_1(x_1x_2,y_1,y_2)=(1+ϵ)^{\mathrm{\Delta }_\varphi }(1+ϵ)^{\mathrm{\Delta }_\varphi }G(x_1^{}x_2^{},y_1^{},y_2^{}).$$
(9)
An infinitesimal special conformal transformation along x is
$`x^{}`$ $`=`$ $`x+ϵ(x^2y^2)`$
$`y^{}`$ $`=`$ $`y+2ϵxy,`$ (10)
and under such a transformation we have
$$G_1(x_1x_2,y_1,y_2)=(1+2ϵx_1)^{\mathrm{\Delta }_\varphi }(1+2ϵx_2)^{\mathrm{\Delta }_\varphi }G(x_1^{}x_2^{},y_1^{},y_2^{}).$$
(11)
By expanding equations (9) and (11) to the first order in $`ϵ`$, we arrive at :
$`u{\displaystyle \frac{G_1}{u}}+y_1{\displaystyle \frac{G_1}{y_1}}+y_2{\displaystyle \frac{G_1}{y_2}}+2\mathrm{\Delta }_\varphi G_1`$ $`=`$ $`0`$
$`(y_1^2y_2^2){\displaystyle \frac{G_1}{u}}+u\left(y_1{\displaystyle \frac{G_1}{y_1}}y_2{\displaystyle \frac{G_1}{y_2}}\right)`$ $`=`$ $`0`$ (12)
in which $`u=x_1x_2`$. The first equation states that $`G_1`$ is a homogeneous function of dimension $`2\mathrm{\Delta }_\varphi `$:
$$G_1=\frac{1}{(u)^{2\mathrm{\Delta }_\varphi }}g_1(\alpha ,\beta )$$
(13)
where $`\alpha =y_1/u`$ and $`\beta =y_2/u`$. Substituting this in the second line of equation (12) one finds
$$\left[\alpha +\frac{\alpha }{\alpha ^2\beta ^2}\right]\frac{g_1}{\alpha }+\left[\beta +\frac{\beta }{\beta ^2\alpha ^2}\right]\frac{g_1}{\beta }+2\mathrm{\Delta }_\varphi g_1=0.$$
(14)
The general solution of this equation is:
$$g_1(\alpha ,\beta )=\frac{1}{(\alpha \beta )^{\mathrm{\Delta }_\varphi }}h_1\left(\frac{1+(\alpha \beta )^2}{\alpha \beta }\right).$$
(15)
So the two-point correlation function is:
$$\mathrm{\Phi }(𝐫_1)\mathrm{\Phi }(𝐫_2)=\frac{1}{(y_1y_2)^{\mathrm{\Delta }_\varphi }}h_1\left(\frac{(x_1x_2)^2+(y_1y_2)^2}{y_1y_2}\right).$$
(16)
We have two other two-point functions, $`G_2(𝐫_1,𝐫_2)=\mathrm{\Phi }(𝐫_1)\mathrm{\Psi }(𝐫_2)`$ and $`G_3(𝐫_1,𝐫_2)=\mathrm{\Psi }(𝐫_1)\mathrm{\Psi }(𝐫_2)`$. We can follow similar steps for these two and the result is
$`G_2(𝐫_1,𝐫_2)`$ $`=`$ $`{\displaystyle \frac{1}{(y_1y_2)^{\mathrm{\Delta }_\varphi }}}\left[h_2(\eta )d\mathrm{log}y_2h_1(\eta )\right]`$
$`G_3(𝐫_1,𝐫_2)`$ $`=`$ $`{\displaystyle \frac{1}{(y_1y_2)^{\mathrm{\Delta }_\varphi }}}\left[h_3(\eta )d\mathrm{log}y_1y_2h_2(\eta )+d^2\mathrm{log}y_1\mathrm{log}y_2h_1(\eta )\right]`$ (17)
where $`\eta =[(x_1x_2)^2+(y_1y_2)^2]/y_1y_2`$.
Far from the boundary, the effect of boundary becomes negligible and we must recover the bulk two-point functions:
$`\mathrm{\Phi }(𝐫_1)\mathrm{\Phi }(𝐫_2)`$ $`=`$ $`0`$
$`\mathrm{\Phi }(𝐫_1)\mathrm{\Psi }(𝐫_2)`$ $`=`$ $`{\displaystyle \frac{a}{r^{2\mathrm{\Delta }_\varphi }}}`$
$`\mathrm{\Psi }(𝐫_1)\mathrm{\Psi }(𝐫_2)`$ $`=`$ $`{\displaystyle \frac{1}{r^{2\mathrm{\Delta }_\varphi }}}(bda\mathrm{log}r)`$ (18)
where $`r=|𝐫_1𝐫_2|`$ and $`a,b`$ are arbitrary constants. These equations were first derived by in two dimensions and were generalized to d-dimensions by .
To go far from the boundary one must let $`y_1`$ and $`y_2`$ tend to infinity keeping $`y_1y_2`$ and $`x_1x_2`$ finite. This means letting $`\eta `$ tend to zero . So we can find $`h_1`$ , $`h_2`$ , $`h_3`$ in the limit $`\eta 0`$ :
$`h_1(\eta )`$ $`=`$ $`{\displaystyle \frac{1}{\eta ^{\mathrm{\Delta }_\varphi }}}\left({\displaystyle \frac{4\frac{a}{d}}{\mathrm{log}\eta }}+{\displaystyle \frac{C_1}{(\mathrm{log}\eta )^2}}+\mathrm{}\right)`$
$`h_2(\eta )`$ $`=`$ $`{\displaystyle \frac{1}{\eta ^{\mathrm{\Delta }_\varphi }}}\left(a+{\displaystyle \frac{C_2}{(\mathrm{log}\eta )}}+\mathrm{}\right)`$
$`h_3(\eta )`$ $`=`$ $`{\displaystyle \frac{1}{\eta ^{\mathrm{\Delta }_\varphi }}}\left(bdC_2{\displaystyle \frac{d^2}{4}}C_1+\mathrm{}\right)`$ (19)
Here $`C_1`$ and $`C_2`$ are arbitrary constants. On the other hand the behaviour of these functions when $`\eta `$ tends to infinity, determines the surface behaviour of correlation functions (To investigate the surface behaviour one must let $`x_1x_2`$ tend to infinity while keeping $`y_1`$ and $`y_2`$ finite).
To find the surface exponents one should know the behaviour of these functions in this limit . This requires a knowledge of the differential equations governing $`h(\eta )`$ which requires details of the structure of conformal field theory.
## 3 Two Dimensions
In two dimensions, however, conformal group is an infinite dimensional group and any analytic function from the plane to itself is a conformal transformation . In LCFT’s, under a conformal transformation $`zw(z)`$ , $`\overline{z}w(\overline{z})`$ primary fields of a Jordanian cell, $`\mathrm{\Phi }_i`$’s, transform as :
$$\mathrm{\Phi }_i(z,\overline{z})\mathrm{\Phi }_i(z,\overline{z})+\left[\alpha ^{}(z)\mathrm{\Delta }_i^j+\delta _i^j\alpha (z)\frac{}{z}+\overline{\alpha ^{}(z)}\overline{\mathrm{\Delta }}_i^j+\delta _i^j\overline{\alpha (z)}\frac{}{\overline{z}}\right]\mathrm{\Phi }_j(z,\overline{z})$$
(20)
The effect of a small transformation may be expressed in terms of correlation functions of the fields with the energy-momentum tensor. In complex coordinates, there are two non-zero components, namely $`T(z)=T_{zz}(z)`$ and $`\overline{T}(\overline{z})=T_{\overline{z}\overline{z}}(\overline{z})`$. In terms of these two we have:
$`{\displaystyle \frac{1}{2\pi i}}{\displaystyle _c}𝑑z\alpha (z)T(z)\mathrm{\Phi }_{i_1}(z_1,\overline{z_1})\mathrm{}{\displaystyle \frac{1}{2\pi i}}{\displaystyle _c}𝑑\overline{z}\overline{\alpha (z)}\overline{T}(\overline{z})\mathrm{\Phi }_{i_1}(z_1,\overline{z}_1)\mathrm{}`$
$`={\displaystyle \underset{k}{}}{\displaystyle \underset{j_k}{}}\left[\alpha ^{}(z_k)\mathrm{\Delta }_{i_k}^{j_k}+\delta _{i_k}^{j_k}\alpha (z_k){\displaystyle \frac{}{z_k}}+\overline{\alpha ^{}(z_k)}\overline{\mathrm{\Delta }}_{i_k}^{j_k}+\delta _{i_k}^{j_k}\overline{\alpha (z_k)}{\displaystyle \frac{}{\overline{z}_k}}\right]\mathrm{\Phi }_{j_1}(z_1,\overline{z}_1)\mathrm{}`$ (21)
where $`c`$ is an arbitrary contour containing all the points $`z_k`$ and $`k`$ is summed over the fields in the correlators and $`j_k`$ is summed over the Jordanian cell containing $`\mathrm{\Phi }_{i_k}`$. In the absence of a boundary, $`\alpha (z)`$ is an arbitrary analytic function, thus $`\alpha `$ and $`\overline{\alpha }`$ can be assumed independent. As a consequence the $`z`$ and $`\overline{z}`$ dependence in equation (21) separates:
$$\frac{1}{2\pi i}_c𝑑z\alpha (z)T(z)\mathrm{\Phi }_{i_1}(z_1,\overline{z}_1)\mathrm{}=\underset{k}{}\underset{j_k}{}\left[\alpha ^{}(z_k)\mathrm{\Delta }_{i_k}^{j_k}+\delta _{i_k}^{j_k}\alpha (z_k)\frac{}{z_k}\right]\mathrm{\Phi }_{j_1}(z_1,\overline{z}_1)\mathrm{}$$
(22)
with a similar equation for $`\overline{T}\mathrm{\Phi }\mathrm{}`$. Using Cauchy’s theorem the correlation function $`T\mathrm{\Phi }\mathrm{}`$ can be expressed in terms of linear differential operators on $`\mathrm{\Phi }\mathrm{}`$ :
$$T(z)\mathrm{\Phi }_{i_1}(z_1,\overline{z}_1)\mathrm{}=\underset{k}{}\underset{j_k}{}\left[\frac{\mathrm{\Delta }_{i_k}^{j_k}}{(zz_k)^2}+\frac{\delta _{i_k}^{j_k}}{(zz_k)}\frac{}{z_k}\right]\mathrm{\Phi }_{j_1}(z_1,\overline{z}_1)\mathrm{}.$$
(23)
In the presence of a boundary, transformations are restricted to those which leave the boundary invariant, therefore the function $`\alpha `$ and $`\overline{\alpha }`$ are no longer independent. Let us take as boundary the real axis, thus the operators $`\mathrm{\Phi }_{i_k}`$ are defined in the upper half-plane $`Im(z)>0`$. The transformations which leave this boundary invariant, are the real analytic functions, i.e $`\overline{\alpha (z)}=\alpha (\overline{z})`$. As result of this the separation of $`z`$ and $`\overline{z}`$ can not take place.
However, as Cardy has shown , we may extend the definition of $`T(z)`$ into the lower half-plane by:
$$T(z):=\overline{T}(z)Im(z)<0$$
(24)
and relable $`\overline{z}_k=z_k^{}`$. Changing variable $`\overline{z}z`$ in the second integral of equation (21) and using $`\overline{\alpha (z)}=\alpha (\overline{z})`$ , it becomes:
$`{\displaystyle \frac{1}{2\pi i}}{\displaystyle _c}𝑑z\alpha (z)T(z)\mathrm{\Phi }_{i_1}(z_1,z_1^{})\mathrm{}+{\displaystyle \frac{1}{2\pi i}}{\displaystyle _{\overline{c}}}𝑑z\alpha (z)T(z)\mathrm{\Phi }_{i_1}(z_1,z_1^{})\mathrm{}`$
$`={\displaystyle \underset{k}{}}{\displaystyle \underset{j_k}{}}\left[\alpha ^{}(z_k)\mathrm{\Delta }_{i_k}^{j_k}+\delta _{i_k}^{j_k}\alpha (z_k){\displaystyle \frac{}{z_k}}+\alpha ^{}(z_k^{})\overline{\mathrm{\Delta }}_{i_k}^{j_k}+\alpha (z_k^{}){\displaystyle \frac{}{z_k^{}}}\right]\mathrm{\Phi }_{j_1}(z_1,z_1^{})\mathrm{},`$ (25)
where $`\overline{c}`$ is a contour in the lower half-plane.
The left-hand side of equation (25) can be written as one integral around a large contour containing all the points $`z_k`$ and $`z_k^{}`$ if we have $`T=\overline{T}`$ on the boundary. This condition is equivalent to the condition $`T_{xy}=0`$ in Cartesian coordinates which means that there is no flux of energy across the boundary. Now that the left-hand side of equation (25) is an integral, one can use Cauchy’s theorem to get:
$`T(z)\mathrm{\Phi }_{i_1}(z_1,z_{}^{}{}_{1}{}^{})\mathrm{}=`$
$`{\displaystyle \underset{k}{}}{\displaystyle \underset{j_k}{}}\left[{\displaystyle \frac{\mathrm{\Delta }_{j_k}^{i_k}}{(zz_k)^2}}+{\displaystyle \frac{\delta _{j_k}^{i_k}}{(zz_k)}}{\displaystyle \frac{}{z_k}}+{\displaystyle \frac{\overline{\mathrm{\Delta }}_{i_k}^{j_k}}{(zz_k^{})^2}}+{\displaystyle \frac{\delta _{i_k}^{j_k}}{(zz_k^{})}}{\displaystyle \frac{}{z_k^{}}}\right]\mathrm{\Phi }_{j_1}(z_1,z_1^{})\mathrm{}.`$ (26)
Comparing equations (26) and (23), we observe that in presence of a boundary, the correlation function $`\mathrm{\Phi }_{i_1}(z_1,\overline{z}_1)\mathrm{}\mathrm{\Phi }_{i_{2n}}(z_{2n},\overline{z}_{2n})`$ regarded as a function of $`(z_1,\mathrm{},z_n,\overline{z}_1,\mathrm{},\overline{z}_n)`$ satisfies the same differential equations as that of the bulk correlation function $`\mathrm{\Phi }_{i_1}(z_1,\overline{z}_1)\mathrm{}\mathrm{\Phi }_{i_{2n}}(z_{2n},\overline{z}_{2n})`$ as a function of $`z_1,\mathrm{},z_{2n}`$. Note that the conformal dimension of fields $`\mathrm{\Phi }_{i_{n+1}}`$ to $`\mathrm{\Phi }_{i_{2n}}`$ is $`\overline{\mathrm{\Delta }}`$.
As an example we calculate the one point function of such fields. Consider a $`2\times 2`$ scalar Jordanian cell, composed of the fields $`\mathrm{\Phi }`$ and $`\mathrm{\Psi }`$ ($`\mathrm{\Psi }`$ is the logarithmic partner). Conformal invariance of correlation functions means that acting with the set $`L_0,L_{\pm 1}`$ on the correlation functions yields zero. Solving the equations obtained, one arrives at:
$`\mathrm{\Phi }(z,\overline{z})`$ $`=`$ $`{\displaystyle \frac{c}{(z\overline{z})^{2\mathrm{\Delta }}}},`$
$`\mathrm{\Psi }(z,\overline{z})`$ $`=`$ $`{\displaystyle \frac{1}{(z\overline{z})^{2\mathrm{\Delta }}}}\left[c^{}2c\mathrm{log}(z\overline{z})\right].`$ (27)
which are the same as two-point functions without the boundary. Also note that these results are consistent with the results obtained for general dimensions in the previous section (equation (8)).
Now consider the two-point correlation functions near the boundary. Again the differential equations which are satisfied by these two-point correlation functions are the same as the equations satisfied by four-point correlation function in the bulk and so the result is:
$`\mathrm{\Phi }(z_1)\mathrm{\Phi }(z_2)`$ $`=`$ $`u^{2\mathrm{\Delta }}f_1(v)`$
$`\mathrm{\Phi }(z_1)\mathrm{\Psi }(z_2)`$ $`=`$ $`u^{2\mathrm{\Delta }}\left(f_2(v)2\mathrm{log}(z_2\overline{z_2})f_1(v)\right)`$
$`\mathrm{\Psi }(z_1)\mathrm{\Psi }(z_2)=u^{2\mathrm{\Delta }}(f_3(v)2\mathrm{log}\left[(z_1\overline{z}_1)(z_2\overline{z}_2)\right]f_2(v)`$
$`+4\mathrm{log}(z_1\overline{z}_1)\mathrm{log}(z_2\overline{z}_2)f_1(v))`$ (28)
where $`u=(z_1\overline{z}_1)(z_2\overline{z}_2)/(z_1z_2)(\overline{z}_1\overline{z}_2)(z_1\overline{z}_2)(\overline{z}_1z_2)`$ and
$`v=(z_1z_2)(\overline{z}_1\overline{z}_2)/(z_1\overline{z}_1)(z_2\overline{z}_2)`$ and $`f_1,f_2,f_3`$ are arbitrary functions (Again compare these results with the equations (16) and (17) derived for general dimensions).
If a singular vector is found in such a theory, one can find some differential equations which are satisfied by $`f_1,f_2,f_3`$ (usually they are hypergeometric ones) and hence the correlation functions will be determined completely. Finding the correlation functions one can investigate the surface behaviour of the theory. In a successing paper we will consider such a theory.
## 4 Application to MHD
The incompressible two-dimensional magnetohydrodynamic system has two independent dynamical variables, the velocity stream function $`\varphi `$ and the magnetic-flux function $`\psi `$. They obey the equations
$`{\displaystyle \frac{w}{t}}`$ $`=`$ $`ϵ_{\alpha \beta }_\alpha \varphi _\beta w+ϵ_{\alpha \beta }_\alpha \psi _\beta J+\mu ^2w`$
$`{\displaystyle \frac{\psi }{t}}`$ $`=`$ $`ϵ_{\alpha \beta }_\alpha \varphi _\beta \psi +\eta J,`$ (29)
where the vorticity $`w=^2\varphi `$ and the current $`J=^2\psi `$ and $`\mu `$ and $`\eta `$ are viscosity and molecular resistivity. The velocity and magnetic field are given in terms of $`\varphi `$ and $`\psi `$:
$`V_\alpha `$ $`=`$ $`ϵ_{\alpha \beta }_\beta \varphi `$
$`B_\alpha `$ $`=`$ $`ϵ_{\alpha \beta }_\beta \psi ,`$ (30)
where $`ϵ_{\alpha \beta }`$ is the totally antisymmetric tensor with $`ϵ_{12}=1`$. It has been argued that the Alf’ven effect implies that $`\varphi `$ and $`\psi `$ should have equal scaling dimension which naturally leads to LCFT’s .
We consider this system near a boundary and calculate the mean values of velocity and magnetic field. As we have derived in the last section the one-point function of the fields $`\varphi `$ and $`\psi `$ are given by equation (27), so for velocity we have:
$`V_x(x,y)`$ $`=`$ $`_y\varphi (x,y)={\displaystyle \frac{2\mathrm{\Delta }C}{y^{2\mathrm{\Delta }+1}}}`$
$`V_y(x,y)`$ $`=`$ $`_x\varphi (x,y)=0,`$ (31)
and for magnetic field:
$`B_x(x,y)`$ $`=`$ $`_y\psi (x,y)={\displaystyle \frac{2\mathrm{\Delta }}{y^{2\mathrm{\Delta }}}}\left[(C^{}+2C)2C\mathrm{log}y\right]`$
$`B_y(x,y)`$ $`=`$ $`_x\psi (x,y)=0.`$ (32)
A specific model is proposed by Rahimi-Tabar and Rouhani with $`\mathrm{\Delta }=\frac{5}{7}`$ . This theory seems to be unphysical at first sight, because V and B grow large far from boundary. However, they acquire physical meaning when this model is regularized, for example by attaching a value to the velocity field at the boundary.
Other boundaries such as strip and circle can be readily investigated by proper transformations. For example for the strip geometry with size $`L`$ one obtains:
$`\varphi (x,y)`$ $`=`$ $`\left({\displaystyle \frac{\pi }{L}}\right)^{2\mathrm{\Delta }}{\displaystyle \frac{C}{(\mathrm{sin}\frac{\pi }{L}y)^{2\mathrm{\Delta }}}}`$
$`\psi (x,y)`$ $`=`$ $`\left({\displaystyle \frac{\pi }{L}}\right)^{2\mathrm{\Delta }}{\displaystyle \frac{1}{(\mathrm{sin}\frac{\pi }{L}y)^{2\mathrm{\Delta }}}}\left(C^{}+2C\mathrm{log}{\displaystyle \frac{\pi }{L}}2C\mathrm{log}\mathrm{sin}({\displaystyle \frac{\pi }{L}}y)\right).`$ (33)
Further development of LCFT, such as complete calculation of the four point functions, is necessary before some interesting questions such as the possible set of surface critical indices can be determined. Work in this direction is in progress.
Acknowledgement
We would like to thank M. R. Rahimi-Tabar and M. Sa’addat for helpful comments and advice on this paper. |
warning/0002/hep-th0002016.html | ar5iv | text | # 1 Introduction
## 1 Introduction
This series of lectures describes the matrix model of M-theory, also known as M(atrix) Theory. Matrix theory is a supersymmetric quantum mechanics theory with matrix degrees of freedom. It has been known for over a decade that matrix theory arises as a regularization of the 11D supermembrane theory in light-front gauge. It was conjectured in 1996 that when the size of the matrices is taken to infinity this theory gives a microscopic second-quantized description of M-theory in light-front coordinates.
These lectures focus on some basic aspects of matrix theory. We begin by describing in some detail the two alternative definitions of the theory in terms of a quantized and regularized supermembrane theory and as a compactification of M-theory on a lightlike circle. Given these definitions of the theory, we then focus on the question of whether the physics of M-theory and 11-dimensional supergravity can be described constructively using finite size matrices. We show that all the objects of M-theory, including the supergraviton, membrane and 5-brane can be constructed explicitly from configurations of matrices, although these results are not yet complete in the case of the 5-brane. We then turn to the gravitational interactions between these objects, and review what is known about the connection between perturbative calculations in the matrix quantum mechanics theory and supergravity interactions. In the last part of the lectures, some discussion is given of how the matrix theory formalism may be generalized to describe compact or curved space-times.
Previous reviews of matrix theory and related work have appeared in .
In Section 2 we show how matrix theory can be derived from the light-front quantization of the supermembrane theory in 11 dimensions. We discuss in Section 3 the conjecture of Banks, Fischler, Shenker and Susskind that matrix theory describes light-front M-theory in flat space, and we review an argument of Seiberg and Sen showing that finite $`N`$ matrix theory describes the discrete light-cone quantization (DLCQ) of M-theory. In Section 4 we show how the objects of M-theory (the supergraviton, supermembrane and M5-brane) can be described in terms of matrix theory degrees of freedom. Section 5 reviews what is known about the interactions between these objects. We discuss the problem of reproducing N-body interactions in 11D classical supergravity from matrix theory, beginning with two-body interactions in the linearized theory and then discussing many-body interactions and nonlinear terms as well as quantum corrections to the supergravity theory. Section 6 contains a discussion of the problems of formulating matrix theory on a compact or curved background geometry. Finally, we conclude in section 7 with a summary of the current state of affairs and the outlook for the future of this theory.
Even if in the long run matrix theory turns out not to be the most useful description of M-theory, there are many features of this theory which make it well worth studying. It is the simplest example of a quantum supersymmetric gauge theory which seems to correspond to a theory of gravity in a fixed background in some limit. It is the only known example of a well-defined quantum theory which has been shown explicitly to give rise to long-range interactions which agree with gravity at the linearized level and which also contain some nonlinearity. Finally, it provides simple examples of many of the remarkable connections between D-brane physics and gauge theory, giving intuition which may be applicable to a wide variety of situations in string theory and M-theory.
## 2 Matrix theory from the quantized supermembrane
In this section we show that supersymmetric matrix quantum mechanics arises naturally as a regularization of the supermembrane action in 11 dimensions. We begin our discussion with some motivational remarks.
In retrospect, the supermembrane is a natural place to begin when trying to construct a microscopic description of M-theory. There are several distinct 10-dimensional supersymmetric theories of gravity. These theories are well-defined classically but, as with all theories of gravity, are difficult to quantize directly. Each of these theories has a bosonic antisymmetric 2-form tensor field $`B_{\mu \nu }`$. This field is analogous to the 1-form field $`A_\mu `$ of electromagnetism, but carries an extra space-time index. Each of these 10D supergravity theories admits a classical stringlike black hole solution which is “electrically” charged under the 2-form field, in the sense that the two-dimensional world-volume $`\mathrm{\Sigma }`$ of the string couples to the $`B`$ field through a term
$$_\mathrm{\Sigma }B_{\mu \nu }ϵ^{ab}(_aX^\mu )(_bX^\nu ).$$
where $`X^\mu `$ are the embedding functions of the string world-volume in 10 dimensions. This is the higher-dimensional analog of the usual coupling of a charged particle to a gauge field through $`A_\mu \dot{X}^\mu `$.
The quantization of strings in 10-dimensional background geometries can be carried out consistently in only a limited number of ways. These constructions lead to the perturbative descriptions of the five superstring theories known as the type I, IIA, IIB and heterotic $`E_8\times E_8`$ and $`SO(32)`$ theories. These quantum superstring theories are first-quantized from the point of view of the target space—that is, a state in the string Hilbert space corresponds to a single particle-like state in the target space consisting of a single string. Although the quantized string spectrum naturally contains states corresponding to quanta of the supergravity fields (including the NS-NS field $`B_{\mu \nu }`$), it is not possible to give a simple description in terms of the string Hilbert space for extended objects such as D-branes and the NS 5-brane. These objects are essentially nonperturbative phenomena in the superstring theories.
One of the most important developments in the last few years has been the discovery of a network of duality symmetries which relates the five superstring theories to each other and to 11-dimensional supergravity. Of these six theories, the quantized superstring gives a microscopic description of the five 10-dimensional theories. It has been hypothesized that there is a microscopic 11-dimensional theory, dubbed M-theory, underlying this structure which reduces in the low-energy limit to 11D supergravity . To date, however, a precise description of this theory is lacking. Such a theory cannot be described by a quantized string since there is no antisymmetric 2-form field in the 11D supergravity multiplet and hence no stringlike solution of the gravity equations. The 11D supergravity theory contains, however, an antisymmetric 3-form field $`A_{IJK}`$, and the classical theory admits membrane-like solutions which couple electrically to this field. It is easy to imagine that a microscopic description of M-theory might be found by quantizing this supermembrane. This idea was explored extensively in the 80’s, when it was first realized that a consistent classical theory of a supermembrane could be realized in 11 dimensions. At that time, while no satisfactory covariant quantization of the membrane theory was found, it was shown that the supermembrane could be quantized in light-front coordinates. In fact, an elegant regularization of this theory was suggested by Goldstone and Hoppe in 1982. They showed that for the bosonic membrane the regularized quantum theory is a simple quantum-mechanical theory of $`N\times N`$ matrices which leads to the membrane theory in the large $`N`$ limit. This approach was generalized to the supermembrane by de Wit, Hoppe and Nicolai , who showed that the regularized supermembrane theory is precisely the supersymmetric matrix quantum mechanics now known as Matrix Theory. A remarkable feature of the quantum supermembrane theory is that unlike the quantized string theories, the membrane theory automatically gives a second quantized theory from the point of view of the target space. This issue will be discussed in more detail in Section 2.
In this section we describe in some detail how matrix theory arises from the quantization of the supermembrane. In 2.1 we review how the bosonic string is quantized in the light-front formalism. This will be a useful reference point for our discussion of membrane quantization. In 2.2 we describe the theory of the relativistic bosonic membrane in flat space. The light-front description of this theory is discussed in 2.3, and the matrix regularization of the theory is described in 2.4. In 2.5 we discuss briefly the description of the bosonic membrane moving in a general background geometry. In 2.6 we extend the discussion to the supermembrane. We discuss the supermembrane in an arbitrary background geometry. We discuss the $`\kappa `$-symmetry of the supermembrane theory which leads, even at the classical level, to the condition that the background geometry satisfies the classical 11D supergravity equations of motion. The matrix theory Hamiltonian is derived from the regularized supermembrane theory. The problem of finding a covariant membrane quantization is discussed in 2.7.
The material in this section roughly follows the original papers . Note, however, that the original derivation of the matrix quantum mechanics theory was done in the Nambu-Goto-type membrane formalism, while we use here the Polyakov-type approach. We only consider closed membranes in the discussion here; little is known about the open membrane which must end on the M-theory 5-brane, but it would be very interesting to generalize the discussion here to the open membrane.
### 2.1 Review of light-front string
We begin with a brief review of the bosonic string. This will be a useful model to compare with in our discussion of the supermembrane.
The Nambu-Goto action for the relativistic bosonic string moving in a flat background space-time is
$$S=T_sd^2\sigma \sqrt{deth_{ab}}$$
(2.1)
where $`T_s=1/(2\pi \alpha ^{})`$ and
$$h_{ab}=_aX^\mu _bX_\mu .$$
(2.2)
It is convenient to use the Polyakov formalism in which an auxiliary world-sheet metric $`\gamma `$ is introduced
$$S=\frac{1}{4\pi \alpha ^{}}d^2\sigma \sqrt{\gamma }\gamma ^{ab}_aX^\mu _bX_\mu $$
(2.3)
Solving the equation of motion for $`\gamma _{ab}`$ leads to
$$\gamma _{ab}=h_{ab}=_aX^\mu _bX_\mu $$
(2.4)
and replacing this in (2.3) gives (2.1).
The action (2.3) is simplified by going to the gauge
$$\gamma _{ab}=\eta _{ab}.$$
(2.5)
In this gauge we simply have the free field action
$$S=\frac{1}{4\pi \alpha ^{}}d^2\sigma \eta ^{ab}_aX^\mu _bX_\mu .$$
(2.6)
The fields $`X^\mu `$ satisfy the equation of motion $`\mathrm{}X^\mu =0`$ and are subject to the auxiliary Virasoro constraints
$`\dot{X}^\mu (X_\mu )`$ $`=`$ $`0`$ (2.7)
$`\dot{X}^\mu \dot{X}_\mu `$ $`=`$ $`(X^\mu )(X_\mu )`$
(we denote $`\tau `$ derivatives by a dot and $`\sigma `$ derivatives by $``$). Because this is a free theory it is fairly straightforward to quantize. The approaches to quantizing this theory include the BRST and light-front formalisms. The Virasoro constraints can be explicitly solved in light-front gauge
$$X^+(\tau ,\sigma )=x^++p^+\tau .$$
(2.8)
In the classical theory we have
$`\dot{X}^{}`$ $`=`$ $`{\displaystyle \frac{1}{2p^+}}\left(\dot{X}^i\dot{X}^i+X^iX^i\right)`$ (2.9)
$`X^{}`$ $`=`$ $`{\displaystyle \frac{1}{p^+}}\dot{X}^iX^i`$
The transverse degrees of freedom $`X^i`$ have Fourier modes with the commutation relations of simple harmonic oscillators. These are straightforward to quantize. The string spectrum is then given by the usual mass-shell condition
$$M^2=2p^+p^{}p^ip^i=\frac{1}{\alpha ^{}}(Na)$$
(2.10)
### 2.2 The bosonic membrane theory
We now discuss the relativistic bosonic membrane moving in an arbitrary number $`D`$ of space-time dimensions. The story begins in a very similar fashion to the relativistic string. We want to use a Nambu-Goto-style action
$$S=Td^3\sigma \sqrt{deth_{\alpha \beta }}$$
(2.11)
where $`T`$ is the membrane tension
$$T=\frac{1}{(2\pi )^2l_p^3}$$
(2.12)
and
$$h_{\alpha \beta }=_\alpha X^\mu _\beta X_\mu $$
(2.13)
is the pullback of the metric to the three-dimensional membrane world-volume, with coordinates $`\sigma _\alpha ,\alpha \{0,1,2\}`$. We will use the notation $`\tau =\sigma _0`$ and use indices $`a,b,\mathrm{}`$ to describe “spatial” indices $`a\{1,2\}`$ on the membrane world-volume.
We again wish to use a Polyakov-type formalism in which an auxiliary world-sheet metric $`\gamma _{\alpha \beta }`$ is introduced
$$S=\frac{T}{2}d^3\sigma \sqrt{\gamma }\left(\gamma ^{\alpha \beta }_\alpha X^\mu _\beta X_\mu 1\right).$$
(2.14)
The need for the extra “cosmological” term arises from the absence of scale invariance in the theory. Computing the equations of motion from varying $`\gamma _{\alpha \beta }`$, and using $`\delta \sqrt{\gamma }=\frac{1}{2}\sqrt{\gamma }\gamma ^{\alpha \beta }\delta \gamma _{\alpha \beta },\delta \gamma ^{ϵ\varphi }=\gamma ^{\alpha ϵ}\gamma ^{\varphi \beta }\delta \gamma _{\alpha \beta }`$, we get
$$\gamma ^{\alpha \gamma }\gamma ^{\beta \delta }h_{\gamma \delta }+\frac{1}{2}\gamma ^{\alpha \beta }t\frac{1}{2}\gamma ^{\alpha \beta }=0$$
(2.15)
where $`t=\gamma ^{\alpha \beta }h_{\alpha \beta }`$. Lowering all indices gives
$$\frac{1}{2}\gamma _{\alpha \beta }(t1)=h_{\alpha \beta }$$
(2.16)
or
$$\gamma _{\alpha \beta }=\frac{2h_{\alpha \beta }}{t1}.$$
(2.17)
Contracting indices gives
$$3=\frac{2t}{t1}$$
(2.18)
so $`t=3`$ and
$$\gamma _{\alpha \beta }=h_{\alpha \beta }=_\alpha X^\mu _\beta X_\mu .$$
(2.19)
Replacing this in (2.14) again gives (2.11). The equation of motion which arises from varying $`X^\mu `$ is
$$_\alpha \left(\sqrt{\gamma }\gamma ^{\alpha \beta }_\beta X^\mu \right)=0.$$
(2.20)
To follow the procedure we used for the bosonic string theory, we would now like to use the symmetries of the theory to gauge-fix the metric $`\gamma _{\alpha \beta }`$. Unfortunately, whereas for the string we had three components of the metric and three continuous symmetries (two diffeomorphism symmetries and a scale symmetry), for the membrane we have six independent metric components and only three diffeomorphism symmetries. We can use these symmetries to fix the components $`\gamma _{0\alpha }`$ of the metric to be
$`\gamma _{0a}`$ $`=`$ $`0`$ (2.21)
$`\gamma _{00}`$ $`=`$ $`{\displaystyle \frac{4}{\nu ^2}}\overline{h}{\displaystyle \frac{4}{\nu ^2}}deth_{ab}`$
where $`\nu `$ is a constant whose normalization has been chosen to make the later matrix interpretation transparent. Once we have chosen this gauge, no further components of the metric $`\gamma _{ab}`$ can be fixed. This gauge can only be chosen when the membrane world-volume is of the form $`\mathrm{\Sigma }\times R`$ where $`\mathrm{\Sigma }`$ is a Riemann surface of fixed topology. The membrane action becomes
$$S=\frac{T\nu }{4}d^3\sigma \left(\dot{X}^\mu \dot{X}_\mu \frac{4}{\nu ^2}\overline{h}\right).$$
(2.22)
It is natural to rewrite this theory in terms of a canonical Poisson bracket on the membrane at constant $`\tau `$ where $`\{f,g\}ϵ^{ab}_af_bg`$ with $`ϵ^{12}=1`$. We will assume that the coordinates $`\sigma `$ are chosen so that with respect to the symplectic form associated to this canonical Poisson bracket the volume of the Riemann surface $`\mathrm{\Sigma }`$ is $`d^2\sigma =4\pi `$. In terms of this metric we have the handy formulae
$`\overline{h}=deth_{ab}`$ $`=`$ $`_1X^\mu _1X^\mu _2X^\nu _2X^\nu _1X^\mu _2X^\mu _1X^\nu _2X^\nu `$ (2.23)
$`=`$ $`{\displaystyle \frac{1}{2}}\{X^\mu ,X^\nu \}\{X_\mu ,X_\nu \}`$
$`_a(\overline{h}h^{ab}_bX^\mu )`$ $`=`$ $`\{\{X^\mu ,X^\nu \},X_\nu \}`$ (2.24)
$`\overline{h}h^{ab}_aX^\mu _bX^\nu `$ $`=`$ $`\{X^\mu ,X^\lambda \}\{X_\lambda ,X^\nu \}`$ (2.25)
In terms of the Poisson bracket, the membrane action becomes
$$S=\frac{T\nu }{4}d^3\sigma \left(\dot{X}^\mu \dot{X}_\mu \frac{2}{\nu ^2}\{X^\mu ,X^\nu \}\{X_\mu ,X_\nu \}\right).$$
(2.26)
The equations of motion for the fields $`X^\mu `$ are
$`\ddot{X}^\mu `$ $`=`$ $`{\displaystyle \frac{4}{\nu ^2}}_a\left(\overline{h}h^{ab}_bX^\mu \right)`$ (2.27)
$`=`$ $`{\displaystyle \frac{4}{\nu ^2}}\{\{X^\mu ,X^\nu \},X_\nu \}`$
The auxiliary constraints on the system are
$`\dot{X}^\mu \dot{X}_\mu `$ $`=`$ $`{\displaystyle \frac{4}{\nu ^2}}\overline{h}`$ (2.28)
$`=`$ $`{\displaystyle \frac{2}{\nu ^2}}\{X^\mu ,X^\nu \}\{X_\mu ,X_\nu \}`$
and
$$\dot{X}^\mu _aX_\mu =0.$$
(2.29)
It follows directly from (2.29) that
$$\{\dot{X}^\mu ,X_\mu \}=0.$$
(2.30)
We have thus expressed the bosonic membrane theory as a constrained Hamiltonian system. The degrees of freedom are $`D`$ functions $`X^\mu `$ on the 3-dimensional world-volume of a membrane which has topology $`\mathrm{\Sigma }\times R`$ where $`\mathrm{\Sigma }`$ is a Riemann surface. This theory is still completely covariant. It is difficult to quantize, however, because of the constraints and the nonlinearity of the equations of motion. The direct quantization of this covariant theory will be discussed further in Section 2.7.
### 2.3 The light-front bosonic membrane
As we did for the bosonic string, we now consider the theory in light-front coordinates
$$X^\pm =(X^0\pm X^{D1})/\sqrt{2}.$$
(2.31)
Just as in the case of the string, the constraints (2.28,2.29) can be explicitly solved in light-front gauge
$$X^+(\tau ,\sigma _1,\sigma _2)=\tau .$$
(2.32)
We have
$`\dot{X}^{}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\dot{X}^i\dot{X}^i+{\displaystyle \frac{2\overline{h}}{\nu ^2}}`$ (2.33)
$`=`$ $`{\displaystyle \frac{1}{2}}\dot{X}^i\dot{X}^i+{\displaystyle \frac{1}{\nu ^2}}\{X^i,X^j\}\{X^i,X^j\}`$
$`_aX^{}`$ $`=`$ $`\dot{X}^i_aX^i`$
We can go to a Hamiltonian formalism by computing the canonically conjugate momentum densities.
$`P^+`$ $`=`$ $`{\displaystyle \frac{\delta }{\delta (\dot{X}^{})}}={\displaystyle \frac{\nu T}{2}}`$ (2.34)
$`P^i`$ $`=`$ $`{\displaystyle \frac{\delta }{\delta (\dot{X}^i)}}={\displaystyle \frac{\nu T}{2}}\dot{X}^i`$
The total momentum in the direction $`P^+`$ is then
$$p^+=d^2\sigma P^+=2\pi \nu T.$$
(2.35)
The Hamiltonian of the theory is given by
$`H`$ $`=`$ $`{\displaystyle d^2\sigma \left(P^i\dot{X}^iP^+\dot{X}^{}\right)}`$
$`=`$ $`{\displaystyle \frac{\nu T}{4}}{\displaystyle d^2\sigma \left(\dot{X}^i\dot{X}^i+\frac{4\overline{h}}{\nu ^2}\right)}`$
$`=`$ $`{\displaystyle \frac{\nu T}{4}}{\displaystyle d^2\sigma \left(\dot{X}^i\dot{X}^i+\frac{2}{\nu ^2}\{X^i,X^j\}\{X^i,X^j\}\right)}.`$
The only remaining constraint is that the transverse degrees of freedom must satisfy
$$\{\dot{X}^i,X^i\}=0$$
(2.37)
This theory has a residual invariance under time-independent area-preserving diffeomorphisms. Such diffeomorphisms do not change the symplectic form and thus manifestly leave the Hamiltonian (2.3)
We now have a Hamiltonian formalism for the light-front membrane theory. Unfortunately, this theory is still rather difficult to quantize. Unlike string theory, where the equations of motion are linear in this formalism, for the membrane the equations of motion (2.27) are nonlinear and difficult to solve.
### 2.4 Matrix regularization
In 1982 a remarkably clever regularization of the light-front membrane theory was found by Goldstone and Hoppe in the case where the surface $`\mathrm{\Sigma }`$ is a sphere $`S^2`$ . According to this regularization procedure, functions on the membrane surface are mapped to finite size $`N\times N`$ matrices. Just as in the quantization of a classical mechanical system defined in terms of a Poisson brackets, the Poisson bracket appearing in the membrane theory is replaced in the matrix regularization of the theory by a matrix commutator.
The matrix regularization of the theory can be generalized to membranes of arbitrary topology, but is perhaps most easily understood by considering the case discussed in , where the membrane has the topology of a sphere $`S^2`$ for all values of $`\tau `$. In this case the world-sheet of the membrane surface at fixed time can be described by a unit sphere with a rotationally invariant canonical symplectic form. Functions on this membrane can be described in terms of functions of the three Cartesian coordinates $`\xi _1,\xi _2,\xi _3`$ on the unit sphere satisfying
$$\xi _1^2+\xi _2^2+\xi _3^2=1.$$
(2.38)
The Poisson brackets of these functions are given by
$$\{\xi _A,\xi _B\}=ϵ_{ABC}\xi _C.$$
This is essentially the same algebraic structure as that defined by the commutation relations of the generators of $`SU(2)`$. It is therefore natural to associate these coordinate functions on $`S^2`$ with the matrices generating $`SU(2)`$ in the $`N`$-dimensional representation. In terms of the conventions we are using here, when the normalization constant $`\nu `$ is integral, the correct correspondence is
$$\xi _A\frac{2}{N}J_A$$
where $`J_1,J_2,J_3`$ are generators of the $`N`$-dimensional representation of $`SU(2)`$ with $`N=\nu `$, satisfying the commutation relations
$$i[J_A,J_B]=ϵ_{ABC}J_C.$$
In general, any function on the membrane can be expanded as a sum of spherical harmonics
$$f(\xi _1,\xi _2,\xi _3)=\underset{l,m}{}c_{lm}y_{lm}(\xi _1,\xi _2,\xi _3)$$
(2.39)
The spherical harmonics can in turn be written as a sum of monomials in the coordinate functions:
$$y_{lm}(\xi _1,\xi _2,\xi _3)=\underset{k}{}t_{A_1\mathrm{}A_l}^{(lm)}\xi _{A_1}\mathrm{}\xi _{A_l}$$
where the coefficients $`t_{A_1\mathrm{}A_l}^{(lm)}`$ are symmetric and traceless (because $`\xi _A\xi _A=1`$). Using the above correspondence, a matrix approximation to each of the spherical harmonics with $`l<N`$ can be constructed, which we denote by $`Y`$.
$$Y_{lm}=\left(\frac{2}{N}\right)^lt_{A_1\mathrm{}A_l}^{(lm)}J_{A_1}\mathrm{}J_{A_l}$$
(2.40)
For a fixed value of $`N`$ only the spherical harmonics with $`l<N`$ can be constructed because higher order monomials in the generators $`J_A`$ do not generate linearly independent matrices. Note that the number of independent matrix entries is precisely equal to the number of independent spherical harmonic coefficients which can be determined for fixed $`N`$
$$N^2=\underset{l=0}{\overset{N1}{}}(2l+1)$$
(2.41)
The matrix approximations (2.40) of the spherical harmonics can be used to construct matrix approximations to an arbitrary function of the form (2.39)
$$F=\underset{l<N,m}{}c_{lm}Y_{lm}$$
(2.42)
The Poisson bracket in the membrane theory is replaced in the matrix regularized theory with the matrix commutator according to the prescription
$$\{f,g\}\frac{iN}{2}[F,G].$$
(2.43)
Similarly, an integral over the membrane at fixed $`\tau `$ is replaced by a matrix trace through
$$\frac{1}{4\pi }d^2\sigma f\frac{1}{N}\mathrm{Tr}F$$
(2.44)
The Poisson bracket of a pair of spherical harmonics takes the form
$$\{y_{lm},y_{l^{}m^{}}\}=g_{lm,l^{}m^{}}^{l^{\prime \prime }m^{\prime \prime }}y_{l^{\prime \prime }m^{\prime \prime }}.$$
(2.45)
The commutator of a pair of matrix spherical harmonics (2.40) can be written
$$[Y_{lm},Y_{l^{}m^{}}]=G_{lm,l^{}m^{}}^{l^{\prime \prime }m^{\prime \prime }}Y_{l^{\prime \prime }m^{\prime \prime }}.$$
(2.46)
It can be verified that in the large $`N`$ limit the structure constant of these algebras agree
$$\underset{N\mathrm{}}{lim}\frac{iN}{2}G_{lm,l^{}m^{}}^{l^{\prime \prime }m^{\prime \prime }}g_{lm,l^{}m^{}}^{l^{\prime \prime }m^{\prime \prime }}$$
(2.47)
As a result, it can be shown that for any smooth functions on the membrane $`f,g`$ defined in terms of convergent sums of spherical harmonics, with Poisson bracket $`\{f,g\}=h`$ we have
$`\underset{N\mathrm{}}{lim}{\displaystyle \frac{1}{N}}\mathrm{Tr}F`$ $`=`$ $`{\displaystyle \frac{1}{4\pi }}{\displaystyle d^2\sigma f}`$ (2.48)
and it is possible to show that
$`\underset{N\mathrm{}}{lim}(({\displaystyle \frac{iN}{2}})[F,G]H)`$ $`=`$ $`0`$ (2.49)
This last relation is really shorthand for the statement that
$$\underset{N\mathrm{}}{lim}\frac{1}{N}\mathrm{Tr}\left(((\frac{iN}{2})[F,G]H)J\right)=0$$
(2.50)
where $`J`$ is the matrix approximation to any smooth function $`j`$ on the sphere.
We now have a dictionary for transforming between continuum and matrix-regularized quantities. The correspondence is given by
$$\xi _A\frac{2}{N}J_A\{,\}\frac{iN}{2}[,]\frac{1}{4\pi }d^2\sigma \frac{1}{N}\mathrm{Tr}$$
(2.51)
The matrix regularized membrane Hamiltonian is therefore given by
$`H`$ $`=`$ $`(2\pi l_p^3)\mathrm{Tr}\left({\displaystyle \frac{1}{2}}𝐏^i𝐏^i\right){\displaystyle \frac{1}{(2\pi l_p^3)}}\mathrm{Tr}\left({\displaystyle \frac{1}{4}}[𝐗^i,𝐗^j][𝐗^i,𝐗^j]\right)`$ (2.52)
$`=`$ $`{\displaystyle \frac{1}{(2\pi l_p^3)}}\mathrm{Tr}\left({\displaystyle \frac{1}{2}}\dot{𝐗}^i\dot{𝐗}^i{\displaystyle \frac{1}{4}}[𝐗^i,𝐗^j][𝐗^i,𝐗^j]\right).`$
This Hamiltonian gives rise to the matrix equations of motion
$$\ddot{𝐗}^i+[[𝐗^i,𝐗^j],𝐗^j]=0$$
which must be supplemented with the Gauss constraint
$$[\dot{𝐗}{}_{}{}^{i},𝐗^i]=0.$$
(2.53)
This is a classical theory with a finite number of degrees of freedom. The quantization of such a system is straightforward, although solving the quantum theory can in practice be quite tricky. Thus, we have found a well-defined quantum theory describing the matrix regularization of the relativistic membrane theory in light-front coordinates.
There are a number of rather deep mathematical reasons why the matrix regularization of the membrane theory works. One way of looking at this regularization is in terms of the underlying symmetry of the theory. After gauge-fixing, the membrane theory has a residual invariance under the group of time-independent area-preserving diffeomorphisms of the membrane world-sheet. This diffeomorphism group can be described in a natural mathematical way as a limit of the matrix group $`U(N)`$ as $`N\mathrm{}`$. In the discrete theory the area-preserving diffeomorphism symmetry thus is replaced by the $`U(N)`$ matrix symmetry. The matrix regularization can also be viewed in terms of a geometrical quantization of the operators associated with functions on the membrane. From this point of view the matrix membrane is like a “fuzzy” membrane which is discrete yet preserves the $`SU(2)`$ rotational symmetry of the original smooth sphere. This point of view ties into recent developments in noncommutative geometry.
We will not pursue these points of view in any depth here. We will note, however, that from both points of view it is natural to generalize the construction to higher genus surfaces. We discuss the matrix torus explicitly in section 4.2.3.
### 2.5 The bosonic membrane in a general background
So far we have only considered the membrane in a flat background Minkowski geometry. Just as for strings, it is natural to generalize the discussion to a bosonic membrane moving in a general background metric $`g_{\mu \nu }`$ and 3-form field $`A_{\mu \nu \rho }`$. The introduction of a general background metric modifies the Nambu-Goto action by replacing $`h_{\alpha \beta }`$ in (2.13) with
$$h_{\alpha \beta }=_\alpha X^\mu _\beta X^\nu g_{\mu \nu }(X).$$
(2.54)
The membrane couples to the 3-form field as an electrically charged object, giving an additional term to the action of the form $`A_{\alpha \beta \gamma }`$ where $`A_{\alpha \beta \gamma }`$ is the pullback to the world-volume of the membrane of the 3-form field. This gives a total Nambu-Goto-type action for the membrane in a general background of the form
$$S=Td^3\sigma \left(\sqrt{deth_{\alpha \beta }}+6\dot{X}^\mu _1X^\nu _2X^\rho A_{\mu \nu \rho }(X)\right).$$
(2.55)
With an auxiliary world-volume metric, this action becomes
$`S`$ $`=`$ $`{\displaystyle \frac{T}{2}}{\displaystyle }d^3\sigma [\sqrt{\gamma }(\gamma ^{\alpha \beta }_\alpha X^\mu _\beta X^\nu g_{\mu \nu }(X)1)`$
$`+12\dot{X}^\mu _1X^\nu _2X^\rho A_{\mu \nu \rho }(X)]`$
We can gauge fix the action (2.5) using the same gauge (2.21) as in the flat space case. We can then consider carrying out a similar procedure for quantizing the membrane in a general background as we described in the case of the flat background. We will return to this question in section 6.3 when we discuss in more detail the prospects for constructing matrix theory in a general background.
### 2.6 The supermembrane
Now let us turn our attention to the supermembrane. In order to make contact with M-theory, and indeed to make the membrane theory well-behaved it is necessary to add supersymmetry to the theory. Supersymmetric membrane theories can be constructed classically in dimensions 4, 5, 7 and 11. These theories have different degrees of supersymmetry, with 2, 4, 8 and 16 independent supersymmetric generators respectively. It is believed that all the supermembrane theories other than the 11D maximally supersymmetric theory suffer from anomalies in the Lorentz algebra. Thus, just as $`D=10`$ is the natural dimension for the superstring, $`D=11`$ is the natural dimension for the supermembrane.
The formalism for describing the supermembrane is rather technically complicated. We will not use most of this formalism in the rest of these lectures, so we restrict ourselves here to a fairly concise discussion of the structure of the supersymmetric theory. The reader not interested in the details of how the supersymmetric form of matrix theory is derived may wish to skip directly to the result of this analysis, the supersymmetric matrix theory Hamiltonian (2.97), on first reading. In Section 2.6.1 we describe using superfield notation the supermembrane action in a general background and its symmetries. We discuss in particular the fact that the $`\kappa `$-symmetry of the theory at the classical level guarantees already that the background geometry satisfies the equations of motion of 11D supergravity. In 2.6.2 we describe in more explicit form the supermembrane action in a flat background. We describe the light-front form of the theory in 2.6.3, where we show how the regularized theory gives precisely the Hamiltonian of the supersymmetric matrix theory.
#### 2.6.1 The supermembrane action
In this section we describe the supermembrane action in an arbitrary background and its symmetries. In particular, we describe the $`\kappa `$-symmetry of the theory, which implies that the background obeys the classical equations of 11D supergravity. For further details see the original paper of Bergshoeff, Sezgin and Townsend or the review paper of Duff .
The standard NSR description of the superstring gives a theory which is fairly straightforward to quantize. This formalism can be used in a straightforward fashion to describe the spectra of the five superstring theories. One disadvantage of this formalism is that the target space supersymmetry of the theory is difficult to show explicitly. There is another formalism, known as the Green-Schwarz formalism (, reviewed in ), in which the target space supersymmetry of the theory is much more clear. In the Green-Schwarz formalism additional Grassmann degrees of freedom are introduced which transform as space-time fermions but as world-sheet vectors. These correspond to space-time superspace coordinates for the string. The Green-Schwarz superstring action does not have a standard world-sheet supersymmetry (it can’t, since there are no world-sheet fermions). The theory does, however, have a novel type of supersymmetry known as a $`\kappa `$-symmetry. The existence of the $`\kappa `$-symmetry in the classical Green-Schwarz string theory already implies that the theory is restricted to $`D=3,4,6`$ or 10. This is already a much stronger restriction than can be gleaned from classical superstring with world-sheet supersymmetry.
Unlike the superstring, there is no known way of formulating the supermembrane in a world-volume supersymmetric fashion (although there has been some recent progress in this direction, for further references see ). A Green-Schwarz formulation of the supermembrane in a general background was first found by Bergshoeff, Sezgin and Townsend . We now review this construction.
We consider an 11-dimensional target space with a general metric $`g_{\mu \nu }`$ described by an elfbein $`e_\mu ^a`$, and an arbitrary background gravitino field $`\psi _\mu `$ and 3-form field $`A_{\mu \nu \rho }`$. In superspace notation we describe the space as having 11 bosonic coordinates $`X^\mu `$ and 32 anticommuting fermionic coordinates $`\theta ^{\dot{\alpha }}`$. These coordinates are combined into a single superspace coordinate
$$Z^M=(X^\mu ,\theta ^{\dot{\alpha }})$$
(2.57)
where $`M`$ is an index with 43 possible values. (Space-time spinor indices $`\dot{\alpha },\dot{\beta },\mathrm{}`$ will carry a dot in this section to distinguish them from world-volume coordinate indices $`\alpha ,\beta ,\mathrm{}`$). In superspace the elfbein becomes a 43-bein $`E_M^A`$, with $`A=(a,\varphi )`$. There is also an antisymmetric superspace 3-form field $`B_{MNP}`$. The superspace formulation of 11D supergravity is written in terms of these two fields. The identification of the superspace degrees of freedom with the component fields $`e_\mu ^a,\psi _\mu `$ and $`A_{\mu \nu \rho }`$ is quite subtle, and involves a careful analysis of the supersymmetry transformations in both formalisms as well as gauge choices. At leading order in $`\theta `$ the component fields are identified through
$`E_\mu ^a`$ $`=`$ $`e_\mu ^a+𝒪(\theta )`$
$`E_\mu ^\varphi `$ $`=`$ $`\psi _\mu ^\varphi +𝒪(\theta )`$ (2.58)
$`B_{\mu \nu \rho }`$ $`=`$ $`A_{\mu \nu \rho }+𝒪(\theta )`$
The identifications of $`E_M^A`$ and $`B_{MNP}`$ in terms of component fields through order $`\theta ^2`$ has only recently been determined . The identification beyond this order has not been determined explicitly.
In terms of these superspace fields, the supermembrane action in a general background is given by
$$S=\frac{T}{2}d^3\sigma \left[\sqrt{\gamma }\left(\gamma ^{\alpha \beta }\mathrm{\Pi }_\alpha ^a\mathrm{\Pi }_\beta ^b\eta _{ab}1\right)+ϵ^{\alpha \beta \gamma }\mathrm{\Pi }_\alpha ^A\mathrm{\Pi }_\beta ^B\mathrm{\Pi }_\gamma ^CB_{ABC}\right]$$
(2.59)
where $`\mathrm{\Pi }_\alpha ^A`$ are the components of the pullback of the 43-bein to the membrane world-volume
$$\mathrm{\Pi }_\alpha ^A=_\alpha Z^ME_M^A$$
(2.60)
and $`B_{ABC}`$ is defined implicitly through
$$B_{MNP}=E_M^AE_N^BE_P^CB_{ABC}$$
(2.61)
The action (2.59) is very closely related to the superspace formulation of the Green-Schwarz action. The superstring action differs in that it has no cosmological term and that the antisymmetric field is a superspace 2-form field.
Let us now review the symmetries of the action (2.59). This action has global symmetries corresponding to space-time super diffeomorphisms, gauge transformations and discrete symmetries, as well as the local symmetries of world-volume diffeomorphisms and $`\kappa `$ symmetry.
Super diffeomorphisms: Under a super diffeomorphism of the target space generated by a super vector field $`\xi ^M`$ the coordinate fields, 43-bein and 3-form field transform under
$`\delta Z^M`$ $`=`$ $`\xi ^M`$
$`\delta E_M^A`$ $`=`$ $`\xi ^N_NE_M^A+_M\xi ^NE_N^A`$ (2.62)
$`\delta B_{MNP}`$ $`=`$ $`\xi ^Q_QB_{MNP}+(_M\xi ^Q)B_{QNP}(_N\xi ^Q)B_{MQP}+(_P\xi ^Q)B_{MNQ}`$
Super gauge transformations: This global symmetry transforms the 3-form superfield by
$$\delta B_{MNP}=_M\mathrm{\Sigma }_{NP}_N\mathrm{\Sigma }_{MP}+_P\mathrm{\Sigma }_{MN}.$$
(2.63)
Discrete symmetries: There is also a discrete symmetry $`Z_2`$ corresponding to taking
$$B_{MNP}B_{MNP}$$
(2.64)
and performing a space-time reflection on a single coordinate.
World-volume diffeomorphisms: Under a world-volume diffeomorphism generated by the vector field $`\eta ^\alpha `$ the fields transform by
$$\delta Z^M=\eta ^\alpha _\alpha Z^M$$
(2.65)
$`\kappa `$-symmetry: The most interesting symmetry of the theory is the fermionic $`\kappa `$-symmetry. The parameter $`\kappa ^\psi `$ is taken to be an anticommuting world-volume scalar which transforms as a space-time 32-component spinor. Under this symmetry the coordinate fields transform under
$`\delta Z^ME_M^a`$ $`=`$ $`0`$ (2.66)
$`\delta Z^ME_M^\varphi `$ $`=`$ $`(1+\mathrm{\Gamma })_\psi ^\varphi \kappa ^\psi `$
where
$$\mathrm{\Gamma }=\frac{1}{6\sqrt{\gamma }}ϵ^{\alpha \beta \gamma }\mathrm{\Pi }_\alpha ^a\mathrm{\Pi }_\beta ^b\mathrm{\Pi }_\gamma ^c\mathrm{\Gamma }_{abc}.$$
(2.67)
The $`\kappa `$-symmetry of the theory has a number of interesting features. For one thing, it can be shown that (2.66) is only a symmetry of the theory when the background fields $`E_M^a,B_{MNP}`$ obey the equations of motion of the classical 11D supergravity theory. Thus, 11D supergravity emerges from the membrane theory even at the classical level. For the details of this analysis, see . This situation is similar to that which arises in the Green-Schwarz formulation of the superstring theories. In the Green-Schwarz formalism there is a local $`\kappa `$-symmetry on the string world-sheet only when the backgrounds satisfy the supergravity equations of motion.
Another interesting aspect of the $`\kappa `$-symmetry arises from the algebraic fact that
$$\mathrm{\Gamma }^2=1.$$
(2.68)
This implies that $`(1+\mathrm{\Gamma })`$ is a projection operator. We can thus use $`\kappa `$-symmetry to gauge away half of the fermionic degrees of freedom $`\theta ^{\dot{\alpha }}`$. This reduces the number of propagating fermionic degrees of freedom to 8. This is also the number of propagating bosonic degrees of freedom, as can be seen by going to a static gauge where $`X^{0,1,2}`$ are identified with $`\tau ,\sigma _{1,2}`$ so that only the 8 transverse directions appear as propagating degrees of freedom.
In general, gauge-fixing the $`\kappa `$-symmetry in any particular way will break the Lorentz invariance of the theory. This makes it quite difficult to find any way of quantizing the theory without breaking Lorentz symmetry. This situation is again analogous to the Green-Schwarz superstring theory, where fixing of $`\kappa `$-symmetry also breaks Lorentz invariance and no covariant quantization scheme is known.
#### 2.6.2 The supermembrane in flat space
To make the connection with matrix theory, we now restrict attention to a flat Minkowski background space-time with vanishing 3-form field $`A_{\mu \nu \rho }`$. We will return to a discussion of general backgrounds in section 6.3.
In flat space the 43-bein becomes
$`E_M^a`$ $`=`$ $`(\delta _\mu ^a,(\mathrm{\Gamma }^a)_{\dot{\alpha }\dot{\beta }}\theta ^{\dot{\beta }})`$ (2.69)
$`E_M^\varphi `$ $`=`$ $`(0,\delta _{\dot{\alpha }}^\varphi )`$
The super 4-form field strength $`H_{MNPQ}`$ has as its only nonvanishing components
$$H_{ab\varphi \psi }=\frac{1}{3}(\mathrm{\Gamma }_{ab})_{\varphi \psi }.$$
(2.70)
From this and the definition $`H=dB`$ it is possible to derive the components of the super 3-form field $`B_{MNP}`$
$`B_{\mu \nu \rho }`$ $`=`$ $`0`$
$`B_{\mu \nu \dot{\alpha }}`$ $`=`$ $`{\displaystyle \frac{1}{6}}(\mathrm{\Gamma }_{\mu \nu }\theta )_{\dot{\alpha }}`$ (2.71)
$`B_{\mu \dot{\alpha }\dot{\beta }}`$ $`=`$ $`{\displaystyle \frac{1}{6}}(\mathrm{\Gamma }_{\mu \nu }\theta )_{(\dot{\alpha }}(\mathrm{\Gamma }^\nu \theta )_{\dot{\beta })}`$
$`B_{\dot{\alpha }\dot{\beta }\dot{\gamma }}`$ $`=`$ $`{\displaystyle \frac{1}{6}}(\mathrm{\Gamma }_{\mu \nu }\theta )_{(\dot{\alpha }}(\mathrm{\Gamma }^\mu \theta )_{\dot{\beta }}(\mathrm{\Gamma }^\nu \theta )_{\dot{\gamma })}`$
From (2.60) it follows that
$$\mathrm{\Pi }_\alpha ^\mu =_\alpha X^\mu +\overline{\theta }\mathrm{\Gamma }^\mu _\alpha \theta .$$
(2.72)
The membrane action (2.59) reduces in flat space to
$`S`$ $`=`$ $`{\displaystyle \frac{T}{2}}{\displaystyle }d^3\sigma \{\sqrt{\gamma }(\gamma ^{\alpha \beta }\mathrm{\Pi }_\alpha ^\mu \mathrm{\Pi }_\beta ^\nu \eta _{\mu \nu }1)`$
$`ϵ^{\alpha \beta \gamma }[{\displaystyle \frac{1}{2}}_\alpha X^\mu (_\beta X^\nu +\overline{\theta }\mathrm{\Gamma }^\nu _\beta \theta )`$
$`+{\displaystyle \frac{1}{6}}(\overline{\theta }\mathrm{\Gamma }^\mu _\alpha \theta )(\overline{\theta }\mathrm{\Gamma }^\nu _\beta \theta )]\overline{\theta }\mathrm{\Gamma }_{\mu \nu }_\gamma \theta \}`$
The extra Wess-Zumino type terms which appear in this action are rather non-obvious from the point of view of the flat space-time theory, although they have arisen naturally in the superspace formalism. These are analogous to terms in the Green-Schwarz superstring action which were originally found by imposing $`\kappa `$-symmetry on the theory. The equation of motion for $`\gamma `$ as usual sets $`\gamma _{\alpha \beta }`$ to be the pullback of the metric
$$\gamma _{\alpha \beta }=\mathrm{\Pi }_\alpha ^\mu \mathrm{\Pi }_\beta ^\nu \eta _{\mu \nu }$$
(2.75)
The action (2.6.2) has the target space supersymmetry
$`\delta \theta `$ $`=`$ $`ϵ`$ (2.76)
$`\delta X^\mu `$ $`=`$ $`\overline{ϵ}\mathrm{\Gamma }^\mu \theta `$
This transformation leaves $`\mathrm{\Pi }_\alpha ^\mu `$ invariant. The fact that it leaves the action invariant follows from the identity
$$\overline{\psi }_{[1}\mathrm{\Gamma }^\mu \psi _2\overline{\psi }_3\mathrm{\Gamma }_{\mu \nu }\psi _{4]}=0$$
(2.77)
which holds in 11 dimensions (as well as in dimensions 4, 5 and 7). The relation (2.77) is also necessary to show that the action is $`\kappa `$-symmetric. This relation is analogous to the relation $`\overline{ϵ}\mathrm{\Gamma }_\mu \psi _{[1}\psi _2\mathrm{\Gamma }^\mu \psi _{3]}=0`$ which holds in 3, 4, 6 and 10 dimensions and which is necessary for the supersymmetry and $`\kappa `$-symmetry of the Green-Schwarz superstring action.
#### 2.6.3 The quantum supermembrane and supersymmetric matrix theory
We now go to light-front gauge. As usual we define
$$X^\pm =(X^0\pm X^{D1})/\sqrt{2}.$$
(2.78)
We write the $`32\times 32`$ $`\mathrm{\Gamma }`$ matrices in the block forms
$`\mathrm{\Gamma }^+`$ $`=`$ $`\left(\begin{array}{cc}0& 0\\ \sqrt{2}i\text{1}\text{ }\text{1}_{16}& 0\end{array}\right)`$ (2.81)
$`\mathrm{\Gamma }^{}`$ $`=`$ $`\left(\begin{array}{cc}0& \sqrt{2}i\text{1}\text{ }\text{1}_{16}\\ 0& 0\end{array}\right)`$ (2.84)
$`\mathrm{\Gamma }_i`$ $`=`$ $`\left(\begin{array}{cc}\gamma ^i& 0\\ 0& \gamma ^i\end{array}\right)`$ (2.87)
where $`\gamma ^i`$ are $`16\times 16`$ Euclidean $`SO(9)`$ gamma matrices.
In addition to setting the gauge
$$X^+=\tau $$
(2.88)
We can also use $`\kappa `$-symmetry to fix
$$\mathrm{\Gamma }^+\theta =0$$
(2.89)
From the above form of the matrices $`\mathrm{\Gamma }^\mu `$ it is clear that this projects onto the 16 Grassmann degrees of freedom $`(0,\theta )`$, and that as a consequence all expressions of the forms
$$\overline{\theta }\mathrm{\Gamma }^\mu _\alpha \theta ,\mu $$
(2.90)
and
$$\overline{\theta }\mathrm{\Gamma }^{ij}_\alpha \theta \mathrm{or}\overline{\theta }\mathrm{\Gamma }^{+\mu }_\alpha \theta $$
(2.91)
must vanish in this gauge. This simplifies the theory in this gauge considerably. First, we have
$$\mathrm{\Pi }_\alpha ^\mu =_\alpha X^\mu ,\mu $$
(2.92)
Second, we find that the terms on the second line of (2.6.2) simplify to
$$\overline{\theta }\mathrm{\Gamma }_{+i}\{X^i,\theta \}$$
(2.93)
Solving for the derivatives $`_\gamma X^{}`$ as in (2.33) we get
$`\dot{X}^{}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\dot{X}^i\dot{X}^i+{\displaystyle \frac{1}{\nu ^2}}\{X^i,X^j\}\{X^i,X^j\}+\overline{\theta }\mathrm{\Gamma }_+\dot{\theta }`$ (2.94)
$`=`$ $`\mathrm{\Pi }_0^{}+\overline{\theta }\mathrm{\Gamma }_+\dot{\theta }`$
$`_aX^{}`$ $`=`$ $`\dot{X}^i_aX^i+\overline{\theta }\mathrm{\Gamma }_+_a\theta `$ (2.95)
$`=`$ $`\mathrm{\Pi }_a^{}+\overline{\theta }\mathrm{\Gamma }_+_a\theta `$
Combining these observations, we find that the light-front supermembrane Hamiltonian becomes
$$H=\frac{\nu T}{4}d^2\sigma \left(\dot{X}^i\dot{X}^i+\frac{2}{\nu ^2}\{X^i,X^j\}\{X^i,X^j\}\frac{2}{\nu }\theta ^T\gamma _i\{X^i,\theta \}\right)$$
(2.96)
where $`\theta `$ is a 16-component Majorana spinor of $`SO(9)`$.
It is straightforward to apply the matrix regularization procedure discussed in section 2.4 to this Hamiltonian. This gives the supersymmetric form of matrix theory
$$H=\frac{1}{(2\pi l_p^3)}\mathrm{Tr}\left(\frac{1}{2}\dot{𝐗}^i\dot{𝐗}^i\frac{1}{4}[𝐗^i,𝐗^j][𝐗^i,𝐗^j]+\frac{1}{2}\theta ^T\gamma _i[𝐗^i,\theta ]\right).$$
(2.97)
This is the matrix quantum mechanics theory which will play a central role in these lectures. This theory was derived in from the regularized supermembrane action, but had been previously found and studied as a particularly simple example of a supersymmetric theory with gauge symmetry .
### 2.7 Covariant membrane quantization
It is natural to think of generalizing the matrix regularization approach to the covariant formulation of the bosonic and supersymmetric membrane theories (2.26) and (2.6.2). Some progress was made in this direction by Fujikawa and Okuyama in . For the bosonic membrane it is straightforward to implement the matrix regularization procedure. The only catch is that the BRST charge needed to implement the gauge-fixing procedure cannot be simply expressed in terms of the Poisson bracket on the membrane. For the supermembrane, there is a more serious complication related to the $`\kappa `$-symmetry of the theory. Essentially, as mentioned above, any gauge-fixing of the $`\kappa `$-symmetry will break the 11-dimensional Lorentz invariance of the theory. This is the same difficulty as one encounters when trying to construct a covariant quantization of the Green-Schwarz superstring. The approach taken in the second paper of is to fix the $`\kappa `$-symmetry in a way which breaks the 32 of SO(10, 1) into $`16_R+16_L`$ of SO(9, 1). Thus, they end up with a matrix formulation of a theory with explicit SO(9, 1) Lorentz symmetry. Although this theory does not have the desired complete SO(10, 1) Lorentz symmetry of M-theory, there are many questions which might be addressed by this theory with limited Lorentz invariance. It would be interesting to study the quantum mechanics of this alternative matrix formulation of M-theory in further detail.
## 3 The BFSS conjecture
As we have already discussed, the fact that the light-front supermembrane theory can be regularized and described as a supersymmetric quantum mechanics theory has been known for over a decade. At the time that this theory was first developed, however, it was believed that the quantum supermembrane theory suffered from instabilities which would make the low-energy interpretation as a theory of quantized gravity impossible. In 1996 the supersymmetric Yang-Mills quantum mechanics theory was brought back into currency as a candidate for a microscopic description of an 11-dimensional quantum mechanical theory containing gravity by Banks, Fischler, Shenker and Susskind (BFSS). The BFSS proposal, which quickly became known as the “Matrix Theory Conjecture” was motivated not by the quantum supermembrane theory, but by considering the low-energy theory of a system of many D0-branes as a partonic description of light-front M-theory.
In this section we discuss the apparent instability of the quantized membrane theory and the BFSS conjecture. We describe the membrane instability in subsection 3.1. We give a brief introduction to M-theory in section 3.2, and describe the BFSS conjecture in subsection 3.3. In subsection 3.4 we describe the resolution of the apparent instability of the membrane theory by an interpretation in terms of a second-quantized gravity theory. Finally, in subsection 3.5 we review an argument due to Seiberg and Sen which shows that matrix theory should be equivalent to a discrete light-front quantization of M-theory, even at finite $`N`$.
### 3.1 Membrane “instability”
At the time that de Wit, Hoppe and Nicolai wrote the paper showing that the regularized supermembrane theory could be described in terms of supersymmetric matrix quantum mechanics, the general hope seemed to be that the quantized supermembrane theory would have a discrete spectrum of states. In string theory the spectrum of states in the Hilbert space of the string can be put into one-to-one correspondence with elementary particle-like states in the target space. The facts that the massless particle spectrum contains a graviton and that there is a mass gap separating the massless states from massive excitations are crucial for this interpretation. For the supermembrane theory, however, the spectrum does not seem to have these properties. This can be seen in both the classical and quantum membrane theories. In this section we discuss this apparent difficulty with the membrane theory, which was first described in detail in .
The simplest way to see the instability of the membrane theory is to consider a classical bosonic membrane whose energy is proportional to the area of the membrane times a constant tension. Such a membrane can have long narrow spikes at very low cost in energy (See Figure 1). If the spike is roughly cylindrical and has a radius $`r`$ and length $`L`$ then the energy is $`2\pi TrL`$. For a spike with very large $`L`$ but a small radius $`r1/TL`$ the energy cost is small but the spike is very long. This indicates that a quantum membrane will tend to have many fluctuations of this type, making it difficult to think of the membrane as a single pointlike object. Note that the quantum string theory does not have this problem since a long spike in a string always has energy proportional to the length of the string.
In the matrix regularized version of the membrane theory, this instability appears as a set of flat directions in the classical theory. For example, if we have a pair of matrices with nonzero entries of the form
$$X^1=\left(\begin{array}{cc}x& 0\\ 0& 0\end{array}\right)X^2=\left(\begin{array}{cc}0& y\\ y& 0\end{array}\right)$$
(3.1)
then a potential term $`[X^1,X^2]^2`$ corresponds to a term $`x^2y^2`$. If either $`x=0`$ or $`y=0`$ then the other variable is unconstrained, giving flat directions in the moduli space of solutions to the classical equations of motion. This corresponds classically to a marginal instability in the matrix theory with $`N>1`$. (Note that in the previous section we distinguished matrices $`𝐗^i`$ from related functions $`X^i`$ by using bold font for matrices. We will henceforth drop this font distinction as long as the difference can easily be distinguished from context.)
In the quantum bosonic membrane theory, the apparent instability from the flat directions is cured because of the 0-modes of off-diagonal degrees of freedom. In the above example, for instance, if $`x`$ takes a large value then $`y`$ corresponds to a harmonic oscillator degree of freedom with a large mass. The zero point energy of this oscillator becomes larger as $`x`$ increases, giving an effective confining potential which removes the flat directions of the classical theory. This would seem to resolve the instability problem. Indeed, in the matrix regularized quantum bosonic membrane theory, there is a discrete spectrum of energy levels for the system of $`N\times N`$ matrices.
When we consider the supersymmetric theory, on the other hand, the problem reemerges. The zero point energies of the fermionic oscillators associated with the extra Grassmann degrees of freedom in the supersymmetric theory conspire to precisely cancel the zero point energies of the bosonic oscillators. This cancellation gives rise to a continuous spectrum in the supersymmetric matrix theory. This result was formally proven by de Wit, Lüscher and Nicolai in . They showed that for any $`ϵ>0`$ and any energy $`E[0,\mathrm{})`$ there exists a state $`\psi `$ in the $`N=2`$ maximally supersymmetric matrix model which is normalizable ($`|\psi |^2=\psi ^2=1`$) and which has
$$(HE)\psi ^2<ϵ.$$
This implies that the spectrum of the supersymmetric matrix quantum mechanics theory is continuous<sup>1</sup><sup>1</sup>1Note that did not resolve the question of whether a state existed with identically vanishing energy $`H=0`$. This question was not resolved until the much later work of Sethi and Stern showed that such a marginally bound state does indeed exist in the maximally supersymmetric theory. This result indicated that it would not be possible to have a simple interpretation of the states of the theory in terms of a discrete particle spectrum. After this work there was little further development on the supersymmetric matrix quantum mechanics theory as a theory of membranes or gravity until almost a decade later.
### 3.2 M-theory
The concept of M-theory has played a fairly central role in the development of the web of duality symmetries which relate the five string theories to each other and to supergravity . M-theory is a conjectured eleven-dimensional theory whose low-energy limit corresponds to 11D supergravity. Although there are difficulties with constructing a quantum version of 11D supergravity, it is a well-defined classical theory with the following field content :
$`e_I^a`$: a vielbein field (bosonic, with 44 components)
$`\psi _I`$: a Majorana fermion gravitino (fermionic, with 128 components)
$`A_{IJK}`$: a 3-form potential (bosonic, with 84 components).
In addition to being a local theory of gravity with an extra 3-form potential field, M-theory also contains extended objects. These consist of a two-dimensional supermembrane and a 5-brane, which couple electrically and magnetically to the 3-form field.
One way of defining M-theory is as the strong coupling limit of the type IIA string. The IIA string theory is taken to be equivalent to M-theory compactified on a circle $`S^1`$, where the radius of compactification $`R`$ of the circle in direction 11 is related to the string coupling $`g`$ through $`R=g^{2/3}l_p=gl_s`$, where $`l_p`$ and $`l_s=\sqrt{\alpha ^{}}`$ are the M-theory Planck length and the string length respectively. The decompactification limit $`R\mathrm{}`$ corresponds then to the strong coupling limit of the IIA string theory. (Note that we will always take the eleven dimensions of M-theory to be labeled $`0,1,\mathrm{},8,9,11`$; capitalized roman indices $`I,J,\mathrm{}`$ denote 11-dimensional indices).
Given this relationship between compactified M-theory and IIA string theory, a correspondence can be constructed between various objects in the two theories. For example, the Kaluza-Klein photon associated with the components $`g_{\mu 11}`$ of the 11D metric tensor can be associated with the R-R gauge field $`A_\mu `$ in IIA string theory. The only object which is charged under this R-R gauge field in IIA string theory is the D0-brane; thus, the D0-brane can be associated with a supergraviton with momentum $`p_{11}`$ in the compactified direction. The membrane and 5-brane of M-theory can be associated with different IIA objects depending on whether or not they are wrapped around the compactified direction; the correspondence between various M-theory and IIA objects is given in Table 1.
### 3.3 The BFSS conjecture
In 1996, motivated by recent work on D-branes and string dualities, Banks, Fischler, Shenker and Susskind (BFSS) proposed that the large $`N`$ limit of the supersymmetric matrix quantum mechanics model (2.97) should describe all of M-theory in a light-front coordinate system . Although this conjecture fits neatly into the framework of the quantized membrane theory, the starting point of BFSS was to consider M-theory compactified on a circle $`S^1`$, with a large momentum in the compact direction. As we have just discussed, when M-theory is compactified on $`S^1`$ the corresponding theory in 10D is the type IIA string theory, and the quanta corresponding to momentum in the compact direction are the D0-branes of the IIA theory. In the limits where the radius of compactification $`R`$ and the compact momentum $`p_{11}`$ are both taken to be large, this correspondence relates M-theory in the “infinite momentum frame” (IMF) to the nonrelativistic theory of many D0-branes in type IIA string theory.
The low-energy Lagrangian for a system of many type IIA D0-branes is the matrix quantum mechanics Lagrangian arising from the dimensional reduction to 0 + 1 dimensions of the 10D super Yang-Mills Lagrangian
$$=\frac{1}{2gl_s}\mathrm{Tr}\left[\dot{X}^a\dot{X}^a+\frac{1}{2}[X^a,X^b]^2+\theta ^T(i\dot{\theta }\mathrm{\Gamma }_a[X^a,\theta ])\right]$$
(3.2)
(the gauge has been fixed to $`A_0=0`$.) The corresponding Hamiltonian is
$$H=\frac{1}{2gl_s}\mathrm{Tr}\left(\dot{X}^i\dot{X}^i\frac{1}{2}[X^i,X^j][X^i,X^j]+\theta ^T\gamma _i[X^i,\theta ]\right).$$
(3.3)
Using the relations $`R=g^{2/3}l_{11}=gl_s`$, we see that in string units ($`2\pi l_s^2=1`$) we can replace $`gl_s=R=2\pi l_{11}^3`$. So the Hamiltonian (3.3) arising in the matrix quantum mechanics picture is in fact precisely equivalent to the matrix membrane Hamiltonian (2.97). This connection and its possible physical significance was first pointed out by Townsend . The matrix theory Hamiltonian is often written, following BFSS, in the form
$`H`$ $`=`$ $`{\displaystyle \frac{R}{2}}\mathrm{Tr}\left(P^iP^i{\displaystyle \frac{1}{2}}[X^i,X^j][X^i,X^j]+\theta ^T\gamma _i[X^i,\theta ]\right)`$ (3.4)
where we have rescaled $`X/g^{1/3}X`$ and written the Hamiltonian in Planck units $`l_{11}=1`$.
The original BFSS conjecture was made in the context of the large $`N`$ theory. It was later argued by Susskind that the finite $`N`$ matrix quantum mechanics theory should be equivalent to the discrete light-front quantized (DLCQ) sector of M-theory with $`N`$ units of compact momentum . We describe in section (3.5) below an argument due to Seiberg and Sen which makes this connection more precise and which justifies the use of the low-energy D0-brane action in the BFSS conjecture.
While the BFSS conjecture was based on a different framework from the matrix quantization of the supermembrane theory we have discussed above, the fact that the membrane naturally appears as a coherent state in the matrix quantum mechanics theory was a substantial piece of additional evidence given by BFSS for the validity of their conjecture. Two additional pieces of evidence were given by BFSS which extended their conjecture beyond the previous work on the matrix membrane theory.
One important point made by BFSS is that the Hilbert space of the matrix quantum mechanics theory naturally contains multiple particle states. This observation, which we discuss in more detail in the following subsection, resolves the problem of the continuous spectrum discussed above. Another piece of evidence given by BFSS for their conjecture is the fact that quantum effects in matrix theory give rise to long-range interactions between a pair of gravitational quanta (D0-branes) which have precisely the correct form expected from light-front supergravity. This result was first shown by a calculation of Douglas, Kabat, Pouliot and Shenker ; we will discuss this result and its generalization to more general matrix theory interactions in Section 5.
### 3.4 Matrix theory as a second quantized theory
The classical equations of motion for a bosonic matrix configuration with the Hamiltonian (2.97) are (up to an overall constant)
$$\ddot{X}^i=[[X^i,X^j],X^j].$$
(3.5)
If we consider a block-diagonal set of matrices
$$X^i=\left(\begin{array}{cc}\widehat{X}^i& 0\\ 0& \stackrel{~}{X}^i\end{array}\right)$$
with first time derivatives $`\dot{X}^i`$ which are also of block-diagonal form, then the classical equations of motion for the blocks are separable
$`\ddot{\widehat{X}}^i`$ $`=`$ $`[[\widehat{X}^i,\widehat{X}^j],\widehat{X}^j]`$
$`\ddot{\stackrel{~}{X}}^i`$ $`=`$ $`[[\stackrel{~}{X}^i,\stackrel{~}{X}^j],\stackrel{~}{X}^j]`$
If we think of each of these blocks as describing a matrix theory object with center of mass
$$\widehat{x}^i=\frac{1}{\widehat{N}}\mathrm{Tr}\widehat{X}^i$$
$$\stackrel{~}{x}^i=\frac{1}{\stackrel{~}{N}}\mathrm{Tr}\stackrel{~}{X}^i$$
then we have two objects obeying classically independent equations of motion (See Figure 2). It is straightforward to generalize this construction to a block-diagonal matrix configuration describing $`k`$ classically independent objects. This gives a simple indication of how matrix theory can encode, even in finite $`N`$ matrices, a configuration of multiple objects. In this sense it is natural to think of matrix theory as a second quantized theory from the point of view of the target space.
Given the realization that matrix theory should describe a second quantized theory, the puzzle discussed above regarding the continuous spectrum of the theory is easily resolved. If there is a state in matrix theory corresponding to a single graviton of M-theory (as we will discuss in more detail in section 4.1) with $`H=0`$ which is roughly a localized state, then by taking two such states to have a large separation and a small relative velocity $`v`$ it should be possible to construct a two-body state with an arbitrarily small total energy. Since the D0-branes of the IIA theory correspond to gravitons in M-theory with a single unit of longitudinal momentum, we would therefore naturally expect to have a continuous spectrum of energies even in the theory with $`N=2`$. This resolves the puzzle found by de Wit, Lüscher and Nicolai in a very pleasing fashion, which suggests that matrix theory is perhaps even more powerful than string theory, which only gives a first-quantized theory in the target space.
The second quantized nature of matrix theory can also be seen naturally in the continuous membrane theory. Recall that the instability of membrane theory appears in the classical theory of a continuous membrane when we consider the possibility of long thin spikes of negligible energy, as discussed in section 3.1. In a similar fashion, it is possible for a classical smooth membrane of fixed topology to be mapped to a configuration in the target space which looks like a system of multiple distinct macroscopic membranes connected by infinitesimal tubes of negligible energy (See Figure 3). In the limit where the tubes become very small, their effect on the classical dynamics of the multiple membrane configuration disappears and we effectively have a system of multiple independent membranes moving in the target space. At the classical level, the sum of the genera of the membranes in the target space must be equal to or smaller than the genus of the single world-sheet membrane, but when quantum effects are included handles can be added to the membrane as well as removed . These considerations seem to indicate that any consistent quantum theory which contains a continuous membrane in its effective low-energy theory must contain configurations with arbitrary membrane topology and must therefore be a “second quantized” theory from the point of view of the target space.
### 3.5 Matrix theory and DLCQ M-theory
A theory which has been compactified on a lightlike circle can be viewed as a limit of a theory compactified on a spacelike circle where the size of the spacelike circle becomes vanishingly small in the limit. This point of view was used by Seiberg and Sen in to argue that light-front compactified M-theory is described through such a limiting process by the low-energy Lagrangian for many D0-branes, and hence by matrix theory. In this section we go through this argument in detail.
Consider a space-time which has been compactified on a lightlike circle by identifying
$$\left(\begin{array}{c}x\\ t\end{array}\right)\left(\begin{array}{c}xR/\sqrt{2}\\ t+R/\sqrt{2}\end{array}\right)$$
(3.6)
This theory has a quantized momentum in the compact direction
$$P^+=\frac{N}{R}$$
(3.7)
The compactification (3.6) can be described as a limit of a family of spacelike compactifications
$$\left(\begin{array}{c}x\\ t\end{array}\right)\left(\begin{array}{c}x\sqrt{R^2/2+R_s^2}\\ t+R/\sqrt{2}\end{array}\right)$$
(3.8)
parameterized by the size $`R_s0`$ of the spacelike circle, which is taken to vanish in the limit.
The system satisfying (3.8) is related through a boost to a system with the identification
$$\left(\begin{array}{c}x^{}\\ t^{}\end{array}\right)\left(\begin{array}{c}x^{}R_s\\ t^{}\end{array}\right)$$
(3.9)
where
$$\left(\begin{array}{c}x^{}\\ t^{}\end{array}\right)=\left(\begin{array}{cc}\frac{1}{\sqrt{1\beta ^2}}& \frac{\beta }{\sqrt{1\beta ^2}}\\ \frac{\beta }{\sqrt{1\beta ^2}}& \frac{1}{\sqrt{1\beta ^2}}\end{array}\right)\left(\begin{array}{c}x\\ t\end{array}\right)$$
(3.10)
The boost parameter $`\beta `$ is given by
$$\beta =\frac{1}{\sqrt{1+\frac{2R_s^2}{R^2}}}1\frac{R_s^2}{R^2}.$$
(3.11)
In the context of matrix theory we are interested in understanding M-theory compactified on a lightlike circle. This is related through the above limiting process to a family of spacelike compactifications of M-theory, which we know can be identified with the IIA string theory. At first glance, it may seem that the limit we are considering here is difficult to analyze from the IIA point of view. The IIA string coupling and string length are related to the compactification radius and 11D Planck length as usual by
$`g`$ $`=`$ $`({\displaystyle \frac{R_s}{l_{11}}})^{3/2}`$ (3.12)
$`l_s^2`$ $`=`$ $`{\displaystyle \frac{l_{11}^3}{R_s}}`$
Thus, in the limit $`R_s\mathrm{}`$ the string coupling $`g`$ becomes small as desired; the string length $`l_s`$, however, goes to $`\mathrm{}`$. Since $`l_s^2=\alpha ^{}`$, this corresponds to a limit of vanishing string tension. Such a limiting theory is very complicated and would not seem to provide a useful alternative description of the theory.
Let us consider, however, how the energy of the states we are interested in behaves in the class of limiting theories with spacelike compactification. If we want to describe the behavior of a state which has light-front energy $`P^{}`$ and compact momentum $`P^+=N/R`$ then the spatial momentum in the theory with spatial $`R_s`$ compactification is $`P^{}=N/R_s`$. The energy in the spatially compactified theory is
$$E^{}=N/R_s+\mathrm{\Delta }E,$$
(3.13)
where $`\mathrm{\Delta }E`$ has the energy scale we are interested in understanding. The term $`N/R_s`$ in the energy is simply the mass-energy of the $`N`$ D0-branes which correspond to the momentum in the compactified M-theory direction. Relating back to the near lightlike compactified theory we have
$`\left(\begin{array}{c}P\\ E\end{array}\right)`$ $`=`$ $`\left(\begin{array}{cc}\frac{1}{\sqrt{1\beta ^2}}& \frac{\beta }{\sqrt{1\beta ^2}}\\ \frac{\beta }{\sqrt{1\beta ^2}}& \frac{1}{\sqrt{1\beta ^2}}\end{array}\right)\left(\begin{array}{c}P^{}\\ E^{}\end{array}\right)`$ (3.20)
so
$$P^{}=\frac{1}{\sqrt{2}}(EP)=\frac{1}{\sqrt{2}}\frac{1+\beta }{\sqrt{1\beta ^2}}\mathrm{\Delta }E\frac{R}{R_s}\mathrm{\Delta }E$$
(3.21)
As a result we see that the energy $`\mathrm{\Delta }E`$ of the IIA configuration needed to approximate the light-front energy $`P^{}`$ is given by
$$\mathrm{\Delta }EP^{}\frac{R_s}{R}$$
(3.22)
We know that the string scale $`1/l_s`$ becomes small as $`R_s0`$. We can compare the energy scale of interest to this string scale, however, and find
$$\frac{\mathrm{\Delta }E}{(1/l_s)}=\frac{P^{}}{R}R_sl_s=\frac{P^{}}{R}\sqrt{R_sl_{11}^3}$$
(3.23)
This ratio vanishes in the limit $`R_s0`$, which implies that although the string scale vanishes, the energy scale of interest is smaller still. Thus, it is reasonable to study the lightlike compactification through a limit of spatial compactifications in this fashion.
To make the correspondence between the light-front compactified theory and the spatially compactified limiting theories more transparent, we perform a change of units to a new Planck length $`\stackrel{~}{l}_{11}`$ in the spatially compactified theories in such a way that the energy of the states of interest is independent of $`R_s`$. For this condition to hold we must have
$$\mathrm{\Delta }E\stackrel{~}{l}_{11}=P^{}\frac{R_sl_{11}^2}{R\stackrel{~}{l}_{11}}$$
(3.24)
where $`E,R`$ and $`P^{}`$ are independent of $`R_s`$ and all units have been explicitly included. This requires us to keep the quantity
$$\frac{R_s}{\stackrel{~}{l}_{11}^2}$$
(3.25)
fixed in the limiting process. Thus, in the limit $`\stackrel{~}{l}_{11}0`$.
We can now summarize the discussion with the following story: to describe the sector of M-theory corresponding to light-front compactification on a circle of radius $`R`$ with light-front momentum $`P^+=N/R`$ we may consider the limit $`R_s0`$ of a family of IIA configurations with $`N`$ D0-branes where the string coupling and string length
$`\stackrel{~}{g}`$ $`=`$ $`(R_s/\stackrel{~}{l}_{11})^{3/2}0`$
$`\stackrel{~}{l}_s`$ $`=`$ $`\sqrt{\stackrel{~}{l}_{11}^3/R_s}0`$ (3.26)
are defined in terms of a Planck length $`\stackrel{~}{l}_{11}`$ and compactification length $`R_s`$ which satisfy
$$R_s/\stackrel{~}{l}_{11}^2=R/l_{11}^2$$
(3.27)
All transverse directions scale normally through
$$\stackrel{~}{x}^i/\stackrel{~}{l}_{11}=x^i/l_{11}$$
(3.28)
To give a very concrete example of how this limiting process works, let us consider a system with a single unit of longitudinal momentum
$$P^+=\frac{1}{R}$$
(3.29)
We know that in the corresponding IIA theory, we have a single D0-brane whose Lagrangian has the Born-Infeld form
$$=\frac{1}{\stackrel{~}{g}\stackrel{~}{l}_s}\sqrt{1\dot{\stackrel{~}{x}}^i\dot{\stackrel{~}{x}}^i}$$
(3.30)
Expanding the square root we have
$$=\frac{1}{\stackrel{~}{g}\stackrel{~}{l}_s}\left(1\frac{1}{2}\dot{\stackrel{~}{x}}^i\dot{\stackrel{~}{x}}^i+𝒪(\dot{\stackrel{~}{x}}^4)\right).$$
(3.31)
Replacing $`\stackrel{~}{g}\stackrel{~}{l}_sR_s`$ and $`\stackrel{~}{x}x\stackrel{~}{l}_{11}/l_{11}`$ gives
$$=\frac{1}{R_s}+\frac{1}{2R}\dot{x}^i\dot{x}^i+𝒪(R_s/R).$$
(3.32)
Thus, we see that all the higher order terms in the Born-Infeld action vanish in the $`R_s0`$ limit. The leading term is the D0-brane energy $`1/R_s`$ which we subtract to compare to the M-theory light-front energy $`P^{}`$. Although we do not know the full form of the nonabelian Born-Infeld action describing $`N`$ D0-branes in IIA, it is clear that an analogous argument shows that all terms in this action other than those in the nonrelativistic supersymmetric matrix theory action will vanish in the limit $`R_s0`$.
This argument apparently demonstrates that matrix theory gives a complete description of the dynamics of DLCQ M-theory. There are several caveats which should be taken into account, however, with respect to this discussion. First, in order for this argument to be correct, it is necessary that there exists a well-defined theory with the properties expected of M-theory, and that there exist a well-defined IIA string theory which arises as the compactification of M-theory. Neither of these statements is at this point definitely established. Thus, this argument must be taken as contingent upon the definition of these theories. Second, although we know that 11D supergravity arises as the low-energy limit of M-theory, this argument does not necessarily indicate that matrix theory describes DLCQ supergravity in the low-energy limit. It may be that to make the connection to supergravity it is necessary to deal with subtleties of the large $`N`$ limit.
In the remainder of these lectures we will discuss some more explicit approaches to connecting matrix theory with supergravity. In particular, we will see how far it is possible to go in demonstrating that 11D supergravity arises from calculations in the finite $`N`$ version of matrix theory, which is a completely well-defined theory. In the last sections we will return to a more general discussion of the status of matrix theory.
## 4 M-theory objects from matrix theory
In this section we discuss how the matrix theory degrees of freedom can be used to construct the various objects of M-theory: the supergraviton, supermembrane and 5-brane. We discuss each of these objects in turn in subsections 4.1, 4.2, 4.3, after which we give a general discussion of the structure of extended objects and their charges in subsection 4.4
### 4.1 Supergravitons
Since in DLCQ M-theory there should be a pointlike state corresponding to a longitudinal graviton with $`p^+=N/R`$ and arbitrary transverse momentum $`p^i`$, we expect from the massless condition $`m^2=p^Ip_I=0`$ that such an object will have matrix theory energy
$$E=\frac{p_i^2}{2p^+}$$
We discuss such states first classically and then in the quantum theory.
#### 4.1.1 Classical supergravitons
The classical matrix theory potential is $`[X^i,X^j]^2`$, from which we have the classical equations of motion
$$\ddot{X}^i=[[X^i,X^j],X^j].$$
One simple class of solutions to these equations of motion can be found when the matrices minimize the potential at all times and therefore all commute. Such solutions are of the form
$$X^i=\left(\begin{array}{cccc}x_1^i+v_1^it& 0& 0& \mathrm{}\\ 0& x_2^i+v_2^it& \mathrm{}& 0\\ 0& \mathrm{}& \mathrm{}& 0\\ \mathrm{}& 0& 0& x_N^i+v_N^it\end{array}\right)$$
This corresponds to a classical $`N`$-graviton solution, where each graviton has
$$p_a^+=1/Rp_a^i=v_a^i/RE_a=v_a^2/(2R)=(p_a^i)^2/2p^+$$
A single classical graviton with $`p^+=N/R`$ can be formed by setting
$$x_1^i=\mathrm{}=x_N^i,v_1^i=\mathrm{}=v_N^i$$
so that the trajectories of all the components are identical. Although this may seem like a very simple model for a graviton, it is precisely such matrix configurations which are used as a background in most computations of quantum effects in matrix theory corresponding to gravitational interactions, as will be discussed further in the following.
#### 4.1.2 Quantum supergravitons
The picture of a supergraviton in quantum matrix theory is somewhat more subtle than the simple classical picture just discussed. Let us first consider the case of a single supergraviton with $`p^+=1/R`$. This corresponds to the U(1) case of the super Yang-Mills quantum mechanics theory. The Hamiltonian is simply
$$H=\frac{1}{2R}\dot{X}^2$$
since all commutators vanish in this theory. The bosonic part of the theory is simply a free nonrelativistic particle. In the fermionic sector there are 16 spinor variables with anticommutation relations
$$\{\theta _\alpha ,\theta _\beta \}=\delta _{\alpha \beta }.$$
By using the standard trick of writing these as 8 fermion creation and annihilation operators
$$\theta _i^\pm =\frac{1}{\sqrt{2}}(\theta _i\pm \theta _{i+8})\mathrm{\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}1}i8$$
we see that the Hilbert space for the fermions is a standard fermion Fock space of dimension $`2^8=256`$. Indeed, this is precisely the number of states needed to represent all the polarization states of the graviton (44), the antisymmetric 3-tensor field (84) and the gravitino (128). For details of how the polarization states are represented in terms of the fermionic Fock space, see .
The case when $`N>1`$ is much more subtle. We can factor out the overall $`U(1)`$ so that every state in the $`SU(N)`$ quantum mechanics theory has 256 corresponding states in the full theory. For the matrix theory conjecture to be correct, as BFSS pointed out, it should then be the case that for every $`N`$ there exists a unique threshold bound state in the $`SU(N)`$ theory with $`H=0`$. As mentioned before, no definitive answer as to the existence of such a state was given in the early work on matrix theory. This result was finally proven by Sethi and Stern for $`N=2`$ in . Progress towards proving the result for arbitrary values of $`N`$ was made in , and the result for a general gauge group was given in (see also ).
### 4.2 Membranes
In this section we discuss the description of M-theory membranes in terms of the matrix quantum mechanics degrees of freedom. It is clear from the derivation of matrix theory as a regularized supermembrane theory that there must be matrix configurations which in the large $`N`$ limit give arbitrarily good descriptions of any membrane configuration. It is instructive, however, to study in detail the structure of such membrane configurations. In subsection 4.2.1 we discuss the significance of the matrix representation of membranes in the language of type IIA D0-branes. In subsection 4.2.2 we discuss in some detail how a spherical membrane can be very accurately described by matrices even with small values of $`N`$. In subsection 4.2.3 we discuss higher genus matrix membranes. In subsection 4.2.4 we discuss noncompact matrix membranes, and finally in subsection 4.2.5 we discuss M-theory membranes which are wrapped on the longitudinal direction and appear as strings in the IIA theory.
#### 4.2.1 D2-branes from D0-branes
As we have mentioned, it is clear from inverting the matrix membrane regularization procedure that smooth membranes can be approximated by finite size matrices. This construction may seem less natural in the language of type IIA string theory, where it corresponds to a construction of a IIA D2-brane out of the degrees of freedom describing a system of $`N`$ D0-branes. In fact, however, the fact that this construction is possible is simply the T-dual of the familiar statement that D0-branes are described by the magnetic flux of the gauge field living on a set of $`N`$ D2-branes . Both of these statements can in turn be seen by performing T-duality on a diagonally wrapped D1-brane on a 2-torus.
To see this explicitly, consider a set of $`N`$ D2-branes on a torus $`T^2`$ with $`k`$ units of magnetic flux
$$\frac{1}{2\pi }F=k$$
(4.1)
Under a T-duality transformation on one direction of the torus, the gauge field component $`A_2`$ is replaced by an infinite matrix
$$X^2=i_2+A_2$$
representing a transverse scalar field for a set of $`N`$ D1-branes living on the dual torus $`(T^2)^{}`$. These matrices are infinite because they contain information about winding strings connecting the infinite number of copies of each brane which live on the infinite covering space of the dual torus. (This construction is described in more detail in section 6.1.) This T-dual configuration corresponds to a single D-string which is diagonally wound $`N`$ times around the $`x^1`$ direction and $`k`$ times around the $`x^2`$ direction; this can be seen from the fact that the T-dual of (4.1) is $`_1X^2=(kL_2^{}/NL_1)\text{1}\text{ }\text{1}`$. Since under T-duality in the $`x^2`$ direction a D1-brane wrapped in the $`x^2`$ direction becomes a D0-brane, we can identify the flux (4.1) with $`k`$ D0-branes in the original theory.
Further T-dualizing in the direction $`x^1`$, we replace
$$X^1=i_1+A_1.$$
where $`X^1,X^2`$ are now infinite matrices describing transverse fields of a system of $`N`$ D0-branes on the dual torus $`(T^2)^{}`$. When the normalization constants are treated carefully, the flux condition (4.1) now becomes the condition on the D0-brane matrices
$$\mathrm{Tr}[X^1,X^2]=\frac{iAk}{2\pi }$$
(4.2)
where $`A`$ is the area of the dual torus. Since the T-dual in the $`x^1`$ direction of the D-string wrapped in the $`x^2`$ direction is a D2-brane, we interpret $`k`$ in (4.2) as the D2-brane charge of a system of $`N`$ D0-branes.
This construction can be interpreted more generally, so that in general a pair of matrices $`X^a,X^b`$ describing a D0-brane configuration satisfying
$$\mathrm{Tr}[X^a,X^b]=\frac{iA}{2\pi }$$
(4.3)
should be interpreted as giving rise to a piece of a D2-brane of area $`A`$. Of course, for finite matrices the trace of the commutator must vanish. This is simply a consequence of the fact that the net D2-brane charge of any compact object must vanish. However, not only is it possible to have a nonzero membrane charge when the matrices are infinite, but it is also possible to treat (4.3) as a local expression by restricting the trace to a subset of the diagonal elements. We will see a specific example of this in the next subsection. The local relation (4.3) will also be useful in constructing higher moments of the membrane charge, which can be nonzero even for finite size configurations, as we shall discuss later.
#### 4.2.2 Spherical membranes
One extremely simple example of a membrane configuration which can be approximated very well even at finite $`N`$ by simple matrix configurations is the symmetric spherical membrane . Imagine that we wish to construct a membrane embedded in an isotropic sphere
$$x_1^2+x_2^2+x_3^2=r^2$$
in the first three dimensions of $`R^{11}`$. The embedding functions for such a continuous membrane can be written as linear functions
$$X^i=r\xi ^i\mathrm{\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}1}i3$$
of the three Euclidean coordinates $`\xi ^i`$ on the spherical world-volume. Using the matrix-membrane correspondence (2.51) we see that the matrix approximation to this membrane will be given by the $`N\times N`$ matrices
$$𝐗^i=\frac{2r}{N}J^i\mathrm{\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}1}i3$$
(4.4)
where $`J^i`$ are the generators of $`SU(2)`$ in the $`N`$-dimensional representation.
It is quite interesting to see how many of the geometrical and physical properties of the sphere can be extracted from the algebraic structure of these matrices, even for small values of $`N`$. We list here some of these properties.
i) Spherical locus: The matrices (4.4) satisfy
$$𝐗_1^2+𝐗_2^2+𝐗_3^2=\frac{4r^2}{N^2}C_2(N)\text{1}\text{ }\text{1}=r^2(11/N^2)\text{1}\text{ }\text{1}$$
where $`C_2(N)=(N^21)/4`$ is the quadratic Casimir of $`SU(2)`$ in the $`N`$-dimensional representation. This shows that the D0-branes are in a noncommutative sense “localized” on a sphere of radius $`r+𝒪(1/N^2)`$.
ii) Rotational invariance: The matrices (4.4) satisfy
$$R_{ij}𝐗_j=U(R)𝐗_iU(R^1)$$
where $`RSO(3)`$ and $`U(R)`$ is the $`N`$-dimensional representation of $`R`$. Thus, the spherical matrix configuration is rotationally invariant up to a gauge transformation.
iii) Spectrum: The matrix $`𝐗^3=2rJ_3/N`$ (as well as the other matrices) has a spectrum of eigenvalues which are uniformly distributed in the interval $`[r,r]`$. This is precisely the correct distribution if we imagine a perfectly symmetric sphere with D0-branes distributed uniformly on its surface and project this distribution onto a single axis.
iv) Local membrane charge: As discussed above, the expression (4.3) gives an area for a piece of a membrane described by a pair of matrices. We can use this formula to check the interpretation of the matrix sphere. We do this by computing the membrane charge in the 1-2 plane of the half of the configuration with eigenvalues $`X^3>0`$. This should correspond to the projected area of the “upper hemisphere” of the sphere. We compute
$$A_h=2\pi i\mathrm{Tr}_{1/2}[𝐗^1,𝐗^2]$$
where the trace is restricted to the set of eigenvalues where $`X^3>0`$ in the standard representation. This is possible since $`[𝐗^1,𝐗^2]𝐗^3`$. We find
$$A_h=2\pi \frac{4}{N^2}r^2\mathrm{Tr}_{1/2}J_3=\pi r^2(1+𝒪(1/N^2))$$
thus, we find precisely the expected area of the projected hemisphere.
v) Energy: In M-theory we expect the tension energy of a (momentarily) stationary membrane sphere to be
$$e=\frac{4\pi r^2}{(2\pi )^2l_{11}^3}=\frac{r^2}{\pi l_{11}^3}$$
Using $`p^Ip_I=e^2`$ we see that the light-front energy should be
$$E=\frac{e^2}{2p^+}$$
(4.5)
in 11D Planck units. Let us compute the matrix membrane energy. It is given by
$$E=\frac{1}{4R}[𝐗^i,𝐗^j]^2=\frac{2r^4}{NR}+𝒪(N^3)$$
in string units. This is easily seen to agree with (4.5).
It is also straightforward to verify that the equations of motion for the membrane are correctly reproduced in matrix theory.
Thus, we see that many of the geometrical and physical properties of the membrane can be extracted from algebraic information about the structure of the appropriate membrane configuration. The discussion we have carried out here has only applied to the simple case of the rotationally invariant spherically embedded membrane. It is straightforward to extend the discussion to a membrane of spherical topology and arbitrary shape, however, simply by using the matrix-membrane correspondence (2.51) to construct matrices approximating an arbitrary smooth spherical membrane. We now turn to the question of membranes with non-spherical topology.
#### 4.2.3 Higher genus membranes
So far we have only discussed membranes of spherical topology. It is possible to describe compact membranes of arbitrary genus by generalizing this construction, although an explicit construction is only known for the sphere and torus. In this section we give a brief description of the matrix torus, following the work of Fairlie, Fletcher and Zachos .
We consider a torus defined by two coordinates $`x_1,x_2[0,2\pi ]`$ with symplectic form $`\omega _{ij}=ϵ_{ij}/\pi `$ corresponding to a total volume $`d^2x\omega =4\pi `$ as in the case of the sphere discussed in section 2.4. As in the case of the sphere we wish to find a map from functions on the torus to matrices which is compatible with the correspondence
$$\{,\}\frac{iN}{2}[,]\frac{1}{4\pi ^2}d^2x\frac{1}{N}\mathrm{Tr}$$
(4.6)
A natural (complex) basis for the functions on $`T^2`$ is given by the Fourier modes
$$y_{nm}(x_1,x_2)=e^{inx_1+imx_2}$$
(4.7)
The real functions on $`T^2`$ are given by the linear combinations
$$\frac{1}{2}\left(y_{nm}+y_{nm}\right),\frac{i}{2}\left(y_{nm}y_{nm}\right).$$
(4.8)
The Poisson bracket algebra of the functions $`y_{nm}`$ is
$$\{y_{nm},y_{n^{}m^{}}\}=\pi (nm^{}mn^{})y_{n+n^{},m+m^{}}$$
(4.9)
To describe the matrix approximations for these functions we use the ’t Hooft matrices
$$U=\left(\begin{array}{ccccc}1& & & & \\ & q& & & \\ & & q^2& & \\ & & & \mathrm{}& \\ & & & & q^{N1}\end{array}\right)$$
(4.10)
and
$$V=\left(\begin{array}{ccccc}& 1& & & \\ & & 1& & \\ & & & \mathrm{}& \\ & & & & 1\\ 1& & & & \end{array}\right)$$
(4.11)
where
$$q=e^{\frac{2\pi i}{N}}.$$
(4.12)
The matrices $`U,V`$ satisfy
$$UV=q^1VU.$$
(4.13)
In terms of these matrices we can define
$$Y_{nm}=q^{nm/2}U^nV^m=q^{nm/2}V^mU^n$$
(4.14)
so that the matrix approximation to an arbitrary function
$$f(x_1,x_2)=\underset{n,m}{}c_{nm}y_{nm}(x_1,x_2)$$
(4.15)
is given by
$$F=\underset{n,m}{}c_{nm}Y_{nm}.$$
(4.16)
By computing
$`[Y_{nm},Y_{n^{}m^{}}]`$ $`=`$ $`(q^{(mn^{}nm^{})/2}q^{(nm^{}mn^{})/2})Y_{n+n^{},m+m^{}}`$
$``$ $`{\displaystyle \frac{2\pi i}{N}}(mn^{}nm^{})Y_{n+n^{},m+m^{}}`$
We see that for fixed $`n,m,n^{},m^{}`$ in the large $`N`$ limit the matrix commutation relations correctly reproduce (4.9) just as in the case of the sphere.
As a concrete example let us consider embedding a torus into $`R^4R^9`$ so that the membrane fills the locus of points satisfying
$$X_1^2+X_2^2=r^2X_3^2+X_4^2=s^2.$$
(4.17)
Such a membrane configuration can be realized through the following matrices
$`𝐗_1`$ $`=`$ $`{\displaystyle \frac{r}{2}}(U+U^{})`$
$`𝐗_2`$ $`=`$ $`{\displaystyle \frac{ir}{2}}(UU^{})`$ (4.18)
$`𝐗_3`$ $`=`$ $`{\displaystyle \frac{s}{2}}(V+V^{})`$
$`𝐗_4`$ $`=`$ $`{\displaystyle \frac{is}{2}}(VV^{})`$
It is straightforward to check that this matrix configuration has geometrical properties analogous to those of the matrix membrane sphere discussed in the previous subsection. In particular, the equation (4.17) is satisfied identically as a matrix equation. Note, however that this configuration is not gauge invariant under $`U(1)`$ rotations in the 12 and 34 planes—only under a $`Z_N`$ subgroup of each of these $`U(1)`$’s.
#### 4.2.4 Infinite membranes
So far we have discussed compact membranes, which can be described in terms of finite-size $`N\times N`$ matrices. In the large $`N`$ limit it is also possible to construct membranes with infinite spatial extent. The matrices $`X^i`$ describing such configurations are infinite-dimensional matrices which correspond to operators on a Hilbert space. Infinite membranes are of particular interest because they can be BPS states which solve the classical equations of motion of matrix theory. Extended compact membranes cannot be static solutions of the equations of motion since their membrane tension always causes them to contract and oscillate, as in the case of the spherical membrane.
The simplest infinite membrane is the flat planar membrane corresponding in IIA theory to an infinite D2-brane. This solution can be found by looking at the limit of the spherical membrane at large radius. It is simpler, however, to simply directly construct the solution by regularizing the flat membrane of M-theory. As in the other cases we have studied, we wish to quantize the Poisson bracket algebra of functions on the brane. Functions on the infinite membrane can be described in terms of two coordinates $`x_1,x_2`$ with a symplectic form $`\omega _{ij}=ϵ_{ij}`$ giving a Poisson bracket
$$\{f(x_1,x_2),g(x_1,x_2)\}=_1f_2g_1g_2f.$$
(4.19)
This algebra of functions can be “quantized” to the algebra of operators generated by $`Q,P`$ satisfying
$$[Q,P]=\frac{iϵ^2}{2\pi }\text{1}\text{ }\text{1}$$
(4.20)
where $`ϵ`$ is a constant parameter. As usual in the quantization process there are operator-ordering ambiguities which must be resolved in determining a general map from functions expressed as polynomials in $`x_1,x_2`$ to operators expressed as polynomials of $`Q,P`$.
This gives a map from functions on $`R^2`$ to operators which allows us to describe fluctuations around a flat membrane geometry with a single unit of $`P^+=1/R`$ in each region of area $`ϵ^2`$ on the membrane. Configurations of this type were discussed in the original BFSS paper and their existence used as additional evidence for the validity of their conjecture. Note that this configuration only makes sense in the large $`N`$ limit.
In addition to the flat membrane solution there are other infinite membranes which are static solutions of M-theory in flat space. In particular, there are BPS solutions corresponding to membranes which are holomorphically embedded in $`C^4=R^8R^9`$. These are static solutions of the membrane equations of motion. Finding a matrix theory description of such membranes is possible but involves some somewhat subtle issues related to choosing a regularization which preserves the complex structure of the brane. The details of this construction for a general holomorphic membrane are discussed in .
#### 4.2.5 Wrapped membranes as matrix strings
So far we have discussed M-theory membranes which are unwrapped in the longitudinal direction and which therefore appear as D2-branes in the IIA language of matrix theory. It is also possible to describe wrapped M-theory membranes which correspond to strings in the IIA picture. The charge in matrix theory which measures the number of strings present is proportional to
$$\frac{i}{R}\mathrm{Tr}\left([X^i,X^j]\dot{X}^j+[[X^i,\theta ^{\dot{\alpha }}],\theta ^{\dot{\alpha }}]\right)$$
(4.21)
This result can be understood in several ways. It was found in as a central charge in the matrix theory SUSY algebra corresponding to string charge; we will discuss this algebra further in the subsection 4.4. An intuitive way of understanding why (4.21) measures string charge is by a T-duality argument analogous to that used in 4.2.1 to derive the D2-brane charge of a system of D0-branes. If we compactify on a 2-torus in the $`i`$ and $`j`$ directions, the D0-branes become D2-branes and the bosonic part of (4.21) becomes
$$\frac{1}{R}F^{ij}F_{j0}.$$
(4.22)
This is the part of the energy-momentum tensor usually referred to as the Poynting vector in the 4D theory, and which describes momentum in the $`i`$ direction. Such momentum is of course T-dual to string winding in the original picture, so we understand the identification of the original charge (4.21) as counting fundamental IIA strings corresponding to wound M-theory membranes. Configurations with nonzero values of this charge were considered by Imamura in .
To realize a classical configuration in matrix theory which contains fundamental strings it is clear from the form of the charge that we need to construct configuration with local membrane charge extended in a pair of directions $`X^i,X^j`$ and to give the D0-branes velocity in the $`X^j`$ direction. For example, we could consider an infinite planar membrane (as discussed in the previous subsection) sliding along itself according to the equation
$`X^1`$ $`=`$ $`Q+t\text{1}\text{ }\text{1}`$ (4.23)
$`X^2`$ $`=`$ $`P`$ (4.24)
This corresponds to an M-theory membrane which has a projection onto the $`X^1,X^2`$ plane and which wraps around the compact direction as a periodic function of $`X^1`$ so that the IIA system contains a D2-brane with infinite strings extended in the $`X^2`$ direction since
$$\dot{X}^1[X^1,X^2]\text{1}\text{ }\text{1}.$$
(4.25)
Another example of a matrix theory system containing fundamental strings can be constructed by spinning the torus from (4.18) in the 12 plane to stabilize it. This gives the system some fundamental strings wrapped around the 34 circle. By taking the radius $`r`$ to be very small we can construct a configuration of a single fundamental string wrapped in a circle of radius $`s`$. As $`s\mathrm{}`$ this becomes an infinite fundamental string.
It is interesting to note that there is no classical matrix theory solution corresponding to a classical string which is truly 1-dimensional and has no local membrane charge. This follows from the appearance of the commutator $`[X^i,X^j]`$ in the string charge, which vanishes unless the matrices describe a configuration with at least two dimensions of spatial extent. We can come very close to a 1-dimensional classical string configuration by considering a one-dimensional array of D0-branes at equal intervals on the $`X^1`$ axis
$$X^1=a\left(\begin{array}{ccccc}\mathrm{}& \mathrm{}& \mathrm{}& & \\ \mathrm{}& 1& 0& \mathrm{}& \\ \mathrm{}& 0& 0& 0& \mathrm{}\\ \mathrm{}& \mathrm{}& 0& 1& \mathrm{}\\ & & \mathrm{}& \mathrm{}& \mathrm{}\end{array}\right)$$
(4.26)
We can now construct an excitation of the off-diagonal elements of $`X^2`$ corresponding to a string threading through the line of D0-branes
$$X^2=b\left(\begin{array}{ccccc}\mathrm{}& \mathrm{}& \mathrm{}& & \\ \mathrm{}& 0& e^{i\omega t}& \mathrm{}& \\ \mathrm{}& e^{i\omega t}& 0& e^{i\omega t}& \mathrm{}\\ \mathrm{}& \mathrm{}& e^{i\omega t}& 0& \mathrm{}\\ & & \mathrm{}& \mathrm{}& \mathrm{}\end{array}\right)$$
(4.27)
where $`\omega =a`$. In the classical theory, this configuration can have arbitrary string charge. If the mode (4.27) is quantized then the string charge is quantized in the correct units. This string is almost 1-dimensional but has a small additional extent in the $`X^2`$ direction corresponding to the extra dimension of the M-theory membrane. From the M-theory point of view this extra dimension must appear because the membrane cannot have momentum in a direction parallel to its direction of extension since it has no internal degrees of freedom. Thus, the momentum in the compact direction represented by the D0-branes must appear on the membrane as a fluctuation in some transverse direction.
### 4.3 5-branes
The M-theory 5-brane can appear in two possible guises in type IIA string theory. If the 5-brane is wrapped around the compact direction it becomes a D4-brane in the IIA theory, while if it is unwrapped it appears as an NS 5-brane. We will refer to these two configurations as “longitudinal” and “transverse” 5-branes in matrix theory. We begin by discussing the transverse 5-brane.
A priori, one might think that it should be possible to see both types of 5-branes in matrix theory. Several calculations, however, indicate that the transverse 5-brane does not carry a conserved charge which can be described in terms of the matrix degrees of freedom. In principle, if this charge existed we would expect it to appear both in the supersymmetry algebra of matrix theory (discussed in the next subsection) and in the set of supergravity currents whose interactions are described by perturbative matrix theory calculations (discussed in section 5.1.2). In fact, no charge or current with the proper tensor structure for a transverse 5-brane appears in either of these calculations.
One way of understanding this apparent puzzle is by comparing to the situation for D-branes in light-front string theory . Due to the Virasoro constraints, strings in the light-front formalism must have Neumann boundary conditions in both the light-front directions $`X^+,X^{}`$. Thus, in light-front string theory there are no transverse D-branes which can be used as boundary conditions for the string. A similar situation holds for membranes in M-theory, which can end on M5-branes. The boundary conditions on the bosonic membrane fields which can be derived from the action (2.22) state that
$$(\overline{h}h^{ab}_bX^i)\delta X^i=0$$
(4.28)
Combined with the Virasoro-type constraint
$$_aX^{}=\dot{X}^i_aX^i$$
(4.29)
we find that, just as in the string theory case, membranes must have Neumann boundary conditions in the light-front directions.
These considerations would seem to lead to the conclusion that transverse 5-branes simply cannot be constructed in matrix theory. On the other hand, it was argued in that there may be a way to construct a transverse 5-brane using S-duality, at least when the theory has been compactified on a 3-torus. To construct an infinite extended transverse 5-brane in this fashion would require performing an S-duality on $`(3+1)`$-dimensional $`𝒩=4`$ supersymmetric Yang-Mills theory with gauge group $`U(\mathrm{})`$, which is a poorly understood procedure to say the least. In , however, a finite size transverse 5-brane with geometry $`T^3\times S^2`$ was constructed using S-duality of the four-dimensional $`U(N)`$ with finite $`N`$. Furthermore, it was shown that this object couples correctly to the supergravity fields even in the absence of an explicit transverse 5-brane charge. This seems to indicate that transverse 5-branes in matrix theory can be constructed locally, but that they are essentially solitonic objects and do not carry independent conserved quantum numbers. It would be nice to have a more explicit construction of a general class of such finite size transverse 5-branes, particular in the noncompact version of matrix theory.
We now turn to the wrapped, or “longitudinal”, M5-brane which we will refer to as the “L5-brane”. This object appears as a D4-brane in the IIA theory. An infinite D4-brane was considered as a matrix theory background in by including extra fields corresponding to strings stretching between the D0-branes of matrix theory and the background D4-brane. As in the case of the membrane, however, we would like to find a way to explicitly describe a dynamical L5-brane using the matrix degrees of freedom. Just as for the D2-brane, it may be surprising that a D4-brane can be constructed from a configuration of D0-branes. This can be seen from the same type of T-duality argument we used for the D2-brane in 4.2.1. By putting D4-branes and D0-branes on a torus $`T^4`$ we find that the charge-volume relation analogous to (4.2) for a D4-brane is
$$\mathrm{Tr}ϵ_{ijkl}X^iX^jX^kX^l=\frac{V}{2\pi ^2}$$
(4.30)
This is the T-dual of the instanton number in a 4D gauge theory which measures D0-brane charge on D4-branes.
Unlike the case of the membrane, there is no general theory describing an arbitrary L5-brane geometry in matrix theory language. In fact, the only L5-brane configurations which have been explicitly constructed to date are those corresponding to the highly symmetric geometries $`S^4,CP^2`$ and $`R^4`$. We now make a few brief comments about these configurations.
The L5-brane with isotropic $`S^4`$ geometry is similar in many ways to the membrane with $`S^2`$ geometry discussed in section 4.2.2. There are a number of unusual features of the $`S^4`$ system, however, which deserve mention. For full details of the construction see .
A rotationally invariant spherical L5-brane can only be constructed for those values of $`N`$ which are of the form
$$N=\frac{(n+1)(n+2)(n+3)}{3}$$
(4.31)
where $`n`$ is integral. For $`N`$ of this form we define the configuration by
$$X_i=\frac{r}{n}G_i,i\{1,\mathrm{},5\}.$$
(4.32)
where $`G_i`$ are the generators of the $`n`$-fold symmetric tensor product representation of the five four-dimensional Euclidean gamma matrices $`\mathrm{\Gamma }_i`$ satisfying $`\mathrm{\Gamma }_i\mathrm{\Gamma }_j+\mathrm{\Gamma }_j\mathrm{\Gamma }_i=2\delta _{ij}`$
$$G_i^{(n)}=\left(\mathrm{\Gamma }_i\text{1}\text{ }\text{1}\mathrm{}\text{1}\text{ }\text{1}+\text{1}\text{ }\text{1}\mathrm{\Gamma }_i\text{1}\text{ }\text{1}\mathrm{}\text{1}\text{ }\text{1}+\mathrm{}+\text{1}\text{ }\text{1}\mathrm{}\text{1}\text{ }\text{1}\mathrm{\Gamma }_i\right)_\mathrm{S}$$
where the subscript $`S`$ indicates that only the completely symmetric representation is used. For any $`n`$ this configuration has the geometrical properties expected of $`n`$ superimposed L5-branes contained in the locus of points describing a 4-sphere. As for the spherical membrane discussed in 4.2.2 the configuration is confined to the appropriate spherical locus
$$X_1^2+X_2^2+X_3^2+X_4^2+X_5^2r^2\text{1}\text{ }\text{1}.$$
(4.33)
The configuration is symmetric under $`SO(5)`$ and has the correct spectrum and the local D4-brane charge of $`n`$ spherical branes. The energy and equations of motion of this system agree with those expected from M-theory.
Although the system can only be defined in a completely symmetric fashion for certain values of $`n,N`$, this does not seem like a fundamental issue. This constraint is a consequence of the imposition of exact rotational symmetry on the system. It may be that for large and arbitrary $`N`$ it is possible to construct a very good approximation to a spherical L5-brane which breaks rotational invariance to a very small degree. A more fundamental problem, however, is that there is no obvious way of including small fluctuations of the membrane geometry around the perfectly isotropic sphere in a systematic way. In the case of the membrane, we know that for any particular geometry the fluctuations around that geometry can be encoded into matrices which form an arbitrarily good approximation to a smooth fluctuation through the procedure of replacing functions described in terms of an orthonormal basis by appropriate matrix analogues. In the case of the L5-brane we have no such procedure. In fact, there seems to be an obstacle to including all degrees of freedom corresponding to local fluctuations of the brane. It is natural to speculate by analogy with the membrane case that arbitrary fluctuations should be encoded in symmetric polynomials in the matrices $`G_i`$. It can be shown, however, that this is not possible. This geometry has been discussed in a related context in the noncommutative geometry literature as a noncommutative version of $`S^4`$. There also, it was found that not all functions on the sphere could be consistently quantized.
As for the infinite membrane, the infinite L5-brane with geometry of a flat $`R^4R^9`$ can be viewed as a local limit of a large spherical geometry or it can be constructed directly. We need to find a set of operators $`X^{14}`$ on some Hilbert space satisfying
$$ϵ_{ijkl}X^iX^jX^kX^l=\frac{ϵ^4}{2\pi ^2}\text{1}\text{ }\text{1}.$$
(4.34)
Such a configuration can be constructed using matrices which are tensor products of the form $`\text{1}\text{ }\text{1}Q,P`$ and $`Q,P\text{1}\text{ }\text{1}`$. This gives a “stack of D2-branes” solution with D2-brane charge as well as D4-brane charge . It is also possible to construct a configuration with no D2-brane charge by identifying $`X^a`$ with the components of the covariant derivative operator for an instanton on $`S^4`$
$$X^i=i^i+A_i.$$
(4.35)
This construction is known as the Banks-Casher instanton . Just as for the spherical L5-brane, it is not known how to construct small fluctuations of the membrane geometry around any of these flat solutions.
The only other known configuration of an L5-brane in matrix theory corresponds to a brane with geometry $`CP^2`$. This configuration was constructed by Nair and Randjbar-Daemi as a particular example of a coset space $`G/H`$ with $`G=SU(3)`$ and $`H=U(2)`$ . They choose the matrices
$$X_i=\frac{rt_i}{\sqrt{N}}$$
(4.36)
where $`t_i`$ are generators spanning $`𝐠/𝐡`$ in a particular representation of $`SU(3)`$. The geometry defined in this fashion seems to be in some ways better behaved than the $`S^4`$ geometry. For one thing, configurations of a single brane with arbitrarily large $`N`$ can be constructed. Furthermore, it seems to be possible to include all local fluctuations as symmetric functions of the matrices $`t_i`$. This configuration is also somewhat confusing, however, as it extends in only four spatial dimensions, which makes the geometrical interpretation somewhat unclear.
Clearly there are many aspects of the L5-brane in matrix theory which are not understood. The principal outstanding problem is to find a systematic way of describing an arbitrary L5-brane geometry including its fluctuations. One approach to this might be to find a way of regularizing the world-volume theory of an M5-brane in a fashion similar to the matrix regularization of the supermembrane. It is also possible that understanding the structure of noncommutative 4-manifolds might help clarify this question. This is one of many places where noncommutative geometry seems to tie in closely with matrix theory. We will discuss other such connections with noncommutative geometry later in these lectures.
### 4.4 Extended objects from matrices
We have seen that not only pointlike graviton states, but also objects extended in one, two, and four transverse directions can be constructed from matrix degrees of freedom. In this subsection we make some general comments about the appearance of these extended objects and their structure.
One systematic way of understanding the conserved charges associated with the longitudinal and transverse membrane and the longitudinal 5-brane in matrix theory arises from considering the supersymmetry algebra of the theory. The 11-dimensional supersymmetry algebra takes the form
$$\{Q_\alpha ,Q_\beta \}P^I(\gamma _I)_{\alpha \beta }+Z^{I_1I_2}(\gamma _{I_1I_2})_{\alpha \beta }+Z^{I_1\mathrm{}I_5}(\gamma _{I_1\mathrm{}I_5})_{\alpha \beta }$$
(4.37)
where the central terms $`Z`$ correspond to 2-brane and 5-brane charges. The supersymmetry algebra of Matrix theory was explicitly computed by Banks, Seiberg and Shenker . Similar calculations had been performed previously ; however, in these earlier analyses terms such as $`\mathrm{Tr}[X^i,X^j]`$ and $`\mathrm{Tr}X^{[i}X^jX^kX^{l]}`$ were dropped since they vanish for finite $`N`$. The full supersymmetry algebra of the theory takes the schematic form
$$\{Q,Q\}P^I+z^i+z^{ij}+z^{ijkl},$$
(4.38)
as we would expect for the light-front supersymmetry algebra corresponding to (4.37). The charge
$$z^ii\mathrm{Tr}\left(\{P^j,[X^i,X^j]\}+[[X^i,\theta ^\alpha ],\theta ^\alpha ]\right)$$
(4.39)
corresponds to longitudinal membranes (strings), the charge
$$z^{ij}i\mathrm{Tr}[X^i,X^j]$$
(4.40)
corresponds to transverse membranes and
$$z^{ijkl}\mathrm{Tr}X^{[i}X^jX^kX^{l]}$$
(4.41)
corresponds to longitudinal 5-brane charge. For all the extended objects we have described in the preceding subsections, these results agree with the charges we motivated by T-duality arguments.
Note that the charges of all the extended objects in the theory vanish when the matrix size $`N`$ is finite. Physically, this corresponds to the fact that any finite-size configuration of strings, 2-branes and 4-branes must have net charges which vanish.
Another approach to understanding the charges associated with the extended objects of matrix theory arises from the study of the coupling of these objects to supergravity fields, which we will discuss in the next section. From this point of view, perturbative matrix theory calculations can be used to determine not only the conserved charges of the theory, but also the higher multipole moments of all the components of the supercurrent describing the matrix configuration. For example , the multipole moments of the membrane charge $`z^{ij}=2\pi i\mathrm{Tr}[X^i,X^j]`$ can be written in terms of the matrix moments
$$z^{ij(k_1\mathrm{}k_n)}=2\pi i\mathrm{STr}\left([X^i,X^j]X^{k_1}\mathrm{}X^{k_n}\right)$$
(4.42)
which are the matrix analogues of the moments
$$d^2\sigma \{X^i,X^j\}X^{k_1}\mathrm{}X^{k_n}$$
(4.43)
for the continuous membrane. The symbol $`\mathrm{STr}`$ indicates a symmetrized trace, wherein the trace is averaged over all possible orderings of the terms $`[X^i,X^j]`$ and $`X^{k_\nu }`$ appearing inside the trace. This corresponds to a particular ordering prescription in applying the matrix-membrane correspondence to (4.43). There is no a priori justification for this ordering prescription, but it is a consequence of explicit calculations of interactions between general matrix theory objects as described in the next section. The same prescription can be used to define the multipole moments of the longitudinal membrane and 5-brane charges.
Although as we have mentioned, the conserved charges in matrix theory corresponding to extended objects all vanish at finite $`N`$, the same is not true of the higher moments of these charges. For example, the isotropic spherical matrix membrane configuration discussed in section 4.2.2 has nonvanishing membrane dipole moments
$`z^{12(3)}=z^{23(1)}=z^{31(2)}`$ $`=`$ $`2\pi i\mathrm{Tr}\left([X^1,X^2]X^3\right)`$ (4.44)
$`=`$ $`{\displaystyle \frac{4\pi r^3}{3}}(11/N^2)`$
which agrees with the membrane dipole moment $`4\pi r^3/3`$ of the smooth spherical membrane up to terms of order $`1/N^2`$. Using the multipole moments of a fixed matrix configuration we can essentially reproduce the complete spatial dependence of the matter configuration to which the matrices correspond. This higher moment structure describing higher-dimensional extended objects through lower-dimensional objects is very general, and has a precise analog in describing the supercurrents and charges of Dirichlet $`(p+2k)`$-branes in terms of the world-volume theory of a system of D$`p`$-branes . This structure has many possible applications to D-brane physics as well as to matrix theory. For example, it was recently pointed out by Myers that putting a system of D$`p`$-branes in a constant background $`(p+4)`$-form flux will produce a dielectric effect in which spherical bubbles of D$`(p+2)`$-branes will be formed with dipole moments which screen the background field.
## 5 Interactions in matrix theory
In this section we discuss interactions in matrix theory between block matrices describing general time-dependent matrix theory configurations which may include gravitons, membranes and 5-branes. We begin by reviewing the perturbative Yang-Mills formalism in background field gauge. This formalism can be used to carry out loop calculations in matrix theory, giving results which can be related to supergravity interactions. We carry out two explicit examples of this calculation at one-loop order: first for a pair of 0-branes with relative velocity $`v`$, following , then for the leading order term in the interaction between an arbitrary pair of bosonic background configurations, following . Following these examples, we summarize the extent to which perturbative Yang-Mills calculations of this kind have been shown to agree with classical supergravity. At the level of linearized supergravity, it has been found that there is an infinite series of terms in the one-loop matrix theory effective potential which precisely reproduce all tree-level supergravity interactions arising from the exchange of a single graviton, 3-form quantum or gravitino. There is limited information about the extent to which nonlinear supergravity effects are reproduced by higher-loop matrix theory calculations, however. While it has been shown that the nonlinear structure of 3-graviton scattering is correctly reproduced by a two-loop matrix theory calculation, there is not a clear picture of what should be expected beyond this. We discuss these results and how they are related to supersymmetric nonrenormalization theorems which protect some terms in the perturbative Yang-Mills expansion from higher-loop corrections.
In this section we primarily focus on the problem of deriving classical 11-dimensional supergravity from matrix theory. A very interesting, but more difficult, question is whether matrix theory can also successfully reproduce string/M-theory corrections to classical supergravity. The first such corrections would be $`^4`$ corrections to the Einstein-Hilbert action. Some work has been done investigating the question of whether these terms can be seen in matrix theory . While more work needs to be done in this direction, the results of indicate that the perturbative loop expansion in matrix theory probably does not correctly reproduce quantum effects in M-theory. The most likely explanation for this discrepancy is that such terms are not subject to nonrenormalization theorems, and are only reproduced in the large $`N`$ limit. We discuss these issues again briefly in the last section.
In subsection 5.1 we describe two-body interactions in matrix theory, and in subsection 5.2 we discuss interactions between more than two objects. Section 5.3 contains a brief discussion of interactions involving longitudinal momentum transfer, which correspond to nonperturbative processes in matrix theory.
### 5.1 Two-body interactions
The background field formalism for describing matrix theory interactions between block matrices which are widely separated in eigenvalue space was first used by Douglas, Kabat, Pouliot and Shenker in to describe interactions between a pair of D0-branes in type IIA string theory moving with relative velocity $`v`$. In this subsection we discuss their result and the generalization to general bosonic background configurations. The matrix theory Lagrangian is
$$=\frac{1}{2R}\mathrm{Tr}\left[D_0X^iD_0X^i+\frac{1}{2}[X^i,X^j]^2+\theta ^T(i\dot{\theta }\gamma _i[X^i,\theta ])\right]$$
(5.1)
where
$$D_0X^i=_tX^ii[A,X^i].$$
(5.2)
We wish to expand each of the matrix theory fields around a classical background. We will assume here for simplicity that the background has a vanishing gauge field and vanishing fermionic fields. For a discussion of the general situation with background fermionic fields as well as bosonic fields see . We expand the bosonic field in terms of a background plus a fluctuation
$`X^i`$ $`=`$ $`B^i+Y^i.`$
We choose the background field gauge
$$D_\mu ^{\mathrm{bg}}A^\mu =_tAi[B^i,X^i]=0.$$
(5.3)
This gauge can be implemented by adding a term $`(D_\mu ^{\mathrm{bg}}A^\mu )^2`$ to the action and including the appropriate ghosts. The nice feature of this gauge is that the terms quadratic in the bosonic fluctuations simplify to the form
$$\dot{Y}^i\dot{Y}^i[B^i,Y^j]^2[B^i,B^j][Y^i,Y^j]$$
(5.4)
The complete gauge-fixed action including ghosts is written in Euclidean time $`\tau =it`$ as
$$S=S_0+S_2+S_3+S_4$$
(5.5)
where
$`S_0`$ $`=`$ $`{\displaystyle \frac{1}{2R}}{\displaystyle 𝑑\tau \mathrm{Tr}\left[_\tau B^i_\tau B^i+\frac{1}{2}[B^i,B^j]^2\right]}`$
$`S_2`$ $`=`$ $`{\displaystyle \frac{1}{2R}}{\displaystyle }d\tau \mathrm{Tr}[_\tau Y^i_\tau Y^i[B^i,Y^j][B^i,Y^j][B^i,B^j][Y^i,Y^j]`$
$`+_\tau A_\tau A[B^i,A][B^i,A]2i\dot{B}^i[A,Y^i]`$
$`+_\tau \overline{C}_\tau C[B^i,\overline{C}][B^i,C]+\theta ^T\dot{\theta }\theta ^T\gamma _i[B^i,\theta ]]`$
and where $`S_3`$ and $`S_4`$ contain terms cubic and quartic in the fluctuations $`Y^i,A,C,\theta `$. Note that we have taken $`AiA`$ as appropriate for the Euclidean formulation.
We now wish to use this gauge-fixed action to compute the effective potential governing the interaction between a pair of matrix theory objects. In general, to calculate the interaction potential to arbitrary order it is necessary to include the terms $`S_3`$ and $`S_4`$ in the action. The propagators for each of the fields can be computed from the quadratic term $`S_2`$. A systematic diagrammatic expansion will then yield the effective potential to arbitrary high order. We begin our discussion of matrix theory interactions, however, with the simplest case: the interaction of two objects at leading order in the inverse separation distance. In 5.1.1 we discuss the simplest case of this situation, the scattering of a pair of gravitons. In 5.1.2 we discuss the situation of two general matrix theory objects, giving an explicit calculation for the leading term in the case where both objects are purely bosonic. After working out these explicit examples we review what is known about the scattering of a general pair of matrix theory objects to arbitrary order in section 5.1.3. We review the special case of a pair of gravitons in section 5.1.4. We discuss the N-body problem in 5.2.
#### 5.1.1 Two graviton interactions at leading order
As we have discussed in 4.1.1, a classical background describing a pair of gravitons with relative velocity $`v`$ and impact parameter $`b`$ (and no polarization information) is given in the center of mass frame by
$`B^1`$ $`=`$ $`{\displaystyle \frac{i}{2}}\left(\begin{array}{cc}v\tau & 0\\ 0& v\tau \end{array}\right)`$ (5.9)
$`B^2`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(\begin{array}{cc}b& 0\\ 0& b\end{array}\right)`$ (5.12)
$`B^i`$ $`=`$ $`0,i>2`$ (5.13)
Inserting these backgrounds into (5.1) we see that at a fixed value of time the Lagrangian at quadratic order for the 10 complex bosonic off-diagonal components of $`A`$ and $`Y^i`$ is that of a system of 10 harmonic oscillators with mass matrix
$$(\mathrm{\Omega }_b)^2=\left(\begin{array}{ccccc}r^2& 2iv& 0& \mathrm{}& 0\\ 2iv& r^2& 0& \mathrm{}& 0\\ 0& 0& r^2& \mathrm{}& \mathrm{}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& 0\\ 0& 0& \mathrm{}& 0& r^2\end{array}\right)$$
(5.14)
where $`r^2=b^2+(vt)^2`$ is the instantaneous separation between the gravitons.
There are two complex off-diagonal ghosts with $`\mathrm{\Omega }^2=r^2`$.
There are 16 fermionic oscillators with a mass-squared matrix
$$(\mathrm{\Omega }_f)^2=r^2\text{1}\text{ }\text{1}_{16\times 16}+v\gamma _1$$
(5.15)
This matrix can be found by writing
$$P^{}P=^2+(\mathrm{\Omega }_f)^2$$
(5.16)
where
$$P=v\tau \gamma _1b\gamma _2$$
(5.17)
To perform a completely general calculation of the two-body effective interaction potential to all orders in $`1/r`$ it is necessary to perform a diagrammatic expansion using the exact propagator for the bosonic and fermionic fields. For example, the bosonic propagator satisfying
$$(^2+b^2+v^2\tau ^2)\mathrm{\Delta }_B(\tau ,\tau ^{}|b^2+v^2\tau ^2)=\delta (\tau \tau ^{})$$
(5.18)
is given by the expression
$`\mathrm{\Delta }_B(\tau ,\tau ^{}|b^2+v^2\tau ^2)`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}dse^{b^2s}\sqrt{{\displaystyle \frac{v}{2\pi \mathrm{sinh}2sv}}}\times `$
$`\mathrm{exp}\left({\displaystyle \frac{v}{2\mathrm{sinh}2sv}}((\tau ^2+\tau ^2)\mathrm{cosh}2sv2\tau \tau ^{})\right).`$
In general, even for a simple 2-graviton calculation there is a fair amount of algebra involved in extracting the effective potential using propagators of the form (5.1.1). If, however, we are only interested in calculating the leading term in the long-range interaction potential we can simplify the calculation by making the quasi-static assumption that all the oscillator frequencies $`\omega `$ of interest are large compared to the ratio $`v/r`$ of the relative velocity divided by the separation scale. In this regime, we can make the approximation that all the oscillators stay in their ground state over the interaction time-scale, so that the effective potential between the two objects is simply given by the sum of the ground-state energies of the boson, ghost and fermion oscillators
$$V_{\mathrm{qs}}=\underset{b}{}\omega _b\underset{g}{}\omega _g\frac{1}{2}\underset{f}{}\omega _f.$$
(5.20)
Note that the bosonic and ghost oscillators are complex so that no factor of 1/2 is included.
In the situation of two-graviton scattering we can therefore calculate the effective potential by diagonalizing the frequency matrices $`\mathrm{\Omega }_b,\mathrm{\Omega }_g`$ and $`\mathrm{\Omega }_f`$. We find that the bosonic oscillators have frequencies
$`\omega _b`$ $`=`$ $`r\mathrm{with}\mathrm{multiplicity}8`$
$`\omega _b`$ $`=`$ $`\sqrt{r^2\pm 2v}\mathrm{with}\mathrm{multiplicity}1\mathrm{each}.`$
The 2 ghosts have frequencies
$$\omega _g=r,$$
(5.21)
and the 16 fermions have frequencies
$$\omega _f=\sqrt{r^2\pm v}\mathrm{with}\mathrm{multiplicity}8\mathrm{each}.$$
(5.22)
The effective potential for a two-graviton system with instantaneous relative velocity $`v`$ and separation $`r`$ is thus given by the leading term in a $`1/r`$ expansion of the expression
$$V=\sqrt{r^2+2v}+\sqrt{r^22v}+6r4\sqrt{r^2+v}+4\sqrt{r^2v}.$$
(5.23)
Expanding in $`v/r^2`$ we see that the terms of order $`r,v/r,v^2/r^3`$ and $`v^3/r^5`$ all cancel. The leading term is
$$V=\frac{15}{16}\frac{v^4}{r^7}+𝒪(\frac{v^6}{r^{11}})$$
(5.24)
This potential was first computed by Douglas, Kabat, Pouliot and Shenker . This result agrees with the leading term in the effective potential between two gravitons with $`P^+=1/R`$ in light-front 11D supergravity. We will discuss the supergravity side of this calculation in more detail in the following section, where we generalize this calculation to an arbitrary pair of matrix theory objects.
#### 5.1.2 General 2-body systems and linearized supergravity at leading order
We now generalize the discussion to an arbitrary pair of matrix theory objects, which are described by a block-diagonal background
$$B^i=\left(\begin{array}{cc}\widehat{X}^i& 0\\ 0& \stackrel{~}{X}^i\end{array}\right)$$
(5.25)
where $`\widehat{X}^i`$ and $`\stackrel{~}{X}^i`$ are $`\widehat{N}\times \widehat{N}`$ and $`\stackrel{~}{N}\times \stackrel{~}{N}`$ matrices describing the two objects. The separation distance between the objects, which we will use as an expansion parameter, is given by
$$r^i=\frac{1}{\widehat{N}}\mathrm{Tr}\widehat{X}^i\frac{1}{\stackrel{~}{N}}\mathrm{Tr}\stackrel{~}{X}^i$$
(5.26)
There are $`\widehat{N}\stackrel{~}{N}`$ independent complex off-diagonal components of the fluctuation matrices $`Y^i`$. We will find it useful to treat these components as an $`\widehat{N}\stackrel{~}{N}`$-component vector $`Z^i`$. We now construct a $`\widehat{N}\stackrel{~}{N}\times \widehat{N}\stackrel{~}{N}`$ matrix which acts on the $`Z^i`$ vectors
$$K_i\widehat{X}_i\text{1}\text{ }\text{1}_{\stackrel{~}{N}\times \stackrel{~}{N}}\text{1}\text{ }\text{1}_{N\times N}\stackrel{~}{X}_i^T.$$
(5.27)
It is convenient to extract the centers of mass explicitly so that $`K^i`$ can be rewritten as
$$K^i=r^i\text{1}\text{ }\text{1}+\overline{K}^i$$
(5.28)
where $`\overline{K}^i`$ is of order $`1`$ in terms of the separation scale $`r`$. The matrices $`K`$ encode the adjoint action of the background $`B`$ on the fluctuations $`Y`$ so that we can extract the part of $`[B,Y]`$ depending on the off-diagonal fields $`Z`$ through
$$[B^i,Y^j]K^iZ^j.$$
(5.29)
This formalism allows us to write the quadratic terms from (5.1) in the action for the off-diagonal fields in a simple form
$`\dot{Y}^i\dot{Y}^i[B^i,Y^j][B^i,Y^j][B^i,B^j][Y^i,Y^j]`$ (5.30)
$`\dot{Z}_i^{}\dot{Z}^iZ_j^{}K^iK_iZ^j2Z_i^{}[K_i,K_j]Z^j`$
Performing a similar operation for the terms quadratic in fluctuations of the $`A`$ field, we find that the full frequency-squared matrices for the bosonic, ghost and fermionic fields can be written
$`\mathrm{\Omega }_b^2`$ $`=`$ $`K^2\text{1}\text{ }\text{1}_{10\times 10}2iF_{\mu \nu }`$
$`\mathrm{\Omega }_g^2`$ $`=`$ $`K^2\text{1}\text{ }\text{1}_{2\times 2}`$ (5.31)
$`\mathrm{\Omega }_f^2`$ $`=`$ $`K^2\text{1}\text{ }\text{1}_{16\times 16}iF_{\mu \nu }\gamma ^\mu \gamma ^\nu `$
where $`\gamma ^0=\text{1}\text{ }\text{1}`$ and the field strength matrix $`F_{\mu \nu }`$ is given by
$`F_{0i}`$ $`=`$ $`\dot{K}^i`$ (5.32)
$`F_{ij}`$ $`=`$ $`i[K^i,K^j]`$
Note that each of the frequencies has a leading term $`r`$ and subleading terms of order 1. Expanding the frequency matrices in powers of $`1/r`$ we find that for a completely arbitrary pair of objects described by the background matrices $`\widehat{X}^i`$ and $`\stackrel{~}{X}^i`$ the potential vanishes to order $`1/r^6`$. At order $`1/r^7`$ we find that the potential is
$`V_{\mathrm{leading}}`$ $`=`$ $`\mathrm{Tr}\left(\mathrm{\Omega }_b\right){\displaystyle \frac{1}{2}}\mathrm{Tr}\left(\mathrm{\Omega }_f\right)2\mathrm{T}\mathrm{r}\left(\mathrm{\Omega }_g\right)`$ (5.33)
$`=`$ $`{\displaystyle \frac{5}{128r^7}}\mathrm{STr}`$ (5.34)
where
$$=24F^\mu {}_{\nu }{}^{}F_{}^{\nu }{}_{\lambda }{}^{}F_{}^{\lambda }{}_{\sigma }{}^{}F_{}^{\sigma }{}_{\mu }{}^{}6F_{\mu \nu }F^{\mu \nu }F_{\lambda \sigma }F^{\lambda \sigma }$$
(5.35)
and $`\mathrm{STr}`$ indicates that the trace is symmetrized over all possible orderings of $`F`$’s in the product $`F^4`$.
From the definition (5.27) it is clear that the field strength $`F_{\mu \nu }`$ decomposes into a piece from each of the two objects
$$F_{\mu \nu }=\widehat{F}_{\mu \nu }\stackrel{~}{F}_{\mu \nu }$$
(5.36)
where $`\widehat{F}_{\mu \nu }`$ and $`\stackrel{~}{F}_{\mu \nu }`$ are defined through (5.32) in terms of $`\widehat{X}`$ and $`\stackrel{~}{X}`$. We can therefore decompose the potential $`V_{\mathrm{leading}}`$ into a sum of terms which are written as products of a function of $`\widehat{X}`$ and a function of $`\stackrel{~}{X}`$, where the terms can be grouped according to the number of Lorentz indices contracted between the two objects. With some algebra, we can write this potential as
$`V_{\mathrm{leading}}`$ $`=`$ $`V_{\mathrm{gravity}}+V_{\mathrm{electric}}+V_{\mathrm{magnetic}}`$ (5.37)
$`V_{\mathrm{gravity}}`$ $`=`$ $`{\displaystyle \frac{15R^2}{4r^7}}(\widehat{𝒯}^{IJ}\stackrel{~}{𝒯}_{IJ}{\displaystyle \frac{1}{9}}\widehat{𝒯}^I{}_{I}{}^{}\stackrel{~}{𝒯}_{}^{J}{}_{J}{}^{})`$ (5.38)
$`V_{\mathrm{electric}}`$ $`=`$ $`{\displaystyle \frac{45R^2}{r^7}}\widehat{𝒥}^{IJK}\stackrel{~}{𝒥}_{IJK}`$ (5.39)
$`V_{\mathrm{magnetic}}`$ $`=`$ $`{\displaystyle \frac{45R^2}{r^7}}\widehat{}^{+ijkl}\stackrel{~}{}^{+ijkl}`$ (5.40)
This is, as we shall discuss shortly, precisely the form of the interactions we expect to see from 11D supergravity in light-front coordinates, where $`𝒯,𝒥`$ and $``$ play the role of the (integrated) stress tensor, membrane current and 5-brane current of the two objects. The quantities appearing in this decomposition are defined as follows.
$`𝒯^{IJ}`$ is a symmetric tensor with components
$`𝒯^{}`$ $`=`$ $`{\displaystyle \frac{1}{R}}\mathrm{STr}{\displaystyle \frac{}{96}}`$ (5.41)
$`𝒯^i`$ $`=`$ $`{\displaystyle \frac{1}{R}}\mathrm{STr}\left({\displaystyle \frac{1}{2}}\dot{X}^i\dot{X}^j\dot{X}^j+{\displaystyle \frac{1}{4}}\dot{X}^iF^{jk}F^{jk}+F^{ij}F^{jk}\dot{X}^k\right)`$
$`𝒯^+`$ $`=`$ $`{\displaystyle \frac{1}{R}}\mathrm{STr}\left({\displaystyle \frac{1}{2}}\dot{X}^i\dot{X}^i+{\displaystyle \frac{1}{4}}F^{ij}F^{ij}\right)`$
$`𝒯^{ij}`$ $`=`$ $`{\displaystyle \frac{1}{R}}\mathrm{STr}\left(\dot{X}^i\dot{X}^j+F^{ik}F^{kj}\right)`$
$`𝒯^{+i}`$ $`=`$ $`{\displaystyle \frac{1}{R}}\mathrm{STr}\dot{X}^i`$
$`𝒯^{++}`$ $`=`$ $`{\displaystyle \frac{N}{R}}`$
$`𝒥^{IJK}`$ is a totally antisymmetric tensor with components
$`𝒥^{ij}`$ $`=`$ $`{\displaystyle \frac{1}{6R}}\mathrm{STr}(\dot{X}^i\dot{X}^kF^{kj}\dot{X}^j\dot{X}^kF^{ki}{\displaystyle \frac{1}{2}}\dot{X}^k\dot{X}^kF^{ij}`$
$`+{\displaystyle \frac{1}{4}}F^{ij}F^{kl}F^{kl}+F^{ik}F^{kl}F^{lj})`$
$`𝒥^{+i}`$ $`=`$ $`{\displaystyle \frac{1}{6R}}\mathrm{STr}\left(F^{ij}\dot{X}^j\right)`$
$`𝒥^{ijk}`$ $`=`$ $`{\displaystyle \frac{1}{6R}}\mathrm{STr}\left(\dot{X}^iF^{jk}+\dot{X}^jF^{ki}+\dot{X}^kF^{ij}\right)`$
$`𝒥^{+ij}`$ $`=`$ $`{\displaystyle \frac{1}{6R}}\mathrm{STr}F^{ij}`$
Note that we retain some quantities — in particular $`𝒥^{+i}`$ and $`𝒥^{+ij}`$ — which vanish at finite $`N`$ (by the Gauss constraint and antisymmetry of $`F^{ij}`$, respectively). These terms represent BPS charges (for longitudinal and transverse membranes) which are only present in the large $`N`$ limit. We define higher moments of these terms below which can be non-vanishing at finite $`N`$.
$`^{IJKLMN}`$ is a totally antisymmetric tensor with
$$^{+ijkl}=\frac{1}{12R}\mathrm{STr}\left(F^{ij}F^{kl}+F^{ik}F^{lj}+F^{il}F^{jk}\right).$$
(5.43)
At finite $`N`$ this vanishes by the Jacobi identity, but we shall retain it because it represents the charge of a longitudinal 5-brane. The other components of $`^{IJKLMN}`$ do not appear in the Matrix potential. In principle, we expect another component of the 5-brane current, $`^{ijklm}`$, to be well-defined. This term arises from a moving longitudinal 5-brane. This term does not appear in the 2-body interaction formula because it would couple to the transverse 5-brane charge $`^{+ijklm}`$ which, as we have discussed, vanishes in light-front coordinates. The component $`^{ijklm}`$ can, however, be determined from the conservation of the 5-brane current, and is given by
$$^{ijklm}=\frac{5}{4R}\mathrm{STr}\left(\dot{X}^{[i}F^{jk}F^{lm]}\right).$$
(5.44)
Let us compare the interaction potential (5.37) with the leading long-range interaction between two objects in 11D light-front compactified supergravity. The scalar propagator in 11D is
$$\text{ }\text{ }\text{ }\text{ }\text{ }{}_{}{}^{1}(x)=\frac{1}{2\pi R}\underset{n}{}\frac{dk^{}d^{\mathrm{\hspace{0.17em}9}}k_{}}{(2\pi )^{10}}\frac{e^{i\frac{n}{R}x^{}ik^{}x^++ik_{}x_{}}}{2\frac{n}{R}k^{}k_{}^2}$$
(5.45)
where $`n`$ counts the number of units of longitudinal momentum $`k^+`$. To compare the leading term in the long-distance potential with matrix theory we only need to extract the term associated with $`n=0`$, corresponding to interactions mediated by exchange of a supergraviton with no longitudinal momentum.
$$\text{ }\text{ }\text{ }\text{ }\text{ }{}_{}{}^{1}(xy)=\frac{1}{2\pi R}\delta (x^+y^+)\frac{15}{32\pi ^4|x_{}y_{}|^7}$$
(5.46)
Note that the exchange of quanta with zero longitudinal momentum gives rise to interactions that are instantaneous in light-front time, as recently emphasized in . This is precisely the type of instantaneous interaction that we find at one loop in Matrix theory. Such action-at-a-distance potentials are allowed by the Galilean invariance manifest in the light-front formalism.
The graviton propagator can be written in terms of this scalar propagator as
$$D_{\mathrm{graviton}}^{IJ,KL}=2\kappa ^2(\eta ^{IK}\eta ^{JL}+\eta ^{IL}\eta ^{JK}\frac{2}{9}\eta ^{IJ}\eta ^{KL})\text{ }\text{ }\text{ }\text{ }\text{ }{}_{}{}^{1}(xy)$$
(5.47)
where $`2\kappa ^2=(2\pi )^5R^3`$ in string units. The effective supergravity interaction between two objects having stress tensors $`\widehat{T}_{IJ}`$ and $`\stackrel{~}{T}_{KL}`$ can then be expressed as
$$S=\frac{1}{4}d^{11}xd^{11}y\widehat{T}_{IJ}(x)D_{\mathrm{graviton}}^{IJ,KL}(xy)\stackrel{~}{T}_{KL}(y)$$
(5.48)
This interaction has a leading term of precisely the form (5.38) if we define $`𝒯^{IJ}`$ to be the integrated component of the stress tensor
$$𝒯^{IJ}𝑑x^{}d^{\mathrm{\hspace{0.17em}9}}x_{}T^{IJ}(x).$$
(5.49)
It is straightforward to show in a similar fashion that (5.39) and (5.40) are precisely the forms of the leading supergravity interaction between membrane currents and 5-brane currents of a pair of objects.
We can calculate the components of the source currents (5.41), (5.1.2) and (5.43) for all the matrix theory objects we have discussed: the supergraviton, the membrane and the L5-brane. For all these objects the currents have the expected values, at least to order $`1/N^2`$. For example, the stress tensor of a graviton can be written in the form
$$𝒯^{IJ}=\frac{p^Ip^J}{p^+}$$
(5.50)
where
$$p^+=N/R,p^i=p^+\dot{x}^i,p^{}=p_{}^2/2p^+$$
(5.51)
The stress tensor and membrane current of the membrane can be computed in the continuum membrane theory from the action (2.5) for the bosonic membrane in a general background. Using the matrix-membrane correspondence (2.51) it is possible to show that the matrix definitions above are compatible with the expressions for the stress tensor and membrane current of the continuum membrane, although the matrix expressions are not uniquely determined by this correspondence.
We have thus shown that to leading order in the separation distance the interaction between any pair of objects in supergravity is precisely reproduced by one-loop quantum effects in matrix theory. We have only shown this explicitly in the case of a pair of bosonic backgrounds, following . The more general case where fermionic background fields are included is discussed in . In the following sections we discuss what is known about the extension of these results to higher order in $`1/r`$ and to interactions of more than two distinct objects.
#### 5.1.3 General 2-body interactions
In the previous subsections we have considered only the leading $`1/r^7`$ terms in the 2-body interaction potential. If we consider all possible Feynman diagrams which might contribute to higher-order terms, it is straightforward to demonstrate by power counting that the complete potential can be written as a sum of terms of the form
$$V=\underset{n,k,l,m,p}{}V_{n,k,l,m,p,\alpha }R^{n1}\frac{X^lD^pF^k\psi ^{2m}}{r^{3n+2k+l+3m+p4}}.$$
(5.52)
where $`n`$ counts the number of loops in the relevant diagrams and $`\psi `$ describes the fermionic background fields. Each $`D`$ either indicates a time derivative or a commutator with an $`X`$, as in $`\psi [X,\psi ]`$. The summation over the index $`\alpha `$ indicates a sum over many possible index contractions for every combination of $`F`$’s, $`X`$’s and $`D`$’s and $`\mathrm{\Gamma }`$ matrices between the $`\psi `$’s.
For a completely general pair of objects, only terms in the one-loop effective action have been understood in terms of supergravity. At one-loop order, when the fields are taken on-shell by imposing the matrix theory equations of motion, all terms with $`k+m+p<4`$ which have been calculated vanish. All terms with $`k+m+p=4`$ which have been calculated have $`mp`$ and can be written in the form
$$V_{1,4mp,l,m,p,\alpha }\frac{X^lF^{(4mp)}\psi ^{2(mp)}(\psi D\psi )^p}{r^{7+mp+l}}.$$
(5.53)
In this expression, the grouping of $`\psi `$ terms indicates the contraction of spinor indices—in general, the terms can be ordered in an arbitrary fashion when considered as $`U(N)`$ matrices. The terms (5.53) have been explicitly determined for $`m<2`$ in , where they were shown to precisely correspond to multipole interaction terms in linearized supergravity. We now briefly describe some of those terms which have been interpreted in this fashion
$`m=p=0,k=4,l=0:`$ These are the leading $`1/r^7`$ terms in the interaction potential between a pair of purely bosonic objects discussed above. They are precisely equivalent to the leading term in the supergravity potential between a pair of objects with appropriate integrated stress tensors, membrane currents and 5-brane currents.
$`m=p=0,k=4,l>0:`$ This infinite set of terms was shown in to be equivalent to the higher-order terms in the linearized supergravity potential arising from higher moments of the bosonic parts of the stress tensor, membrane current and 5-brane current. The simplest example (discussed in ) is the term of the form $`F^4X/r^8`$ which appears in the case of a graviton moving in the long-range gravitational field of a matrix theory object with angular momentum
$$J^{ij}=T^{+i(j)}T^{+j(i)}$$
(5.54)
where the first moment of the matrix theory stress tensor component $`T^{+i}`$ is defined through (as discussed in subsection 4.4)
$$T^{+i(j)}=\frac{1}{R}\mathrm{Tr}\left(\dot{X}^iX^j\right)$$
(5.55)
In it was shown that terms of the general form $`F^4X^l/r^{7+l}`$ can describe higher-moment membrane-5-brane and D0-brane-D6-brane interactions as well as membrane-membrane and 5-brane-5-brane interactions, generalizing previous results in .
$`m=1,p=1,k=2,l0:`$ The terms of the form
$$\frac{F^2(\psi D\psi )X^l}{r^{7+l}}$$
(5.56)
correspond again to leading and higher-moment interactions in linearized supergravity, where now the components of the (integrated) gravity currents have contributions from the fermionic backgrounds as well as the bosonic backgrounds. These terms are also related to linearized supercurrent interactions arising from single gravitino exchange, as discussed in .
$`m=2,p=2,k=0,l0:`$ The terms of the form
$$\frac{(\psi D\psi )(\psi D\psi )X^l}{r^{7+l}}$$
(5.57)
correspond, just like the terms (5.56), to fermionic contributions to the linearized supergravity interaction arising from fermion contributions to the integrated supergravity currents.
$`m=1,p=0,k=3,l0:`$ The terms of the form
$$\frac{F^3\psi \psi X^l}{r^{8+l}}$$
(5.58)
have a similar interpretation to the terms (5.56). In these terms, however, the dipole moments of the currents have nontrivial fermionic contributions in which no derivatives act on the fermions . The simplest example of this is the spin contribution to the matrix theory angular momentum
$$J_{\mathrm{fermion}}^{ij}=\frac{1}{4R}\mathrm{Tr}\left(\psi \gamma ^{ij}\psi \right)$$
(5.59)
This contribution was first noted in the context of spinning gravitons in . This angular momentum term couples to the component $`T^iF^3`$ of the matrix theory stress-energy tensor through terms of the form $`\widehat{J}^{ij}\stackrel{~}{T}^ir^j/r^9`$.
$`m>1,k=4m,l0:`$ The terms of the form
$$\frac{F^2\psi ^4X^l}{r^{9+l}},\frac{F\psi ^6X^l}{r^{10+l}},\frac{\psi ^8X^l}{r^{11+l}},\frac{(\psi D\psi )F\psi ^2X^l}{r^{8+l}},\frac{(\psi D\psi )\psi ^4X^l}{r^{9+l}}$$
(5.60)
have not been completely calculated or related to supergravity interactions, although as we will discuss in the following section these terms are known and agree with supergravity interactions in the special case $`N=2`$. From the structure which has already been understood it seems most likely that these terms arise from fermion contributions to the higher multipole moments of the supergravity currents, and that these terms will also agree with the corresponding supergravity interactions.
This is all that is known about the 2-body interaction for a completely general (and not necessarily supersymmetric) pair of matrix theory objects. To summarize, it has been shown that all terms of the form $`F^k(\psi D\psi )^p\psi ^{2(mp)}`$ with $`k+p+m=4`$ correspond to supergravity interactions, at least for the terms with $`m<2`$. It seems likely that this correspondence persists for the remaining values of $`m>1`$, but the higher order fermionic contributions to the multipole moments of the supergravity currents have not yet been calculated for a general matrix theory object. It is likely that all these terms are protected by a supersymmetric nonrenormalization theorem of the type discussed in the following section. This has not yet been proven, but might follow from arguments similar to those in . The only other known results are for a pair of gravitons, which we now review.
#### 5.1.4 General two-graviton interactions
In the case of a pair of gravitons, the general interaction potential (5.52) simplifies to
$$V=\underset{n=0}{\overset{\mathrm{}}{}}\underset{k=0}{\overset{\mathrm{}}{}}\underset{m=0}{\overset{4}{}}V_{n,k,m}R^{n1}\frac{v^k\psi ^{2m}}{r^{3n+2k+3m4}}$$
(5.61)
The sum over $`m`$ is finite since in the $`SU(2)`$ theory all terms with fermions can be described in terms of a product of 2, 4, 6 or 8 $`\psi `$’s. The leading terms for each value of $`m`$ have been computed using the one-loop approach, and agree with supergravity. The sum of these terms is (see and references therein for further details)
$`V_{(1)}`$ $`=`$ $`{\displaystyle \frac{15}{16}}[v^4+2v^2v_iD^{ij}_j+2v_iv_jD^{ik}D^{jl}_k_l`$
$`+{\displaystyle \frac{4}{9}}v_iD^{ij}D^{km}D^{lm}_j_k_l+{\displaystyle \frac{2}{63}}D^{in}D^{jn}D^{km}D^{lm}_i_j_k_l]{\displaystyle \frac{1}{r^7}}`$
where
$$D^{ij}=\psi \gamma ^{ij}\psi $$
(5.63)
The term with a single $`D`$ proportional to $`1/r^8`$ arises from the spin angular momentum term described in (5.59).
No further checks have been made on the matrix theory/supergravity correspondence for terms with nontrivial fermion backgrounds. Simplifying to the spinless case, the complete effective potential (5.61) simplifies still further to
$$V=\underset{n,k}{}V_{n,k}R^{n1}\frac{v^k}{r^{3n+2k4}}.$$
(5.64)
Following , we write these terms in matrix form
$$\begin{array}{cccccccccc}V& =& \frac{1}{R}V_{0,2}v^2& & & & & & & \\ & & +& V_{1,4}\frac{v^4}{r^7}& +& V_{1,6}\frac{v^6}{r^{11}}& +& V_{1,8}\frac{v^8}{r^{15}}& +& \mathrm{}\\ & & +& RV_{2,4}\frac{v^4}{r^{10}}& +& RV_{2,6}\frac{v^6}{r^{14}}& +& RV_{2,8}\frac{v^8}{r^{18}}& +& \mathrm{}\\ & & +& R^2V_{3,4}\frac{v^4}{r^{13}}& +& R^2V_{3,6}\frac{v^6}{r^{17}}& +& R^2V_{3,8}\frac{v^8}{r^{21}}& +& \mathrm{}\\ & & +& \mathrm{}& +& \mathrm{}& +& \mathrm{}& +& \mathrm{}\end{array}$$
(5.65)
where each row gives the contribution at fixed loop order. We will now give a brief review of what is known about these coefficients. First, let us note that in Planck units this potential is (restoring factors of $`\alpha ^{}=l_{11}^3/R`$ by dimensional analysis)
$$V=\underset{n,k}{}V_{n,k}\frac{l_{11}^{3n+3k6}}{R^{k1}}\frac{v^k}{r^{3n+2k4}}.$$
(5.66)
Since the gravitational coupling constant is $`\kappa ^2=2^7\pi ^8l_{11}^9`$ we only expect terms with
$$n+k2(\mathrm{mod}\mathrm{\hspace{0.33em}3})$$
(5.67)
to correspond with classical supergravity interactions, since all terms in the classical theory have integral powers of $`\kappa `$. Of the terms explicitly shown in (5.65) only the diagonal terms satisfy this criterion. By including factors of $`\widehat{N}`$ and $`\stackrel{~}{N}`$ for semi-classical graviton states with finite momentum $`P^+`$ and comparing to supergravity, one can argue that the terms on the diagonal are precisely those which should correspond to classical supergravity. The terms beneath the diagonal should vanish for a naive agreement with supergravity at finite $`N`$, while the terms above the diagonal correspond to quantum gravity corrections. It was argued in that the sum of diagonal terms corresponding to the effective classical supergravity potential between two gravitons should be given by
$$V_{\mathrm{classical}}=\frac{2}{15R^2}\left(1\sqrt{1\frac{15R}{2}v^2}\right).$$
(5.68)
Now let us discuss the individual terms in (5.65). As we have discussed, the one-loop analysis gives a term
$$V_{1,4}=\frac{15}{16}$$
(5.69)
which agrees with linearized supergravity. The analysis of DKPS can be extended to the remaining one-loop terms. The next one-loop term vanishes
$$V_{1,6}=0.$$
(5.70)
Some efforts have been made to relate the higher order terms $`V_{1,8},\mathrm{}`$ to quantum effects in 11D supergravity, but so far this interpretation is not clear. For some discussion of this issue see and references therein. The term
$$V_{2,4}=0$$
(5.71)
was computed by Becker and Becker . As expected, this term vanishes. The term
$$V_{2,6}=\frac{225}{32}$$
(5.72)
was computed in . This term agrees with the expansion of (5.68). A general expression for the two-loop effective potential given by the second line of (5.65) was given in .
It was shown in by Paban, Sethi and Stern that there can be no higher-loop corrections to the $`v^4`$ and $`v^6`$ terms on the diagonal. Their demonstration of these results follows from a consideration of the terms with the maximal number of fermions which are related to the $`v^4`$ and $`v^6`$ terms by supersymmetry. For example, this is the $`\psi ^8/r^{11}`$ term in the case $`v^4`$. They show that the fermionic terms are uniquely determined by supersymmetry, and that this in turn uniquely fixes the form of the bosonic terms proportional to $`v^4`$ and $`v^6`$ (see also for more about the case of $`v^4`$). Thus, they have shown that that
$$V_{(n>1),4}=V_{(n>2),6}=0.$$
(5.73)
This nonrenormalization theorem was originally conjectured by BFSS in analogy to similar known theorems for higher-dimensional theories .
This completes our summary of what is known about 2-body interactions in matrix theory. The complete set of known terms is given by
$$\begin{array}{ccccccccccc}V& =& \frac{1}{2R}v^2& & & & & & & & \\ & & & +& \frac{15}{16}\frac{v^4}{r^7}& +& 0& +& (\mathrm{known})& & \\ & & & +& 0& +& \frac{225}{32}R\frac{v^6}{r^{14}}& +& (\mathrm{known})& & \\ & & & +& 0& +& 0& +& \mathrm{?}& +& \mathrm{}\\ & & & & & & & +& \mathrm{}& +& \mathrm{}\end{array}$$
(5.74)
It has been proposed that for arbitrary $`N`$ the analogues of the higher-loop diagonal terms should naturally take the form of a supersymmetric Born-Infeld type action , which would give rise in the case $`N=2`$ to a sum of the form (5.68). There is as yet, however, no proof of this statement beyond two loops. One particular obstacle to calculating the higher-loop terms in this series is that it is necessary to integrate over loops containing propagators of massless fields. These propagators can give rise to subtle infrared problems with the calculation. Some of these difficulties can be avoided by trying to reproduce higher-order supergravity interactions from interactions of more than two objects in matrix theory, the subject to which we will turn in section 5.2.
#### 5.1.5 The Equivalence Principle in matrix theory
We have seen that the form of the linearized theory of 11D supergravity is precisely reproduced by a one-loop calculation in matrix theory. This equivalence follows provided that the expressions in (5.41-5.43), as well as the higher moments of these expressions and related expressions for the fermion components of the supercurrent are interpreted as definitions of the stress tensor, membrane current and other supercurrent components of a given matrix theory object. It is perhaps somewhat surprising given that this correspondence holds exactly at finite $`N`$ to observe that Einstein’s Equivalence Principle breaks down at finite $`N`$, even in the linearized theory .
The Equivalence Principle essentially states that given a background gravitational field produced by some source matter configuration, any two objects which are small compared to the scale of variation in the metric and which have the same initial space-time velocity vector $`\dot{x}^I`$ will follow identical trajectories through space-time. This follows from the fact that objects which are moving in the influence of a gravitational field follow geodesics in space-time. Of course, this result is only valid if the objects are not influenced by any other fields in the theory such as an electromagnetic 1-form or 3-form field.
To see a simple example of a case where the equivalence principle is violated in matrix theory, consider a source at the origin consisting of a single graviton with $`p^+=\stackrel{~}{N}/R`$ and $`\stackrel{~}{v}^i=0`$. This source produces a long-range gravitational field and no 3-form or gravitino field. Now consider a probe object at a large distance $`r`$. We take the probe to be a small membrane sphere, initially stationary, of radius $`r_0`$ and with longitudinal momentum $`p^+=N/R`$. It is straightforward to calculate the energy $`p^{}`$ of the membrane; we find that the 11-momentum of the membrane has the light-front components
$$p^+=N/Rp^i=0p^{}=\frac{8r_0^4}{RN^3}c_2.$$
The initial velocity of the membrane is then
$$\dot{x}^+=1\dot{x}^i=0\dot{x}^{}=\frac{p^{}}{p^+}=8r_0^4\frac{c_2}{N^4}.$$
According to the equivalence principle, any two membrane spheres with different values of $`r_0`$ but the same value of $`\dot{x}^{}=r_0^4c_2/N^4`$ should experience precisely the same acceleration. Using the general formula for the 2-body interaction potential in matrix theory, however, it is straightforward to calculate
$$\ddot{x}^i=\frac{R}{N}\frac{V_{\mathrm{matrix}}}{x^i}=1680R\stackrel{~}{N}\frac{x^i}{|x|^9}\frac{r_0^8}{N^8}\left(c_2^2\frac{1}{3}c_2\right).$$
The leading term in an expansion in $`1/N`$ of this acceleration is indeed a function of $`\dot{x}^{}`$. Thus, in the large $`N`$ limit the equivalence principle is indeed satisfied. The subleading term, however, has a different dependence on $`r_0`$ and $`N`$. Thus, the equivalence principle is not satisfied at finite $`N`$.
This result implies that even if finite $`N`$ matrix theory is equivalent to DLCQ M-theory, this theory does not seem to be related to a smooth theory of Einstein-Hilbert gravity, even on a compact space and with restrictions on longitudinal momentum. This is not a problem if one only takes seriously the large $`N`$ version of the conjecture. If one wishes to make sense of the finite $`N`$ theory in terms of some theory with a reasonable classical limit, however, it may be necessary to consider some new ideas for what this theory may be. It is tempting to think that the theory at finite $`N`$ might be some sort of theory of classical gravity on a noncommutative space. Since the equivalence principle in the form we have been using it depends upon the geodesic equations, which are defined only on a smooth commutative space, it is natural to imagine that this principle might have to be corrected at finite $`N`$ when the space has nontrivial noncommutative structure.
### 5.2 The N-body problem
So far we have seen that in general the linearized theory of supergravity is correctly reproduced by an infinite series of terms arising from one-loop calculations in matrix theory. We have also discussed 2-loop calculations of two-graviton interactions which seem to agree with supergravity. If matrix theory is truly to reproduce all of classical supergravity, however, it must reproduce all the nonlinear effects of the fully covariant gravitational theory. The easiest way to study these nonlinearities is to consider N-body interaction processes. For example, following let us consider a probe body at position $`r_3`$ in the long-range gravitational field produced by a pair of bodies at positions $`r_1=0,r_2r_3`$. We can consider a perturbative expansion of Einstein’s equations. At leading order we have the linearized theory which gives a long-range field satisfying (schematically, dropping indices)
$$^2hT$$
where $`T`$ is a matter source. At the next order we have
$$^2h+h^2h+(h)^2T+Th,$$
which we can rewrite in the form
$$^2hT+Th+h^2h+(h)^2$$
(5.75)
The action of a probe object in the long-range field produced by objects 1 and 2 can be written in a double expansion in the inverse separations $`r_3`$ and $`r_2`$ as
$$T_3h_{12}\frac{T_3(T_1+T_2)}{r_3^7}+\frac{T_3T_{(12)}}{r_3^7r_2^7}+\mathrm{}$$
(5.76)
where $`T_{(12)}`$ is an interaction term contributing through the quadratic terms on the RHS of (5.75). On the matrix theory side, an apparently analogous calculation can be performed by first doing the one-loop calculation we have already described to find the linearized interaction between the 3rd object and the 1-2 system, and then doing a further loop integration to evaluate the quantum corrections to the long-range field generated by the first two sources, giving an expression of the form
$$\frac{T_3T_{1+2}}{r_3^7}.$$
We expect quantum corrections to the expectation value of the schematic form
$$T_{1+2}T_1+T_2+\frac{T_{(12)}}{r_2^7}+\mathrm{}$$
which roughly conforms to the structure expected from (5.76). Thus, in principle, it seems like it should be possible to make a correspondence between the double power series expansions computed in the two theories, given the results of the one-loop expansion for a completely general pair of objects such as was calculated in . Indeed, a simple subset of terms were shown to correspond in this way in . The terms considered in that paper were the terms in the 3-graviton interaction potential proportional to $`v_3^4/r_3^7`$. Considering the form discussed above for the components of the matrix stress tensor, it is clear that such terms only arise in the part of the interaction potential corresponding to
$$\frac{v_3^4T_{1+2}^{++}}{r_3^7}.$$
(5.77)
But the stress tensor component
$$T_{1+2}^{++}=\frac{N_1+N_2}{R}$$
is a constant which suffers no quantum corrections in matrix theory. This is a conserved charge: the total longitudinal momentum of the 1-2 system, and is responsible for the long-range component $`h^{++}`$ of the metric. It is therefore easy to see that this term is correctly reproduced by matrix theory. The terms corresponding to other powers of $`v`$ are more complicated, however, as the relevant components of $`T_{1+2}`$ are corrected by quantum effects.
In addition to the practical complications of the calculation, there are several conceptual subtleties in using the approach we have just described to making a concrete correspondence between the matrix theory and supergravity descriptions of a general 3-body interaction process. The first subtlety arises, as was pointed out by Okawa and Yoneya in , from the fact that the complete gravity action is not simply the probe-source term (5.76), but also contains terms cubic in the gravitational field $`(h^3)`$. These terms have a more complicated structure than the simple probe-source terms considered above, and it is more complicated to relate them to the results of the matrix theory calculation. The second subtlety which arises is that the precise choice of gauge made in the matrix theory calculation has a very strong impact on the form of the expressions found in the resulting effective action. Of course, for any physical quantity such as an S-matrix element, the result of a complete calculation will be independent of gauge choice. Nonetheless, to compare terms in the fashion we are suggested here will require a careful choice of gauge in matrix theory to match the appropriate gauge chosen in the supergravity theory. From this point of view, it is somewhat remarkable that in the calculation of the leading-order terms the natural gauge choices in the two theories (background field gauge in matrix theory and linearized gauge in supergravity) give rise to results which can be easily compared.
In any case, one might hope to navigate through these complications in the general 3-body problem, although this clearly would involve a substantial amount of work. In a very impressive pair of papers by Okawa and Yoneya (see also the more recent work ), the full S-matrix calculation was carried out for the interaction between 3 gravitons in both matrix theory and in supergravity, and it was shown that there was a precise agreement between all terms. Unlike other work on this problem, Okawa and Yoneya did not use the double expansion to simplify the problem but simply carried out the complete calculation.
One would naturally like to extend these results beyond the 3-body problem to the general N-body problem. The hierarchy of scales leading to the double expansion discussed above can be generalized, so that one has a different scale for each distance in the problem. This organizes the large number of terms in the N-body interaction into a more manageable structure. To date, however there has been very little work done on the problem of understanding higher order nonlinearities in the theory beyond those involved in the 3-body problem.
A very intriguing paper by Dine, Echols and Gray attempts to find a matrix-supergravity correspondence for some special terms in the general N-body interaction potential. Although they find that some terms agree, they also find some terms which appear in the matrix theory potential which have the wrong scaling behavior to correspond to supergravity terms. We briefly describe these terms here in the language we have been using of stress tensor components.
For a 3-graviton system the term (5.77) is associated with an infinite series of higher-moment terms, as described in subsection 5.1.3 and in more detail in . The first of these higher moment terms is
$$v_3^4T^{++(ij)}_{12}_i_j\frac{1}{r_3^7}$$
(5.78)
This expectation value is given by
$$X^iX^j\frac{\delta _{ij}}{r_2}+\frac{v_2^iv_2^j+\delta ^{ij}v_2^2}{r_2^5}+\mathrm{}$$
The contribution to (5.78) from the first delta function vanishes since $`^2r^7=0`$ away from the origin in the 9-dimensional transverse space. The second term gives rise to a term in the 3-body potential of the form
$$V_a\frac{v_3^4(v_2)^2}{r_2^5}\frac{1}{r_3^7}$$
Dine, Echols and Gray argue that such a term should also be found in supergravity, giving an example of an agreement between two-loop matrix theory and tree level supergravity in the $`U(3)`$ theory at order $`v^6/r^{14}`$. This argument can be repeated by taking a higher moment of this term in a 4-body system
$$V_av_4^4(v_3)^2\frac{1}{r_4^7}X^iX^j_{12}_i_j\frac{1}{r_3^5}$$
This time, however, the first term in the expectation value does not give 0, so that matrix theory predicts a term of the form
$$V_av_4^4\left((v_3)^2\frac{1}{r_4^7}\right)\left(^2\frac{1}{r_3^5}\right)\frac{1}{r_2}$$
As argued by Dine, Echols and Gray, this term has the wrong scaling to correspond to a classical supergravity interaction. Indeed, this term is of the form $`v^6/r^{17}`$, corresponding to a term “below the diagonal”, which is expected to vanish.
The appearance of this term in the matrix theory perturbation series is troubling. It seems to indicate that there may be a breakdown of the correspondence between matrix theory and even classical supergravity. This is the first concrete calculation where the two perturbative expansions have been shown to contain terms which may disagree. On the other hand, there are subtleties in this calculation which may need be resolved. For one thing, there are the issues of gauge choices mentioned above. This calculation implicitly assumes a gauge which may not be appropriate for comparison to the 4-body interaction terms being considered in supergravity. There are also issues of infrared divergences which may lead to unexpected cancellations. In any case, clearly more work is needed to determine whether this indeed represents a breakdown of the relationship between matrix theory and classical supergravity which works so well for lower order terms.
### 5.3 Longitudinal momentum transfer
In this section we have so far concentrated on interactions in matrix theory and supergravity where no longitudinal momentum is transferred from one object to another. A supergravity process in which longitudinal momentum is transferred from one object to another is described in the IIA theory by a process where one or more D0-branes are exchanged between coherent states consisting of clumps of D0-branes. Such processes are exponentially suppressed since the D0-branes are massive, and thus are not relevant for the expansion of the effective potential in terms of $`1/r`$ which we have been discussing. In the matrix theory picture, this type of exponentially suppressed process can only appear from nonperturbative effects. Clearly, however, for a full understanding of interactions in Matrix theory it will be necessary to study processes with longitudinal momentum transfer in detail and to show that they also correspond correctly with processes in supergravity and M-theory. Some progress has been made in this direction. Polchinski and Pouliot have calculated the scattering amplitude for two 2-branes for processes in which a 0-brane is transferred from one 2-brane to the other . In the Yang-Mills picture on the world-volume of the 2-branes, the incoming and outgoing configurations in this calculation are described in terms of an $`U(2)`$ gauge theory with a scalar field taking a VEV which separates the branes. The transfer of a 0-brane corresponds to an instanton-like process where a unit of flux is transferred from one brane to the other. The amplitude for this process was computed by Polchinski and Pouliot and shown to be in agreement with expectations from supergravity. This result suggests that processes involving longitudinal momentum transfer may be correctly described in Matrix theory. It should be noted, however, that the Polchinski-Pouliot calculation is not precisely a calculation of membrane scattering with longitudinal momentum transfer in Matrix theory since it is carried out in the 2-brane gauge theory language. In the T-dual Matrix theory picture the process in question corresponds to a scattering of 0-branes in a toroidally compactified space-time with the transfer of membrane charge. Processes with 0-brane transfer and the relationship between these processes and graviton scattering in matrix theory have been studied further in .
## 6 Matrix theory in a general background
So far we have only discussed matrix theory as a description of M-theory in infinite flat space. In this section we consider the possibility of extending the theory to compact and curved spaces. As a preliminary to the discussion of compactification, we give an explicit description of T-duality in gauge theory language in subsection 6.1. We then discuss the compactification of the theory on tori in subsection 6.2. Following the discussion of matrix theory compactification, we turn in subsection 6.3 to the problem of using matrix theory methods to describe M-theory in a curved background space-time.
### 6.1 T-duality
In this subsection we briefly review how T-duality may be understood from the point of view of super Yang-Mills theory. For more details see .
In string theory, T-duality is a symmetry which relates the type IIA theory compactified on a circle of radius $`R_9`$ with type IIB theory compactified on a circle with dual radius $`\widehat{R}_9=\alpha ^{}/R_9`$. In the perturbative type II string theory, T-duality exchanges winding and momentum modes of the closed string around the compact direction. For open strings, Dirichlet and Neumann boundary conditions are exchanged by T-duality, so that Dirichlet $`p`$-branes are mapped under T-duality to Dirichlet $`(p\pm 1)`$-branes .
It was argued by Witten that the low-energy theory describing a system of $`N`$ parallel D$`p`$-branes in flat space is the dimensional reduction of $`𝒩=1`$, $`(9+1)`$-dimensional super Yang-Mills theory to $`p+1`$ dimensions. In the case of $`N`$ D0-branes, this gives the Lagrangian (3.2). To understand T-duality from the point of view of this low-energy field theory, we consider the simplest case of $`N`$ D0-branes moving in a space which has been compactified in a single direction by identifying
$$x^9x^9+2\pi R^9.$$
(6.1)
To interpret this equivalence in terms of the matrix degrees of freedom of the D0-branes it is natural to pass to the covering space $`R^{9,1}`$, where the $`N`$ D0-branes are each represented by an infinite number of copies labeled by integers $`nZ`$. We can thus describe the dynamics of $`N`$ D0-branes on $`R^{8,1}\times S^1`$ by a set of infinite matrices $`M_{ma,nb}^i`$ where $`a,b\{1,\mathrm{},N\}`$ are $`U(N)`$ indices and $`m,nZ`$ index copies of each D0-brane which differ by translation in the covering space (See Figure 4). In terms of this set of infinite matrices, the quotient condition (6.1) becomes a set of constraints on the allowed matrices which can be written (dropping the $`U(N)`$ indices $`a,b`$) as
$`X_{mn}^i`$ $`=`$ $`X_{(m1)(n1)}^i,i<9`$
$`X_{mn}^9`$ $`=`$ $`X_{(m1)(n1)}^9,mn`$ (6.2)
$`X_{nn}^9`$ $`=`$ $`2\pi R_9\text{1}\text{ }\text{1}+X_{(n1)(n1)}^9.`$
From the structure of the constraints (6.2) it is natural to interpret the matrices $`X_{mn}^i`$ in terms of the $`(nm)`$th Fourier modes of a theory on the dual circle. The infinite matrix $`X^9`$ becomes a covariant derivative operator
$$X^9(2\pi \alpha ^{})(i_9+A_9)$$
(6.3)
in a $`U(N)`$ Yang-Mills theory on the dual torus, while $`X^i`$ for $`i<9`$ becomes an adjoint scalar field. The fermionic fields in the theory can be interpreted similarly.
This gives a precise equivalence between the low-energy world-volume theory of a system of $`N`$ D0-branes on $`S^1`$ and a system of $`N`$ D1-branes on the dual circle. The relationship between winding modes $`X_{mn}^i`$ in the D0-brane theory and modes with $`nm`$ units of momentum in the dual theory corresponds precisely to the mapping from winding to momentum modes in the closed string theory under T-duality.
This argument can easily be generalized to a system of multiple D$`p`$-branes transverse to a torus $`T^d`$, which are equivalent to a system of wrapped D$`(p+d)`$-branes on the dual torus. When we compactify in multiple dimensions, the possibility arises of having a topologically nontrivial gauge field configuration on the dual torus. To discuss this possibility it is useful to use a slightly more abstract language to describe the T-duality.
The constraints (6.2) can be formulated by saying that there exists a translation operator $`U`$ under which the infinite matrices $`X^a`$ transform as
$$UX^aU^1=X^a+\delta ^{a9}2\pi R_9\text{1}\text{ }\text{1}.$$
(6.4)
This relation is satisfied formally by the operators
$$X^9=i_9+A_9,U=e^{2\pi i\widehat{x}^9R_9}$$
(6.5)
which correspond to the solutions discussed above. In this formulation of the quotient theory, the operator $`U`$ generates the group $`\mathrm{\Gamma }=Z`$ of covering space transformations. Generally, when we take a quotient theory of this type, however, there is a more general constraint which can be satisfied. Namely, the translation operator only needs to preserve the state up to a gauge transformation. Thus, we can consider the more general constraint
$$UX^aU^1=\mathrm{\Omega }(X^a+\delta ^{a9}2\pi R_9\text{1}\text{ }\text{1})\mathrm{\Omega }^1.$$
(6.6)
where $`\mathrm{\Omega }U(N)`$ is an arbitrary element of the gauge group. This relation is satisfied formally by
$$X^9=i^9+A_9,U=\mathrm{\Omega }e^{2\pi i\widehat{x}^9R_9}$$
(6.7)
This is precisely the same type of solution as we have above; however, there is the additional feature that the translation operator now includes a nontrivial gauge transformation. On the dual circle $`\widehat{S}^1`$ this corresponds to a gauge theory on a bundle with a nontrivial boundary condition in the compact direction 9.
A similar story occurs when several directions are compact. In this case, however, there is a constraint on the translation operators in the different compact directions. For example, if we have compactified on a 2-torus in dimensions 8 and 9, the generators $`U_8`$ and $`U_9`$ of a general twisted sector must generate a group isomorphic to $`Z^2`$ and therefore must commute. The condition that these generators commute can be related to the condition that the boundary conditions in the dual gauge theory correspond to a well-defined $`U(N)`$ bundle over the dual torus. For compactifications in more than one dimension such boundary conditions can define a topologically nontrivial bundle. It is interesting to note that this construction can even be generalized to situations where the generators $`U_i`$ do not commute. Physically, such a configuration is produced when there is a background NS-NS $`B`$ field. This construction leads to a dual theory which is described by gauge theory on a noncommutative torus . A description of this scenario along the lines of the preceding discussion is given in . The connection between nontrivial background field configurations and noncommutative geometry has been a subject of much recent interest .
### 6.2 Matrix theory on tori
From the discussion in the previous section, it follows that the matrix theory description of M-theory compactified on a torus $`T^d`$ becomes $`(d+1)`$-dimensional super Yang-Mills theory. The argument of Seiberg and Sen in is valid in this situation, so that $`U(N)`$ super Yang-Mills theory on $`(T^d)^{}`$ should describe M-theory compactified on $`T^d`$. When $`d3`$ the quantum super Yang-Mills theory is renormalizable so this is a sensible way to approach the theory. As the dimension of the torus increases, however, the matrix description of the theory develops more and more complications. In general, the super Yang-Mills theory on the $`d`$-torus encodes the full U-duality symmetry group of M-theory on $`T^d`$ in a rather nontrivial fashion.
Compactification of the theory on a two-torus was discussed by Sethi and Susskind . They pointed out that as the $`T^2`$ shrinks, a new dimension appears whose quantized momentum modes correspond to magnetic flux on the $`T^2`$. In the limit where the area of the torus goes to 0, an $`O(8)`$ symmetry appears. This corresponds with the fact that IIB string theory appears as a limit of M-theory on a small 2-torus .
Compactification of the theory on a three-torus was discussed in . In this case, M-theory on $`T^3`$ is equivalent to $`(3+1)`$-dimensional super Yang-Mills theory on a torus. This theory is conformal and finite. M-theory on $`T^3`$ has a special type of T-duality symmetry under which all three dimensions of the torus are inverted. In the matrix description this is encoded in the Montanen-Olive S-duality of the 4D super Yang-Mills theory.
When more than three dimensions are toroidally compactified, the theory undergoes even more remarkable transformations . When compactified on $`T^4`$, the manifest symmetry group of the theory is $`SL(4,Z)`$. The expected U-duality group of M-theory compactified on $`T^4`$ is $`SL(5,Z)`$, however. It was pointed out by Rozali that the U-duality group can be completed by interpreting instantons on $`T^4`$ as momentum states in a fifth compact dimension. This means that Matrix theory on $`T^4`$ is most naturally described in terms of a (5 + 1)-dimensional theory with a chiral $`(2,0)`$ supersymmetry. This unusual $`(2,0)`$ theory with 16 supersymmetries appears to play a crucial role in numerous aspects of the physics of M-theory and 5-branes, and has been studied extensively in recent years.
Compactification on tori of higher dimensions than four continues to lead to more complicated situations, particularly when one gets to $`T^6`$, when the matrix theory description seems to be as complicated as the original M-theory. A significant amount of literature has been produced on this subject, to which the reader is referred to further details (see for reviews and further references). Despite the complexity of $`T^6`$ compactification, however, it was suggested by Kachru, Lawrence and Silverstein that compactification of Matrix theory on a more general Calabi-Yau 3-fold might actually lead to a simpler theory than that resulting from compactification on $`T^6`$. If this speculation is correct and a more explicit description of the theory on a Calabi-Yau compactification could be found, it might make matrix theory a possible approach for studying realistic 4D phenomenology.
### 6.3 Matrix theory in curved backgrounds
We now consider matrix theory in a space which is infinite but may be curved or have other nontrivial background fields. We would like to generalize the matrix theory action to one which includes a general supergravity background given by a metric tensor, 3-form field, and gravitino field which together satisfy the equations of motion of 11D supergravity. This issue has been discussed in . In it was argued that light-front M-theory on an arbitrary compact or non-compact manifold should be reproduced by the low-energy D0-brane action on the same compact manifold; no explicit description of this low-energy theory was given, however. In an explicit prescription was given for the first few terms of a matrix theory action on a general Kähler 3-fold which agreed with a general set of axioms proposed in . In and , however, it was argued that no finite $`N`$ matrix theory action could correctly reproduce physics on a large K3 surface. We review here an explicit proposal for a formulation of matrix theory in an arbitrary background geometry originally presented in .
If we assume that matrix theory is a correct description of M-theory around a flat background, then there is a large class of curved backgrounds for which we know it is possible to construct a matrix theory action for $`N\times N`$ matrices. This is the class of backgrounds which can be produced as long-range fields produced by some other supergravity matter configuration with a known description in matrix theory. Imagine that a background metric $`g_{IJ}=\eta _{IJ}+h_{IJ}`$, a 3-form field $`A_{IJK}`$ and a gravitino field $`\psi _I`$ of light-front compactified 11-dimensional supergravity can be produced by a matter configuration described in matrix theory by matrices $`\stackrel{~}{X}^i`$. Then the matrix theory action describing $`N\times N`$ matrices $`X^i`$ in this background should be precisely the effective action found by considering the block-diagonal matrix configuration
$$X^i=\left[\begin{array}{cc}X^i& 0\\ 0& \stackrel{~}{X}^i\end{array}\right]$$
(and a similar fermion configuration) and integrating out the off-diagonal fields as well as fluctuations around the background $`\stackrel{~}{X}`$.
From the results found in , we know that for weak background fields, the first few terms in an expansion of this effective action in the background metric are given by
$`S_{\mathrm{eff}}`$ $`=`$ $`S_{\mathrm{matrix}}+{\displaystyle 𝑑xT^{IJ}(x)h_{IJ}(x)}+\mathrm{}`$
$`=`$ $`S_{\mathrm{matrix}}+{\displaystyle 𝑑x^+\{T^{IJ}h_{IJ}(0)+T^{IJ(i)}_ih_{IJ}(0)+\mathrm{}\}}+\mathrm{}`$
where $`T^{IJ(\mathrm{})}`$ are the moments of the matrix theory stress-energy tensor, and there are analogous terms for the coupling of the membrane, 5-brane and fermionic components of the supercurrent to $`A_{IJK}`$ and $`S_I`$. If the standard formulation of matrix theory in a flat background is correct, the absence of corrections to the long-range $`1/r^7`$ potential around an arbitrary matrix theory object up to at least order $`1/r^{11}`$ implies that this formulation must be correct at least up to terms of order $`^4h`$ and $`h^2`$.
As we have derived it, this formulation of the effective action is only valid for certain background geometries which can be produced by well-defined matrix theory configurations. It is natural, however, to suppose that this result can be generalized to an arbitrary background. Thus, it is proposed in that up to nonlinear terms in the background, the general form of the matrix theory action in an arbitrary but weak background is given by
$`S_{\mathrm{weak}}`$ $`=`$ $`{\displaystyle }d\tau {\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{i_1,\mathrm{},i_n}{}}{\displaystyle \frac{1}{n!}}(T^{IJ(i_1\mathrm{}i_n)}_{i_1}\mathrm{}_{i_n}h_{IJ}`$
$`+J^{IJK(i_1\mathrm{}i_n)}_{i_1}\mathrm{}_{i_n}A_{IJK}`$
$`+M^{IJKLMN(i_1\mathrm{}i_n)}_{i_1}\mathrm{}_{i_n}A_{IJKLMN}^D`$
$`+\mathrm{fermion}\mathrm{terms})`$
Let us make several comments about this action. First, this formulation is only appropriate for backgrounds with no explicit $`x^{}`$ dependence, as we do not understand how to encode higher modes in the compact direction in the components of the supergravity currents. Second, note that the coupling to $`A^D`$ is free of ambiguity since the net 5-brane charge must vanish for any finite matrices, so that only first and higher derivatives of $`A^D`$ appear in the action. Third, note that though we only have explicit expressions for the fermion terms in the zeroeth and some of the first moments of $`T`$, we may in principle generalize the calculations of to determine all the fermionic contributions from higher order terms in the one loop matrix theory potential.
The linearized couplings in the action (6.3) are motivated by the results of one-loop calculations in matrix theory. In principle, it may be possible to extend the formulation of matrix theory in weak background fields to higher order by performing general higher-loop calculations in matrix theory. For example, a complete description of the 2-loop interaction in matrix theory between an arbitrary 3 background configurations would suggest the form of the coupling between one object considered as a probe and the quadratic terms in the background produced by the other pair of objects. Generally, knowing the full $`n`$-loop interaction between $`n+1`$ matrix theory objects would suggest the $`n`$th order coupling of the matrix degrees of freedom to the background fields. Unfortunately, as we have discussed such calculations are rather complicated. In addition to the technical difficulties of doing the general 2-loop calculation, there are subtleties related to the gauge choice and possible infrared divergences. Furthermore, finite $`N`$ calculations will only help us to learn the higher-order couplings to the background if the results of these calculations are protected by supersymmetric nonrenormalization theorems, and as we have discussed there is no strong reason to believe that such nonrenormalization theorems hold for the general $`n`$-loop $`SU(N)`$ calculation. Thus, to write a completely general coupling of matrix theory to a nontrivial supergravity background, it is probably necessary to find a new general principle, such as a matrix version of the principle of coordinate invariance.
Another approach which one might take to define matrix theory in a general background geometry is to follow the original derivation of matrix theory as a regularized membrane theory, but to include a general background geometry instead of a flat background as was used in . The superspace formulation of a supermembrane theory in a general 11D supergravity background was given in . In principle, it should be possible to simply apply the matrix regularization procedure to this theory to derive matrix theory in a general background geometry. Unfortunately, however, the connection between superspace fields and component fields is not well-understood in this theory. Until recently, in fact, the explicit expressions for the superspace fields were only known up to first order in the component fermion fields $`\theta `$ . In , this analysis was extended to quadratic order in $`\theta `$ with the goal of finding an explicit formulation of the supermembrane in general backgrounds in terms of component fields, to which the matrix regulation procedure could be applied to generate a general background formulation of matrix theory. These results can be compared with the proposal just described for the linear couplings to the background. The two formulations seem to be completely compatible , although extra terms appear in the matrix theory action which cannot be predicted from the form of the continuous membrane theory.
In , Douglas proposed that any formulation of matrix theory in a curved background should satisfy a number of axioms. All these axioms are satisfied in a straightforward fashion by the proposal in , except one: this exception is the axiom that states that a pair of D0-branes at points $`x^i`$ and $`y^i`$ should correspond to diagonal $`2\times 2`$ matrices where the masses of the off-diagonal fields should be equal to the geodesic distance between the points $`x^i`$ and $`y^i`$ in the given background metric. In it was shown that the linearized terms in the action (6.3) are consistent with this condition and that the linear variation in geodesic distance between a pair of D0-branes is correctly reproduced by coupling the matrix theory stress tensor to the background metric through a combinatorial identity which follows from the particular ordering implied by the symmetrized trace form of the multipole moments of the stress tensor. The fact that this condition can be satisfied at linear order provides hope that it might be possible to extend the action to all orders in a consistent way. In , it was indeed shown by Douglas, Kato and Ooguri that a set of some higher order terms for the action on a Ricci-flat Kähler manifold can be found which are consistent with the geodesic length condition, but these authors also found that this condition did not uniquely determine most of the terms in the action so that a more general principle is still needed to construct the action to all orders.
We synopsize the discussion in this section as follows: (6.3) seems to be a consistent proposal for the linearized couplings between matrix theory and weak supergravity background fields. The expressions for the higher moments of the supergravity currents which couple to the derivatives of the background fields are known up to terms quadratic in the fermions, and the remaining terms can be found from a one-loop matrix theory computation. This proposal can be generalized to $`m`$th order in the background fields, where matrix expressions are needed for quantities which can be determined from an $`m`$-loop matrix theory calculation. Whether these terms can be calculated and sensibly organized into higher-order couplings of matrix theory to background fields depends on whether higher-loop matrix theory results are protected by supersymmetric nonrenormalization theorems. It is worth emphasizing that the definitions of the matrix theory currents we have described here depend upon gauge choices for the propagating supergravity fields. For a given gauge choice, the theory is only defined for backgrounds compatible with the gauge condition. Making the appropriate gauge choices represents another obstacle to carrying out this analysis to higher order.
## 7 Outlook
We conclude with a brief review of the connection between matrix theory and M-theory, and a short discussion of the current state of affairs and the outlook for further developments in matrix theory.
We have discussed two complementary ways of thinking about matrix theory: first as a quantized regularized theory of a supermembrane, which naturally describes a second-quantized theory of objects moving in an 11-dimensional target space, and second as the DLCQ of M-theory which is equivalent to a simple limit of type IIA string theory through the Seiberg-Sen limiting argument.
Using matrix degrees of freedom, it is possible to describe pointlike objects which have many of the physical properties of supergravitons. It is also possible to use the matrix degrees of freedom to describe extended objects which behave like the supermembrane and 5-brane of M-theory. For supergravitons and membranes this story seems fairly complete; for 5-branes, only a few very special geometries have been described in matrix language, and a complete description of dynamical (longitudinal) 5-branes, even at the classical level, is still lacking.
As we have discussed, to date all perturbative calculations except the 3-loop calculation of Dine, Echols and Gray indicate that matrix theory correctly reproduces classical 11D supergravity. It has been suggested that the agreement between the theories at 1-loop and 2-loop orders is essentially an accident of supersymmetry, however there is little understanding of how to interpret or organize higher-loop terms. There is also very little understanding at this point of how quantum corrections to the supergravity theory can be understood in terms of matrix theory, although there is evidence that quantum gravity effects are not reproduced by perturbative calculations in matrix theory but will require a better understanding of the large $`N`$ limit of the theory.
At this point there are essentially 4 possible scenarios for the validity of the matrix theory conjecture:
i) Matrix theory is correct, and DLCQ supergravity is reproduced at finite $`N`$ by perturbative matrix theory calculations.
ii) Matrix theory is correct in the large $`N`$ limit, and noncompact supergravity is reproduced by a naive large $`N`$ limit of the standard perturbative matrix theory calculations.
iii) Matrix theory is correct in the large $`N`$ limit, but to connect it with supergravity, even at the classical level, it is necessary to deal with subtleties in the large $`N`$ limit. (i.e., there are problems with the standard perturbative analysis at higher order)
iv) Matrix theory is simply wrong, and further terms need to be added to the dimensionally reduced super Yang-Mills action to find agreement with M-theory even in the large $`N`$ limit.
Now let us examine the evidence:
$``$ The breakdown of the Equivalence Principle seems incompatible with (i), but compatible with all other possibilities.
$``$ If the result of Dine, Echols and Gray in is correct, and has been correctly interpreted, clearly (i) and (ii) are not possible. The fact that the methods of Paban, Sethi and Stern for proving nonrenormalization theorems in the $`SU(2)`$ theory break down for $`SU(3)`$ at two loops and at higher loop order also hints that (ii) may not be correct.
$``$ The analysis of Seiberg and Sen seems to indicate that one of the possibilities (i)-(iii) should hold.
It seems that (iii) is the most likely possibility, given this limited evidence. There are several issues which are extremely important in understanding how this problem will be resolved. The first is the issue of Lorentz invariance. If a theory contains linearized gravity and is Lorentz invariant, then it is well known that it must be either the complete generally covariant gravity theory or just the pure linearized theory. Since we know that matrix theory has some nontrivial nonlinear structure which reproduces part of the nonlinearity of supergravity, it would seem that the conjecture must be valid if and only if the theory is Lorentz invariant. Unfortunately, so far there is no complete understanding of whether the quantum theory is Lorentz invariant (classical Lorentz invariance was demonstrated in ). It was suggested by Lowe in that the problems found in might be related to a breakdown of Lorentz invariance and that in fact extra terms must be added to the theory to restore this invariance; this would lead to possibility (iv) above.
Another critical issue in understanding how the perturbative matrix theory calculations should be interpreted is the issue of the order of limits. In the perturbative calculations discussed here we have assumed that the longitudinal momentum parameter $`N`$ is fixed for each of the objects we are taking as a background, and we have then taken the limit of large separations between each of the objects. Since the size of the wavefunction describing a given matrix theory object will depend on $`N`$ but not on the separation from a distant object, this gives a systematic approximation scheme in which the bound state and wavefunction effects for each of the bodies can be ignored in the perturbative analysis. If we really are interested in the large $`N`$ theory, however, the correct order of limits to take is the opposite. We should fix a separation distance $`r`$ and then take the large $`N`$ limit. Unfortunately, in this limit we have no systematic approximation scheme. The wavefunctions for each of the objects overlap significantly as the size of the objects grows. Indeed, it was argued recently by Polchinski that the size of the bound state wavefunction of $`N`$ D0-branes will grow at least as fast as $`N^{1/3}`$. As emphasized by Susskind in , this overlap of wavefunctions makes the theory very difficult to analyze. Indeed, if possibility (iii) above is correct, it may be very difficult to use matrix theory to reproduce all the nonlinear structure of classical supergravity, let alone to derive new results about quantum supergravity. On the other hand, it may be that whatever mechanism allows the one-loop and two-loop matrix theory results to correctly reproduce the first few terms in supergravity and to evade the problem of wavefunction overlap may persist at higher orders. Indeed, one of the must important outstanding questions regarding matrix theory is to understand precisely which terms in the naive perturbative expansion of the quantum mechanics will agree with classical supergravity, and more importantly, why these terms agree. As mentioned in the last section, one of the other main outstanding problems in matrix theory is understanding how the matrix quantum mechanics theory behaves when M-theory is compactified on a curved manifold. In order to use matrix theory to make new statements about corrections to classical supergravity in phenomenologically interesting models such as M-theory on compact 7-manifolds or orbifolds, it will be necessary to solve both of these problems. In each case, a certain amount of luck will be needed for it to be possible to probe physically interesting questions using existing computational techniques.
In these lectures we have focused on understanding some basic aspects of matrix theory: the definitions of the theory in terms of the membrane and DLCQ of M-theory, and the construction of the objects and supergravity interactions of M-theory using matrix degrees of freedom. We conclude with a few brief words about some of the topics we have not discussed.
In addition to the matrix model of M-theory, there have been numerous related models suggested in the literature in the last few years. Some of these which have received particular attention are the $`(0+0)`$-dimensional matrix model of IIB string theory suggested by Ishibashi, Kawai, Kitazawa and Tsuchiya , the $`(1+1)`$-dimensional matrix string theory of Dijkgraaf, Verlinde and Verlinde and the family of AdS/CFT conjectures proposed by Maldacena . All these proposals relate a particular limit of string theory or M-theory in a fixed background to a field theory. Many connections between these models have been made, and in fact most of these proposals are related by a duality symmetry to the matrix theory we have discussed here. A fundamental question at this point, however, is how we may move away from a fixed background and discuss questions of cosmological significance.
Even within the framework of the matrix model of M-theory we have discussed in these lectures, there are many very interesting directions and particular applications which have been pursued which we did not have time to review here in any detail. These include questions about black holes in matrix theory (see, e.g., and references therein), higher dimensional compactifications and the matrix model of the (2, 0) theory which arises upon compactification on $`T^4`$ (, see for a review and further references), the detailed structure of the $`N=2`$ bound state (see e.g., and references therein), and many other directions of recent research.
In closing, it seems that matrix theory has achieved something which just a few years ago would have been deemed virtually impossible to accomplish in such a simple fashion: it gives a well-defined framework for M-theory and quantum gravity which reduces any problem, at least in light-front coordinates, to a computation which can in principle be defined and fed into a computer. Thus, in some sense this may be the first concrete answer to the problem of finding a consistent theory of quantum gravity. Unfortunately, even though this theory is a simple quantum mechanics theory, and not even a field theory, it is computationally intractable at this point to ask many of the really interesting questions about M-theory using this model. It is clearly a very interesting problem to try to find better ways of doing interesting M-theory calculations using the matrix model. But even if matrix theory is never able to give us a computational handle on some of the subtle aspects of M-theory, it certainly has given us a new perspective on how to think about a microscopic theory of quantum gravity. One of the most interesting aspects of the matrix picture is the appearance of dynamical higher-dimensional extended objects from a system of ostensibly pointlike degrees of freedom, as discussed in Section 4. It seems likely that this feature of matrix theory may play a key role in future attempts to describe a more covariant or background-independent microscopic model for M-theory, string theory or quantum gravity.
## Acknowledgments
I would like to thank L. Thorlacius and the other organizers of the “Quantum Gravity” NATO Advanced Study Institute for putting together an excellent summer school and for inviting me to participate. Thanks to the students and other lecturers for providing a stimulating atmosphere and for many interesting discussions and questions. I would also like to thank N. Prezas and M. Van Raamsdonk for pointing out errors in a preliminary version of these lecture notes and suggesting constructive changes.
This work is supported in part by the A. P. Sloan Foundation and in part by the DOE through contract #DE-FC02-94ER40818.
Some of the material in these lectures was previously presented as part of a course at MIT: physics 8.871, Fall 1998. Some of this material was also previously presented in lectures at the Korean Institute for Advanced Study in October 1998 and at the Komaba ’99 workshop in Tokyo, 1999. |
warning/0002/hep-th0002088.html | ar5iv | text | # Noether Currents for Bosonic Branes
## 1 Introduction
A brane is a relativistic extended object propagating in a background spacetime, usually treated as fixed, of some given dimension $`N`$. The spacetime trajectory of the brane is a timelike worldvolume of some lower dimension $`D`$: for a relativistic string $`D=2`$; for a domain wall, or membrane, $`D=3`$, and so on. The dynamics of the brane is described by some local action constructed using scalars that characterize the geometry of the worldvolume. The possibilities are limited by worldvolume reparameterization invariance, ambient spacetime diffeomorphism invariance, and, when $`D<N1`$, invariance under rotations in the $`ND`$ dimensional plane normal to the worldvolume. The simplest action of this kind is the Dirac-Nambu-Goto \[DNG\] action, proportional to the volume swept out by the brane in the course of its evolution. This action depends only on first order derivatives of the field variables, the embedding functions that define the worldvolume. Originally proposed by Dirac in the context of an extensible model for the electron ($`D=3`$) , it was later exploited by Nambu and Goto for the case of a relativistic string ($`D=2`$) to model hadronic matter . It constitutes the point of departure in the construction of modern string theory . Higher dimensional ($`D>1`$) DNG branes also play an increasingly important role as solitonic solutions of this theory (see e.g. ).
Despite significant developments in string theory, it would appear fair to claim that a comprehensive understanding of the dynamics of relativistic extended objects is still lacking. In particular, there are many contexts where the DNG action is clearly inadequate, and it has been necesssary to consider a more general action that includes higher-order geometrical scalars which depend on the curvature of the worldvolume; possibly even its derivatives. These additions do modify the dynamics substantially. The lowest order correction to the dynamics of a purely geometrical object is described by an action quadratic in the extrinsic curvature. For example, in a phenomenological approach, the addition of such a ‘rigidity’ term to the DNG action was proposed in the eighties by Polyakov, and independently by Kleinert, as an improved effective action for QCD . Earlier still, hamiltonians of this form were studied in condensed matter physics as models to describe the mechanical properties of lipid membranes . Higher-order corrections can also arise in systematic approximations. For example, in the context of the physics of topological defects, curvature terms are induced by considering an expansion in the thickness of the defect . The addition of such terms typically associates an energy penalty with the formation of a spike. These corrections thus become important on short distance scales serving to smooth out the curvature singularities which will arise in the course of the DNG dynamics. In Ref. , an effective action of this form was obtained explicitly for domain walls by integrating out the microscopic degrees of freedom of the underlying field theory which vary rapidly on the length scale of the wall. In the case of strings it has been argued that the lowest order effective action which describes the behavior in the neighborhood of a cusp must contain terms which are of fourth order in the extrinsic curvature of the worldvolume .
Another, different, source of higher-order terms comes from supersymmetric branes. They arise once an effective action is obtained by integrating out the fermionic degrees of freedom .
In this paper, we explore the structure of the conservation laws associated with the Poincaré invariance of a local brane action propagating in Minkowski spacetime. These occur as a consequence of the induced internal symmetries on the field theory described by the embedding of the brane worldvolume in spacetime.
At one level, the determination of the associated conserved quantities is straightforward. One can write the action explicitly in terms of the embedding functions, perhaps in terms of a natural parametrization, if available. Noether’s theorem then generates the conserved momentum and angular momentum. This procedure is perfectly adequate for a DNG object where the lagrangian is a function of the intrinsic geometry alone (see, for example, and for a DNG string, or for higher dimensional DNG objects.) If, however, the geometrical action involves either the intrinsic or the extrinsic curvature, so that the action depends on second derivatives of the embedding functions or higher, the manifest covariance of the geometrical action we start out with gets mutilated in this straighforward approach, and the conserved quantities do not possess any obvious geometrical form.
There are several ways to remedy this shortcoming. One way exploits the elegant formalism developed by Iyer and Wald for diffeomorphism invariant theories . This involves grouping higher derivative terms into totally symmetric combinations, and obtaining the general expressions for the Noether charges. For the case of branes, however, this procedure does not result in the natural geometric quantities of the worldvolume, and the translation is ackward. An alternative strategy is to perform a hamiltonian analysis. This is certainly necessary if one is interested in the canonical quantization of the theory. However, since the extrinsic curvature of a brane can be considered, roughly, as a generalization to higher dimensional objects of the acceleration of a relativistic particle, one is dealing with the complicated hamiltonian analysis of a higher derivative theory. Another approach which is spacetime manifestly covariant was pioneered by Carter (see e.g. and references therein). It focuses on the derivation of the stress-energy tensor for the theory, and the conserved quantities are constructed by contraction with the appropriate background Killing vector fields. This approach is certainly convenient when external fields are present, but, in our opinion, it does not take full advantage of the natural geometric structures on the worldvolume.
In this paper we develop an independent approach, close in spirit to the one developed by Carter, but differing in our emphasis on the worldvolume geometry. We exploit the geometrical formalism introduced by two of the authors in , tailored to the worldvolume of the extended object, to describe deformations of the worldvolume to express the variations in the worldvolume geometry induced by a Poincaré transformation in a covariant way. This will permit us to write down the corresponding conservation laws in a geometrical form adapted to the worldvolume.
Following Carter, we express the equations of motion in terms of the conservation of the linear momentum. However, we go one step further by separating these equations into two sets, obtained by projecting the local form of the conservation law onto and normal to the worldvolume. In doing this, we necessarily dismantle the worldvolume divergence appearing in the conservation law, which one might consider a step in the wrong direction. The advantage, however, is that the $`ND`$ normal projections encode completely the worldvolume diffeomorphism invariant content of the former equations. The associated gauge redundancy in the conservation laws is captured in the $`D`$ tangential projections or Bianchi identities. When we cast the conservation law this way, the distinct dynamical roles played by the tangential and normal projections of the linear momentum density in the theory also become explicit. It is then an easy matter to identify characteristics of the projections associated with a specific theory, such as, for example, the conditions under which a theory is conformally invariant.
We demonstrate how this structure can be exploited in perturbation theory. In particular, we exploit the linearization of the conservation law to provide a novel derivation of the linearized equations of motion, which, in addition to its technical merits, also throws light on the underlying structure.
We also examine angular momentum conservation. The structure of the spin for a DNG action has a very special form. We demonstrate explicitly how new structures enter when actions of higher order are considered. We point out that the Regge inequality for string theory continues to hold along extremal solutions of certain higher order theories.
An important issue in the path integral formulation of the quantum theory is the identification of appropriate boundary conditions to be imposed on the initial and final configurations in the variational principle. These boundary conditions are largely a matter of choice in the classical theory. An example is provided by the brane Einstein-Hilbert action which involves terms linear in second derivatives of the embedding functions, so that the equations of motion, in common with DNG theory, is also second order in derivatives. A spurious surface term occurs in the variational principle which signals problems in the path integral quantization of the theory. We provide a geometric approach to analysing this problem which is guided by the analogous problem in General Relativity .
It should be mentioned that various special cases the Noether currents of particular branes have been analyzed in the literature, using different approaches. For example, Letelier has considered a DNG membrane ($`D=2`$) propagating in four-dimensional Minkowski spacetime, with the addition of a term porportional to the worldvolume intrinsic curvature, the brane Einstein-Hilbert action . A relativistic string with arbitrary curvature corrections has been studied by Boisseau and Letelier . The more general case of a brane action at most quadratic in the extrinsic curvature in an arbitrary background spacetime is the subject of a detailed analysis by Carter in . For an alternative treatment of extrinsic curvature actions which uses the language of differential forms, see Hartley and Tucker in . We find complete agreement with their treatment, once details of notation have been taken into account.
As is well known, the expressions for the momenta themselves are ambigous. One is always free to add a total divergence to the lagrangian, which contributes to the momenta without affecting the dynamics. Moreover, one can add a term that is identically conserved to the momentum densities. We will keep in mind this freedom, and we will have occasion to use it for example in the specific case of the brane Einstein-Hilbert action, but we will not develop the general formalism to treat it in full generality. For the interested reader, this subject and its consequences on the global structure of the space of solutions of a relativistic theory is developed in detail e.g. by Anderson and Torre in .
This paper is organized as follows. In the next section, we define the linear and angular momentum for a general brane. In Sect. 3, we show how the equations of motion can be expressed in terms of the conservation of the linear momentum. The familiar case of a DNG brane is the subject of Sect. 4, which serves as a useful illustration of the formalism. In Sect. 5, we move on to the general case of a rigid lagrangian depending on the extrinsic curvature of the worldvolume. We obtain both the linear and angular momentum for this class of theories. In order to make our treatment more concrete we specialize to low order specific examples in Sect. 6, including the Einstein-Hilbert action for a brane. In Sect. 7, we provide a novel approach to the linearized equations of motion, in terms of the linearization of the linear momentum density. The introduction of a surface term of the York type in the brane Einstein-Hilbert action is the subject of Sect. 8. We conclude in Sect. 9 with a brief discussion. Finally, in an appendix, we sketch the extension of the formalism to lagrangians that depend on derivatives of the extrinsic curvature.
## 2 Conservation laws
Let us consider a brane, of dimension $`D1`$, propagating in a fixed background Minkowski spacetime of dimension $`N`$. For simplicity, we suppose that this object is either infinite in extent, or closed. In the case of an infinite object, we will assume in the following that appropriate fall-off conditions are chosen on the fields, so that the formal expressions we derive actually exist. The extension of the formalism to objects with finite (timelike) boundaries, or with loaded edges, relevant for the treatment of hybrid branes, such as a string with monopoles at its edges, or a wall bounded by strings (see e.g. ), is straightforward, and will not be considered in this paper.
The domain of integration is the timelike worldvolume $`m`$, given by the embedding functions
$$x^\mu =Y^\mu (\xi ^a),$$
(1)
where $`x^\mu `$ are local coordinates in the ambient Minkowski spacetime, and $`\xi ^a`$ local coordinates for the worldvolume ($`\mu ,\nu ,\mathrm{}=0,1,\mathrm{}N1`$, and $`a,b,\mathrm{}=0,1\mathrm{},D1`$). The worldvolume is described by the evolution of a $`D1`$-dimensional brane between fixed initial and final (spacelike) configurations.
The initial and final configurations are described completely by two spacelike hypersurfaces on the worldvolume $`m`$, $`\mathrm{\Sigma }_{(i)}`$ and $`\mathrm{\Sigma }_{(f)}`$. For definiteness, we parametrize these hypersurfaces by the embedding functions,
$`\xi ^a`$ $`=`$ $`X_{(i)}^a(u^A),`$ (2)
$`\xi ^a`$ $`=`$ $`X_{(f)}^a(u^A),`$ (3)
respectively, where $`u^A`$ are local coordinates for the spacelike hypersurfaces ($`A,B,\mathrm{}=1,\mathrm{},D1`$).
It turns out that it is also convenient to consider them as embedded directly in Minkowski spacetime, via map composition, so that e.g. with (1) and (2) we have
$$x^\mu =X_{(i)}^\mu (u^A)=e^\mu {}_{a}{}^{}X_{(i)}^{a}(u^A),$$
(4)
where we denote by
$$e^\mu {}_{a}{}^{}=\frac{Y^\mu }{\xi ^a},$$
(5)
the $`D`$ tangent vectors to the worldvolume $`m`$. We will also use
$$ϵ^\mu {}_{A}{}^{}=\frac{X^\mu }{u^A},$$
(6)
to denote the $`D1`$ tangent vectors to $`\mathrm{\Sigma }_{(i)}`$ or $`\mathrm{\Sigma }_{(f)}`$ as embedded in spacetime.
The union of these two hypersurfaces is the boundary $`m`$ of the worldvolume $`m`$, with the understanding that the natural orientation of the initial hypersurface $`\mathrm{\Sigma }_{(i)}`$ is opposite to that of $`\mathrm{\Sigma }_{(f)}`$. We consider only oriented branes.
We now make an infinitesimal deformation of the embedding functions for the worldvolume $`m`$,
$$Y^\mu (\xi )Y^\mu (\xi )+\delta Y^\mu (\xi ).$$
(7)
This displacement can be seen as a diffeomorphism of the ambient spacetime, and it will induce a deformation of the worldvolume geometry.
We decompose an arbitrary infinitesimal deformation of the embedding $`\delta Y^\mu `$ into its parts tangential and normal to the worldvolume,
$$\delta Y^\mu =\mathrm{\Phi }^ae^\mu {}_{a}{}^{}+\mathrm{\Phi }^in^\mu {}_{i}{}^{},$$
(8)
where the tangent vectors $`e^\mu _a`$ have been defined earlier, by Eq. (5), and $`n^\mu _i`$ are the $`ND`$ vectors normal to the worldvolume $`m`$ ($`i,j,\mathrm{}=1,2,\mathrm{},ND`$). These are defined by
$$\eta _{\mu \nu }e^\mu {}_{a}{}^{}n_{}^{\nu }{}_{i}{}^{}=0,$$
where $`\eta _{\mu \nu }`$ is the Minkowski metric, with only one minus sign, which will be used to raise and lower spacetime indices. We choose to normalize the normal vector fields as
$$\eta _{\mu \nu }n^\mu {}_{i}{}^{}n_{}^{\nu }{}_{j}{}^{}=\delta _{ij},$$
where $`\delta _{ij}`$ is the Kronecker delta.
The worldvolume geometry is described completely by the induced metric $`\gamma _{ab}`$, the extrinsic curvature $`K_{ab}^i`$, and the extrinsic twist $`\omega _a^{ij}`$, when the appropriate integrability conditions are satisfied. The induced metric is defined by
$$\gamma _{ab}=e^\mu {}_{a}{}^{}e_{}^{\nu }{}_{b}{}^{}\eta _{\mu \nu }^{}.$$
(9)
This metric, together with its inverse $`\gamma ^{ab}`$, will be used to lower and raise worldvolume indices. The quantity $`K_{ab}{}_{}{}^{i}=K_{ba}^i`$ is the extrinsic curvature along the $`i`$-th normal vector field $`n^\mu _i`$,
$$K_{ab}{}_{}{}^{i}=n_\mu {}_{}{}^{i}_{a}^{}e^\mu {}_{b}{}^{}.$$
(10)
The extrinsic twist $`\omega _a^{ij}`$, defined by
$$\omega _a{}_{}{}^{ij}=\eta _{\mu \nu }n^{\mu j}_an^{\nu i},$$
(11)
is the connection associated with covariance under normal rotations.
To evaluate the variation of these quantities under a worldvolume deformation, we will exploit the covariant formalism describing deformations of the worldvolume geometry developed in Ref.. In this approach, the effect of the deformation on geometrical tensors is covariant not only with respect to reparameterizations of the worldvolume, but also with respect to local rotations of the normals $`n^{\mu i}`$.
We consider a local action, depending on the embedding functions $`Y^\mu `$, which is both invariant under worldvolume reparametrization, and under rotations of the normals,
$$S[Y]=_m\sqrt{\gamma }L$$
(12)
(For convenience of notation, we have absorbed the worldvolume differential $`d^D\xi `$ into the integral sign. We will do likewise for the integrals over the boundary $`m`$. ) The lagrangian $`L`$ is constructed locally from the geometry of the worldvolume $`m`$
$$L=L(\gamma _{ab},K_{ab}{}_{}{}^{i},\stackrel{~}{}_aK_{bc}{}_{}{}^{i},\mathrm{}),$$
(13)
where we denote by $`\stackrel{~}{}_a`$ the covariant derivative under rotation of the normal vector fields introduced in . For an arbitrary normal vector $`\mathrm{\Psi }^i`$, it is defined by
$$\stackrel{~}{}_a\mathrm{\Psi }^i=_a\mathrm{\Psi }^i\omega _a{}_{}{}^{i}{}_{j}{}^{}\mathrm{\Psi }_{}^{j},$$
(14)
where $`_a`$ denotes the (torsionless) worldvolume covariant derivative compatible with $`\gamma _{ab}`$.
The infinitesimal variation of the action which is induced by a worldvolume deformation can always be decomposed into its tangential and normal parts,
$$\delta S=\delta _{}S+\delta _{}S.$$
(15)
Away from the boundary, the tangential deformation can be identified with a diffeomorphism of $`m`$, since $`\delta _{}S`$ is a boundary term. Let us examine explicitly how this occurs. We note that $`\delta _{}f=\mathrm{\Phi }^a_af`$ for any scalar function $`f(\xi )`$ defined on the worldvolume, $`m`$. In addition, under a tangential deformation, the induced metric on $`m`$ tranforms as a Lie derivative,
$$\delta _{}\gamma _{ab}=_a\mathrm{\Phi }_b+_b\mathrm{\Phi }_a.$$
(16)
Thus
$$\delta _{}\sqrt{\gamma }=\sqrt{\gamma }_a\mathrm{\Phi }^a.$$
(17)
For an arbitrary variation of the action, we have
$$\delta S=_m\left\{\left(\delta \sqrt{\gamma }\right)L+\sqrt{\gamma }\left(\delta L\right)\right\}.$$
(18)
A tangential deformation of the worldvolume thus always results in a pure divergence,
$$\delta _{}S=_m\sqrt{\gamma }_a\left(L\mathrm{\Phi }^a\right)=_m\sqrt{h}L\eta _a\mathrm{\Phi }^a,$$
(19)
where we have used Stokes theorem in the second equality, $`h`$ is the determinant of the metric $`h_{AB}`$ induced on $`m`$ by the embeddings described by Eq. (4),
$$h_{AB}=ϵ^\mu {}_{A}{}^{}ϵ_{}^{\nu }{}_{B}{}^{}\eta _{\mu \nu }^{},$$
(20)
and $`\eta ^a`$ is the unit timelike normal on $`m`$ pointing into $`m`$, i.e. the N-velocity of an observer sitting on the hypersurface $`\mathrm{\Sigma }`$. A diffeomorphism of $`m`$ can only move its boundary.
The action is stationary with respect to tangential deformations of the worldvolume with a vanishing normal component, $`\eta _a\mathrm{\Phi }^a=0`$ on $`\mathrm{\Sigma }_{(i)}`$ and $`\mathrm{\Sigma }_{(f)}`$. The remaining components of $`\mathrm{\Phi }^a`$ may range freely. These are precisely the components that generate a diffeomorphism of the spacelike configurations: the initial and final configurations are fixed, but not the coordinates chosen to describe these configurations.
Whereas the tangential variation of the action is simple, the normal variation is, in general, non-trivial. The normal deformation can always be cast in the form
$$\delta _{}S=_m\sqrt{\gamma }[_i(L)\mathrm{\Phi }^i+_a\mathrm{\Pi }^a{}_{i}{}^{}[\mathrm{\Phi }^i]],$$
(21)
i.e. as a worldvolume part, and a pure divergence. Here $`_i(L)`$ is the Euler-Lagrange derivative of $`L`$ projected onto the normals $`n^\mu _i`$ to the worldvolume; $`\mathrm{\Pi }^a_i`$ is a linear differential operator defined on $`m`$ which arises when the worldvolume gradients of the normal deformation in the worldvolume bulk are confined to a pure divergence, using integration by parts. The argument of the operator $`\mathrm{\Pi }^a_i`$ is indicated within the square bracket. We can use the divergence theorem in the second term, to obtain,
$$\delta _{}S=_m\sqrt{\gamma }_i(L)\mathrm{\Phi }^i+_m\sqrt{h}\eta _a\mathrm{\Pi }^a{}_{i}{}^{}[\mathrm{\Phi }^i].$$
(22)
When the classical equations of motion are satisfied,
$$_i(L)=0,$$
(23)
the action is stationary with respect to normal deformations of $`m`$. A well posed variational problem requires the vanishing of the boundary terms on the initial and final configurations. The variational principle restricted to normal deformations gives the classical dynamics.
So far, we have considered arbitrary deformations of the embedding functions. Let us now specialize to an infinitesimal Poincaré transformation,
$$\delta Y^\mu =ϵ^\mu +\omega ^\mu {}_{\nu }{}^{}Y_{}^{\nu },$$
(24)
where $`ϵ^\mu `$ is an infinitesimal constant translation, and an infinitesimal Lorentz transformation is given by $`\mathrm{\Lambda }^\mu {}_{\nu }{}^{}=\delta _\nu ^\mu +\omega ^\mu _\nu `$ , with $`\omega _{\mu \nu }=\omega _{\nu \mu }`$. We decompose $`\delta Y^\mu `$ according to Eq.(8). For an infinitesimal spacetime translation $`\delta Y^\mu =ϵ^\mu `$, we have
$$\mathrm{\Phi }_i=n^\mu {}_{i}{}^{}ϵ_{\mu }^{},\mathrm{\Phi }_a=e^\mu {}_{a}{}^{}ϵ_{\mu }^{}.$$
Substituting this into Eqs.(19) and (21), and summing, the variation of the action associated with a spacetime translation can be cast in the form,
$$\delta S=ϵ_\mu _m\left[\sqrt{\gamma }^i(L)n^\mu {}_{i}{}^{}+_a𝒫^{a\mu }\right].$$
(25)
The worldvolume vector density of weight one $`𝒫^{a\mu }`$ is given by
$$𝒫^{a\mu }=\sqrt{\gamma }(\mathrm{\Pi }^a{}_{i}{}^{}[n^{\mu i}]+Le^{\mu a}).$$
(26)
This expression for the variation of the action plays a central role in what follows. While tangential deformations do not participate in the variational derivation of the equations of motion, we see that they do contribute in an essential way to the construction of conserved quantities.
The total boundary contribution associated with a translation, using the divergence theorem, is
$$\delta S=ϵ_\mu _m\eta _a𝒫^{a\mu },$$
(27)
with the understanding that $`𝒫^{a\mu }`$ is to be evaluated at the boundary. This integral is well defined since $`\eta _a𝒫^{a\mu }`$ is a density of weight one when evaluated there.
Similarly, for an infinitesimal Lorentz transformation, we have
$$\mathrm{\Phi }_i=\omega _{\mu \nu }n^\mu {}_{i}{}^{}Y_{}^{\nu },\mathrm{\Phi }_a=\omega _{\mu \nu }e^\mu {}_{a}{}^{}Y_{}^{\nu }.$$
and in this case the variation of the action associated with a Lorentz transformation reduces to
$$\delta S=\omega _{\mu \nu }_m\left[\sqrt{\gamma }^i(L)n^\mu {}_{i}{}^{}Y_{}^{\nu }+_a^{a\mu \nu }\right],$$
(28)
where the worldvolume vector density of weight one $`^{a\mu \nu }`$ is given by
$$^{a\mu \nu }=\frac{1}{2}\sqrt{\gamma }[\mathrm{\Pi }^a{}_{i}{}^{}[n^{\mu i}Y^\nu ]+Le^{\mu a}Y^\nu (\mu \nu )].$$
(29)
The boundary contribution to $`\delta S`$ gives,
$$\delta S=\omega _{\mu \nu }_m\eta _a^{a\mu \nu }.$$
(30)
In our derivation of the expressions for $`𝒫^{a\mu }`$ and $`^{a\mu \nu }`$ it is important to emphasize that we did not enforce any boundary conditions on the induced variations. Indeed, it would be an error to attempt to enforce the boundary conditions which are appropriate for the variational derivation of the equations of motion.
Let us now suppose that the equations of motion $`_i(L)=0`$ are satisfied. We have
$$\delta S=ϵ_\mu \left[P^\mu (\mathrm{\Sigma }_{(f)})P^\mu (\mathrm{\Sigma }_{(i)})\right]+\omega _{\mu \nu }\left[M^{\mu \nu }(\mathrm{\Sigma }_{(f)})M^{\mu \nu }(\mathrm{\Sigma }_{(i)})\right].$$
(31)
The linear momentum $`P^\mu (\mathrm{\Sigma })`$ of the spatial hypersurfaces is defined by
$$P^\mu (\mathrm{\Sigma })=_\mathrm{\Sigma }\eta _a𝒫^{a\mu }.$$
(32)
We identify therefore $`𝒫^{a\mu }`$ as the linear momentum density. From Eq. (31) on we depart from our earlier convention, understanding the unit normal $`\eta ^a`$ to be future pointing when referring to a spacelike hypersurface on $`m`$.
We emphasize that $`P^\mu (\mathrm{\Sigma })`$ is a quantity associated with the spacelike hypersurface $`\mathrm{\Sigma }`$. In fact, it is possible to express it purely in terms of the geometry of $`\mathrm{\Sigma }`$ itself. We refrain from showing it explicitly, since the investment in additional formalism is not compensated by a corresponding gain in information. On the other hand, the density $`𝒫^{a\mu }`$ leads a double life: depending on circumstances, it lives either on the worldvolume $`m`$, or on its boundary $`m`$, as exemplified by Eqs. (25) and (27).
We define the angular momentum $`M^{\mu \nu }(\mathrm{\Sigma })`$ by
$$M^{\mu \nu }(\mathrm{\Sigma })=_\mathrm{\Sigma }\eta _a^{a\mu \nu },$$
(33)
which identifies $`^{a\mu \nu }`$ as the angular momentum density.
It is useful to express $`^{a\mu \nu }`$ in the alternative form
$$^{a\mu \nu }=\frac{1}{2}[𝒫^{a\mu }X^\nu +\pi ^a{}_{i}{}^{}[n^{\mu i}Y^\nu ](\mu \nu )],$$
(34)
where we introduce
$$\pi ^a{}_{i}{}^{}[n^{\mu i}Y^\nu ]=\sqrt{\gamma }(\mathrm{\Pi }^a{}_{i}{}^{}[n^{\mu i}Y^\nu ]n^{\mu i}\mathrm{\Pi }^a{}_{i}{}^{}[X^\nu ]).$$
(35)
The antisymmetric part of $`\pi ^a{}_{i}{}^{}[n^{i\mu }X^\nu ]`$ denotes that part of $`^{a\mu \nu }`$ which is not determined completely by the linear momentum density, $`𝒫^{a\mu }`$. As we show below, this is precisely the part that is interesting in higher derivative theories.
If the action is Poincaré invariant, so that $`\delta S=0`$, we have
$$P^\mu (\mathrm{\Sigma }_{(f)})=P^\mu (\mathrm{\Sigma }_{(i)})$$
and
$$M^{\mu \nu }(\mathrm{\Sigma }_{(f)})=M^{\mu \nu }(\mathrm{\Sigma }_{(i)}).$$
However, $`\mathrm{\Sigma }_{(i)}`$ and $`\mathrm{\Sigma }_{(f)}`$ are arbitrarily chosen initial and final configurations of the brane. Thus $`P^\mu (\mathrm{\Sigma })`$ and $`M^{\mu \nu }(\mathrm{\Sigma })`$ are both independent of the spacelike hypersurface $`\mathrm{\Sigma }`$. In this extremely broad sense, both $`P^\mu `$ and $`M^{\mu \nu }`$ are constants of the motion.
Our treatment so far has been entirely general. We have not indicated how to evaluate either the Euler-Lagrange derivatives of $`L`$, or the differential operator, $`\mathrm{\Pi }^a_i`$. We will consider some concrete examples, beginning with the most elementary, in Sect. 4.
## 3 Equations of motion from momentum conservation
Before we consider explicit examples, in this section we show how the equations of motion can be expressed in terms of the conservation of the linear momentum.
From the variation of the action under a spacetime translation, Eq. (25), when the action is invariant under translations, so that $`\delta S=0`$ on the left hand side, we have
$$\sqrt{\gamma }^i(L)n^\mu {}_{i}{}^{}=_a𝒫^{a\mu }.$$
(36)
This says that $`_a𝒫^{a\mu }`$ is normal to the worldvolume, and the equations of motion imply the conservation of the linear momentum density,
$$_a𝒫^{a\mu }=0,$$
(37)
and vice versa. That the equations of motions can be restated in terms of the conservation of linear momenta should not come as a surprise. In fact, it is a special case of the fact that the equations of motion can be expressed in terms of the conservation of the stress-energy tensor (for a relativistic string, see e.g. ). As we mentioned in the introduction, this is the approach adopted by Carter in his treatment of brane dynamics (see e.g. Ref. , and references therein.).
However, the form (37) is not the most useful expression of the conservation law, as it involves the mixed spacetime-worldvolume density, $`𝒫^{a\mu }`$; it does not isolate its non-trivial part. It is possible, however, to express the equations of motion in purely worldvolume terms. First, we decompose the spacetime vector density $`𝒫^{a\mu }`$ into its tangential and normal parts,
$$𝒫^{a\mu }=𝒫^{ab}e^\mu {}_{b}{}^{}+𝒫^{ai}n^\mu {}_{i}{}^{}.$$
(38)
Note that, in general, the worldvolume tensor field density $`𝒫^{ab}`$ will not be symmetric in its indices.
The worldvolume covariant divergence of $`𝒫^{a\mu }`$ gives
$$_a𝒫^{a\mu }=(_a𝒫^{ab}+K^b{}_{ai}{}^{}𝒫_{}^{ai})e^\mu {}_{b}{}^{}+(\stackrel{~}{}_a𝒫^{ai}K_{ab}{}_{}{}^{i}𝒫_{}^{ab})n^\mu {}_{i}{}^{},$$
(39)
where we have made use of the Gauss-Weingarten equations for the worldvolume $`m`$, (see e.g. ),
$`_ae^\mu _b`$ $`=`$ $`K_{ab}{}_{}{}^{i}n_{}^{\mu }{}_{i}{}^{},`$ (40)
$`\stackrel{~}{}_an^\mu _i`$ $`=`$ $`K_{abi}e^{\mu b}.`$ (41)
The worldvolume projections of the expression (36), using Eq. (39), are therefore given by
$`\stackrel{~}{}_a𝒫^{ai}K_{ab}{}_{}{}^{i}𝒫_{}^{ab}`$ $`=`$ $`\sqrt{\gamma }^i(L).`$ (42)
$`_a𝒫^{ab}+K^b{}_{ai}{}^{}𝒫_{}^{ai}`$ $`=`$ $`0,`$ (43)
The first equation expresses the Euler-Lagrange derivative in a divergence form, so that the equations of motion take the form
$$\stackrel{~}{}_a𝒫^{ai}K_{ab}{}_{}{}^{i}𝒫_{}^{ab}=0.$$
(44)
This equation is like a Gauss law for a $`O(ND)`$ Yang-Mills “electric field” $`𝒫^{ai}`$, with “source” $`K_{ab}{}_{}{}^{i}𝒫_{}^{ab}`$. Note that only the symmetric part of $`𝒫^{ab}`$ enters the equations of motion.
The Bianchi identity (43) is a non-obvious integrability condition required for the existence of solutions of (42). Here, also the antisymmetric part of $`𝒫^{ab}`$ contributes.
We will exploit these expressions to examine the different role played by the tangential and normal parts of the momentum density. Moreover, as shown below in Sect. 7, they allow for a novel approach to the linearization of the equations of motion.
We mentioned earlier that $`𝒫^{ab}`$ need not be symmetric. However, the conservation of angular momentum, $`_a^{a\mu \nu }=0`$, requires the anti-symmetric part to vanish.
To conclude this section, let us point out that, in the variational principle, one usually keeps volume terms; boundary terms, which contribute to the momenta, are thrown away. The approach via conservation of momenta to the equations of motion points to a complementary strategy: keep only boundary terms, since they also are sufficient to reconstruct the dynamics.
## 4 Dirac-Nambu-Goto action
In order to illustrate the general formalism developed in the previous sections, let us begin with the familiar case of a DNG brane. The DNG action for a relativistic extended object is
$$S_{(0)}=\mu _m\sqrt{\gamma },$$
(45)
where the constant $`\mu `$ is the brane tension. This is the simplest action one can write down for such an object. It depends only on the intrinsic geometry of the worldvolume.
The normal deformation of this action is given by
$$\delta _{}S_{(0)}=\mu _m\sqrt{\gamma }K_i\mathrm{\Phi }^i,$$
(46)
where we use the familiar expression relating the Lie derivative along the normals of the volume element of $`m`$ to its mean extrinsic curvature $`K^i=\gamma ^{ab}K_{ab}^i`$,
$$\delta _{}\sqrt{\gamma }=\sqrt{\gamma }K^i\mathrm{\Phi }_i.$$
(47)
The tangential deformation of this action is simply given by Eq. (19), with $`L=\mu `$.
In this geometrical language, the equations of motion are given by the vanishing of the mean extrinsic curvature,
$$\mu K^i=0.$$
(48)
The worldvolumes that extremize the DNG action are extremal timelike surfaces. This is a system of $`ND`$ second-order hyperbolic partial differential equations for the embedding functions, $`Y^\mu (\xi )`$. The appropriate boundary conditions in the variational principle are $`\eta _a\mathrm{\Phi }^a=0`$ on $`m`$. Both $`\mathrm{\Phi }^i`$ and $`\mathrm{\Phi }_{}^a=\mathrm{\Phi }^a\eta _b\mathrm{\Phi }^b\eta ^a`$ are arbitrary. We only need to fix the initial and final hypersurface geometries. Anything more is superfluous.
To bring the equations of motion into a more familiar form, using the Gauss-Weingarten equations (40), we have that
$$K^i=n_\mu ^i\mathrm{\Delta }Y^\mu ,$$
(49)
where $`\mathrm{\Delta }`$ is the worldvolume d’Alembert operator. The tangential projections of $`\mathrm{\Delta }Y^\mu `$ vanish identically. We can now peel Eq.(49), and its tangential counterpart to recover the familiar harmonicity condition,
$$\mu \mathrm{\Delta }Y^\mu =0.$$
(50)
Note that in this model there is no surface term arising from the normal variation, so that we have that the operator introduced in Eq. (21) vanishes identically, $`\mathrm{\Pi }^a{}_{i}{}^{}[\mathrm{\Phi }^i]=0`$. This is a feature which is unique to the DNG action.
The invariance of the DNG action under Poincaré transformations gives the linear momentum density
$$𝒫^{a\mu }=\mu \sqrt{\gamma }e^{\mu a},$$
(51)
so that the total momentum $`P^\mu (\mathrm{\Sigma })`$ is given by
$$P^\mu (\mathrm{\Sigma })=\mu _\mathrm{\Sigma }\sqrt{h}\eta ^\mu ,$$
(52)
where $`\eta ^\mu =\eta ^ae^\mu _a`$ is the unit velocity vector at a given point on $`\mathrm{\Sigma }`$. The momentum density $`𝒫^{a\mu }`$ is not only tangent to the worldsheet, it also lies parallel to the tangent vector, $`e^\mu _a`$. In this sense, extremal surfaces, like geodesics, are self-parallel.
The linear momentum of a DNG brane is defined directly in terms of initial data on the spacelike hypersurface $`\mathrm{\Sigma }`$, without explicit reference to the worldvolume that will be generated by these initial data. This expression for the linear momentum of a DNG brane generalizes the expression for a free relativistic massive particle, $`P^\mu =mU^\mu `$, with $`m`$ its mass, and $`U^\mu `$ its unit velocity.
The worldvolume projections of $`𝒫^{a\mu }`$ are, respectively,
$`𝒫^{ab}`$ $`=`$ $`\mu \sqrt{\gamma }\gamma ^{ab},`$ (53)
$`𝒫^{ai}`$ $`=`$ $`0.`$ (54)
The vanishing of the normal part was to be expected — the DNG equations of motion are of second order in the embedding functions. Moreover, the tangential projection is explicitly symmetric. Substituting these projections into the worldvolume projections of the linear momentum conservation equation, Eqs. (42) and (43), we find that the first reproduces the extremal dynamics, $`\mu K^i=0`$, whereas the second is satisfied identically. We also note that $`𝒫^{ab}`$ is scale invariant, under $`\gamma _{ab}\mathrm{\Omega }^2\gamma _{ab}`$, if and only if $`D=2`$, i.e. for a relativistic string. This is a consequence of the scale invariance of the DNG string, which becomes manifest in the Polyakov formulation of the theory .
The angular momentum density is simply
$$^{a\mu \nu }=𝒫^{a[\mu }X^{\nu ]}.$$
(55)
When the equations of motion hold, this is automatically conserved. The total angular momentum $`M^{\mu \nu }(\mathrm{\Sigma })`$ is given by
$$M^{\mu \nu }(\mathrm{\Sigma })=\frac{\mu }{2}_\mathrm{\Sigma }\sqrt{h}\{\eta ^\mu X^\nu (\mu \nu )\}.$$
(56)
Having covered the familiar case of a DNG brane, in the next section, we move on to a less trivial class of applications.
## 5 Extrinsic curvature actions
In this section, we consider the simplest actions involving extrinsic curvature, described by some lagrangian $`L=L(\gamma ^{ab},K_{ab}{}_{}{}^{i})`$. The more general case of a lagrangian that depends on derivatives of the extrinsic curvature as well can be treated along the same lines, and we will consider it in an Appendix. The possibilities are limited by the requirement that the lagrangian must transform as a scalar under diffeomorphisms of the worldvolume, and under rotations of the normals. In general, such theories will involve derivatives of the embedding functions higher than first, since the extrinsic curvature generalizes to a brane the acceleration of a point particle.
For a hypersurface, $`D=N1`$, a lagrangian proportional to (odd powers of) $`K`$ is admissible. For arbitrary co-dimension, however, the lowest order lagrangian invariant under normal rotations is quadratic in $`K_{ab}^i`$,
$$L=\mu +\alpha _1K^iK_i+\alpha _2K_{ab}{}_{}{}^{i}K_{}^{ab}{}_{i}{}^{},$$
(57)
where $`\alpha _1`$ and $`\alpha _2`$ are constants that measure the rigidity of the brane.
The worldvolume scalars $`K^iK_i`$ and $`K_{ab}{}_{}{}^{i}K_{}^{ab}_i`$ are not independent. The completely contracted Gauss-Codazzi equations in a flat spacetime background relates their difference to the worldvolume scalar curvature $``$, with
$$=K^iK_iK^{abi}K_{abi}.$$
(58)
When $`D=2`$, for a (closed) relativistic string, the action constructed from $``$ is a topological invariant, as follows from the Gauss-Bonnet theorem, so that the actions determined by the two quadratics are locally equivalent.
As before, the tangential variation of the action is straightforward, see Eq. (19). We exploit the chain rule to write down the normal variation of the lagrangian in terms of the variations of its arguments. We then have
$$\delta _{}S=_m\sqrt{\gamma }\{K^i\mathrm{\Phi }_iL+L^{ab}{}_{i}{}^{}(\stackrel{~}{\delta }_{}K_{ab}{}_{}{}^{i})+L_{ab}\left(\delta _{}\gamma ^{ab}\right)\},$$
(59)
where we have used Eq. (47) for the first term, and we have defined,
$`L^{ab}{}_{i}{}^{}={\displaystyle \frac{L}{K_{ab}^i}}=L^{ba}{}_{i}{}^{},`$ (60)
$`L_{ab}={\displaystyle \frac{L}{\gamma ^{ab}}}=L_{ba}.`$ (61)
We now exploit the results of Ref., specialized to a flat background, to express the normal variation of the inverse induced metric and the extrinsic curvature as,
$`\delta _{}\gamma ^{ab}`$ $`=`$ $`2K^{ab}{}_{i}{}^{}\mathrm{\Phi }_{}^{i},`$ (62)
$`\stackrel{~}{\delta }_{}K_{ab}^i`$ $`=`$ $`\stackrel{~}{}_a\stackrel{~}{}_b\mathrm{\Phi }^i+K_{ac}{}_{}{}^{i}K_{}^{c}{}_{bj}{}^{}\mathrm{\Phi }_{}^{j}.`$ (63)
The deformation operator $`\stackrel{~}{\delta }_{}`$ acting on $`K_{ab}^i`$ is constructed analogously to $`\stackrel{~}{}_a`$, and it involves a deformation connection $`\gamma _{ij}=\gamma _{ji}`$, so that when acting e.g. on a vector under normal rotations, it is defined by
$$\stackrel{~}{\delta }_{}A^i=\delta _{}A^i\gamma _i{}_{}{}^{j}A_{j}^{}.$$
(64)
This refinement is necessary to ensure covariance of the deformation under normal rotations, but coinciding with $`\delta _{}`$ when considering the deformation of normal rotation scalars. (For more detail the reader may refer to .)
We remove the hessian of $`\mathrm{\Phi }^i`$ appearing in the second term on the right hand side of Eq.(59), when we insert Eq. (63). The result is
$`\delta _{}S`$ $`=`$ $`{\displaystyle _m}\sqrt{\gamma }\{K^iL2K^{ab}{}_{}{}^{i}L_{ab}^{}\stackrel{~}{}_a\stackrel{~}{}_bL^{ab}{}_{i}{}^{}K_{aci}K_b{}_{}{}^{cj}L_{}^{ab}{}_{j}{}^{}\}\mathrm{\Phi }^i`$ (65)
$`+`$ $`{\displaystyle _m}\sqrt{\gamma }_a\{L^{ab}{}_{i}{}^{}\stackrel{~}{}_{b}^{}\mathrm{\Phi }^i+(\stackrel{~}{}_bL^{ab}{}_{i}{}^{})\mathrm{\Phi }^i\}.`$
We identify the Euler-Lagrange derivative as
$$_i(L)=\stackrel{~}{}_a\stackrel{~}{}_bL^{ab}{}_{i}{}^{}+L^{ab}{}_{j}{}^{}K_{ac}^{}{}_{}{}^{j}K_{}^{c}{}_{bi}{}^{}+LK_i2L_{ab}K^{ab}{}_{i}{}^{}.$$
(66)
Generically, the Euler-Lagrange equations $`^i=0`$ are of second order in derivatives of $`K_{ab}^i`$, so they are of fourth order in derivatives of the embedding functions $`Y^\mu `$.
Appropriate boundary conditions in the variational principle are that $`\mathrm{\Phi }^i=0`$, in order to cancel the last term in Eq. (65), and, since this implies already the vanishing of the derivative of $`\mathrm{\Phi }^i`$ along $`\mathrm{\Sigma }`$, in order to cancel the next to last term in Eq. (65), we need only to require $`\eta ^a\stackrel{~}{}_a\mathrm{\Phi }^i=0`$ on $`m`$, independently of $`L^{ab}_i`$.
We identify the operator $`\mathrm{\Pi }^a_i`$ introduced in Eq.(21) as the covariant ‘wronskian’:
$$\mathrm{\Pi }^a{}_{i}{}^{}[\mathrm{\Phi }^i]=L^{ab}{}_{i}{}^{}\stackrel{~}{}_{b}^{}\mathrm{\Phi }^i+(\stackrel{~}{}_bL^{ab}{}_{i}{}^{})\mathrm{\Phi }^i.$$
(67)
In particular, for a $`\mathrm{\Phi }^i`$ which corresponds to a background translation,
$$ϵ_\mu \mathrm{\Pi }^a{}_{i}{}^{}[n^{\mu i}]=ϵ_\mu [L^{ab}{}_{i}{}^{}K_{b}^{}{}_{}{}^{ci}e_{}^{\mu }{}_{c}{}^{}+(\stackrel{~}{}_bL^{ab}{}_{i}{}^{})n^{\mu i}].$$
(68)
where we exploit the Gauss-Weingarten equation (41) to simplify the first term. We thus have from Eq. (26), the general expression for the linear momentum density,
$$𝒫^{a\mu }=\sqrt{\gamma }\left[(L\gamma ^{ab}L_i^{ac}K^b{}_{ci}{}^{})e^\mu {}_{b}{}^{}+(\stackrel{~}{}_bL^{ab}{}_{i}{}^{})n^{\mu i}\right].$$
(69)
Unlike the DNG case, in general, the momentum density $`𝒫^{a\mu }`$ now possesses a component normal to the worldvolume.
We can write the tangential part in a different way. First we separate the quantity $`L_i^{ac}K^b_{ci}`$ in its symmetric and anti-symmetric parts. It is easy to show that the symmetric part is given by
$$L^{ca}{}_{i}{}^{}K_{}^{b}{}_{c}{}^{}{}_{}{}^{i}+L^{cb}{}_{i}{}^{}K_{}^{a}{}_{c}{}^{}{}_{}{}^{i}=2L^{ab}.$$
(70)
Then, we can rewrite the tangential part of $`𝒫^{a\mu }`$ in the form
$$𝒫^{ab}=\sqrt{\gamma }\left(L\gamma ^{ab}L^{ab}\right)+𝒬^{ab},$$
(71)
where we isolate the anti-symmetric part of $`𝒫^{ab}`$,
$$𝒬^{ab}=\sqrt{\gamma }(L^{ca}{}_{i}{}^{}K_{}^{b}{}_{ci}{}^{}L^{cb}{}_{i}{}^{}K_{}^{a}{}_{ci}{}^{})=0,$$
(72)
which vanishes identically for the geometrical actions we consider.
In order to check in a non-trivial case that the equations of motion can be expressed in the form (44), we can substitute the normal and tangential projections into the left hand side of Eq.(42), we obtain the equations of motion in the form (66). Moreover, we can check that Eq.(43), corresponding to a worldvolume translation, is in fact an identity.
Inspection of the form (42) of the equations of motion, or directly of Eq. (66), shows that the necessary and sufficient condition for the equations of motion to be of second order in derivatives of the embedding functions is simply that the normal component of the linear momentum density vanishes, $`𝒫^{ai}=0`$.
The alternative expression (71) for the tangential part of the linear momentum density is also useful for recovering the result that invariance of the action under scale transformations implies
$$𝒫^{ab}\gamma _{ab}=0.$$
(73)
To see this, consider that, using Eq. (18), the variation of the action under a change of the worldvolume metric is simply
$`\delta S`$ $`=`$ $`{\displaystyle _m}\sqrt{\gamma }\left[L\gamma ^{ab}L^{ab}\right]\delta \gamma _{ab}`$ (74)
$`=`$ $`{\displaystyle _m}𝒫^{ab}\delta \gamma _{ab},`$
The second line follows from the fact that only the symmetric part of $`𝒫^{ab}`$ enters. This identifies the symmetric part of $`𝒫^{ab}`$ as the worldvolume stress-energy tensor. When the induced metric undergoes a scale transformation, scale invariance of the action, $`\delta S=0`$, implies the tracelessness condition (73), and viceversa.
Let us now move on to the angular momentum. The operator appearing in Eq.(35) is
$$\pi ^a{}_{i}{}^{}[n^{i\mu }X^\nu ]=L^{ab}{}_{i}{}^{}n_{}^{\mu i}e^\nu {}_{b}{}^{},$$
(75)
so that the angular momentum density is given by
$$^{a\mu \nu }=\frac{1}{2}[𝒫^{a\mu }X^\nu \sqrt{\gamma }L^{ab}{}_{i}{}^{}n_{}^{\mu i}e^\nu {}_{b}{}^{}(\mu \nu )].$$
(76)
which does not involve derivatives in $`K_{ab}^i`$ other than those already contained in $`𝒫^{a\mu }`$. The second term may be thought of as an effect of the finite width of the worldvolume, when we go beyond the DNG approximation. Its appearance is necessary to ensure conservation of angular momentum. Neither term alone is conserved. To see this, let us assume that the equations of motion hold or, equivalently, that the linear momentum is conserved, $`_a𝒫^{a\mu }=0`$. The divergence of the angular momentum density is
$`_a^{a\mu \nu }={\displaystyle \frac{1}{2}}[𝒫^{a\mu }e^\nu _a`$ $``$ $`\sqrt{\gamma }(\stackrel{~}{}_aL^{ab}{}_{i}{}^{})n^{\mu i}e^\nu {}_{b}{}^{}\sqrt{\gamma }L^{ab}{}_{i}{}^{}K_{a}^{}{}_{}{}^{ci}e_{}^{\mu }{}_{c}{}^{}e_{}^{\nu }_b`$ (77)
$`+`$ $`\sqrt{\gamma }L^{ab}{}_{i}{}^{}K_{abj}^{}n^{\mu i}n^{\nu j}(\mu \nu )],`$
where we have used the Gauss-Weingarten equations (40), (41). In this expression, all of the possible spacetime bivectors appear. However, if we express $`𝒫^{a\mu }`$ in terms of its projections, we obtain
$$_a^{a\mu \nu }=\frac{1}{2}[2𝒬^{ab}e^\mu {}_{a}{}^{}e_{}^{\nu }{}_{b}{}^{}\sqrt{\gamma }L^{ab}{}_{i}{}^{}K_{abj}^{}n^{\mu i}n^{\nu j}(\mu \nu ))].$$
(78)
The two terms are independent, the first is proportional to a bivector parallel to the worlvolume, the second to a bivector normal to it. Therefore the conditions necessary for conservation of angular momentum are
$`𝒬^{ab}`$ $`=`$ $`0,`$ (79)
$`L^{ab}{}_{i}{}^{}K_{abj}^{}L^{ab}{}_{j}{}^{}K_{abi}^{}`$ $`=`$ $`0.`$ (80)
The first condition was to be expected from classical elasticity theory , and, as mentioned above, is satisfied identically. The second condition is a new “thickness” effect associated with co-dimension $`ND>1`$, and it is also satisfied identically for the geometrical actions we consider.
We also note that at this order, $`^{a\mu \nu }`$ does not involve any term proportional to $`n^{\mu i}n^{\nu j}(\mu \nu )`$, which is permitted if $`ND>1`$. Such a term will show up in higher order theories, as we show in Appendix A.
## 6 Extrinsic curvature actions: examples
In order to make our discussion more concrete, in this section we consider some specific examples of the class of theories treated in the previous section. We begin with the simple case of a hypersurface action of the form
$$S_{(1)}=\alpha _0_m\sqrt{\gamma }K.$$
(81)
Since $`L_{ab}=\alpha _0K_{ab}`$ and $`L^{ab}{}_{}{}^{}=\alpha _0\gamma ^{ab}`$ (we use the symbol $``$ to denote the only normal direction), from Eq. (66), we find the Euler-Lagrange derivative in the form
$$_{}=\alpha _0(K^2K^{ab}K_{ab})=\alpha _0,$$
(82)
where we have used the contracted Gauss-Codazzi equation (58). This action is extremized by worldvolumes with vanishing scalar Ricci curvature. It is a topological invariant (the winding number) for a pointlike trajectory. The equations of motion, like the DNG case, are of second order in time derivatives of the embedding functions. Indeed, the total linear momentum is given by
$$P^\mu (\mathrm{\Sigma })=\alpha _0_\mathrm{\Sigma }\sqrt{h}\eta _a(K\gamma ^{ab}K^{ab})e^\mu {}_{b}{}^{},$$
(83)
and we see that the linear momentum density is purely tangential, and manifestly symmetric. Moreover, reading off the quantity $`𝒫^{ab}`$, and substituting in the equations of motion in the form (43), reproduces $`_{}=0`$, with the Euler-Lagrange derivative given by Eq. (82).
Although the addition of this action to the DNG action does not change the order of the equations of motion, it does change the boundary conditions. Now we need to require $`\eta ^a\stackrel{~}{}_a\mathrm{\Phi }^i=0`$. We will discuss this issue in more detail below, in the special case of a Einstein-Hilbert brane action.
The total angular momentum is
$$M^{\mu \nu }(\mathrm{\Sigma })=\frac{1}{2}_\mathrm{\Sigma }\eta _a[𝒫^{a\mu }X^\nu \alpha _0\sqrt{h}n^\mu e^{\nu a}(\mu \nu )].$$
(84)
The second term is proportional to the normal bivector $`n^\mu \eta ^\nu n^\nu \eta ^\mu `$. In the equation for angular momentum conservation, $`_a^{a\mu \nu }=0`$, both terms are conserved separately. The first because of the vanishing of $`𝒫^{ai}`$ and the symmetry of $`𝒫^{ab}`$. The second is conserved automatically, because the condition (80) is vacuus. Alternatively, this can be checked directly using the Gauss-Weingarten equations (40), (41), specialized to the case of a hypersurface.
Let us now examine the case of an action quadratic in the extrinsic curvature. In particular, let us consider first,
$$S_{(2,a)}=\alpha _1_m\sqrt{\gamma }K^iK_i.$$
(85)
We have $`L_{ab}=2\alpha _1K^iK_{abi}`$, and $`L^{ab}{}_{i}{}^{}=2\alpha _1K_i\gamma ^{ab}`$. Thus from Eq. (66), the Euler -Lagrangian derivative is
$$_i=2\alpha _1\left[\stackrel{~}{\mathrm{\Delta }}K_i+K^{ab}{}_{i}{}^{}K_{ab}^{}{}_{}{}^{j}K_{j}^{}\frac{1}{2}K^jK_jK_i\right].$$
(86)
For a quadratic action, the Euler-Lagrange derivative is proportional to $`\stackrel{~}{\mathrm{\Delta }}K_i`$ plus some cubic in $`K_{ab}^i`$. We note that the extremal solutions $`K^i=0`$ continue to be solutions of this particular theory .
The linear momentum density takes the form
$$𝒫^{a\mu }=\alpha _1\sqrt{\gamma }\left\{K_i(K^i\gamma ^{ab}2K^{abi})e^\mu {}_{b}{}^{}+2\left(\stackrel{~}{}_aK_i\right)n^{\mu i}\right\}.$$
(87)
Thus, on all extremal solutions of the theory defined by $`S=S_{(0)}+S_{(2,a)}`$, with $`K^i=0`$, the rigidity makes no contribution to the momentum. If the complete action is given by the rigidity term, solutions with $`K^i=0`$ carry no momentum. This was noted by Boisseau and Letelier in for the special case of a relativistic string.
The scale invariance of this action in the case of a relativistic string ($`D=2`$) can be checked by verifying that $`𝒫^{ab}\gamma _{ab}=0`$.
The normal projection of the linear momentum density is
$$𝒫^{ai}=2\alpha _1\sqrt{\gamma }\stackrel{~}{}^aK^i.$$
(88)
Therefore, if the mean curvature is constant, $`K^i=`$ const., so that $`𝒫^{ai}=0`$, then the equations of motion are of second order in derivatives of the embedding functions.
The angular momentum density is given by
$$^{a\mu \nu }=\frac{1}{2}[𝒫^{a\mu }X^\nu 2\sqrt{\gamma }K_in^{\mu i}e^{\nu a}(\mu \nu )].$$
(89)
The second term is proportional to the bivector normal to $`\mathrm{\Sigma }`$. Both conditions (79), (80) are identically satisfied, so that angular momentum is conserved.
The total angular momentum $`M^{\mu \nu }`$ also vanishes when $`K^i=0`$. For a DNG string, the well-known Regge inequality bounds the spin by the mass. We conclude that for a rigid string, the inequality continues to hold for all DNG solutions .
We have seen already that for strings the two quadratics are not independent so that the $`K^iK_i`$ theory is, in fact, the unique theory quadratic in extrinsic curvature. In general, for arbitrary $`D`$, however we also have the other possibility,
$$S_{(2,b)}=\alpha _2_m\sqrt{\gamma }K_{ab}{}_{}{}^{i}K_{}^{ab}{}_{i}{}^{}.$$
(90)
We find easily that $`L_{ab}=2\alpha _2K^c{}_{ai}{}^{}K_{bc}^{}^i`$ and $`L^{ab}{}_{i}{}^{}=2\alpha _2K^{ab}_i`$. From Eq. (66), the Euler-Lagrange derivative is
$$_i=2\alpha _2\left(\stackrel{~}{}_a\stackrel{~}{}_bK^{ab}{}_{i}{}^{}+K^{ac}{}_{j}{}^{}K_{cb}^{j}K_a{}_{}{}^{b}{}_{i}{}^{}\frac{1}{2}K_{ab}{}_{}{}^{j}K_{}^{ab}{}_{j}{}^{}K_{i}^{}\right).$$
(91)
This expression for the Euler-Lagrange derivative is perfectly legitimate. On the other hand, we can offer an alternative expression, closer to the one given for the Euler-Lagrange derivative of the other quadratic action, Eq. (86). We use the once contracted Gauss-Codazzi equation,
$$_{ab}K^iK_{abi}+K_{ac}{}_{}{}^{i}K_{}^{c}{}_{bi}{}^{}=0,$$
(92)
and the contracted Codazzi-Mainardi integrability condition,
$$\stackrel{~}{}_a(K^{abi}\gamma ^{ab}K^i)=0,$$
(93)
to reduce the Euler-Lagrange derivative (91) to the form
$$_i=2\alpha _1(\stackrel{~}{\mathrm{\Delta }}K_i+K^{ab}{}_{i}{}^{}K_{ab}^{}{}_{}{}^{j}K_{j}^{}\frac{1}{2}K^jK_jK_i𝒢_{ab}K^{ab}{}_{i}{}^{}),$$
(94)
where $`𝒢_{ab}=_{ab}(1/2)\gamma _{ab}`$ is the worldvolume Einstein tensor. This alternative expression is identical to Eq. (86), except for the addition of the last term, which vanishes identically for a string, $`D=2`$.
The linear momentum density for this theory is
$$𝒫^{a\mu }=\alpha _2\sqrt{\gamma }[(K_{cd}{}_{}{}^{i}K_{}^{cd}{}_{i}{}^{}\gamma _{}^{ab}2K^a{}_{ci}{}^{}K_{}^{bci})e^\mu {}_{b}{}^{}+2(\stackrel{~}{}_bK^{abi})n^\mu {}_{i}{}^{}].$$
(95)
We can use the Codazzi-Mainardi equation (93) to simplify the second term, and the contractions of the Gauss-Codazzi equation (58), (92), to write the linear momentum density in the alternative form,
$$𝒫^{a\mu }=\alpha _2\sqrt{\gamma }[(K_iK^i\gamma ^{ab}2K_iK^{abi}2𝒢^{ab})e^\mu {}_{b}{}^{}+2(\stackrel{~}{}^aK^i)n^\mu {}_{i}{}^{}].$$
(96)
We note that also for this theory, for constant mean curvature solutions, $`K^i=`$ const., the equations of motion are of second order in derivatives of the embedding functions. Moreover, for $`D=2`$, the theory is scale invariant, since $`𝒢_{ab}=0`$, identically, so that $`𝒫^{ab}\gamma _{ab}=0`$.
The total angular momentum is
$$M^{\mu \nu }(\mathrm{\Sigma })=\frac{1}{2}_m\eta _a[𝒫^{a\mu }X^\nu +2\alpha _2\sqrt{h}K^{abi}n^\mu {}_{i}{}^{}e_{}^{\nu }{}_{b}{}^{}(\mu \nu )].$$
(97)
Whereas in the cases considered previously the second term, independent of the linear momentum density, involved only the bivector normal to $`\mathrm{\Sigma }`$, here it includes also the bivector normal-tangential to $`\mathrm{\Sigma }`$.
As a last example of an action quadratic in the extrinsic curvature, we consider now the Einstein-Hilbert action,
$$S_{EH}=\beta _m\sqrt{\gamma }.$$
(98)
We have,
$`L_{ab}`$ $`=`$ $`2\beta (K_iK_{ab}{}_{}{}^{i}K_{aci}K_b{}_{}{}^{ci})=2\beta _{ab},`$
$`L^{ab}_i`$ $`=`$ $`2\beta (K_i\gamma ^{ab}K^{ab}{}_{i}{}^{}).`$
Note that the Codazzi-Mainardi equation (93) implies
$$\stackrel{~}{}_bL^{ab}{}_{i}{}^{}=0.$$
This kills the boundary term proportional to $`\mathrm{\Phi }_i`$.
To calculate the variation of this action, we can use the contracted Gauss-Codazzi equation for a flat space-time background, Eq.(58), and the results of the previous section. We obtain,
$`\delta S_{EH}`$ $`=`$ $`\beta {\displaystyle _m}\sqrt{\gamma }\left\{2K^{abi}𝒢_{ab}\mathrm{\Phi }_i\right\}`$ (99)
$`+`$ $`\beta {\displaystyle _m}\sqrt{h}\eta _a\left\{2\left[K^{abi}\gamma ^{ab}K^i\right]\stackrel{~}{}_b\mathrm{\Phi }_i+\mathrm{\Phi }^a\right\},`$
where $`𝒢_{ab}`$ is the worldvolume Einstein tensor, $`𝒢_{ab}=_{ab}\frac{1}{2}\gamma _{ab}`$. It follows at once that the Euler-Lagrange derivative of the Einstein-Hilbert action is given by
$$_i=2\beta K^{ab}{}_{i}{}^{}𝒢_{ab}^{}.$$
(100)
When $`D=2`$, for a string, the Einstein tensor vanishes identically, $`𝒢_{ab}=0`$ and $`_i=0`$ identically. This can also be checked by considering the difference of the Euler-Lagrange derivatives for the two actions quadratic in the curvature, Eqs. (86), (94).
In general, a DNG action with an Einstein-Hilbert correction satisfies the Euler-Lagrange equations,
$$K^{ab}{}_{i}{}^{}(\mu \gamma _{ab}+2\beta 𝒢_{ab})=0.$$
(101)
In particular, any embedded Einstein manifold, satisfying $`𝒢_{ab}=(\mu /2\beta )\gamma _{ab}`$ is a solution.
We note that an Einstein-Hilbert addition to the DNG action does not change the order of the equations of motion. They remain second order. The action, however, is no longer stationary under the same boundary conditions as the DNG theory, even when $`D=2`$. We now need to require that $`\eta ^a\stackrel{~}{}_a\mathrm{\Phi }^i=0`$. Technically, we trace this necessity to the fact that the action contains a term linear in the second derivative of the induced metric. In general, this spells trouble in a functional integral approach to the quantum theory. The resolution is, of course, well known: introduce a surface correction to the action to cancel these offending surface terms. The essence of the variational problem is encountered in the elementary description of a free particle by the lagrangian, $`L(x,\ddot{x})=(1/2)x\ddot{x}`$ which is equivalent to the lagrangian $`(1/2)\dot{x}^2`$ with the addition of a boundary term to the corresponding action. This term is simply $`(1/2)x\dot{x}`$. To do this in a covariant way for the problem at hand is a little harder. We are guided by the analogous problem in general relativity. This issue is treated in Sect. 8.
Here we examine the conserved quantitied associated with the unadorned Einstein-Hilbert action. With or without surface modifications, Noether’s theorem applies so long as the action possesses the required invariance. Noether’s theorem is independent of the particular boundary conditions that are invoked to determine the Euler-Lagrange equations.
The Einstein-Hilbert contribution to the momentum density is given by
$`𝒫^{a\mu }`$ $`=`$ $`\beta \sqrt{\gamma }[2\left(K_{ab}{}_{}{}^{i}\gamma _{ab}K^i\right)K_i^{bc}+]e^\mu _c`$ (102)
$`=`$ $`2\beta \sqrt{\gamma }𝒢^{ab}e^\mu {}_{b}{}^{},`$
where we have exploited the once contracted Gauss-Codazzi relation, to obtain the second line. We note that the vanishing of this momentum density gives the vacuum Einstein field equations for the worldvolume. If one is interested in an embedding formulation of GR in a flat background spacetime, this observation isolates the relevant subspace of the space of solutions, solving the problem posed by Regge and Teitelboim in . We plan to develop the consequences of this observation elsewhere.
The angular momentum density $`^{a\mu \nu }`$ decomposes into a sum of two separately conserved quantities:
$$^{a\mu \nu }=\frac{1}{2}[𝒫^{a\mu }X^\nu +2\beta \sqrt{\gamma }(K^{abi}\gamma ^{ab}K^i)n^\mu {}_{i}{}^{}e_{}^{\nu }{}_{b}{}^{}(\mu \nu )].$$
(103)
Again, both conditions (79) and (80) are identically satisfied, so the total angular momentum is conserved.
In the case of a string with worldvolume dimension $`D=2`$, the Einstein-Hilbert lagrangian is a topological invariant and does not contribute to the equations of motion. We would hope that the corresponding contributions to the conserved quantities also vanish. As expected $`𝒫^{a\mu }=0`$; however, whereas the first term vanishes in $`^{a\mu \nu }`$, the second does not. However, as we have seen, this non-vanishing part is conserved kinematically, it is divergence free off shell. One is always at liberty to add such a tensor to produce a ‘new’ conservation law. The addition of the appropriate boundary term, as shown below in Sect. 8, removes this extra kinematic term.
## 7 Linearized equations of motion from momenta
In this section, we consider the second order variation of the action, and we show how the linearized equations of motion can be expressed in terms of the normal variation of the linear momentum density $`𝒫^{a\mu }`$. It turns out that this approach to the linearized equations of motion is more economical than a direct approach, which considers the variation of the Euler-Lagrange derivative, $`\delta ^i(L)`$ . The virtue of this approach is that it involves the variation of geometrical quantities of lower order in derivatives of the embedding functions.
In Sect. 3, we have shown explicitly how to express the equations of motion in terms of the conservation of the appropriate projections of the momenta. The key equation we used, for an infinitesimal translation, is Eq. (25). We consider only its normal variation for the moment. We obtain
$$\delta _{}(\delta S)=ϵ_\mu \{_m\sqrt{\gamma }\left[K^j\mathrm{\Phi }_j^in^\mu {}_{i}{}^{}+(\stackrel{~}{\delta }_{}^i)n^\mu {}_{i}{}^{}^i(\stackrel{~}{}_a\mathrm{\Phi }_i)e^{\mu a}\right]+_m_a\delta _{}𝒫^{a\mu }\}.$$
(104)
Here we have used Eq. (47), and the second normal deformation Gauss-Weingarten equation ,
$$\stackrel{~}{\delta }_{}n^\mu {}_{i}{}^{}=(\stackrel{~}{}_a\mathrm{\Phi }^i)e^{\mu a}.$$
(105)
In the last term, we have exploited the fact that, since $`𝒫^{a\mu }`$ is a vector density of weight one, variation and covariant differentiation commute, $`\delta _a𝒫^{a\mu }=_a\delta 𝒫^{a\mu }`$.
When the action is invariant under translations, $`\delta S=0`$, this equation relates the normal variation of the Euler-Lagrange derivative to the divergence of the normal variation of the linear momentum density.
Following the same strategy we used for the equations of motion, let us decompose $`𝒫^{a\mu }`$ into its worldvolume projections,
$$\delta _{}𝒫^{a\mu }=[\delta _{}𝒫]^{ab}e^\mu {}_{b}{}^{}+[\delta _{}𝒫]^{ai}n^\mu {}_{i}{}^{},$$
(106)
where we use the brackets notation in order to distinguish the projection of the variation from the variation of the projection, since in general they are different, e.g. $`[\delta _{}𝒫]^{ab}\delta _{}𝒫^{ab}`$. In fact, this is the reason that makes the approach described here more efficient than a direct approach.
A straightforward calculation gives that when the action is invariant under translations, so that $`\delta S=0`$ on the left hand side of Eq. (104), the worldvolume projections of Eq. (104) are
$`\sqrt{\gamma }(\stackrel{~}{\delta }_{}^i)`$ $`=`$ $`\sqrt{\gamma }K^j\mathrm{\Phi }_j^i\stackrel{~}{}_a[\delta _{}𝒫]^{ai}+[\delta _{}𝒫]^{ab}K_{ab}{}_{}{}^{i},`$ (107)
$`\sqrt{\gamma }^i(\stackrel{~}{}^b\mathrm{\Phi }_i)`$ $`=`$ $`\stackrel{~}{}_a[\delta _{}𝒫]^{ab}+[\delta _{}𝒫]^{ai}K_a{}_{}{}^{b}{}_{i}{}^{}.`$ (108)
This latter equation is merely an identity, as expected from reparametrization invariance. The first equation can be used to express the linearized equations of motion, about a solution of the equations of motion, in the form
$$\stackrel{~}{}_a[\delta _{}𝒫]^{ai}[\delta _{}𝒫]^{ab}K_{ab}{}_{}{}^{i}=0.$$
(109)
Note the apparent similarity with the equations of motion expressed in the form (44). The usefulness of this expression is that in general, the variation of the momentum density is easier to calculate than the variation of the Euler-Lagrange derivative.
The variation of the total linear momentum is then simply
$$\delta _{}P^\mu (\mathrm{\Sigma })=_m\eta _a\delta _{}𝒫^{a\mu }.$$
(110)
Again, it should be emphasized that this is much easier than a direct variation of $`P^\mu (\mathrm{\Sigma })`$.
In order to illustrate this approach, and to give a simple application, let us consider a DNG object. The linear momentum density is given by Eq. (51). We can use Eqs. (62), (47), together with the first Gauss-Weingarten deformation equation (see Ref. )
$$\delta _{}e^\mu {}_{a}{}^{}=K_{ab}{}_{}{}^{i}\mathrm{\Phi }_{i}^{}e^{\mu b}+(\stackrel{~}{}_a\mathrm{\Phi }^i)n^\mu {}_{i}{}^{},$$
(111)
so that
$$\delta _{}𝒫^{a\mu }=\mu \sqrt{\gamma }[(K_i\gamma ^{ab}K^{ab}{}_{i}{}^{})\mathrm{\Phi }^ie^\mu {}_{b}{}^{}+(\stackrel{~}{}^a\mathrm{\Phi }^i)n^\mu {}_{i}{}^{}].$$
(112)
We read off the projections,
$$[\delta _{}𝒫]^{ab}=\mu \sqrt{\gamma }(K_i\gamma ^{ab}K^{ab}{}_{i}{}^{})\mathrm{\Phi }^i,$$
(113)
$$[\delta _{}𝒫]^{ai}=\mu \sqrt{\gamma }\stackrel{~}{}^a\mathrm{\Phi }^i.$$
(114)
Recalling that, for a DNG object, $`_i=\mu K_i`$, and substituting in Eq. (107), we find
$`(\stackrel{~}{\delta }_{}K^i)`$ $`=`$ $`K^j\mathrm{\Phi }_jK^i\stackrel{~}{}_a(\stackrel{~}{}^a\mathrm{\Phi }^i)+(K_j\gamma ^{ab}K^{ab}{}_{j}{}^{})\mathrm{\Phi }^jK_{ab}^i`$ (115)
$`=`$ $`\stackrel{~}{}_a(\stackrel{~}{}^a\mathrm{\Phi }^i)K_{ab}{}_{}{}^{i}K_{}^{ab}{}_{j}{}^{}\mathrm{\Phi }_{}^{j},`$
which agrees with the expression derived e.g. in Ref. (see also ).
Let us now confirm that Eq. (108) is in fact an identity. Substitution of the projections gives
$$K^i(\stackrel{~}{}^b\mathrm{\Phi }_i)=\stackrel{~}{}_a[(K_i\gamma ^{ab}K^{ab}{}_{i}{}^{})\mathrm{\Phi }^i]+(\stackrel{~}{}^a\mathrm{\Phi }^i)K_a{}_{}{}^{b}{}_{i}{}^{},$$
and this expression is identically zero, as follows from the Codazzi-Mainardi integrability condition, Eq. (93).
This example renders transparent the advantages of this approach. Rather than dealing with the variation of the extrinsic curvature, all one needs here is the variation of the intrinsic geometry of the worldvolume.
What about the second variation parallel to the worldvolume, $`\delta _{}(\delta S)`$? As expected from reparametrization invariance, it adds nothing new. We show this explicitly in Appendix B for a DNG brane.
## 8 Einstein-Hilbert action with surface term
In this section, we consider the addition of a surface term to the Einstein-Hilbert brane action we have briefly described at the end of Sect. 6,
$$S=S_{EH}+S_m,$$
(116)
where the first term is defined in Eq. (98), and
$$S_m=2\beta _m\sqrt{h}\kappa ,$$
(117)
and $`\kappa `$ is the trace of the extrinsic curvature $`\kappa _{AB}`$ of the spacelike boundary provided by $`\mathrm{\Sigma }_{(i)}`$ and $`\mathrm{\Sigma }_{(f)}`$ embedded in $`m`$.
An analogous term is added to the Einstein-Hilbert action in general relativity in order that the variational principle applied to the action yield the Einstein equations in a bounded region subject to the boundary condition that the metric induced on the boundary is fixed and no more . Technically, the variation of the boundary term precisely cancels normal derivatives of the variation of the metric tensor on the boundary, which occur in the variation of the Einstein Hilbert action. In general relativity this term is diffeomorphism invariant. Here, it is also Poincaré invariant. Interestingly enough, the particle analogue mentioned in Sect. 6 is incomplete in this respect. Whereas the original lagrangian $`L_0(x,\dot{x})=(1/2)\dot{x}^2`$ is invariant under translations, the lagrangian $`L(x,\ddot{x})=(1/2)x\ddot{x}`$ is not and so Noether’s theorem cannot be applied directly to it. The boundary term is necessary to restore the translation invariance of the problem.
The variation of Eq.(117) poses new technical difficulties. In general relativity, the dynamical variables to be varied are the spacetime metric coefficients. In the present context, we need to vary the embedding functions describing the worldvolume. Fortunately, the relevant formalism has been developed in Refs. , for the case of a timelike boundary. Its adaptation to the case of interest here requires only some minor sign modifications.
The problem is simplified by treating the boundary of the worldvolume as two embedded spacelike surfaces in the background Minkowski space-time described by Eq. (4). We denote the normals to $`m`$ in spacetime by $`m^\mu _I`$ ($`I,J,\mathrm{}=0,1,2,\mathrm{},ND`$). It is also convenient to exploit a normal basis which is adapted to the worldvolume $`m`$: we choose the basis $`m^{\mu I}=\{\eta ^\mu ,n^\mu {}_{i}{}^{}\}`$, supplementing the normals $`n^\mu _i`$ to $`m`$ in spacetime with $`\eta ^\mu =e^\mu {}_{a}{}^{}\eta _{}^{a}`$, the normal to the boundary which is tangent to the worldvolume. Let $`L_{AB}^I`$ represent the extrinsic curvatures associated with the embedding $`X`$. With respect to the adapted basis, $`L_{AB}{}_{}{}^{i}=K_{ab}{}_{}{}^{i}ϵ_{}^{a}{}_{A}{}^{}ϵ_{}^{b}{}_{B}{}^{}=K_{AB}^i`$, and $`L_{AB}{}_{}{}^{0}=\kappa _{AB}`$. (For details see .
Let us examine an arbitrary deformation of $`m`$ in spacetime. If $`m`$ is closed (as we assume), we need only consider a normal deformation of $`m`$. Let us first expand
$$\delta _\overline{}X^\mu =\mathrm{\Phi }^Im^\mu {}_{I}{}^{}=\psi \eta ^\mu +\varphi ^in^\mu {}_{i}{}^{}.$$
(118)
We use the symbol $`\overline{}`$ to distinguish this normal variation, which includes a variation along $`\eta ^\mu `$, from the worldvolume normal variation used earlier. We have also expressed the variation with respect to the adapted basis, so that $`\mathrm{\Psi }^0=\psi ,\mathrm{\Psi }^i=\psi ^i`$.
We now exploit the formalism developed in Ref. to express the induced normal variations of $`h_{AB}`$ and $`L_{AB}^I`$ as follows:
$`\delta _\overline{}h_{AB}`$ $`=`$ $`2L_{AB}{}_{}{}^{I}\mathrm{\Phi }_{I}^{},`$ (119)
$`\delta _\overline{}L_{AB}^I`$ $`=`$ $`\widehat{𝒟}_A\widehat{𝒟}_B\mathrm{\Phi }^I+L_A{}_{}{}^{CI}L_{CBJ}^{}\mathrm{\Phi }^J+\widehat{\gamma }_J^IL_{AB}{}_{}{}^{J}.`$ (120)
Here $`𝒟_A`$ is the covariant derivative compatible with $`h_{AB}`$, and $`\widehat{𝒟}_A`$ its extension that is covariant under rotations of normals to $`m`$ in spacetime. We need to compute the variation for the component $`I=0`$ in this last expression. The boundary term is not covariant under rotations of the normals, $`m^\mu _I`$: $`\kappa _{AB}`$ is the extrinsic curvature which corresponds to $`\eta ^\mu `$, there is no rotational arbitrariness here. For this reason, the appropriate deformation operator we apply to $`\kappa `$ is $`\delta _\overline{}`$ and not the manifestly rotationally covariant object $`\stackrel{~}{\delta }_\overline{}`$ . This implies the necessity to restore the deformation connection $`\widehat{\gamma }^I_J`$ on the right hand side of Eq.(120). This would spell disaster if we needed to evaluate all of $`\widehat{\gamma }_{IJ}`$. However, we only require $`\widehat{\gamma }_{0i}`$ which is well defined on $`m`$. We then find
$`\delta _\overline{}\kappa `$ $`=`$ $`𝒟_A𝒟^A\psi +\left(K^{Ai}K_{Ai}\kappa ^{AB}\kappa _{AB}h^{AB}K_{AB}{}_{}{}^{i}K_{abi}^{}\eta ^a\eta ^b\right)\psi `$ (121)
$``$ $`𝒟_A\left(K^{iA}\varphi _i\right)+(\eta ^ah^{AB}K_{AB}{}_{}{}^{i}ϵ^{aA}K_A{}_{}{}^{i})\stackrel{~}{}_a\varphi _i\kappa ^{AB}K_{AB}{}_{}{}^{i}\varphi _{i}^{}.`$
To arrive at this expression we have used the fact that $`\widehat{\gamma }^{0i}=\eta ^a\eta ^bK_{ab}^i\psi +\eta ^a\stackrel{~}{}_a\varphi ^i`$ and we have introduced the notation, $`K_A{}_{}{}^{i}=\eta ^aϵ^b{}_{A}{}^{}K_{ab}^{i}`$. Some simplification is possible. Let us begin with the terms linear in $`\psi `$. We can use the contracted Gauss-Codazzi equation to obtain,
$$K^{Ai}K_{Ai}h^{AB}K_{AB}{}_{}{}^{i}K_{ab}^{}{}_{i}{}^{}\eta _{}^{a}\eta ^b=\eta ^a\eta ^b_{ab}.$$
In addition, we have
$$(\eta ^ah^{AB}K_{AB}{}_{}{}^{i}ϵ^{aA}K_A{}_{}{}^{i})\stackrel{~}{}_a\varphi _i=\eta _a(K^i\gamma ^{ab}K^{abi})\stackrel{~}{}_b\varphi _i.$$
Inserting these partial results into Eq. (121), we obtain
$`\delta _\overline{}\kappa `$ $`=`$ $`𝒟_A𝒟^A\psi \left(\kappa ^{AB}\kappa _{AB}+\eta ^a\eta ^b_{ab}\right)\psi `$ (122)
$``$ $`𝒟_A\left(K^{iA}\varphi _i\right)+\eta _a\left(K^i\gamma ^{ab}K^{abi}\right)\stackrel{~}{}_b\varphi _i\kappa ^{AB}K_{AB}{}_{}{}^{i}\varphi _{i}^{}.`$
We have now all the ingredients needed for the calculation of the normal variation of the boundary action (117). We neglect total derivatives, and we find
$`\delta _\overline{}S_m`$ $`=`$ $`2\beta {\displaystyle _m}\sqrt{h}\{[\kappa ^2\kappa _{AB}\kappa ^{AB}\eta ^a\eta ^b_{ab}]\psi `$ (123)
$`+`$ $`(h^{AB}\kappa \kappa ^{AB})K_{AB}{}_{}{}^{i}\varphi _{i}^{}+\eta _a(K^i\gamma ^{ab}K^{abi})\stackrel{~}{}_b\varphi _i\}.`$
To establish contact with the bulk variation, we must identify $`\psi =\eta ^a\mathrm{\Phi }_a`$. The total boundary contribution, obtained by summing this expression with the second line of Eq. (102), therefore does not contain any term involving derivatives of the $`\mathrm{\Phi }^i`$.
Let us consider now a translation of the boundary in the background space-time. The linear momentum associated with this contribution is then given by
$$p^\mu (\mathrm{\Sigma })=2\beta _\mathrm{\Sigma }\sqrt{h}\{\eta ^\mu [\kappa _{AB}\kappa ^{AB}\kappa ^2+\eta ^a\eta ^b_{ab}]+\eta ^a_{ab}e^{\mu b}+K_{AB}^i(h^{AB}\kappa \kappa ^{AB})n^\mu {}_{i}{}^{}\}.$$
(124)
It should be noticed that the surface momentum has an additional normal component. The linear momentum for the total action (116) using Eq. (102), becomes
$$P^\mu (\mathrm{\Sigma })+p^\mu (\mathrm{\Sigma })=\beta _\mathrm{\Sigma }\sqrt{h}\{\eta ^\mu [\kappa _{AB}\kappa ^{AB}\kappa ^2+\eta ^a\eta ^b𝒢_{ab}]+K_{AB}^i(h^{AB}\kappa \kappa ^{AB})n^\mu {}_{i}{}^{}\}.$$
(125)
Note that this expression vanishes identically for a string, since for $`D=2`$, we have $`𝒢_{ab}=0`$, and moreover in the degenerate case of a one dimensional boundary, $`k_{AB}=h_{AB}k`$.
Consider now an infinitesimal Lorentz transformation. The boundary contribution to the angular momentum is given by
$$m^{\mu \nu }(\mathrm{\Sigma })=\frac{1}{2}_\mathrm{\Sigma }\{\stackrel{~}{p}^\mu X^\nu +2\beta \sqrt{h}\eta _a(K^i\gamma ^{ab}K^{abi})n^\mu {}_{i}{}^{}e_{}^{\nu }{}_{b}{}^{}(\mu \nu )\},$$
(126)
where the quantity $`\stackrel{~}{p}^\mu `$ in the first term is the integrand appearing in Eq.(124).
The second term is what is needed to cancel the offending term in the angular momentum for the bulk Einstein-Hilbert action, in Eq. (103), so that the angular momentum for the total action is now
$$M^{\mu \nu }(\mathrm{\Sigma })+m^{\mu \nu }(\mathrm{\Sigma })=\frac{1}{2}_\mathrm{\Sigma }[\mathrm{\Pi }^\mu X^\nu (\mu \nu )],$$
(127)
where we denote with $`\mathrm{\Pi }^\mu `$ the integrand of the total linear momentum (125).
In conclusion, the same surface term that is appropriate to lower the order of the boundary conditions, also cancels the kinematic term in the angular momentum density.
## 9 Discussion
In this paper, we have presented a new approach to the derivation of the linear and angular momentum for a brane propagating in Minkowski spacetime based on the worldvolume geometry. We have considered a large class of brane actions, depending on the intrinsic and as well on the extrinsic geometry of the worldvolume, up to a first derivative of the extrinsic curvature. The generalization to a more general action is straightforward, following the guidelines given in the paper.
This analysis may be extended in a straightforward way to treat the corresponding conserved quantities for a brane propagating on a background spacetime possessing Killing vector fields.
A particular simple case we have not discussed explicitly here is the degenerate case of a point object described by a higher order action. We will discuss this case elsewhere.
It would also be interesting to apply this geometrical approach to supersymmetric branes . While the geometry of such objects is well understood, to our knowledge, the geometry of deformations of superembedded surfaces remains to be developed.
Acknowledgements
We have received partial support from CONACyT grant no. 211085-5-0118PE. Our thanks to Charles Torre for useful suggestions.
Appendix A
In this section, we consider the more general case of an action that depends also on first derivatives of the extrinsic curvature,
$$S_{(ho)}[Y]=_m\sqrt{\gamma }L(\gamma ^{ab},K_{ab}{}_{}{}^{i},\stackrel{~}{}_aK_{bc}{}_{}{}^{i}).$$
(128)
We can recycle the results of Sect. 5, for the variation with respect to the first arguments, that will not be repeated here. The new part is
$$\delta _{}S_{(ho)}=_mL^{abc}{}_{i}{}^{}\stackrel{~}{\delta }_{}^{}\stackrel{~}{}_aK_{bc}{}_{}{}^{i},$$
(129)
where we have defined
$$L^{abc}{}_{i}{}^{}=\frac{L}{\stackrel{~}{}_aK_{bc}^i}.$$
(130)
Note that, as a consequence of the Codazzi-Mainardi integrability condition, and the symmetry of $`K_{ab}^i`$, we have the symmetry property
$$\stackrel{~}{}_aK_{bc}{}_{}{}^{i}=\stackrel{~}{}_{(a}K_{bc)}{}_{}{}^{i},$$
(131)
from which it follows that
$$L^{abc}{}_{i}{}^{}=L^{(abc)}{}_{i}{}^{}.$$
(132)
In order to evaluate Eq. (129), we need to commute $`\stackrel{~}{\delta }_{}`$ with $`\stackrel{~}{}_a`$. To do this, we need the following expressions (see also )
$`\delta _{}\gamma _{ab}^c`$ $`=`$ $`\gamma ^{cd}\left[_a\left(K_{bd}{}_{}{}^{i}\mathrm{\Phi }_{i}^{}\right)+_b\left(K_{ad}{}_{}{}^{i}\mathrm{\Phi }_{i}^{}\right)_d\left(K_{ab}{}_{}{}^{i}\mathrm{\Phi }_{i}^{}\right)\right],`$ (133)
$`\delta _{}\omega _a^{ij}`$ $`=`$ $`K_a{}_{}{}^{bi}\stackrel{~}{}_{b}^{}\mathrm{\Phi }^jK_a{}_{}{}^{bj}\stackrel{~}{}_{b}^{}\mathrm{\Phi }^i.`$ (134)
Now, it is a direct computation to obtain that the contribution of (129) to the boundary term of the normal variation of the action is given by,
$`\delta _{}S_{(ho)}`$ $`=`$ $`{\displaystyle _m}\sqrt{h}\eta _a[(\stackrel{~}{}_b\stackrel{~}{}_cL^{abc}{}_{i}{}^{})\mathrm{\Phi }^i+(\stackrel{~}{}_bL^{abc}{}_{i}{}^{})\stackrel{~}{}_c\mathrm{\Phi }^iL^{abc}{}_{i}{}^{}\stackrel{~}{}_{b}^{}\stackrel{~}{}_c\mathrm{\Phi }^i`$ (135)
$``$ $`L^{bcd}{}_{i}{}^{}K_{bc}^{}{}_{}{}^{j}K_{}^{a}{}_{dj}{}^{}\mathrm{\Phi }_{}^{i}+3L^{bcd}{}_{j}{}^{}K_{bci}^{}K_d{}_{}{}^{aj}\mathrm{\Phi }_{}^{i}3L^{abc}{}_{j}{}^{}K_{b}^{}{}_{}{}^{dj}K_{cdi}^{}\mathrm{\Phi }^i].`$
If this seems complicated, well the Euler-Lagrange derivative is worse.
For an infinitesimal translation, we obtain that the total contribution to the momentum density is
$`𝒫^{a\mu }`$ $`=`$ $`\sqrt{\gamma }\{[L_{(ho)}\delta ^a{}_{d}{}^{}L^{abc}{}_{i}{}^{}(\stackrel{~}{}_bK_{cd}{}_{}{}^{i})+(\stackrel{~}{}_bL^{abc}{}_{i}{}^{})K_{cd}{}_{}{}^{i}]e^{\mu d}`$ (136)
$`+`$ $`[(\stackrel{~}{}_b\stackrel{~}{}_cL^{abc}{}_{i}{}^{})L^{bcd}{}_{i}{}^{}K_{bc}^{}{}_{}{}^{j}K_{}^{a}{}_{dj}{}^{}+3L^{bcd}{}_{j}{}^{}K_{bci}^{}K_d^{aj}`$
$``$ $`2L^{abc}{}_{j}{}^{}K_{b}^{}{}_{}{}^{dj}K_{cdi}^{}]n^{\mu i}\}.`$
The contribution to the angular momentum density is
$`^{a\mu \nu }={\displaystyle \frac{1}{2}}\{𝒫^{a\mu }X^\nu `$ $`+`$ $`\sqrt{\gamma }[(\stackrel{~}{}_bL^{abc}{}_{i}{}^{})n^{\mu i}e^\nu {}_{c}{}^{}2L^{abc}{}_{i}{}^{}K_{cd}^{}{}_{}{}^{i}e_{}^{\mu d}e^\nu _b`$
$`+`$ $`L^{abc}{}_{i}{}^{}K_{bc}^{}{}_{}{}^{j}n_{}^{\mu i}n^\nu {}_{j}{}^{}](\mu \nu )\}.`$
We note that in this expression all of the bivectors enter.
Appendix B
In this appendix, we show explicitly that, as stated at the end of Sect. 7, the vanishing of the second variation parallel to the worldsheet, $`\delta _{}(\delta S)`$ results in mere identities.
The parallel variation of Eq. (25) gives
$$\delta _{}(\delta S)=ϵ_\mu _m_a\left[\sqrt{\gamma }\mathrm{\Phi }^a^in^\mu {}_{i}{}^{}+\delta _{}𝒫^{a\mu }\right],$$
(137)
where we have used Eq. (19), valid for any worldvolume scalar density of weight one. We decompose $`\delta _{}𝒫^{a\mu }`$ in terms of its worldvolume projections,
$$\delta _{}𝒫^{a\mu }=[\delta _{}𝒫]^{ab}e^\mu {}_{b}{}^{}+[\delta _{}𝒫]^{ai}n^\mu {}_{i}{}^{}.$$
Substituting into Eq. (137), and using the Gauss-Weingarten equations (40),(41), we find the parallel analogue of Eqs. (107), (108),
$$_a\left(\sqrt{\gamma }\mathrm{\Phi }^a^i\right)=\stackrel{~}{}_a[\delta _{}𝒫]^{ai}+[\delta _{}𝒫]^{ab}K_{ab}{}_{}{}^{i},$$
(138)
$$\sqrt{\gamma }\mathrm{\Phi }^a^iK_a{}_{}{}^{b}{}_{i}{}^{}=\stackrel{~}{}_a[\delta _{}𝒫]^{ab}[\delta _{}𝒫]^{ai}K_a{}_{}{}^{b}{}_{i}{}^{}.$$
(139)
Both these equations are mere identities. Let us confirm this for the special case of a DNG object. Using the parallel deformation Gauss-Weingarten equation (see ),
$$\delta _{}e^\mu {}_{a}{}^{}=(_a\mathrm{\Phi }^b)e^\mu {}_{b}{}^{}+K_{ab}{}_{}{}^{i}\mathrm{\Phi }_{}^{b}n^\mu {}_{i}{}^{},$$
we have that the parallel variation of Eq. (51) gives
$$\delta _{}𝒫^{a\mu }=\mu \sqrt{\gamma }[(_c\mathrm{\Phi }^c\gamma ^{ab}^b\mathrm{\Phi }^a)e^\mu {}_{b}{}^{}K^{abi}\mathrm{\Phi }_bn^\mu {}_{i}{}^{}],$$
so that the worldvolume projections are
$$[\delta _{}𝒫]^{ab}=\mu \sqrt{\gamma }\left(_c\mathrm{\Phi }^c\gamma ^{ab}^b\mathrm{\Phi }^a\right),$$
$$[\delta _{}𝒫]^{ai}=\mu \sqrt{\gamma }K^{abi}\mathrm{\Phi }_b.$$
Substituting into Eq. (138) gives
$$_b(\mathrm{\Phi }^bK^i)=_a[K^{abi}\mathrm{\Phi }_b]+(_c\mathrm{\Phi }^c\gamma ^{ab}^b\mathrm{\Phi }^a)K_{ab}{}_{}{}^{i}.$$
Using the contracted Codazzi-Mainardi integrability condition, Eq. (93), one can easily verify that this equation is satisfied identically. On the other hand, substitution of the projections into Eq. (139) gives
$$\mathrm{\Phi }^aK_iK_a{}_{}{}^{bi}=\mathrm{\Phi }_cK^{aci}K_a{}_{}{}^{b}{}_{i}{}^{}_a(_c\mathrm{\Phi }^c\gamma ^{ab}^b\mathrm{\Phi }^a).$$
This equation can be seen to vanish identically as well. This requires the contracted Gauss-Codazzi equation (58), for the lefthand side together with the first term on the right hand side, and the Ricci identity for the remaining two terms,
$$\left(_a_b_b_a\right)\mathrm{\Phi }^b=_{ab}\mathrm{\Phi }^b.$$ |
warning/0002/astro-ph0002089.html | ar5iv | text | # Relic Neutrinos and Z-Resonance Mechanism for Highest-Energy Cosmic Rays
## Abstract
The origin of the highest-energy cosmic rays remains elusive. The decay of a superheavy particle (X) into an ultra-energetic neutrino which scatters from a relic (anti-)neutrino at the Z-resonance has attractive features. Given the necessary X mass of $`10^{1415}`$ GeV, the required lifetime, $`10^{1516}`$ y, renders model-building a serious challenge but three logical possibilities are considered: (i) X is a Higgs scalar in $`SU(15)`$ belonging to high-rank representation, leading to power-enhanced lifetime; (ii) a global X quantum number has exponentially-suppressed symmetry-breaking by instantons; and (iii) with additional space dimension(s) localisation of X within the real-world brane leads to gaussian decay suppression, the most efficient of the suppression mechanisms considered.
preprint: IFP-778-UNC astro-ph/0002089 February 2000
The confluence of cosmology and particle phenomenology benefits both disciplines and can lead to important new insights.
For protons propagating through the cosmological background radiation there is an energy cut-off, as discussed in e.g. , well-known as the GKZ effect, at an energy of $`E5\times 10^{19}`$ eV. Above this energy, the photoproduction of pions at the 3-3 resonance provides an energy attenuation that prohibits travel over a distance greater than $`50`$ Mpc.
Nevertheless, air showers initiated by a proton (or photon) with energies above the GKZ bound have been observed, and this fact needs explanation.
One possibility is that the origin involves the decay of a superheavy particle as in (for related earlier works, see e.g. ) but that the decay now produces a high energy neutrino which scatters from a relic background neutrino at the $`Z`$ pole and produces the primary. This Z-burst scenario was suggested in and further analysed in .
The kinematics of the neutrino-neutrino collision at the $`Z`$ pole requires an energy $`E_{resonance}=M(Z)^2/(2M(\nu ))`$ and taking $`M(\nu )=0.07eV`$ as suggested by the SuperKamiokande data gives $`M_{resonance}10^{23}eV`$, just as needed to explain the data. This is the most attractive feature of the model.
We first estimate the mass and lifetime of the superheavy particle needed to fit the data. This requires two relationships derived in . Namely the flux of cosmic rays beyond the GKZ cut-off is estimated as:
$`\mathrm{\Phi }_{CR}`$ $`=`$ $`{\displaystyle \frac{C_1}{(4\pi sr)km^2(100y)}}`$ (1)
$`\times `$ $`\left({\displaystyle \frac{N}{10}}\right)\left({\displaystyle \frac{\eta }{0.14}}\right)\left({\displaystyle \frac{\mathrm{\Omega }_X}{0.2}}\right)\left({\displaystyle \frac{h}{0.65}}\right)^2`$ (2)
$`\times `$ $`\left(B_\nu {\displaystyle \frac{10^7t_0}{\tau _X}}\right)\left({\displaystyle \frac{0.07eV}{M(\nu )}}\right)^{3/2}\left({\displaystyle \frac{10^{14}GeV}{M(X)}}\right)^{5/2}`$ (3)
and
$$\frac{\tau _X}{t_0}B_\nu ^1>C_2\times 10^5\left(\frac{\mathrm{\Omega }_X}{0.2}\right)\left(\frac{h}{0.65}\right)^2\left(\frac{10^{14}GeV}{M(X)}\right)^{3/4}$$
(4)
For the numerical dimensionless coefficients we find the values $`C_1=0.33`$ and $`C_2=12.8`$ which we use in the following analysis. The notation is: $`N`$ is the number of protons and photons per annihilation event; $`\eta `$ is the relic neutrino density relative to the present photon number density $`\eta =(n_{\nu ,relic}/n_{\gamma ,0})`$; $`\mathrm{\Omega }_X`$ is the contribution of $`X`$ particles to the energy density, relative to the critical density; $`h`$ is the Hubble constant in units of $`100km/s/Mpc`$; $`B_\nu `$ is the branching ratio of $`X`$ into neutrinos; $`t_0`$ is the age of the universe; $`\tau _X`$ is the lifetime of $`X`$; $`M(\nu )`$ is the neutrino mass; and $`M(X)`$ is the mass of $`X`$.
Assuming central values of all other parameters we plot the allowed region of $`M(X)`$ and $`\tau _X`$ in Figure 1; variations in $`N,\eta ,\mathrm{\Omega }_X,h,B_\nu ,M(\nu )`$ can extend the allowed region but here we need only the order of magnitude estimate.
The value of $`M(X)`$ must certainly exceed $`2E_{resonance}`$ so that a two body decay can lead to a neutrino acquiring enough energy. Higher energies can be red-shifted down to $`E_{resonance}`$ if the progenitor $`X`$ particle is at a red shift $`z>0`$.
The resultant spectrum will cut-off at $`M(X)/2`$ and will be expected to provide a two-component type of overall spectrum, with a dip around $`EE_{GKZ}`$ as can be seen in the data.
Since $`Z`$ decay gives rise to approximately 10 times as many photons as nucleons the model predicts a concomitant number of high-energy photons as cosmic-ray primaries. Because the data is sparse, it is not yet possible to discriminate on this basis, as discussed in ; this is an important prediction of the Z-burst scenario.
From the above analysis we conclude that the required particle properties $`M(X)`$ and $`\tau _X`$ for the hypothetical state X are well defined in order of magnitude. Namely, the mass M(X) should lie between $`10^{14}`$ and $`10^{15}`$ GeV and the lifetime $`\tau (X)`$ should lie between $`10^{16}`$ and $`10^{17}`$ years.
The remainder of the paper will discuss three possible microscopic theories or, better, scenarios for this combination of M(X) and $`\tau (X)`$. We present these three scenarios in what we regard as their increasing appeal, from (i) to (iii).
(i) Power suppression.
The expectation for a particle of this mass is that, unless it is absolutely stable due to some exact conservation law, it will decay exceedingly quickly with a lifetime expected to be $`\tau 10^{24}`$ seconds. Since the required lifetime is larger by some 46 or so orders of magnitude, the longevity is the principal difficulty, as emphasized in .
One viewpoint is that this extraordinary suppression of the decay is an argument against the model, as is the problem, already mentioned, of super-high-energy photons concomitant with the protons.
Let us here take the viewpoint, as discussed in that the photons are a prediction of the model, to be tested in future experiments, rather than a fatal flaw. The data on HECR is probably too sparse to reach any stronger conclusion.
Therefore the only remaining question is longevity.
The first scenario (i) is that considered (in a different model) in . We assume the particle X is a boson and posit a coupling
$`{\displaystyle \frac{g}{M^p}}X_{\beta _1\beta _2\mathrm{}..\beta _n}^{\alpha _1\alpha _2\mathrm{}..\alpha _n}(\overline{\psi }^{\beta _1}\psi _{\alpha _1})(\overline{\psi }^{\beta _2}\psi _{\alpha _2})\mathrm{}\mathrm{}\mathrm{}(\overline{\psi }^{\beta _n}\psi _{\alpha _n})`$
where the power is $`p=3(n1)/2`$. Let us assume that such a coupling is gravity-induced and that M is the reduced Planck mass $`MM(Pl)10^{18}GeV`$. Then one expects the lifetime $`\tau (X)`$ to be of the order of magnitude
$`\tau (X)(10^{24}sec.)\times \left({\displaystyle \frac{M(X)}{M(Pl)}}\right)^{2p}`$
in which the mass ratio is $`\frac{M(X)}{M(Pl)}10^310^4`$. To arrive at a suppression of $`10^{46}`$ thus requires $`2p1215`$ and $`n56`$. Thus the X field must have a high tensorial rank. If this is too much for the reader, skip to scenario (ii).
In the case $`n=2`$ was considered. In the spontaneous breaking of $`SU(15)`$ theory such a tensor appears ”naturally” in the Higgs sector. There is no apparent need for such a high rank as n = 5 or 6, but equally no reason for their absence. The dimensions of such scalar representations in SU(15) are astronomical - even for n=2 the dimension is 14,175 while for n=5 and 6 this becomes respectively 125,846,784 and 1,367,127,216.
It is difficult to believe that such power suppression could be responsible for the longevity. It is a logical possibility which appears highly contrived. Thus exponential or gaussian suppression is more appealing.
(ii) Exponential suppression.
Let us assume that the superheavy particle X carries a conserved quantum number $`Q_X`$ (analagous to baryon number, B) and that in perturbation theory the quantum number is exactly conserved. If there is no open channel which conserves $`Q_X`$ then the state will be absolutely stable.
In the case of B in the standard model, it was first shown in 1976 by ’t Hooft that nonperturbative instanton effects violate conservation and lead to decay of otherwise stable states such as the proton. The resutant rate is typically exponentially suppressed by an exponential of the form $`exp(constant/g^2)`$ where $`g`$ is the gauge coupling constant. Many other examples of such suppression are covered in .
Thus, one scenario is that $`Q_X`$ generates a symmetry of the lagrangian but $`X`$ decays with exponential suppression due to instanton effects. We mention this only for completeness - any quantitative estimation would require many hypotheses.
(iii) Gaussian suppression.
This scenario which is, in our opinion, the most appealing involves the assumption of at least one extra spatial dimension. We will take five space-time dimensions, four space and one time.
Let the coordinates be $`(x_0,𝐱,y)`$ with $`y`$ as the hypothetical extra dimension on which we now focus.
It used to be thought, up to a decade ago, that any such $`y`$ must be compactified at or beyond the GUT scale of $`(10^{16}GeV)^110^{32}m`$ (recall $`1(GeV)^12\times 10^{16}m`$). In 1990, Antoniadis was the first to entertain very much larger compactification scales $`(1TeV)^110^{19}m`$. In 1998 it was pointed out that, although the strong and electroweak interactions of the quarks and leptons need be confined to a region of $`y`$ not exceeding $`10^{19}m`$ (the real-world brane on which we live), the gravitational interaction could be decompactified even out to $`1mm=10^3m`$ without contradicting experimental data, offering the possibility of detecting such additional dimensions by deviation from Newton’s Law of Gravity at millimeter scales.
Models in which the fifth dimension contains a real-world brane and a suitably separated Planck brane which delimits gravitational propagation have been discussed in . Such models are of interest mainly because they suggest how to incorporate gravity in the conformality approach which ab initio describes a flat (gravitationless) space-time.
Let us assume, therefore, that the standard model states are all confined within a real-world brane with a thickness of order $`10^{19}m`$ in the $`y`$ direction. Following (see also and, for the many fold universe, ) a scalar field $`\mathrm{\Phi }(x_\mu ,y)`$ which has a $`\mathrm{\Phi }`$ domain wall in the $`y`$ dimension. In the vicinity of the wall centered by convention at $`y=0`$, and with the normalizations of the field has the value
$$\mathrm{\Phi }(y)=2\mu ^2y$$
(5)
First consider chiral fermions (after all, $`X`$ could be a fermion but we will consider the boson possibility later). In this case we write the five-dimensional action
$$S=d^4x𝑑y\overline{\mathrm{\Psi }}_i[i\gamma _\mu /y_\mu +\mathrm{\Phi }(y)m_i]\mathrm{\Psi }_i+\mathrm{}$$
(6)
The fermion $`\mathrm{\Psi }_i`$ is now localised at $`y=m_i/(2\mu ^2)`$. If the Higgs $`H(y)`$ is unlocalised inside the domain wall then the resulting coupling of $`\mathrm{\Psi }_i`$ to $`\mathrm{\Psi }_j`$ and $`H`$ has a gaussian suppression
$$exp(C(m_im_j)^2)/\mu ^2)$$
(7)
where C is a coefficient of order unity. This is the gaussian overlap of the gaussian tails of the two wave functions. The thickness of the real-world brane is $`\tau (\mu )^1`$ while the separation of the two fermion wave functions is $`\sigma (\mathrm{\Delta }y)_{ij}`$. To obtain the required suppression of $`10^{46}`$ we need $`\tau /\sigma 10`$. For example, if $`\tau 10^{19}m`$, one needs $`\sigma 10^{20}m`$. Clearly only the ratio $`\tau /\sigma `$ matters, but need not be a large number, in order to obtain the necessary suppression.
It is worth remarking that this gaussian suppression of Yukawa couplings by localization in the fifth dimension has a long historical counterpart in the localization of states on orbifold singularities starting with in 1987 and subsequent derivative literature.
The superheavy particle $`X`$ may not be a fermion but a boson, e.g. the Higgs scalar considered in scenario (i) above. Fortunately we can easily extend the localization argument of to the case of a boson and obtain a similar result for the gaussian suppression of the decay amplitude and consequent longevity. We replace Eq.(6) by the following:
$$S=d^4x𝑑y\varphi _i^{}[\mathrm{}m_S^2+\mathrm{\Phi }(y)^2]\varphi _i+\mathrm{}.$$
(8)
and, analagously to the fermion $`\mathrm{\Psi }`$, the boson $`\varphi `$ is localized around $`y=m_S/2\mu ^2`$ when $`\mathrm{\Phi }`$ again has a domain wall centered at $`y=0`$. Identifying $`\varphi `$ with $`X`$ then provides the required longevity by the same mechanism. It would be amusing if the highest-energy cosmic rays provided the first evidence for an extra spatial dimension!
To summarize the Z-burst mechanism for the highest energy cosmic rays, it has two positive features:
1) The resonant energy derived from the $`Z`$ mass and the SuperKamiokande neutrino mass is numerically close to the required energy.
2) The spectrum is predicted to have the two-component shape suggested by the present data.
On the other hand, there are also two apparently negative features:
3) The concomitant high-energy photons are not confirmed by present data. Better data will confirm or refute this important prediction.
4) The longevity of the superheavy particle is such a challenge to microscpic model- building that it may render the model less credible.
It is point (4) which we have attempted to ameliorate in the present article.
Acknowledgments
This work was supported in part by the US Department of Energy under the Grant No. DE-FG02-97ER-41036. We thank Ignatios Antoniadis, Don Ellison, Markus Luty and Tom Weiler for discussions.
Figure 1.
Allowed region of $`M(X)\tau _X`$ from Eq.(3) and Eq.(4) of the text. In the Figure $`a_X=\tau _X(10^7t_0)^1`$ where $`t_0`$ is the age of the universe and $`b_X=M(X)/(10^{14}GeV)`$. Variations in $`N,\eta ,\mathrm{\Omega }_X,h,B_\nu ,M(\nu )`$ can extend the allowed region but we use only such order of magnitude estimates. |
warning/0002/hep-ex0002021.html | ar5iv | text | # Analytical and Numerical Flash-Algorithms for Track Fits
## I Introduction
In many particle physics experiments high-precision alignment can be performed provided there is sufficient data and enough computer (CPU) time alocated. While statistics may suffice for many experiments, CPU time can only be decreased through the use of flash reconstruction routines. For instance in pixel detectors, alignment parameters can be tuned iteratively to satisfy good track reconstruction in all events. However, reconstruction of all tracks in all events, for each iteration in alignment parameter space<sup>*</sup><sup>*</sup>*An example of an analytical alignment procedure (vs. an iterative one) is illustrated in for the SLD End-Cap Čerenkov Ring Imaging Detector., is CPU exhausting and in practice can only be applied to blocks of tracks. If the track fits themselves are also iterative , double-nesting of iteration loops occurs and CPU times reach on the order of months. On the other hand if track fits are semi/analytical, CPU times drop 3-4 orders in magnitude and such an approach becomes feasable.
The paper starts with the basic $`\chi ^2`$ fit, examines the validity of the standard quadratic approximation for $`\chi ^2`$ and presents analytical and semi-analytical track reconstruction algorithms, together with CPU-clocked examplesThe tests were performed on a DEC-ALPHA 878 machine, with EV 5.6 processor at 433 MHz, 640 MB RAM, and running under OSF1 V4.0. Coding was performed in FORTRAN.. The methods were developed gradually from the simple case of straight tracks (magnetic field free environment, or high momentum tracks), to the more demanding case of helical tracks. The Test Beam setup of the ATLAS Pixel Detector MCM-D modules was used as prototype example for the applied methods, although the latter methods are general and applicable to a number of detectors using tracks.
## II General Fits
The general expression for $`\chi ^2`$ in a fit is the sum of the squared normalised residuals over the set of experimental points:
$$\chi ^2\stackrel{\mathrm{def}}{=}\underset{\mathrm{j}=1}{\overset{\mathrm{N}}{}}\left(\frac{\stackrel{}{\mathrm{\Delta }}_\mathrm{j}}{\sigma _{\mathrm{j}_{(\mathrm{\Delta })}}}\right)^2$$
(1)
where the error of the $`j^{th}`$-residual $`\sigma _{j_{(\mathrm{\Delta })}}`$ is related to the direction of the residual itself, $`\stackrel{}{\mathrm{\Delta }}_j`$ through the covariance matrix $`\sigma _j^2`$ :
$$\sigma _{j_{(\mathrm{\Delta })}}^2=\frac{\stackrel{}{\mathrm{\Delta }}_j\sigma _j^2\stackrel{}{\mathrm{\Delta }}_j}{\left(\stackrel{}{\mathrm{\Delta }}_j\right)^2}$$
(2)
leading to:
$$\chi _{exact}^2=\underset{j=1}{\overset{N}{}}\frac{\left(\stackrel{}{\mathrm{\Delta }}_j\right)^2}{\stackrel{}{\mathrm{\Delta }}_j\sigma _j^2\stackrel{}{\mathrm{\Delta }}_j}$$
(3)
More often however, a quadratic$`\stackrel{}{\mathrm{\Delta }}_j\sigma _j^2\stackrel{}{\mathrm{\Delta }}_j`$ = $`\mathrm{\Delta }_x^2/\sigma _x^2+\mathrm{\Delta }_y^2/\sigma _y^2+\mathrm{\Delta }_z^2/\sigma _z^2`$, for $`\sigma _j^2`$ assumed diagonal. approximation of $`\chi _{exact}^2`$ is used:
$$\chi _{approx}^2\stackrel{\mathrm{def}}{=}\underset{\mathrm{j}=1}{\overset{\mathrm{N}}{}}\stackrel{}{\mathrm{\Delta }}_\mathrm{j}\sigma _\mathrm{j}^2\stackrel{}{\mathrm{\Delta }}_\mathrm{j}$$
(4)
with the equivalent residual error:
$$\sigma _{j_{(\mathrm{\Delta })}}^2=\frac{\left(\stackrel{}{\mathrm{\Delta }}_j\right)^2}{\stackrel{}{\mathrm{\Delta }}_j\sigma _j^2\stackrel{}{\mathrm{\Delta }}_j}$$
(5)
If a track impacts a point’s error ellipsoid at $`n=tg(\theta _{incid}`$) with respect to one of the principal axes, then approximating $`|\mathrm{\Delta }_j`$ as perpendicular to the trajectory, the two $`\sigma `$’s can be written as:
$`\sigma _{exact}^2`$ $``$ $`{\displaystyle \frac{\sigma _{xy}^2+n^2\sigma _z^2}{1+n^2}}`$ (6)
$`\sigma _{approx}^2`$ $``$ $`{\displaystyle \frac{1+n^2}{\sigma _{xy}^2+n^2\sigma _z^2}}`$ (8)
both reducing to $`\sigma _{xy}^2`$ or $`\sigma _z^2`$ for tracks impacting along one of the principal axes.
The Test Beam stand for the ATLAS Pixel MCM-D modules is a Telescope setup, with 4 tracking elements (Sirocco strip Detectors) and two slots for Pixel Module evaluation. The setup was mounted on a marble optical-bench, itself placed on a rail that allowed it to be moved into the active area of a 1.4 Tesla Spectrometer Magnet. The latter was used to determine the magnetic field influence (x-shifts) in an environment comparable to the ATLAS Detector. The Sirocco strips were mounted along $`\mathrm{\Delta }z`$ = 1.4 m of beam-line with a tolerance of $`\pm 0.5`$ mm. The strips, 30$`\mu `$m wide, provided a resolution of 4$`\mu `$m in the x- and y-directions, while the Pixels, 50$`\mu `$m $`\times `$ 400$`\mu `$m, a resolution on the order of 14$`\mu `$m $`\times `$ 180$`\mu `$m depending on the cast technology of the chips. The error matrices of the Test Beam stand Telescope points had thus associated ellipsoids with aspect ratios of 1:125, the difference between fitting a 3D-line to these $`\sigma `$’s and one to spherical $`\sigma `$’s being 0.1 $`\mu `$m in the Telescope’s mid-plane.
The difference between $`\chi _{exact}^2`$ and $`\chi _{approx}^2`$ depends strongly on the impact angle of the track onto the individual error ellipsoids:
$$\frac{\mathrm{\Delta }\chi ^2}{\chi _{exact}^2}=\frac{\left(\sigma _{xy}/\sigma _z\right)^2+\left(\sigma _z/\sigma _{xy}\right)^22}{\left(n+1/n\right)^2}$$
(9)
For the Telescope setup the track’s impact angle was on the average 0.15 mrad, with a corresponding $`\mathrm{\Delta }\chi ^2/\chi _{exact}^2`$ on the order of 0.02%. When the tracks impact however at an arbitrary angle, $`\mathrm{\Delta }\chi ^2/\chi _{exact}^2`$ can reach as high as 3000 for the current $`\sigma _z/\sigma _{xy}`$ ratio, even if $`\mathrm{\Delta }\sigma ^2/\sigma _{exact}^2`$ is on the order of unity.
For tracks impacting all points at a constant angle (“stiff”-tracks), $`\mathrm{\Delta }\chi ^2`$ is constant along the trajectory and the two methods yield identical results. If the track is measured however piecewise in two different sub-systems, or it is composed of an ensemble of points with different $`\sigma `$’s (different types of detectors), then even for straight tracks the two solutions differ. For tracks bending in magnetic field, the track’s impact angle changes continuously along the track, $`\mathrm{\Delta }\chi ^2`$ following as:
$$d\left(\frac{\mathrm{\Delta }\chi ^2}{\chi _{exact}^2}\right)/\frac{\mathrm{\Delta }\chi ^2}{\chi _{exact}^2}=\frac{dn}{n}2\frac{1n^2}{1+n^2}$$
(10)
Within the 1.4 m of the Telescope, the 180 GeV/c tracks used bend in the B = 1.4 T magnetic field equivalently to $`\mathrm{\Delta }\chi ^2/\chi _{exact}^2`$ $``$ 0.02% up front and 16% downstream, the approximative method pulling the fit increasingly tighter towards the end - on the order of 0.5 $`\mu `$m per point. The effect is evidently insignificant, both in the Telescope setup, as well as in the real B-physics context of ATLAS, meaning 2 Tesla magnetic field, tracks of momentum greater than 1 GeV/c, and a measured track length of approximately 0.14 m. Over this span the expected point to point change in $`\chi ^2`$ is less than 4 % per %-$`\mathrm{\Delta }\chi ^2`$.
## III Line Fits
In most cases it is possible to interchange the non-linear expression (3) with its quadratic approximation (4), allowing analytical solutions to be given for “stiff-tracks” i.e. \- particle out of magnetic field, weak field with respect to track momentum, or distance travelled small with respect to existing resolution.
The simplest fit is for $`\sigma _j^2=\sigma ^2\mathrm{𝟏}=const.`$ :
$$\chi ^2=\underset{j=1}{\overset{N}{}}\left(\stackrel{}{\mathrm{\Delta }}_j\right)^2=min.$$
(11)
Parametrising the tracks as:
$$\stackrel{}{r}=\stackrel{}{r}_0+\lambda \stackrel{}{n}$$
(12)
with $`\stackrel{}{n}^2=1`$ and $`\stackrel{}{r}_0\stackrel{}{n}=0`$, equation (11) becomes:
$$\chi ^2=\underset{j=1}{\overset{N}{}}\stackrel{}{\delta }_j^2=min.$$
(13)
where $`\stackrel{}{\delta }_j=\stackrel{}{r}_j\stackrel{}{r}_0\lambda _j\stackrel{}{n}`$. The minimum condition implies locally $`\lambda _j=\stackrel{}{r}_j\stackrel{}{n}`$ and globally:
$`\stackrel{}{r}_0`$ $`=`$ $`(\mathrm{𝟏}\stackrel{}{n}\stackrel{}{n})\stackrel{}{r}`$ (14)
$`𝐌\stackrel{}{n}`$ $`=`$ $`\mu _0\stackrel{}{n}`$ (15)
where $`\stackrel{}{r}`$ denotes average over measured points, and $`𝐌=\stackrel{}{r}\stackrel{}{r}\stackrel{}{r}\stackrel{}{r}`$ the spread ellipsoid of points around $`\stackrel{}{r}`$. The 3 eigen-values of $`𝐌`$ represent the length of the track ($`\mu _0`$), and the two transversal variances to the line-fit.
For $`\sigma _j^2=diag(\sigma _x^2,\sigma _y^2,\sigma _z^2)=const.`$ equation (13) holds again, however in normalised form, with $`\stackrel{}{r}_i\sigma ^1\stackrel{}{r}_i`$, $`\stackrel{}{r}_0\sigma ^1\stackrel{}{r}_0`$ and $`\stackrel{}{n}\sigma ^1\stackrel{}{n}`$. Adapted to the geometry of the Telescope this solution has been clocked to 0.033 $`\mu `$s/2D-fit and 2.6 $`\mu `$s/3D-fit. These times are half or less for current 1000-1500 MHz clock machines. CPU-wise this allows the alignment to be performed using blocks of tracks. The intuitive picture of this procedure is aligning two incomplete 3D images: where there is a track, there is a pixel on the image. The more tracks in an alignment loop, the more pixels per image, and the better the alignment. This idealistic picture is cut cold however, by either lack of data, or of CPU power. In most of the contemporary HEP experiments the limiting factor is CPU power, as the reconstruction of the objects used in the alignment (in this case tracks) can be quite CPU costly. Eliminating this problem with flash-algorithms opens the way to high accuracy alignment. The increase in alignment resolution is proportional to $`\sqrt{N}_{block}`$, and the CPU cost to $`N_{block}`$. Holding CPU time fixed, any speed increase can be equivalenced to a resolution increase. Thus a factor of 4000 in speed, would be equivalent to a factor of 63 in resolution. To exemplify this, consider the Test Beam setup, where the tracks impact the Sirocco planes at almost normal incidence. At precisely normal incidence any z-misalignment would be un-noticed. The tracks do impact however at an angle on the order of 0.15 mrad, giving an alignment “lever arm” for $`\mathrm{\Delta }z`$ misalignments on the order of 0.15 $`\mu `$m/mm. With the existing Sirocco resolution ($``$ 4 $`\mu `$m) the resolving power (per track) for $`\mathrm{\Delta }z`$ would be on the order of 27 mm. Increasing the number of tracks per one $`\mathrm{\Delta }z`$ misalignment loop the resolution improves dramatically, as illustrated in figure 1.
Any multi-tile pixel detector benefiting from this high precision alignment method would better perform in physics involving the resolution of vertices. With the increasing energy frontier (LHC, TESLA/NLC) particle ID techniques are replaced with other means of signal identification. Signal “concentrators” that can serve this purpose are the B-mesons. Their identification depends crucially on vertexing resolution<sup>§</sup><sup>§</sup>§For example selecting B-events with high purities and efficiency can be achieved by applying a minimally missing $`\stackrel{}{p}_{}`$ correction to the m<sub>π</sub> evaluated vertex mass . In the ATLAS context $`\stackrel{}{p}_{}`$ would most likely be referenced to the axis of the jet containing the B sub-jet. The method depends crucially on the vertexing accuracy, yielding for instance at SLD B-events with sample purities on the order of 91-99 % and 65-20% corresponding efficiencies..
An improvement to the example considered would be to include in the track fit also the Pixel Demonstrator points. This would mean fitting to points with different error ellipsoids and the imposibility of “absorbing” all $`\sigma `$’s into $`\stackrel{}{r}_0`$ and $`\stackrel{}{n}`$ in an unique way, as done previously. For such $`\sigma _i^2\sigma _j^2`$ cases, the solution is given by a set of self-consistent equations:
$`\lambda _j`$ $`=`$ $`{\displaystyle \frac{\stackrel{}{n}\sigma _j^2(\stackrel{}{r}_j\stackrel{}{r}_0)}{\stackrel{}{n}\sigma _j^2\stackrel{}{n}}}`$ (16)
$`\stackrel{}{r}_0`$ $`=`$ $`\stackrel{}{r}\lambda \stackrel{}{n}`$ (18)
$`\stackrel{}{n}`$ $`=`$ $`{\displaystyle \frac{\lambda \stackrel{}{r}\lambda \stackrel{}{r}}{\lambda ^2\lambda ^2}}`$ (20)
solvable semi-analytically in approximately 3 iterations, starting from the previous analytically exact solution.
The approach used in this paper was to tune the alignmentThe alignment consisted of two separate stages, firstly out of magnetic field, and secondly in magnetic field. In both cases the translational and rotational degrees of freedom were firstly determined for the 4 Sirocco planes. The Pixel Modules were aligned with respect to the Telescope’s Sirocco planes. The alignment in magnetic field required an external marker, which was the setup’s trigger: a PIN Diode. Its change in image with magnetic field, on a “focal” plane, determined the absolute shifts in the tracking elements of the Telescope setup (the 4 Sirocco planes), and by reference, those of the Pixel Modules under evaluation. by looping over the reconstruction of 2000 tracks, the alignment residuals of the Telescope Sirocco plane-1 being shown in figure 2 (top). The Pixel Detectors under evaluation were to first order aligned analytically (exact solution), and subsequently tuned in a fashion similar to that of the Sirocco planes, in order to compensate for the effect of the low $`y`$-resolution on the better $`x`$-resolution. The difference between the Pixel demonstrator measured and Telescope predicted positions of the track hits is shown in figure 2 (middle and bottom) for a set of two Pixel Detectors under test. The identical *side*-resolution of the pixels in the $`x`$ and $`y`$-directions (fit parameter P3) - expected from uniform technology on the chip - gives credit that the alignment conducted to physically tangible results.
## IV Helix Fits
Track fits over small arcs of helices are very sensitive to the fluctuations of the experimental points. The radius of curvature (giving the momentum, or conversely the magnetic field) has consequently a large error, due to the points moving within their error bars. To reduce the errors of such “weak” parameters, fits to collections of tracks - experiencing the same conditions (magnetic field in this case) are used. Iterative fit routines consum however between 90000-300000 $`\mu `$s/3D helix-fit, the previously illustrated method (figure 1) being inapplicable in this case. On the other hand, although desired, flash semi-analytical algorithms for helix-fits are non-trivial to derive. Most of contemporary High Energy Physics experiments however involve large energies and the helical tracks span typically less than 1<sup>o</sup> of an arc (0.2<sup>o</sup> in the case of the Telescope). It is possible then in such cases to perform 3D-helix fits semi-analytically, by perturbatively curving a line-fit into a helix-fit.
The helix equations are:
$`d_t\stackrel{}{p}`$ $`=`$ $`{\displaystyle \frac{ec^2}{E}}\stackrel{}{p}\times \stackrel{}{B}`$ (21)
$`d_t\stackrel{}{r}`$ $`=`$ $`c^2\stackrel{}{p}/E`$ (23)
where $`E`$ is the particle’s energy, $`\stackrel{}{p}`$ its momentum and $`\stackrel{}{B}`$ the magnetic field. The solution to equations (23) is:
$$\stackrel{}{r}=\stackrel{}{r}_0+\lambda \stackrel{}{n}+\frac{\lambda ^2}{2!}f\left(\frac{\lambda }{R}\right)𝐅\stackrel{}{n}+\frac{\lambda ^3}{3!}g\left(\frac{\lambda }{R}\right)𝐆\stackrel{}{n}$$
(24)
where $`\lambda =\stackrel{}{v}_0t=\omega Rt`$ is a “linear” distance travelled by the particle, $`\stackrel{}{\omega }=|e|c^2\stackrel{}{B}/E`$ the helical rotation pulsation, $`\stackrel{}{n}=\stackrel{}{p}_0/p_0`$ the direction of engagement of the particle onto the magnetic field, $`\stackrel{}{n}_B=e\stackrel{}{B}/|e|B`$, $`R=p_0/|e|B`$ a parameter related to the radius of curvature of the helix $`R_{helix}=R\sqrt{1(\stackrel{}{n}\stackrel{}{n}_B)^2}`$, $`f(\zeta )`$ and $`g(\zeta )`$ two functions:
$`f(\zeta )`$ $`=`$ $`{\displaystyle \frac{2!}{\zeta ^2}}(1cos\zeta )\stackrel{\zeta 0}{}\mathrm{\hspace{0.17em}1}`$ (25)
$`g(\zeta )`$ $`=`$ $`{\displaystyle \frac{3!}{\zeta ^3}}(\zeta sin\zeta )\stackrel{\zeta 0}{}\mathrm{\hspace{0.17em}1}`$ (27)
respectively $`𝐅`$ and $`𝐆`$ two tensors:
$`𝐅`$ $`=`$ $`\times \stackrel{}{C}`$ (28)
$`𝐆`$ $`=`$ $`\stackrel{}{C}\stackrel{}{C}\stackrel{}{C}^2\mathrm{𝟏}`$ (29)
that satisfy $`𝐅^{}=𝐅`$, $`𝐆^{}=𝐆`$, $`𝐅𝐆=𝐆𝐅=\stackrel{}{C}^2𝐅`$, $`𝐅^{\mathrm{\hspace{0.17em}2}}=𝐆`$, and $`𝐆^2=\stackrel{}{C}^2𝐆`$. The vector $`\stackrel{}{C}`$ is $`\stackrel{}{n}_B/R`$.
It is evident that for $`R\mathrm{}`$ (or equivalently $`\lambda 0`$), expression (24) reduces to the parametrisation of the line (12) used in performing line fits - which is the requirement for the perturbative approach. In most experiments the third order approximation $`f(\zeta )1`$ and $`g(\zeta )1`$ holds up to the following limiting factors:
* geometric \- the arc of helix should not exceed a length beyond the approximation validity for $`f(\zeta )`$ and $`g(\zeta )`$. This is related to the demanded resolution $`\sigma `$ and the particle’s momentum:
$$p\frac{\lambda }{16(\sigma /\lambda )^{1/3}}5\mathrm{GeV}/\mathrm{c}$$
(30)
where in the above, $`\lambda `$ and $`\sigma `$ are expressed in \[m\] and $`p`$ in \[GeV/c\]. The value for the momentum is for the Telescope setup.
* dE/dx \- the loss of energy along the trajectory determines a “tighter” helix, the deviation:
$$\sigma =\frac{\lambda ^2E}{2\pi p^2c^2}\left(\frac{dE}{dx}\right)$$
(31)
needing to be smaller than 4 times the allowed tolerance in the Pixel plane.
* multiple scattering \- multiple deviations from the direction of flight add up to a displacement of:
$$\sigma 0.6\lambda \theta _{rms}$$
(32)
where $`\lambda `$ is expressed in \[mm\], $`\theta _{rms}`$ in \[mrad\] and $`\sigma `$ in \[$`\mu `$m\]. This should determine an error in the Pixel plane no larger than the allowed tolerance.
The semi-analytical helix fit procedure has 3 steps:
1. Estimation of $`\stackrel{}{n}`$, the engagement direction of the particle onto the magnetic field. This is obtained with small CPU demand via a 3D line flash-fit to the first 3-4 points of the trajectory. The vector $`\stackrel{}{n}`$ is an eigen-vector of $`𝐌=\stackrel{}{r}\stackrel{}{r}\stackrel{}{r}\stackrel{}{r}`$, hence any perturbation $`\delta 𝐌=𝐌_{helix}𝐌_{line}`$ changes it only to second order, and for numerical purposes $`\stackrel{}{n}`$ can be considered constant.
2. Using the $`\stackrel{}{n}`$ found above, the second order term corrections to $`\stackrel{}{r}_0`$ and $`\lambda _i`$ can be estimated:
$`\mathrm{\Delta }\lambda _i`$ $`=`$ $`{\displaystyle \frac{\lambda _i}{R}}(\stackrel{}{r}_i\stackrel{}{r}_0)\stackrel{}{n}\times \stackrel{}{n}_B`$ (33)
$`\mathrm{\Delta }\stackrel{}{r}_0`$ $`=`$ $`\stackrel{}{r}\stackrel{}{r}_0\lambda \stackrel{}{n}{\displaystyle \frac{\lambda ^2}{2R}}\stackrel{}{n}\times \stackrel{}{n}_B`$ (34)
computations again only modestly CPU demanding.
3. Introducing the third order term and using the previously corrected $`(\stackrel{}{n},\stackrel{}{r}_0,\lambda _i)`$, local and global equations for the parameters can be written:
$`(`$ $`\stackrel{}{n}+\lambda _i𝐅\stackrel{}{n}+{\displaystyle \frac{\lambda _i^2}{2}}𝐆\stackrel{}{n})\stackrel{}{\rho }_i=0`$ (35)
$``$ $`(\lambda \mathrm{𝟏}{\displaystyle \frac{\lambda ^2}{2}}𝐅+{\displaystyle \frac{\lambda ^3}{6}}𝐆)\stackrel{}{\rho }=0`$ (36)
$``$ $`\stackrel{}{\rho }=0`$ (37)
where $`\stackrel{}{\rho }_i`$ are the residuals of the points to the fitted curve:
$$\stackrel{}{\rho }_i=\stackrel{}{r}_i+\stackrel{}{r}_0+\lambda _i\stackrel{}{n}+\frac{\lambda _i^2}{2}𝐅\stackrel{}{n}+\frac{\lambda _i^3}{6}𝐆\stackrel{}{n}$$
(38)
Expanding to first order, the corresponding corrections $`(\mathrm{\Delta }\stackrel{}{n},\mathrm{\Delta }\stackrel{}{r}_0,\mathrm{\Delta }\lambda _i)`$ must satisfy:
$`\alpha _i\mathrm{\Delta }\lambda _i+\stackrel{}{a}_i\mathrm{\Delta }\stackrel{}{n}+\stackrel{}{b}_i\mathrm{\Delta }\stackrel{}{r}_0+\beta _i=0`$ (39)
$``$ $`\stackrel{}{a}\mathrm{\Delta }\lambda +\lambda ^2\mathrm{𝟏}\mathrm{\Delta }\stackrel{}{n}+𝐃\mathrm{\Delta }\stackrel{}{r}_0+\stackrel{}{\delta }=0`$ (40)
$``$ $`\stackrel{}{b}\mathrm{\Delta }\lambda +𝐃^{}\mathrm{\Delta }\stackrel{}{n}+\mathrm{𝟏}\mathrm{\Delta }\stackrel{}{r}_0+\stackrel{}{\sigma }=0`$ (41)
where:
$`\alpha _i`$ $`=`$ $`\stackrel{}{n}^2\stackrel{}{r}_i𝐅\stackrel{}{n}+\stackrel{}{r}_0𝐅\stackrel{}{n}`$ (42)
$`\beta _i`$ $`=`$ $`\stackrel{}{r}_0\stackrel{}{n}\stackrel{}{r}_i\stackrel{}{n}+\lambda _i\stackrel{}{n}^2+\lambda _i\stackrel{}{r}_0𝐅\stackrel{}{n}`$ (44)
$`\lambda _i\stackrel{}{r}_i𝐅\stackrel{}{n}+{\displaystyle \frac{1}{2}}\lambda _i^2\stackrel{}{r}_0𝐆\stackrel{}{n}{\displaystyle \frac{1}{2}}\lambda _i^2\stackrel{}{r}_i𝐆\stackrel{}{n}+`$ (46)
$`{\displaystyle \frac{1}{6}}\lambda _i^3\stackrel{}{n}𝐆\stackrel{}{n}`$ (48)
$`\stackrel{}{a}_i`$ $`=`$ $`\stackrel{}{r}_0\stackrel{}{r}_i+2\lambda _i\stackrel{}{n}+\lambda _i𝐅\stackrel{}{r}_i\lambda _i𝐅\stackrel{}{r}_0`$ (50)
$`\stackrel{}{b}_i`$ $`=`$ $`\stackrel{}{n}+\lambda _i𝐅\stackrel{}{n}`$ (52)
$`\stackrel{}{\delta }`$ $`=`$ $`\lambda \stackrel{}{r}_0\lambda \stackrel{}{r}+\lambda ^2\stackrel{}{n}+{\displaystyle \frac{1}{2}}𝐅\lambda ^2\stackrel{}{r}`$ (54)
$`{\displaystyle \frac{1}{2}}\lambda ^2𝐅\stackrel{}{r}_0+{\displaystyle \frac{1}{6}}\lambda ^3𝐆\stackrel{}{r}_0{\displaystyle \frac{1}{6}}𝐆\lambda ^3\stackrel{}{r}+`$ (56)
$`{\displaystyle \frac{1}{12}}\lambda ^4𝐆\stackrel{}{n}`$ (58)
$`\stackrel{}{\sigma }`$ $`=`$ $`\stackrel{}{r}_0\stackrel{}{r}+\lambda \stackrel{}{n}+{\displaystyle \frac{1}{2}}\lambda ^2𝐅\stackrel{}{n}+{\displaystyle \frac{1}{6}}\lambda ^3𝐆\stackrel{}{n}`$ (60)
$`𝐃`$ $`=`$ $`\lambda \mathrm{𝟏}{\displaystyle \frac{1}{2}}\lambda ^2𝐅`$ (62)
By eliminating $`\mathrm{\Delta }\lambda _i=(\beta _i+\stackrel{}{a}_i\mathrm{\Delta }\stackrel{}{n}+\stackrel{}{b}_i\mathrm{\Delta }\stackrel{}{r}_0)/\alpha _i`$ equations (41) become:
$`𝐌\mathrm{\Delta }\stackrel{}{n}+𝐍\mathrm{\Delta }\stackrel{}{r}_0`$ $`=`$ $`\stackrel{}{\tau }`$ (63)
$`𝐍^{}\mathrm{\Delta }\stackrel{}{n}+𝐑\mathrm{\Delta }\stackrel{}{r}_0`$ $`=`$ $`\stackrel{}{\pi }`$ (65)
where:
$`𝐌`$ $`=`$ $`<{\displaystyle \frac{\stackrel{}{a}\stackrel{}{a}}{\alpha }}>\lambda ^2\mathrm{𝟏}`$ (66)
$`𝐍`$ $`=`$ $`<{\displaystyle \frac{\stackrel{}{a}\stackrel{}{b}}{\alpha }}>\lambda \mathrm{𝟏}+{\displaystyle \frac{1}{2}}\lambda ^2𝐅`$ (68)
$`𝐑`$ $`=`$ $`<{\displaystyle \frac{\stackrel{}{b}\stackrel{}{b}}{\alpha }}>\mathrm{𝟏}`$ (70)
and :
$`\stackrel{}{\tau }`$ $`=`$ $`\stackrel{}{\delta }<{\displaystyle \frac{\beta \stackrel{}{a}}{\alpha }}>`$ (71)
$`\stackrel{}{\pi }`$ $`=`$ $`\stackrel{}{\sigma }<{\displaystyle \frac{\beta \stackrel{}{b}}{\alpha }}>`$ (73)
To zero<sup>th</sup> order the three $`𝐌`$, $`𝐍`$ and $`𝐑`$ tensors are proportional to $`(\mathrm{𝟏}\stackrel{}{n}\stackrel{}{n})`$, the non-invertable perpendicular projection to $`\stackrel{}{n}`$, by factors of $`\lambda ^2`$, $`\lambda `$ and 1. Therefore in the numerical approach, the inversion is obtained by decomposing the operators into a part proportional to $`(\mathrm{𝟏}\stackrel{}{n}\stackrel{}{n})`$ and a “remainder”:
$$𝐑=(\mathrm{𝟏}\stackrel{}{n}\stackrel{}{n})(2tr𝐑𝐑_{\stackrel{}{n}\stackrel{}{n}}𝐑_{\stackrel{}{n}\stackrel{}{n}}^{})/4+\mathrm{}$$
(74)
The solution $`(\mathrm{\Delta }\stackrel{}{n},\mathrm{\Delta }\stackrel{}{r}_0,\mathrm{\Delta }\lambda _i)`$ is therefore:
$`\mathrm{\Delta }\stackrel{}{n}`$ $`=`$ $`(\mathrm{𝐑𝐍}^1𝐌𝐍^{})^1(\mathrm{𝐑𝐍}^1\stackrel{}{\tau }\stackrel{}{\pi })`$ (75)
$`\mathrm{\Delta }\stackrel{}{r}_0`$ $`=`$ $`𝐍^1(\stackrel{}{\tau }𝐌\mathrm{\Delta }\stackrel{}{n})`$ (77)
$`\mathrm{\Delta }\lambda _i`$ $`=`$ $`{\displaystyle \frac{1}{\alpha _i}}(\beta _i+\stackrel{}{a}_i\stackrel{}{n}+\stackrel{}{b}_i\mathrm{\Delta }\stackrel{}{r}_0)`$ (79)
All quantities in this section were considered normalised \- *i.e.* $`\sigma ^1\stackrel{}{r}\stackrel{}{r}`$, although for notation simplicity they were written as the quantities themselves.
The CPU demand of the above 3 steps is under 15 $`\mu `$s. For better precision however, the last step can be repeated twice, bringing the 3D-helix fit to 22 $`\mu `$s. This is at least 4000 times faster than any iterative version of the fit.
To complete the fit in all generality, $`dE/dx`$ can be incorporated by perturbatively bending a pure helical track into a $`dE/dx`$-helix. This would allow in principle to infer the track’s mass, although with large errors in certain energy ranges.
Using the fit over blocks of 2000 tracks the alignment in magnetic field was checked and adjusted. The fit was then used on blocks of 10 tracks for a fine scan of the beam energy delivered by the SPS to the Test Beam setup. This was expressed in normal-impact radius of curvature equivalentThe normal-impact radius of curvature equivalent is the radius of curvature of the particle’s track, should the particle impact the field orthogonaly. It is a measure of the particle-momentum, preferred in the context of multiple tracks, when each track impacts the field at a different angle., and it is shown in figure 3 (bottom-right).
## V Conclusions
Analytical methods were shown to have a dramatic impact on the speed of line and helix fits, bringing down CPU usage by 3-4 orders of magnitude. The methods developed were successfully tested in the alignment and data reconstruction of the Test Beam results for the ATLAS Pixel Detector MCM-D modules, and can be of great impact for high precision alignment and reconstruction in all pixel detectors. Such routines could be used for instance by the ATLAS Inner Detector to specify precision $`\stackrel{}{p}_{}`$ corrections to the $`m_\pi `$ evaluated mass of B-vertices, in order to strongly suppress c-decay backgrounds.
## VI Acknowledgements
I am thankful to the High Energy Physics group of the Wuppertal University - in particular to Prof. Dr. K.-H. Becks, for the kind hospitality and facilities provided during completion of this work, as well as to the ATLAS Pixel Detector Collaboration for the opportunity of engaging in the Test Beam activity. I would especially like to thank the Alexander von Humboldt Foundation for support during this period and for the opportunity of better knowing German research and culture. |
warning/0002/astro-ph0002076.html | ar5iv | text | # 1 Introduction
## 1 Introduction
The intensity of Extensive Air Showers (EAS) with fixed shower sizes $`N_e`$ is assumed to decrease exponentially with increasing atmospheric depth of the observation level. This is considered to be due to the absorption of the particles of the EAS cascade following an exponential law
$$N_e(X)=N_e(X_0)exp\left(\frac{XX_0}{\mathrm{\Lambda }}\right),\mathrm{with}XX_0.$$
(1)
$`X_0`$ is a definite initial atmospheric depth after the maximum of the longitudinal development where the number of (charged) particles is $`N_e(X_0)`$ and further decreasing exponentially, $`N_e(X)`$ is the number of particles of the EAS at the slant depth $`X[gcm^2]`$.
The quantity $`\mathrm{\Lambda }`$ controls the attenuation of particles of the individual cascade (size attenuation length). It is related to the inelastic cross sections (to the mean free path length $`\lambda _A`$) of the interaction of the primary cosmic ray particles with air nuclei. The attenuation of the flux intensity of Extensive Air Showers is characterized by a related quantity $`\lambda _N`$ (intensity attenuation length, absorption), which can be directly measured by cosmic rays detector arrays. Thus measurements of the attenuation of the EAS intensity in the atmosphere are considered to be an interesting source of information about hadronic interactions, especially if extended to the ultrahigh energy region expected from the forthcoming LHC and TESLA accelerators. In addition due to the sensitivity of the cross sections to the mass of the primary, alterations of the attenuation length with the energy may be indicative for the variations of the mass composition. Measured results imply tests of the energy dependence of the extrapolated cross sections used for Monte Carlo simulations.
The investigations of the present paper are based on an EAS sample measured 1997-1999 with the MAKET ANI array on Mt. Aragats station (Armenia) and registered for different angles-of-incidence in the zenith angle interval $`\mathrm{\Theta }=045^{}`$. The data basis of the analysis can be enlarged by published data from KASCADE $`(1046gcm^2)`$ and EAS TOP $`(810gcm^2)`$ experiments. Spectra measured by EAS TOP are given in Ref. . Data and zenith angle dependence for KASCADE results are obtained by scanning the spectrum plots communicated by the KASCADE collaboration .
We apply different procedures to deduce the attenuation. First we consider the degradation of the EAS flux with fixed shower size $`N_e`$ with increasing zenith angle i.e. increasing atmospheric thickness of the shower development (characterized by the intensity attenuation length ($`\lambda _N`$) ). Differently the technique of the constant intensity cut (CIC) considers the intensity spectrum of EAS events and relates equal intensities observed at different atmospheric depths.
There is the tacit assumption that the shower size reflects the energy of the primary. The procedure can be refined by using the knee position in the $`N_e`$ spectrum as a bench mark for a well defined energy, so far we may associate the knee phenomenon to a feature of the primary energy spectrum of cosmic radiation.
## 2 Experimental spectra
The experimental basis of the present investigations are measurements of shower size spectra in the knee region and their zenith-angle dependence performed with the MAKET ANI array of the Mt. Aragats Cosmic Ray Station (3200 m a.s.l.) in Armenia. Details of the measurements and the experimental procedures taking into account the detector response are given elsewhere . For a detailed description of the knee region the traditional approximation with two different spectral indices below and above the knee, defining the knee position as intersection of two lines in a logarithmic presentation, appears to be insufficient. Hence a more sophisticated method has been applied with parameterization of the slope of the spectra (see Ref.).
Tab.1 compiles the characteristics of the size spectra measured with the MAKET ANI installation, the changes of the slopes in the knee region ($`\mathrm{\Delta }N_{e_k}`$), expressed by different spectral indices below ($`\gamma _1`$) and above ($`\gamma _2`$) the knee position $`N_{e_k}`$ for the zenith-angle range of $`\mathrm{\Theta }=045^{}`$. For the display and the analysis of the zenith-angle dependence, the size spectra are determined in 5 angular bins of equal $`\mathrm{\Delta }\mathrm{sec}\mathrm{\Theta }`$ widths. The accuracy of the zenith angle determination is estimated to be about $`1.5^{}`$ . A correction due to barometric pressure changes, which lead to small fluctuations of the atmospheric absorption, has not been made. Figure 2 displays the spectra of mean values of each atmospheric depth bin and compares with the results from EAS-TOP and KASCADE experiments.
Following fixed intensities of the experimental spectra (see sect.3.2) the average $`N_e`$ cascade development can be immediately reconstructed as shown in Figure 2. Note that the results in the range of the slant depth observed with the ANI array deviate from the exponential decrease (eq.1). That is an interesting feature which can be revealed more clearly when combining spectra accurately measured on different altitudes. In the present paper we base the formulation of the procedures estimating the attenuation on the exponential decrease (eq.1). It is our interest to explore, if this assumption applied to the ANI and KASCADE data lead to consistent results.
## 3 Procedures for inference of the attenuation length <br>from size spectra
We consider the differential and integral size spectra $`I(N_e,X)`$ and $`I(>N_e,X)`$, respectively. In addition to the basic assumption of exponential attenuation of $`N_e`$ (eq.1) a power-law dependence of the size spectrum
$$I(N_e,X)N_{e}^{}{}_{}{}^{\gamma },$$
(2)
with the spectral index $`\gamma `$ is adopted.
### 3.1 Attenuation of the intensity of fixed $`𝑵_𝒆`$: absorption length
For different fixed values of shower size $`N_e`$, on different depths in the atmosphere or/and different zenith angles of incidence, from measured spectra (see vertical dotted lines on Figure 4) we obtain several values of corresponding intensities from the equivalent depths from 700 till 1280 $`gcm^2`$. Fitting the depth dependence of the intensities by the straight line (in logarithmic scale) according to equation:
$$I(N_e,X)=I(N_e,X_0)exp\left(\frac{XX_0}{\lambda _N}\right)$$
(3)
we obtain the estimate of the absorption length $`\lambda _N`$. The absorption length can be estimated both by integral and differential spectra.
### 3.2 Constant intensity cut
The basic idea of this procedure is to compare the average size of showers which have the same rate (showers per $`m^2ssr`$) in the different bins of the zenith angle of shower incidence and different slant depth, respectively .
Considering two different depths in atmosphere $`X_1`$,$`X_2>X_0`$ the expressions of differential intensities $`I(N_e,X)`$ has the form
$$N_e(X_1)^\gamma exp\left[\left(\gamma 1\right)\frac{X_1X_0}{\mathrm{\Lambda }}\right]=N_e(X_2)^\gamma exp\left[\left(\gamma 1\right)\frac{X_2X_0}{\mathrm{\Lambda }}\right]$$
(4)
With simple transformations we obtain:
$$\mathrm{\Lambda }_{diff}(I)=\frac{\gamma 1}{\gamma }\frac{X_2X_1}{ln\left(\frac{N_e(X_1)}{N_e(X_2)}\right)}$$
(5)
The attenuation lengths, obtained by integral spectra do not depend explicitly on spectral index:
$$\mathrm{\Lambda }_{int}(I)=\frac{X_2X_1}{ln\left(\frac{N_e(X_1)}{N_e(X_2)}\right)}$$
(6)
Practically the estimate of the attenuation length is obtained by fitting the $`N_e`$ dependence on the depth in atmosphere by the straight line according to the equation (1). The sequence of $`N_e`$ values is obtained according to the fixed values of the flux intensity, selected from the interpolation of the differential or integral size spectra.
For each $`N_e`$ value, the slope index $`\gamma `$ used in equation 5, is obtained by averaging over all used slant depths. Selecting equal intensities ($``$ primary energies) corresponding to different shower sizes $`N_e`$ and different depths the value of $`\mathrm{\Lambda }_{diff}(I)`$ is estimated. Intensity values from $`10^9`$ to $`5.10^6`$ were used for CIC method.
### 3.3 Attenuation of the size of the knee
A special variant of the constant intensity cut is to follow the decrease of the shower size at a constant primary energy in the size spectrum. Assuming that the knee phenomenon reflects a feature of the primary flux, the variation shower size at the knee with the zenith angle provides the possibility to extract the attenuation length.
Considering the assigned knee position of the data from various experiments, differences within 30% are noticed for all X-bins.
The knee positions obtained by the differential and integral spectra are a bit shifted to the smaller $`N_e`$ values (see Figure 4). The shift is approximately uniform over all investigated depths interval, therefore the estimates of the attenuation length by the differential and integral size spectra are very close to each other.
### 3.4 The relation between the absorption and attenuation length
We consider the quantity $`I(N_e,X)dN_e`$ \- the number of EAS at the depth $`X`$ which comprise $`N_e`$ to $`N_e+dN_e`$ particles:
$$I(N_e,X)dN_eN_e^\gamma exp\left[\left(\gamma 1\right)\frac{XX_0}{\mathrm{\Lambda }}\right]dN_e$$
(7)
With eq.3 we obtain:
$$\mathrm{\Lambda }_{diff}(N_e)=(\gamma (N_e)1)\lambda _N,$$
(8)
where, $`\gamma (N_e)`$ is the differential size spectra index (here we indicate the $`N_e`$ dependence of the slope index explicitly). For the integral spectra:
$$\mathrm{\Lambda }_{int}(N_e)=\gamma (N_e)\lambda _N,$$
(9)
where, $`\gamma (N_e)`$ is integral size spectra index.
For the evaluation of the inelastic cross section and for comparison of the three methods described above we propose to use the calculated values of the attenuation length $`\mathrm{\Lambda }`$ (instead of using absorption length $`\lambda _N`$). The attenuation of the number of particles in the individual cascade is more directly connected with the characteristics of the strong interaction and is independent from the parameters of the cosmic ray flux incident on the atmosphere. In turn the absorption length, i.e. the attenuation of the CR flux intensity, reflects also characteristics of the primary flux and is dependent on the change of the slope of the spectra.
### 3.5 Estimate of the inelastic cross section
The inelastic cross sections, of the primary nuclei with atmosphere nuclei is related by :
$$\sigma _{Aair}^{inel}(mbarn)=\frac{2.4110^4}{\lambda _A(gcm^2)},$$
(10)
where A denotes the primary nuclei. The quantity $`\lambda _A`$ is the interaction length of the A-nucleus in the atmosphere (note: in some publications the interaction length is denoted by $`\lambda _N`$, where N is primary nuclei, in contrast in this paper N is reserved for the shower size). The interaction length $`\lambda _A`$ is related with the absorption length $`\mathrm{\Lambda }_A`$ by
$$\lambda _A=K(E)\mathrm{\Lambda }_A$$
(11)
The K(E) coefficient reflects peculiarities of the strong interaction model used for simulation. The value of the parameter K has to be determined by simulations of the EAS development in the atmosphere. Such studies require the development of procedures for the selection of EAS initiated by primaries of a definite type (see for example in ).
## 4 Application to the data
The mean values of the attenuation lengths obtained by various methods from data of the ANI and KASCADE installations, as well as for the joint ANI & KASCADE data by the differential ($`\mathrm{\Lambda }_{diff}`$) and integral spectra ($`\mathrm{\Lambda }_{int}`$) are compiled in the Tables 2,3,4.
The alternative estimates of the attenuation length reflect the inherent uncertainties of the methods and the statistical errors, as well as the fluctuations of cascade development in the atmosphere, the
energy dependence of the inelastic cross section and possible changes in mass composition. As obvious in Figure 2, the values corresponding to the minimal equivalent depths of used MAKET ANI data, deviate significantly from the exponential decrease. The observations reflects the flattening of the cascade curve just after the shower maximum in the altitude $`500600gcm^2`$. Due to these features the attenuation lengths calculated by MAKET ANI data appear to be significantly larger than those derived for the KASCADE data (Tables 2, 3).
Therefore, for the combined analysis of the KASCADE and ANI data we omitted the first and the second zenith angle bins of MAKET ANI and calculate the attenuation lengths by the remaining 9 (minimal equivalent depth $`758gcm^2`$) and 8 (minimal equivalent depth $`816gcm^2`$) angular bins. The dependences of estimated values of attenuation length on the shower size and flux intensity for different amount of the angular bins used, are displayed in Figures 6 (note, that higher intensities on the X axes correspond to the lower primary energies) and 6. The attenuation length estimates obtained from the differential and integral spectra agree fairly well. The results of both CIC and recalculation from absorption length agree within the error bars. The results obtained by the ”attenuation of knee position” are larger for MAKET ANI and KASCADE. As pointed out by S. Ostapchenko it is the consequence of the large EAS fluctuations with the tendency to shift the knee position to the lower energies (and correspondingly to higher fluxes) in a way to ”slow down” the cascade curve attenuation.
Well below the shower development maximum starting from $`816gcm^2`$ KASCADE and MAKET ANI data could be fitted with one decay parameter (see Figure 7). There is a concentration of the knee positions on the curve showing the dependence of the attenuation of the flux intensity ($``$ primary energy). In turn, the curve displaying the dependence of the attenuation length on the shower size demonstrates a rather large dispersion of the ”knee positions”. These observations in size and energy scales may be interpreted as an indication of the astrophysical nature of the knee phenomenon.
## 5 Conclusion
Experimental studies of EAS characteristic like the depth of the shower maximum $`X_{max}`$, the elongation rate $`dX_{max}/dlog_{10}E`$ and the attenuation length $`\mathrm{\Lambda }`$ are of particular importance, since they map rather directly basic features of the hadronic interaction. Strictly, however, the interpretation of these quantities in terms of hadronic cross sections cannot bypass the necessity of detailed calculations of the shower development. Nevertheless these type of EAS quantities, if compared with Monte Carlo simulation results, provide stringent tests of the interaction model ingredients of the simulations.
The recent results of various experimental installations are sufficiently accurate to enable relevant studies of this kind, and combining the data from arrays situated on different altitudes (like MAKET ANI and KASCADE) allows a large span in the atmospheric slant depth for reconstructing the development of the charge particle size. In fact such studies, if using a sufficiently large data sample, could be continued in a more detailed manner by separating the muon component and taking into account the deviations from the exponential shape of the cascade decline. The penetrating muon component contributes with smaller attenuation to the development of the considered charged particle component, but hardly with an exponential degrading (according to eq.1). Actually by use of methods in progress to isolate different primary groups (”pure nuclear beams”) of the size spectra , these kind of interaction studies would get of extreme interest.
Acknowledgment
This publication is based on experimental results of the ANI collaboration. The MAKET ANI detector installation has been set up as collaborative project of the Yerevan Physics Institute, Armenia and the Lebedev Institute, Moscow. The continuous contributions and assistance of the Russian colleagues in operating the detector installation and in the data analyses are gratefully acknowledged. In particular, we thank S. Nikolski and V. Romakhin for their encouraging interest and useful discussions.
First perspectives of combined considerations of the KASCADE and ANI experimental data have been discussed in 1998 during the ANI-98 workshop in the cosmic ray observatory station Nor-Amberd of Mt.Aragats (Armenia). The MAKET ANI group would like to thank the German colleagues for stimulating discussions and encouragement, in particular H. Rebel for his numerous valuable comments and interesting suggestions to the topic of this paper. We acknowledge the useful discussions with K.-H. Kampert, H. Klages and R. Glasstetter. The suggestions of S. Ostapchenko are highly appreciated.
The work has been partly supported by the research grant No.96-752 of the Armenian Government and by the ISTC project A116. The assistance of the Maintenance Staff of the Aragats Cosmic Ray Observatory in operating the MAKET ANI installation is highly appreciated. |
warning/0002/hep-th0002179.html | ar5iv | text | # References
1. Introduction
The Casimir effect is one of the most interesting manifestations of nontrivial properties of the vacuum state in quantum field theory , . Since it’s first prediction by Casimir in 1948 this effect is investigated for various cases of boundary geometries and various types of fields. The Casimir effect can be viewed as a polarization of vacuum by boundary conditions. Another type of vacuum polarization arises in the case of external gravitational field. In this paper we shall consider a simple example when these two types of sources for vacuum polarization are present. We investigate the vacuum expectation values of the energy-momentum tensor for conformally coupled scalar field in the standard parallel plate geometry with Dirichlet boundary conditions and on background of planar static domain wall case. It has been shown in and , that the gravitational field of the vacuum domain wall with a source of the form
$$T_\mu ^\nu =\sigma \delta (x)diag(1,\mathrm{\hspace{0.17em}0},\mathrm{\hspace{0.17em}1},\mathrm{\hspace{0.17em}1})$$
(1)
does not correspond to any exact static solution of Einstein equations (on domain wall solutions of Einstein-scalar-field equations see ). However the static solutions can be constructed in presence of an additional background energy-momentum tensor. Such a type solution has been fond in . First we calculate the vacuum expectation values of energy-momentum tensor by using the mode sums. We obtain the result as a direct sum of two terms: boundary term and term which presents the vacuum polarization in the domain wall geometry in the case of absence of boundaries. It is shown that boundary part of total energy between the plates and corresponding pressures on plates are related by standard thermodynamical relation. Then we show that corresponding properties can be obtained by using the conformal properties of the problem.
2. Vacuum expectation values of energy-momentum tensor
In this paper we shall consider the conformally coupled real scalar field $`\varphi `$, which satisfies
$$(\mathrm{}+\frac{1}{6}R)\varphi =0,\mathrm{}=\frac{1}{\sqrt{g}}_\mu (\sqrt{g}g^{\mu \nu }_\nu ),$$
(2)
and propagates on background of gravitational field generated by static domain wall solution from . The corresponding metric has the form
$$ds^2=A^{2\alpha }(dt^2dy^2dz^2)A^{2(\alpha +\gamma +1)}dx^2,A=A(x)1+K|x|,$$
(3)
where $`\alpha >0,K>0`$. The equation (3) describes a planar domain wall with energy-momentum tensor $`T_\mu ^\nu =\delta (x)diag(1,0,1,1)`$ in the background field with Ricci tensor
$$R_\mu ^\nu =\alpha (\gamma 2\alpha )K^2A^{2(\alpha +\gamma )}diag(1,\frac{3\gamma }{\gamma 2\alpha },1,1)$$
(4)
Note that for the energy density of the background to be positive we must have $`\gamma <\alpha /2`$. In what follows as a boundary configuration we shall consider two plates parallel to each other and to domain wall, with $`x`$ coordinates equal to $`x_1`$ and $`x_2`$ (to be definit we shall consider right half space of domain wall geometry $`x_1,x_2>0`$). For the points on plate the scalar field obeys Dirichlet boundary condition
$$\varphi (x=x_1)=\varphi (x=x_2)=0$$
(5)
The quantization of field (2) on background of equation (3) is standard. Let $`\varphi _\alpha ^{(\pm )}(x)`$ be complete set of orthonormalized positive and negative frequency solutions to the field equation (2), obeying boundary conditions (5). The canonical quantization can be done by expanding the general solution of (2) in terms of $`\varphi _\alpha ^{(\pm )}`$,
$$\varphi =\underset{\alpha }{}(\varphi _\alpha ^+a_\alpha +\varphi _\alpha a_\alpha ^{(+)})$$
(6)
and declearing the coeficients $`a_\alpha `$ ,$`a_\alpha ^+`$ as operators satisfying standard commutation relation for bosonic fields. The vacuum state $`|0>`$ is defiend as $`a_\alpha |0>=0`$. This state is different from the vacuum state for domain wall geometry without bondaries, $`|\overline{0}>`$. To investigate effects due to the presence of boundaries we shall consider vacuum expectation values of energy-momentum tensor operator, $`<0|T_{\mu \nu }|0>`$. By substituting the expansion (6) and using the definition of vacuum state it can be easly seen that
$$<0|T_{\mu \nu }|0>=\underset{\alpha }{}T_{\mu \nu }\{\varphi _\alpha ^{(+)},\varphi _\alpha ^{()}\}$$
(7)
Here on the rhs the bilinear form $`T_{\mu \nu }\{\varphi ,\psi \}`$ is determined by the classical energy-momentum tensor for conformally coupled scalar field (see for exampel ). To calculate the vacuum expectation values by (7) we need the explicit form of eigenfunctions $`\varphi _\alpha ^{(\pm )}`$. For this case the metric and boundary conditions are static and translation invariant in directions parallel to the domain wall. It follows from here that the corresponding part has standard plane wave structure:
$$\varphi _\alpha ^\pm =\phi (x)\mathrm{exp}\left[\pm i(k_yy+k_zz\omega t)\right]$$
(8)
The equation for $`\phi (x)`$ is obtained from the field equation (2) and for domain wall metric (3) has the form
$$\phi ^{^{\prime \prime }}(x)+(\gamma +12\alpha )Ksgn(x)A^1\phi ^{^{}}+[\alpha (\alpha \gamma )K^2A^2+k_x^2A^{2(\gamma +1)}]\phi =0$$
(9)
with $`k_{x}^{}{}_{}{}^{2}=\omega ^2k_{t}^{}{}_{}{}^{2}`$, $`k_{t}^{}{}_{}{}^{2}=k_{y}^{}{}_{}{}^{2}+k_{z}^{}{}_{}{}^{2}`$. We shall consider the region between plates. The solution to equation (9) in this region obeying boundary conditions (5) is
$$\phi (x)=constA^\alpha \mathrm{sin}(k_xv),k_x=\frac{n\pi }{a},n=1,2,\mathrm{}$$
(10)
with the relation
$$v(x)=\frac{[(1+K|x_1|)^\gamma (1+K|x|)^\gamma ]}{K\gamma },a=v(x_2)$$
(11)
By using this relations and normalizing the eigenfunctions by standard way one obtains
$$\varphi _\beta ^{(\pm )}(t,\stackrel{}{r})=\frac{A^\alpha }{2\pi \sqrt{\omega a}}\mathrm{sin}(k_xv)e^{\pm i(k_yy+k_zz\omega t)},\beta =(n,k_y,k_z),\omega ^2=\left(\frac{n\pi }{a}\right)^2+k_t^2$$
(12)
Before to start specific calculation with formula (7) it is convenient to present energy-momentum tensor for conformally coupeld scalar field in the form
$$T_{ik}=_i\varphi _k\varphi \frac{1}{6}(_i_k+\frac{1}{2}g_{ik}\mathrm{}+R_{ik})\varphi ^2$$
(13)
By using the equation (7) with this form of energy-momentum tensor and with eigenmodes (12) we receive
$$<0|T_{ik}|0>=<0|_i\varphi _k\varphi |0>\frac{1}{6}(_i_k+\frac{1}{2}g_{ik}\mathrm{}+R_{ik})<0|\varphi ^2|0>.$$
(14)
It is convenient first to calculate the quantity
$$<0|\varphi (x)\varphi (x^{})|0>=$$
$$\frac{A^\alpha (x)A^\alpha (x^{})}{4\pi ^2a}d^2k_t\underset{n=1}{\overset{\mathrm{}}{}}\frac{1}{\omega }\mathrm{sin}(k_xv(x))\mathrm{sin}(k_xv(x^{}))e^{i[k_y(yy^{})+k_z(zz^{})\omega (tt^{})]}$$
(15)
Using Abel-Plana summation formula
$$\underset{n=1}{\overset{\mathrm{}}{}}f(n)=_0^{\mathrm{}}f(x)𝑑x\frac{1}{2}f(0)+i_0^{\mathrm{}}\frac{f(ix)f(ix)}{e^{2\pi x}1}𝑑x$$
(16)
to sum over $`n`$, and after calculating arising integrals one obtains
$$<0|\varphi (x)\varphi (x^{})|0>=\frac{A^\alpha (x)A^\alpha (x^{})}{8a^2}\frac{\mathrm{sinh}(u)}{u}$$
$$\{[\mathrm{cosh}(u)\mathrm{cos}(\frac{\pi }{a}(v(x)v(x^{})))]^1[\mathrm{cosh}(u)\mathrm{cos}(\frac{\pi }{a}(v(x)+v(x^{})))]^1\}$$
(17)
where $`u=\frac{\pi }{a}[(yy^{})^2+(zz^{})^2(tt^{})^2]^{1/2}`$. The vacuum expectation value of energy-momentum tensor may be found now by using the relation
$$<0|T_{ik}(x)|0>=lim_{xx^{}}\widehat{T}_{ik}<0|\varphi (x)\varphi (x^{})|0>$$
(18)
where the form of the second order operator $`\widehat{T}_{ik}`$ is obvious from (14):
$$\widehat{T}_{ik}=_i_k^{^{}}\frac{1}{6}(_i_k+\frac{1}{2}g_{ik}\mathrm{}+R_{ik})$$
(19)
and $`_k^{^{}}=/x_k^{^{}}`$. The vacuum expectation value of course is infinite. This divergencies come from the first term in figure brackets of (17). In this paper we are mainly interested in quantum effects due to the existence of boundaries. Let $`|\overline{0}>`$ be the vacuum state for conformally coupled scalar field in the case of absence of boundaries. Let us consider the difference
$$<T_{ik}^{(b)}(x)>=<0|T_{ik}(x)|0><\overline{0}|T_{ik}(x)|\overline{0}>=$$
$$lim_{xx^{}}\widehat{T}_{ik}[<0|\varphi (x)\varphi (x^{})|0><\overline{0}|\varphi (x)\varphi (x^{})|\overline{0}>]$$
(20)
This quantity describes the boundary contribution to the polarization of the vacuum and is finit. To see this note that the expression for $`<\overline{0}|\varphi (x)\varphi (x^{})|\overline{0}>`$ may be obtained from (17) by taking $`a\mathrm{}`$ and has the following form
$$<\overline{0}|\varphi (x)\varphi (x^{})|\overline{0}>=\frac{A^\alpha (x)A^\alpha (x^{})}{4\pi ^2}\left\{[v(x)v(x^{})]^2+(yy^{})^2+(zz^{})^2(tt^{})^2\right\}$$
(21)
As it can be easily seen the divergences in (17) and (21) cancel in calculating boundary contribution (20). Substitving (17) and (21) into (20) after some calculations and arising the second index one obtains for the region between plates
$$<T_i^{(b)k}>=\frac{\pi ^2A^{4\alpha }}{1440a^4}diag(1,3,1,1),x_1xx_2$$
(22)
and zero outside of this region. Here $`a`$ is defined by equation (11). Note that $`a`$ is different from the proper distance $`a_p`$ between the plates,
$$a_p=_{x_1}^{x_2}g_{11}𝑑x_1=\left[(1+Kx_1)^{\alpha \gamma }(1+Kx_2)^{\alpha \gamma }\right]/k(\alpha +\gamma )$$
(23)
In calculating (22) all terms which come from derivatives $`A/x`$ and $`^2A/x^2`$ and are proportional to $`K`$ and $`K^2`$ are cancelled. As we shall see this is direct consequence of the conformal properties of metric and field under consideration. In the limit of no geravitation, $`K0`$ from (22) we obtain standard Casimir result for parallel plate configuration. In the case of scalar field and $`a=x_2x_1`$. By using this result the vacuum expectation value for total energy-momentum tensor can be written as
$$<0|T_{ik}|0>=<T_{ik}^{(b)}>+<T_{ik}^{(g)}>$$
(24)
where the second summand of right-hand side is the part describing the polarization of scalar vacuum by domain wall geravitational field in the case of absence of boundaries. All divergences are contained in this part. The corresponding regularization can be done by using the standard methods of quantum field theory in curved space-time (see, for example ). Most simply this can be done by using the conformally flatness of the metric (refmetric) (see and below). In this cas the anomalous trace determines the total energy-momentum tensor (see ) and the regular part of purely gravity contribution to (24) is equal to
$$<T_{ik}^{(g)}>=reg<\overline{0}|T_{ik}^{(g)}|\overline{0}>=\frac{\alpha +4\gamma }{2880\pi ^2}\alpha (3\gamma ^22\alpha ^2)A^{4(\alpha +\gamma )}diag(1,\frac{3\alpha }{\alpha +4\gamma },1,1)$$
now from (24) and (S0.Ex4) it follows that regularized total energy-momentum tensor in the region between plates are given by
$$<T_i^k>=<T_i^{(b)k}>+<T_i^{(g)k}>$$
(25)
where boundary, $`<T_i^{(b)k}>`$, and geravitational, $`<T_i^{(g)k}>`$, parts are determined by (22) and (S0.Ex4) respectivly. In the regions $`x<x_1`$ and $`x>x_2`$ the boundary part is zero and only gravitational polarization part remains. The forces acting on plates are determined by boundary part only. The effective pressure created by gravitational part in (25) is equal to
$$p_{g1}=<T_1^{(g)1}>=\frac{\alpha ^2K^4(3\gamma ^22\alpha ^2)}{960\pi ^2a^4}A^{4(\alpha +\gamma )}(x)$$
(26)
and is the same from the both sides of the plates, and hence leads to the zero effective force. Vacuum boundary part pressures acting on plates are
$$p_{b1}^{(1,2)}=p_{b1}(x=x_{1,2})=<T_1^{(b)1}(x=x_{1,2})>=\frac{\pi ^2A^{4\alpha }(x_{1,2})}{480a^4}$$
(27)
and have attractive nature. The boundary part of the total energy between the plates can be found by standard way:
$$E_b=_{x_1}^{x_2}𝑑x𝑑y𝑑z\sqrt{g}<T_0^{(b)0}>=\frac{\pi ^2}{1440a^3}𝑑y𝑑z$$
(28)
where we have used the definition of $`a`$ in accordance with (11). It can be easily seen that total energy (28) and pressures (27) are connected by standard thermodynamical relation
$$p_{b1}(x_1)=\frac{dE}{dV_1}|_{x_2=const}=\frac{dE}{dx_1dydz}A^{4\alpha +\gamma +1}(x_1)|_{x_2=const},dV_1=A^{4\alpha \gamma 1}(x_1)dx_1dydz$$
(29)
and similar relation for $`p_{b1}(x_2)`$.
Only in the case of conformally invariant fields the eigenmodes have simple form (12) in the domain wall geravitational field (3). This is a direct consequnce of the conformally equivalence of the metric (3) to the Minkowskian one. Indeed by the coordinate transformation
$$x=f(X),(1+K|f(X)|)^\gamma =1K\gamma |X|$$
(30)
It can be seen that the metric (3) takes a manifastly conformally flat form .
$$ds^2=(1K\gamma |X|)^{\frac{2\alpha }{\gamma }}(dt^2dX^2dy^2dz^2)$$
(31)
Let $`<T_{ik}^{(M)}>`$ be the regularized standard energy-momentum tensor for a conformally coupled scalar field in the case of parallel plate counfiguration in flat space-time with metric $`\eta ik`$
$$<T_i^{(M)k}>=\frac{\pi ^2}{1440a^4}diag(1,3,1,1)$$
(32)
where $`a=x_2x_1`$ and $`x_i`$ are coordinates for the plates. Using the standard relation between the energy-momentum tensor for conformally coupled situation
$$<T_i^k[\stackrel{~}{g}_{lm}]>=(\frac{\eta }{\stackrel{~}{g}})^{\frac{1}{2}}<T_i^{k(M)}[\eta _{lm}]>\frac{1}{2880\pi ^2}\left[\frac{1}{6}^{(1)}\stackrel{~}{H}_i^k^{(3)}\stackrel{~}{H}_i^k\right]$$
(33)
(the standard notations $`{}_{}{}^{(1,3)}H_{i}^{k}`$ for some combinations of curvature tensor components see ), where tilde notes the quantities in the coordinate system $`(t,X,y,z,)`$ with metric (31). By making the transformation (see (30)) to the initial coordinate system $`(t,x,y,z)`$ from (33) we receive the result (25), where $`a`$ is expressed via the coordinates $`x_1,x_2`$ of plates in system (3) by relation (11).
3. Concluding remarks
In this paper we calculate the Casimir energy for conformally invariant scalar field in the standard parallel plate, on background of planar static domain wall. The boundary conditions over scalar field on the plate are Dirichlet boundary conditions. For calculating the vacuum expectation values of the energy-momentum tensor we use the mode sums method and Abel-Plana summation formula. The result contains two terms, one comes from the boundary conditions and the other one from the effect of gravitation over the vacuum of scalar filed. The quantity which reflects the effect of bondary conditions is finit. If we write this term in flat space-time limit we obtain the standard result of casimir effect for parallel plates. All divergences are in the part which describs the polarization of scalar vacuum by domain wall background in case of absence of boundaries. The effective pressure created by gravitational part is the same for both sides of the plates and hence leads to the zero effective force, but the vacuum boundary part pressures acting on plates is attractive. When the total energy between the plates are found, one can readily see that the total energy and pressures are connected through the standard thermodynamical relations. These results can easily be obtained using the conformal properties of the metric and the scalar field.
Acknowledgment
M. R. Setare would like to thank R. Mansouri for his valuable hints and comments. He also thanks F.H. Jafarpour for his help in preparing the Latex file. |
warning/0002/hep-ph0002154.html | ar5iv | text | # Axion Searches
## 1 Introduction
The axion was postulated nearly two decades ago to explain why the strong interactions conserve $`P`$ and $`CP`$ in spite of the fact that the weak interactions violate those symmetries. Consider the Lagrangian of QCD:
$`_{QCD}={\displaystyle \frac{1}{4}}G_{\mu \nu }^aG^{a\mu \nu }+{\displaystyle \underset{j=1}{\overset{n}{}}}[\overline{q}_j\gamma ^\mu iD_\mu q_j`$
$`(m_jq_{Lj}^{}q_{Rj}+\text{h.c.})]+{\displaystyle \frac{\theta g^2}{32\pi ^2}}G^a_{\mu \nu }\stackrel{~}{G}^{a\mu \nu }.`$ (1)
The last term is a 4-divergence and hence does not contribute in perturbation theory. That term does however contribute through non-perturbative effects associated with QCD instantons . Such effects can make the physics of QCD depend upon the value of $`\theta `$. Using the Adler-Bell-Jackiw anomaly , one can show that $`\theta `$ dependence must be present if none of the current quark masses vanishes. Indeed otherwise QCD would have a $`U_A(1)`$ symmetry and would predict the mass of the $`\eta ^{}`$ pseudo-scalar meson to be less than $`\sqrt{3}m_\pi 240`$ MeV , contrary to observation. Using the anomaly, one can further show that QCD depends upon $`\theta `$ only through the combination of parameters:
$$\overline{\theta }=\theta \text{arg}(m_1,m_2,\mathrm{}m_n).$$
(2)
If $`\overline{\theta }0`$, QCD violates $`P`$ and $`CP`$. The absence of $`P`$ and $`CP`$ violations in the strong interactions therefore places an upper limit upon $`\overline{\theta }`$. The best constraint follows from the experimental bound on the neutron electric dipole moment which yields: $`\overline{\theta }<10^9`$.
The question then is: why is $`\overline{\theta }`$ so small? In the Standard Model of particle interactions, the quark masses originate in the electroweak sector which violates $`P`$ and $`CP`$. There is no reason why the overall phase of the quark mass matrix should exactly match the value of $`\theta `$ from the QCD sector to yield $`\overline{\theta }<10^9`$. In particular, if $`CP`$ violation is introduced in the manner of Kobayashi and Maskawa , the Yukawa couplings that give masses to the quarks are arbitrary complex numbers and hence arg det $`m_q`$ and $`\overline{\theta }`$ are expected to be of order one.
The problem why $`\overline{\theta }<10^9`$ is usually referred to as the “strong $`CP`$ problem”. The existence of an axion solves this problem in a simple manner which is rich in implications for experiment, for astrophysics and for cosmology. There are two alternative solutions however. The first alternative is to set $`m_u=0`$. This removes the $`\theta `$-dependence of QCD and thus solves the strong $`CP`$ problem. The well-known calculation of the pseudo-scalar meson masses in lowest order of chiral perturbation theory yields $`m_u4`$ MeV, which is incompatible with $`m_u=0`$. This calculation also predicts the successful Gell-Mann - Okubo relation among the pseudo-scalar masses squared. It is possible to have $`m_u=0`$ by invoking second order effects . This a reasonable proposition because $`m_s`$ happens to be of order the QCD scale. However, when second order effects are included , the Gell-Mann - Okubo relation is in general violated. Thus the price for having $`m_u=0`$ through higher order effects is that the Gell-Mann - Okubo relation becomes an accident. The second alternative solution to the strong CP problem is to assume that $`CP`$ and/or $`P`$ is spontaneously broken but is otherwise a good symmetry. In this case, $`\overline{\theta }`$ is calculable and may be arranged to be small .
In addition, it is worth emphasizing that the strong $`CP`$ problem need not be solved in the low energy theory. Indeed, as Ellis and Gaillard pointed out, if in the standard model $`\overline{\theta }=0`$ near the Planck scale then $`\overline{\theta }10^9`$ at the QCD scale.
Peccei and Quinn proposed to solve the strong $`CP`$ problem by postulating the existence of a global $`U_{PQ}(1)`$ quasi-symmetry. $`U_{PQ}(1)`$ must be a symmetry of the theory at the classical (i.e., at the Lagrangian) level, it must be broken explicitly by those non-perturbative effects that make the physics of QCD depend upon $`\theta `$, and finally it must be spontaneously broken. The axion is the quasi-Nambu-Goldstone boson associated with the spontaneous breakdown of $`U_{PQ}(1)`$. One can show that, if a $`U_{PQ}(1)`$ quasi-symmetry is present, then
$$\overline{\theta }=\theta arg(m_1\mathrm{}m_n)\frac{a(x)}{f_a},$$
(3)
where $`a(x)`$ is the axion field and $`f_a=v/N`$ is called the axion decay constant. $`v`$ is the vacuum expectation value which spontaneously breaks $`U_{PQ}(1)`$ and $`N`$ is an integer which expresses the color anomaly of $`U_{PQ}(1)`$. Axion models have $`N`$ degenerate vacua . The non-perturbative effects that make QCD depend upon $`\overline{\theta }`$ produce an effective potential $`V(\overline{\theta })`$ whose minimum is at $`\overline{\theta }=0`$. Thus, by postulating an axion, $`\overline{\theta }`$ is allowed to relax to zero dynamically and the strong $`CP`$ problem is solved.
The properties of the axion can be derived using the methods of current algebra . The axion mass is given in terms of $`f_a`$ by
$$m_a0.6eV\frac{10^7GeV}{f_a}.$$
(4)
All the axion couplings are inversely proportional to $`f_a`$. Of particular interest here is the axion coupling to two photons:
$$_{a\gamma \gamma }=g_\gamma \frac{\alpha }{\pi }\frac{a(x)}{f_a}\stackrel{}{E}\stackrel{}{B}$$
(5)
where $`\stackrel{}{E}`$ and $`\stackrel{}{B}`$ are the electric and magnetic fields, $`\alpha `$ is the fine structure constant, and $`g_\gamma `$ is a model-dependent coefficient of order one. $`g_\gamma =0.36`$ in the DFSZ model whereas $`g_\gamma =0.97`$ in the KSVZ model . The coupling of the axion to a spin 1/2 fermion $`f`$ has the form:
$$_{a\overline{f}f}=ig_f\frac{m_f}{v}a\overline{f}\gamma _5f$$
(6)
where $`g_f`$ is a model-dependent coefficient of order one. In the KSVZ model the coupling to electrons is zero at tree level. Models with this property are called ’hadronic’.
## 2 Constraints from laboratory searches and astrophysics
The searches for the axion in high energy and nuclear physics experiments are discussed in the reviews by Kim and by Peccei . A complete list of experiments can be found in the Review of Particle Physics . If the axion is heavier than 1 MeV and decays quickly into $`e^+e^{}`$ (lifetime of order $`10^{11}`$ sec or less), then it is ruled out by negative searches for rare particle decays such as $`\pi ^+a(e^+e^{})e^+\nu _e`$ . The rate of this reaction follows simply from the mixing of $`a`$ and $`\pi ^0`$ and the known decay $`\pi ^+\pi ^0e^+\nu _e`$, and hence it is almost completely model-independent . Alternatively, the axion is long lived ($`10^{11}`$ sec or more). In this case it is severely constrained by negative searches in beam dumps. In a beam dump, axions are produced by many different processes involving independent couplings, such as $`a\pi ^0`$, $`a\eta `$ and $`a\eta ^{}`$ mixing, the axion couplings to two photons and to two gluons, and the couplings to quarks and gluons. It is difficult to calculate these processes with precision at all energies. On the other hand, the many processes add up incoherently and they can not all vanish. Thus one can give a reliable estimate for the total production ($`p+Na+X`$, or $`e+Na+X`$) and interaction ($`a+NX`$) cross-sections. Many such searches have been carried out. They rule out axions with mass larger than about 50 keV.
The astrophysical contraints on the axion are described in detail in the reviews by Turner and Raffelt . Axions are emitted by stars in a variety of processes such as Compton-like scattering ($`\gamma +ea+e`$), axion bremstrahlung ($`e+NN+e+a`$) and the Primakoff process ($`\gamma +NN+a`$). The rate at which a star burns its nuclear fuel is limited by the rate at which it can loose the energy produced. Emission of light weakly coupled bosons, such as the axion, allows a star to radiate energy efficiently because such particles can escape the star at once (without rescattering) whereas photons are emitted only from the stellar surface . Thus the existence of an axion accelerates stellar evolution, which may be inconsistent with observation. The longevity of red giants rules out the mass range 200 keV $`>m_a>0.5`$ eV for hadronic axions. Above 200 keV the axion is too heavy to be copiously emitted in the thermal processes taking in red giants, whereas below 0.5 eV it is too weakly coupled. For axions with a large coupling to electrons \[$`g_e=0(1)`$ in Eq. 6\] the ruled out range can be extended to 200 keV $`>m_a>10^2`$ eV because axion emission through the Compton-like process $`\gamma +ea+e`$ cools the helium core to such an extent as to prevent the onset of helium burning .
Finally the range 2 eV $`>m_a>310^3`$ eV is ruled out by Supernova 1987a . The constraint follows from the fact that the duration of the associated neutrino events in the large underground proton decay detectors is consistent with theoretical expectations based on the premise that the collapsed supernova core cools by emission of neutrinos. If the axion mass is in the above-mentioned range, the core cools instead by axion emission and the neutrino burst is excessively shorthened. The supernova constraint is quite axion model-independent because the axions are emitted by axion bremstrahlung in nucleon-nucleon scattering ($`N+NN+N+a`$) and the relevant couplings follow simply from the mixing of the axion with the $`\pi ^0`$ which is a general feature of axion models.
When the limits from laboratory searches are combined with the astrophysical contraints, all of the axion mass range down to approximately $`310^3`$ eV is ruled out.
## 3 Axion cosmology
The implications of the existence of an axion for the history of the early universe may be briefly described as follows. At a temperature of order $`v`$, a phase transition occurs in which the $`U_{PQ}(1)`$ symmetry becomes spontaneously broken. This is called the PQ phase transition. At these temperatures, the non-perturbative QCD effects which produce the effective potential $`V(\overline{\theta })`$ are suppressed , the axion is massless and all values of $`a(x)`$ are equally likely. Axion strings appear as topological defects. One must distinguish two cases: 1) inflation occurs with reheat temperature higher than the PQ transition temperature (equivalently, for our purposes, inflation does not occur at all) or 2) inflation occurs with reheat temperature less than the PQ transition temperature. In case 2 the axion field gets homogenized by inflation and the axion strings are ’blown’ away.
When the temperature approaches the QCD scale, the potential $`V(\overline{\theta })`$ turns on and the axion acquires mass. There is a critical time, defined by $`m_a(t_1)t_1=1`$, when the axion field starts to oscillate in response to the turn-on of the axion mass. The corresponding temperature $`T_11`$ GeV . The initial amplitude of this oscillation corresponds to how far from zero the axion field is when the axion mass turns on. The axion field oscillations do not dissipate into other forms of energy and hence contribute to the cosmological energy density today . This contribution is called of ‘vacuum realignment’. It is further described below. Note that the vacuum realignment contribution may be accidentally suppressed in case 2 because the homogenized axion field happens to lie close to zero.
In case 1 the axion strings radiate axions from the time of the PQ transition till $`t_1`$ when the axion mass turns on. At $`t_1`$ each string becomes the boundary of $`N`$ domain walls. If $`N=1`$, the network of walls bounded by strings is unstable and decays away. If $`N>1`$ there is a domain wall problem because axion domain walls end up dominating the energy density, resulting in a universe very different from the one observed today. There is a way to avoid this problem by introducing an interaction which slightly lowers one of the $`N`$ vacua with respect to the others. In that case, the lowest vacuum takes over after some time and the domain walls disappear. There is little room in parameter space for that to happen and we will not consider this possibility further here. A detailed discussion is given in Ref. . Henceforth, we assume $`N=1`$.
In case 1 there are three contributions to the axion cosmological energy density. One contribution is from axions that were radiated by axion strings before $`t_1`$. A second contribution is from axions that were produced in the decay of walls bounded by strings after $`t_1`$ . A third contribution is from vacuum realignment .
Let me briefly indicate how the vacuum alignment contribution is evaluated. Before time $`t_1`$, the axion field did not oscillate even once. Soon after $`t_1`$, the axion mass is assumed to change sufficiently slowly that the total number of axions in the oscillations of the axion field is an adiabatic invariant. The number density of axions at time $`t_1`$ is
$$n_a(t_1)\frac{1}{2}m_a(t_1)a^2(t_1)\pi f_a^2\frac{1}{t_1}$$
(7)
where $`f_a`$ is the axion decay constant introduced earlier. In Eq. (7), we used the fact that the axion field $`a(x)`$ is approximately homogeneous on the horizon scale $`t_1`$. Wiggles in $`a(x)`$ which entered the horizon long before $`t_1`$ have been red-shifted away . We also used the fact that the initial departure of $`a(x)`$ from the nearest minimum is of order $`f_a`$. The axions of Eq. (7) are decoupled and non-relativistic. Assuming that the ratio of the axion number density to the entropy density is constant from time $`t_1`$ till today, one finds
$$\mathrm{\Omega }_a\left(\frac{0.610^5\text{ eV}}{m_a}\right)^{\frac{7}{6}}\left(\frac{200\text{ MeV}}{\mathrm{\Lambda }_{QCD}}\right)^{\frac{3}{4}}h^2$$
(8)
for the ratio of the axion energy density to the critical density for closing the universe. $`h`$ is the present Hubble rate in units of 100 km/s.Mpc. The requirement that axions do not overclose the universe implies the constraint $`m_a>610^6`$ eV.
The contribution from axion string decay has been debated over the years. The main issue is the energy spectrum of axions radiated by axion strings. Battye and Shellard have carried out computer simulations of bent strings (i.e. of wiggles on otherwise straight strings) and have concluded that the contribution from string decay is approximately ten times larger than that from vacuum realignment, implying a bound on the axion mass approximately then times more severe, say $`m_a>610^5eV`$ instead of $`m_a>610^6eV`$. My collaborators and I have done simulations of bent strings , of circular string loops and non-circular string loops . We conclude that the string decay contribution is of the same order of magnitude than that from vacuum realignment. Recently, Yamaguchi, Kawasaki and Yokoyama have done computer simulations of a network of strings in an expanding universe, and concluded that the contribution from string decay is approximately three times that of vacuum realignment. The contribution from wall decay has been discussed in detail in ref. . It is probably subdominant compared to the vacuum realignment and string decay constributions.
It should be emphasized that there are many sources of uncertainty in the cosmological axion energy density aside from the uncertainty about the constribution from string decay. The axion energy density may be diluted by the entropy release from heavy particles which decouple before the QCD epoch but decay afterwards , or by the entropy release associated with a first order QCD phase transition. On the other hand, if the QCD phase transition is first order , an abrupt change of the axion mass at the transition may increase $`\mathrm{\Omega }_a`$. If inflation occurs with reheat temperature less than $`T_{PQ}`$, there may be an accidental suppression of $`\mathrm{\Omega }_a`$ because the homogenized axion field happens to lie close to a $`CP`$ conserving minimum. Because the RHS of Eq. (7) is multiplied in this case by a factor of order the square of the initial vacuum misalignment angle $`\frac{a(t_1)}{v}N`$ which is randomly chosen between $`\pi `$ and $`+\pi `$, the probability that $`\mathrm{\Omega }_a`$ is suppressed by a factor $`x`$ is of order $`\sqrt{x}`$. This rule cannot be extended to arbitrarily small $`x`$ however because quantum mechanical fluctuations in the axion field during the epoch of inflation do not allow the suppression to be perfect .
The axions produced when the axion mass turns on during the QCD phase transition are cold dark matter (CDM) because they are non-relativistic from the moment of their first appearance at 1 GeV temperature. Studies of large scale
structure formation support the view that the dominant fraction of dark matter is CDM. Any form of CDM necessarily contributes to galactic halos by falling into the gravitational wells of galaxies. Hence, there is excellent motivation to look for axions as constituent particles of our galactic halo.
Finally, let’s mention that there is a particular kind of clumpiness which affects axion dark matter if there is no inflation after the Peccei-Quinn phase transition. This is due to the fact that the dark matter axions are inhomogeneous with $`\delta \rho /\rho 1`$ over the horizon scale at temperature $`T_1`$ 1 GeV, when they are produced at the start of the QCD phase-transition, combined with the fact that their velocities are so small that they do not erase these inhomogeneities by free-streaming before the time $`t_{eq}`$ of equality between the matter and radiation energy densities when matter perturbations can start to grow. These particular inhomogeneities in the axion dark matter are in the non-linear regime immediately after time $`t_{eq}`$ and thus form clumps, called ‘axion mini-clusters’ . They have mass $`M_{mc}10^{13}M_{}`$ and size $`l_{mc}10^{13}`$ cm.
The various constraints on the axion, from accelerator seaches, astrophysics and cosmology are summarized in Fig. 1.
## 4 The cavity detector of galactic halo axions
An electromagnetic cavity permeated by a strong static magnetic field can be used to detect galactic halo axions . The relevant coupling is given in Eq. (5). Galactic halo axions have velocities $`\beta `$ of order $`10^3`$ and hence their energies $`E_a=m_a+\frac{1}{2}m_a\beta ^2`$ have a spread of order $`10^6`$ above the axion mass. When the frequency $`\omega =2\pi f`$ of a cavity mode equals $`m_a`$, galactic halo axions convert resonantly into quanta of excitation (photons) of that cavity mode. The power from axion $``$ photon conversion on resonance is found to be :
$`P=\left({\displaystyle \frac{\alpha }{\pi }}{\displaystyle \frac{g_\gamma }{f_a}}\right)^2VB_0^2\rho _aC{\displaystyle \frac{1}{m_a}}\text{Min}(Q_L,Q_a)`$
$`=\mathrm{0.5\hspace{0.33em}10}^{26}\text{Watt}\left({\displaystyle \frac{V}{500\text{ liter}}}\right)\left({\displaystyle \frac{B_0}{7\text{ Tesla}}}\right)^2`$
$`C\left({\displaystyle \frac{g_\gamma }{0.36}}\right)^2\left({\displaystyle \frac{\rho _a}{\frac{1}{2}10^{24}\frac{g_r}{\text{cm}^3}}}\right)`$
$`\left({\displaystyle \frac{m_a}{2\pi (\text{GHz})}}\right)\text{Min}(Q_L,Q_a)`$ (9)
where $`V`$ is the volume of the cavity, $`B_0`$ is the magnetic field strength, $`Q_L`$ is its loaded quality factor, $`Q_a=10^6`$ is the ‘quality factor’ of the galactic halo axion signal (i.e. the ratio of their energy to their energy spread), $`\rho _a`$ is the density of galactic halo axions on Earth, and $`C`$ is a mode dependent form factor given by
$$C=\frac{\left|_Vd^3x\stackrel{}{E}_\omega \stackrel{}{B}_0\right|^2}{B_0^2V_Vd^3xϵ|\stackrel{}{E}_\omega |^2}$$
(10)
where $`\stackrel{}{B}_0(\stackrel{}{x})`$ is the static magnetic field, $`\stackrel{}{E}_\omega (\stackrel{}{x})e^{i\omega t}`$ is the oscillating electric field and $`ϵ`$ is the dielectric constant.
Because the axion mass is only known in order of magnitude at best, the cavity must be tunable and a large range of frequencies must be explored seeking a signal. The cavity can be tuned by moving a dielectric rod or metal post inside it. Using Eq. (8), one finds the scanning rate to perform a search with signal to noise ratio $`s/n`$:
$`{\displaystyle \frac{df}{dt}}`$ $`=`$ $`{\displaystyle \frac{12\text{GHz}}{\text{year}}}\left({\displaystyle \frac{4n}{s}}\right)^2\left({\displaystyle \frac{V}{500\text{ liter}}}\right)^2\left({\displaystyle \frac{B_0}{7\text{ Tesla}}}\right)^4`$ (11)
$``$ $`C^2\left({\displaystyle \frac{g_\gamma }{0.36}}\right)^4\left({\displaystyle \frac{\rho _a}{\frac{1}{2}10^{24}\frac{gr}{\text{cm}^3}}}\right)^2\left({\displaystyle \frac{3K}{T_n}}\right)^2`$
$``$ $`\left({\displaystyle \frac{f}{\text{GHz}}}\right)^2{\displaystyle \frac{Q_L}{Q_a}},`$
where $`T_n`$ is the sum of the physical temperature of the cavity plus the electronic noise temperature of the microwave receiver that detects the photons from $`a\gamma `$ conversion. Eq. (11) assumes that $`Q_L<Q_a`$ and that some strategies have been followed which optimize the search rate. The best quality factors attainable at present, using oxygen free copper, are of order $`10^5`$ in the GHz range.
Eq. (11) shows that a galactic halo search with the required sensitivity is feasible with presently available technology, provided the form factor $`C`$ can be kept at values of order one for a wide range of frequencies. For a cylindrical cavity and a homogeneous longitudinal magnetic field, $`C=0.69`$ for the lowest TM mode. The form factors of the other modes are much smaller. The resonant frequency of the lowest TM mode of a cylindrical cavity is $`f`$=115 MHz $`\left(\frac{1m}{R}\right)`$ where $`R`$ is the radius of the cavity. Since $`10^6\text{ eV}=2\pi `$ (242 MHz), a large cylindrical cavity is convenient for searching the low frequency end of the range of interest. To extend the search to high frequencies without sacrifice in volume, one may power-combine many identical cavities which fill up the available volume inside a magnet’s bore . This method allows one to maintain $`C=0(1)`$ at high frequencies, albeit at the cost of increasing engineering complexity as the number of cavities increases.
Pilot experiments were carried out at Brookhaven National Laboratory and at the University of Florida . These experiments used relatively small magnets and hence the limits they placed on the local axion dark matter density are not severe. However they developed the various aspects of the cavity detection technique and demonstrated its feasibility.
Second generation experiments are presently under way at Lawrence Livermore National Laboratory (LLNL) and at Kyoto University . The LLNL experiment is similar in concept to the UF pilot experiment but uses a much larger magnet ($`B_0^2V=12T^2m^3`$). It is well engineered and runs with a near 100% duty cycle . The results from its first year of running are reported in ref. . The exclusion plot is shown in Fig. 2. By definition, $`g_{a\gamma \gamma }=\frac{\alpha }{\pi }\frac{g_\gamma }{f_a}`$. The limits shown assume that the local halo density, estimated to be $`7.510^{25}`$g/cm<sup>3</sup> , is entirely in axions. The experiment has ruled out the hypothesis that 100% of the local halo density is in KSVZ axions with mass in the range shown. Since then, the frequency range 500-800 MHz ($`2.1m_a3.3\mu `$eV) has been searched but the results have not been published yet. Up till now, the experiment has used a single cavity with a variety of dielectric rods and metal posts. However, a four-cavity array will soon be used to search higher frequencies. Ultimately, the LLNL experiment will cover the mass range $`1.310^6`$eV to $`1310^6`$eV at KSVZ sensitivity or better (see below).
A development project is under way to equip the LLNL detector with SQUID microwave receivers. These would replace the HEMT receivers presently in use. The HEMT receivers have noise temperature $`T_n3K`$ . It appears that $`T_n<0.3K`$ will be reached with the SQUIDs . To take advantage of such low electronic noise temperatures, the experiment will have to be equipped with a dilution refrigerator. Also bucking coils must be installed to cancel the static magnetic field at the location of the SQUID. When this development project is completed, the LLNL detector will have sufficient sensitivity to detect DFSZ axions at even a fraction of the local halo density.
The Kyoto experiment exploits resonant $`a\gamma `$ conversion in a cavity permeated by a large static magnetic field, as do the other experiments, but uses a beam of Rydberg atoms to count the photons from $`a\gamma `$ conversion . Single photon counting constitutes a dramatic improvement in microwave detection sensitivity. With HEMT amplifiers one needs to have thousands of $`a\gamma `$ conversions per second and integrate for about 100 sec to find a signal in the noise. With single photon counting, a few $`a\gamma `$ conversions suffice in principle. To build a beam of Rydberg atoms capable of single photon counting is a considerable achievement in itself. In addition, a dilution refrigerator is necessary to cool the cavity down to a temperature ($`10`$ mK) where the thermal photon background is negligible. The projected sensitivity of the Kyoto experiment is sufficient to detect DFSZ axions at even a fraction of the local halo density.
## 5 Other axion searches
There are a number of other techniques which have been used to search for very weakly coupled (so-called ’invisible’) axions. Although these searches have not ruled out parameter space that is not also presently ruled out by the astrophysical limits described above, they do provide completely independent constraints.
### 5.1 Solar axion searches
The conversion of axions to photons in a magnetic field can be used to look for solar axions too . The flux of solar axions on Earth is $`\frac{7.410^{11}}{\mathrm{sec}\mathrm{cm}^2}(\frac{g_\gamma }{0.36})^2(\frac{m_a}{\mathrm{eV}})^2`$ from the Primakoff conversion of thermal photons in the sun . The actual flux may be larger because other processes, such as Compton-like scattering, contribute if the axion has an appreciable coupling to the electron. At any rate the flux is huge compared to what can be produced by man-made processes on Earth and it is cost free. Solar axions have a broad spectrum of energies of order the temperature in the solar core, from one to a few keV.
Since the magnetic field is homogeneous on the length scale set by the axion de Broglie wavelength, the final photon is colinear with the initial axion. The photon and axion also have the same energy assuming the magnetic field is time-independent. The $`a\gamma `$ conversion probability is
$$p=\frac{1}{4}(\frac{\alpha g_\gamma }{\pi f_a})^2(B_0LF(q))^2$$
(12)
if $`B_0(z)=B_0b(z)`$ is the magnetic field transverse to the direction of the colinear axion and photon, $`z`$ is the coordinate along this direction, $`L`$ is the depth over which the magnetic field extends and $`F(q)`$ is the form factor
$$F(q)=\frac{1}{L}_0^L𝑑ze^{iqz}b(z)$$
(13)
where $`q=k_\gamma k_a=E_a\sqrt{E_a^2m_a^2}\frac{m_a^2}{2E_a}`$ is the momentum transfer. If the magnetic field is homogeneous ($`b=1`$), then
$`F(q)`$ $`=`$ $`{\displaystyle \frac{2}{qL}}sin{\displaystyle \frac{qL}{2}}`$ (14)
$``$ $`1\mathrm{for}qL1.`$
For $`qL1`$, the conversion probability goes as $`sin^2(\frac{qL}{2})`$ because the axion and photons oscillate into each other back and forth. The form factor $`F(q)`$ can be improved by filling the conversion region with a gas whose pressure is adjusted in such a way that the plasma frequency, which acts as an effective mass for the photon, equals the axion mass .
Multiplying the flux times the conversion probability, one obtains the event rate:
$$\frac{\mathrm{events}}{\mathrm{time}}\frac{200}{\mathrm{day}}\frac{VL}{\mathrm{meter}^4}F(q)^2(\frac{B_0}{8\mathrm{T}\mathrm{e}\mathrm{s}\mathrm{l}\mathrm{a}})^2(\frac{m_a}{\mathrm{eV}})^4.$$
(15)
The final state photons are soft x-rays which may be detected with good efficiency. There are radioactive backgrounds to worry about however.
The above type of detector is usually referred to as an axion helioscope. If a signal is found due to axions or familons, the detector immediately becomes a marvelous new tool for the study of the solar interior. Experiments were carried out at Brookhaven National Lab. and more recently at the University of Tokyo . The BNL experiment used a stationary Isabelle dipole magnet whose aperture was directed towards the sun at sunset. The total exposure time was of order 15 minutes. The Tokyo experiment uses a superconducting dipole magnet mounted on a altazimuth which tracks the sun. The 95% confidence level upper limit based on a few days of data taking is $`g_{a\gamma \gamma }\alpha g_\gamma /\pi f_a6.010^{10}\mathrm{GeV}^1`$ for $`m_a0.03`$ eV.
F. Avignone and his collaborators have exploited a different method to search for solar axions, namely the coherent Primakoff conversion of axions to photons in a crystal lattice. When the incident angle fulfills the Bragg condition for a given crystalline plane, the rate is enhanced. As the crystal detector turns with the Earth relative to the Sun’s direction, a characteristic diurnal temporal pattern is produced. Using 1.94 kg.yr of data from a Ge detector in Sierra Grande, Argentina, the bound $`g_{a\gamma \gamma }<2.710^9`$ GeV<sup>-1</sup> was obtained, independent of axion mass up to approximately 1 keV.
S. Moriyama proposed looking for monochromatic axions emitted in the deexcitation of $`{}_{}{}^{54}\mathrm{Fe}`$ in the sun. $`{}_{}{}^{54}\mathrm{Fe}`$ has an M1 transition, between the first excited state and the ground state, with excitation energy 14.4 keV. The monochromatic axions can be resonantly absorbed by the same nucleus in the laboratory because the axions are Doppler broadened due to the thermal motion of the axion emitter in the Sun. An experiment of this type was carried out by M. Krčmar et al. .
### 5.2 Laser experiments
Eqs.(12,14) give the conversion probability in a static magnetic field of an axion to a photon of the same energy. The polarization of the photon is parallel to the component of the magnetic field transverse to the direction of motion. The inverse process, conversion of such a photon to an axion, occurs with the same probability $`p`$, of course. K. van Bibber et al. proposed a ’shining light through walls’ experiment in which a laser beam is passed through a long dipole magnet like those used for high-energy physics accelerators. In the field of the magnet, a few of the photons convert to axions. Another dipole magnet is set up in line with the first, behind a wall. Since the axions go through the wall unimpeded, this setup allows one to ’shine light through the wall.’ An experiment of this type was carried out by the RBF collaboration . Compared with a solar axion search, it has the advantage of greater control over experimental parameters. But the signal is much smaller because one pays twice the price of the very small axion-photon conversion rate.
Other types of laser experiments were proposed in which one looks at the effect of the axion on the propagation of light through a magnetic field. If the photon beam is linearly polarized and the polarization direction is at an angle to the direction of the magnetic field, the plane of polarization turns because the component of light polarized parallel to the magnetic field gets depleted whereas the perpendicular component does not. There is an additional effect of birefringence because the component of light polarized parallel to the magnetic field mixes with the axion and hence moves more slowly than in vacuo. Birefringence affects the ellipticity of the polarization as the light travels on. The birefringence associated with the axion is considerably smaller than that due to the box diagram in QED, i.e. an electron running in a loop with four external photon lines. The polarization rotation and the birefringence effects were searched for by the RBF collaboration , which boosted these effects by passing the laser beam hundreds of times in an optical cavity within the magnet. For more recent work, see ref. .
### 5.3 A telescope search
The axion decays to two photons at the rate:
$`\mathrm{\Gamma }(a2\gamma )`$ $`=`$ $`g_\gamma ^2{\displaystyle \frac{\alpha ^2}{64\pi ^3}}{\displaystyle \frac{m_a^3}{f_a^2}}`$ (16)
$`=`$ $`{\displaystyle \frac{g_\gamma ^2}{6.810^{24}\mathrm{sec}}}\left({\displaystyle \frac{m_a}{\mathrm{eV}}}\right)^5.`$
For axions in the 10 eV mass range, the decay rate is comparable to the age of the universe. The dominant contribution to the cosmological energy density of such axions is thermal production . The energy density in thermal axions is proportional to the axion mass and becomes equal to the critical energy density at a mass of order 100 eV, the exact value depending on how many particle species annihilate after the axion decoupled. One can search for relic axions by looking for the monochromatic photons from their decay. Such photons arrive to us from all directions but preferentially from large agglomerations of mass, such as clusters of galaxies. The relative width $`\frac{\mathrm{\Delta }\lambda }{\lambda }`$ of the photon line from axion decay in galactic clusters is of order the virial velocity there, i.e $`10^2`$. By subtracting the spectrum of light from a galactic cluster from that of the nightsky ’off cluster’, one can subtract some of the background. The latter is dominated by lines in the spectrum of airglow . A search of this type was carried out at Kitt Peak National Laboratory by the TSAR collaboration . They placed the limit $`g_{a\gamma \gamma }10^{10}\mathrm{GeV}^1`$ in the range $`3m_a8`$ eV (3100-8300 Å).
## 6 Macroscopic forces mediated by axions
J. Moody and F. Wilczek analyzed the apparent deviations from the $`1/r^2`$ gravitational force law due to the exchange of virtual axions. The coupling of the axion to a spin 1/2 fermion $`f`$ has the general form:
$$_{aff}=g_f\frac{m_f}{v}a\overline{f}(i\gamma _5+\theta _f)f$$
(17)
when allowance is made for the fact that CP is violated in the electroweak interactions. $`g_f`$ is of order one, whereas $`\theta _f`$ is of order $`\overline{\theta }`$, which is of order $`10^{17}`$ in the Standard Model. The second term in Eq. (17) produces a coupling of the axion field to the mass density of a macroscopic collection of non-relativisitc fermions, whereas the first term produces a coupling of the axion field to the spin density of that macroscopic body. The first type of coupling was called ’monopole’, the second ’dipole’. The axion mediated forces between two macroscopic bodies therefore fall into three categories: monopole-monopole, monopole-dipole, and dipole-dipole. The monopole-monopole force is suppressed by two powers of $`\theta _f`$ and therefore very small. The dipole-dipole force has a very large background from ordinary magnetic forces. This background can be suppressed by using superconducting shields but not well enough for the axion mediated contribution to be detected. The monopole-dipole is the least difficult to detect. It is suppressed by only one factor of $`\theta _f`$ and it can be modulated by rotating the spin polarized body. Experiments of this type have been carried out . Unfortunately, their sensitivity is many orders of magnitude short of what is required to see the effect. |
warning/0002/hep-ph0002200.html | ar5iv | text | # Strange quark polarization of the nucleon: A parameter-independent prediction of the chiral potential model
## Abstract
We perform a one-loop calculation of the strange quark polarization ($`\mathrm{\Delta }s`$) of the nucleon in a SU(3) chiral potential model. We find that if the intermediate excited quark states are summed over in a proper way, i.e., summed up to a given energy instead of given radial and orbital quantum numbers, $`\mathrm{\Delta }s`$ turns out to be almost independent of all the model parameters: quark masses, scalar- and vector-potential strengths. The contribution from the quark-antiquark pair creation and annihilation “$`Z`$” diagrams is found to be significant. Our numerical results agree quite reasonably with experiments and lattice QCD calculations.
PACS numbers: 12.39.Ki, 12.39.Fe, 12.39.Pn, 13.88.+e
The intrinsic strangeness content of the nucleon is a key ingredient to understand the structure and dynamics inside baryons. While the experimental investigation of the nucleon spin structure clearly indicates that a strange quark sea exists and is also polarized relative to the nucleon spin, the successes of the naive spin-flavor SU(6) valence quark model in various aspects suggest that the strangeness content should belong to higher order effects for the nucleon. The SU(3) flavor chiral quark model, which couples light quarks to octet pseudoscalar mesons by the requirement of chiral symmetry, provides a natural mechanism for such a perturbative picture: at zeroth order the ground state octet baryons are described by a SU(6) wave function of three valence quarks, and at second order in the quark-meson coupling sea quarks can be generated by emitting a meson from the valence quark. For example, in the nucleon the strange quark can be generated by emitting a $`K^+`$ from the $`u`$ quark or by emitting a $`K^0`$ from the $`d`$ quark (Fig. 1), and hence strange quarks can contribute to the nucleon spin structure.
In this paper we will adopt the standard perturbation theory to calculate the strange quark polarization of the nucleon in the framework of a SU(3) chiral potential model. As will be shown, up to second order the diagrams of Fig. 1 are the only contributions. Hence, a quantitative determination of the strange quark polarization is a clean test for the interaction picture of the chiral quark model, which allows in this case only admixtures of $`K`$ mesons. On the other hand the $`u`$ or $`d`$ quark polarizations obtain contributions from many other diagrams and therefore depend on many more model parameters.
To set up our calculation scheme, we first define the effective Lagrangian for the SU(3) chiral potential model:
$``$ $`=`$ $`\overline{\psi }[i/S(r)\gamma ^0V(r)]\psi `$ (3)
$`{\displaystyle \frac{1}{2F_\pi }}\overline{\psi }[S(r)(\sigma +i\gamma ^5\lambda ^i\varphi _i)+(\sigma +i\gamma ^5\lambda ^i\varphi _i)S(r)]\psi +`$
$`{\displaystyle \frac{1}{2}}(_\mu \sigma )^2+{\displaystyle \frac{1}{2}}(_\mu \varphi _i)^2{\displaystyle \frac{1}{2}}m_\sigma ^2\sigma ^2{\displaystyle \frac{1}{2}}m_i^2\varphi _i^2.`$
The model Lagrangian is derived from the $`\sigma `$ model in which meson fields are introduced to restore chiral symmetry . $`\psi `$ is the quark field with flavor and color indices implied; the scalar term $`S(r)=cr+m`$ represents the the linear scalar confinement potential $`cr`$ and the quark mass matrix $`m`$; $`V(r)=\alpha /r`$ is the Coulomb type vector potential and $`F_\pi `$=93MeV is the pion decay constant. $`\sigma `$ and $`\varphi _i`$ ($`i`$ runs from $`1`$ to $`8`$) are the scalar and pseudoscalar meson fields, respectively and $`\lambda _i`$ are the Gell-Mann matrices. The quark-meson interaction term of Eq.(3) is symmetrized since the mass matrix $`m`$ does not commute with all $`\lambda _i`$ for different quark masses.
The zeroth order quark Hamiltonian is set up as
$$H_q=d^3x\psi ^{}[\stackrel{}{\alpha }\frac{1}{i}\stackrel{}{}+\beta S(r)+V(r)]\psi .$$
(4)
It has discrete eigenstates which are obtained by numerical solution of the Dirac equation with a scalar and vector field . We write the solution as:
$$\psi (x)=\underset{\alpha }{}u_\alpha (x)a_\alpha +\underset{\beta }{}v_\beta (x)b_\beta ^{}.$$
(5)
Eq. (5) forms the basis of our unperturbed wave functions, where quarks are bound permanently by the confinement potential which is included in $`H_q`$. From Eq. (5) we can construct the quark propagator:
$`D(x_1,x_2)`$ $``$ $`0|T\{\psi (x_1),\overline{\psi }(x_2)\}|0`$ (6)
$`=`$ $`\theta (t_1t_2){\displaystyle \underset{\alpha }{}}u_\alpha (x_1)\overline{u_\alpha }(x_2)`$ (8)
$`\theta (t_2t_1){\displaystyle \underset{\beta }{}}v_\beta (x_1)\overline{v_\beta }(x_2),`$
The meson propagator given by Eq. (3) is the free one:
$`\mathrm{\Delta }_{ij}(x_1,x_2)`$ $``$ $`0|T\{\varphi _i(x_1),\varphi _j(x_2)\}|0`$ (9)
$`=`$ $`{\displaystyle \frac{i}{(2\pi )^4}}{\displaystyle d^4q\frac{\delta _{ij}e^{iq(x_1x_2)}}{q^2m_i^2+iϵ}}.`$ (10)
Given the unperturbed basis we can construct any physical quantity up to a desired order in the quark-meson interaction. In the following we are studying the quark contribution of flavor $`q`$ ($`q=u,d,s`$) to the nucleon spin which is defined through
$$\mathrm{\Delta }q=\frac{N|d^3x\overline{\psi _q}\gamma ^3\gamma ^5\psi _q|N}{N|N}.$$
(11)
At zeroth order $`H_q`$ gives the usual SU(6) three-quark states for the nucleon with the single quark wave function $`u_\alpha `$ in the ground state. The zeroth order diagram for the numerator of Eq. (11) is indicated in Fig. 2A, and the denominator by the diagram of Fig. 3A, which is simply unity. Clearly, strange quarks do not contribute at this order.
The corresponding Feynman diagrams which contribute to $`\mathrm{\Delta }s`$ up to second order are shown in Figs. 2 and 3. The denominator $`N|N`$ can be denoted as $`(1+\mathrm{const}./F_\pi ^2)`$, which can be expanded $`(1\mathrm{const}./F_\pi ^2+\mathrm{})`$ and has then to be multiplied with the numerator $`N|d^3x\overline{\psi }\gamma ^3\gamma ^5\psi |N`$. If finally, only terms of order $`1/F_\pi ^2`$ are kept in the product of the normalization and the matrix element of the spin, this has no effect on $`\mathrm{\Delta }s`$, since already the lowest order admixture of $`s`$ quark is proportional to $`1/F_\pi ^2`$.
In the Lagrangian of Eq. (3) the main effect of the nonperturbative quark-gluon interaction is supposed to be included by the scalar and vector potentials. In principle we can also include a residual perturbative gluon piece. This will introduce further modifications on $`\mathrm{\Delta }u`$ and $`\mathrm{\Delta }d`$. However since the perturbative quark-gluon interaction is diagonal in flavor space, it cannot generate strange quark admixtures in second order for the nucleon.
Now we are in a good position to calculate $`\mathrm{\Delta }s`$ for the nucleon: up to second order the only diagrams we need to consider are the subset of the diagrams of Fig. 2, which are given in Fig. 1. For the evaluation, we first give the explicit form for $`u_\alpha (x)`$ and $`v_\beta (x)`$ with: $`u_\alpha (x)=e^{iE_\alpha t}u_\alpha (\stackrel{}{x})\tau _\alpha `$, $`v_\beta (x)=e^{iE_\beta t}v_\beta (\stackrel{}{x})\tau _\beta `$, where $`\tau `$ is the flavor wavefunction and the spatial wavefunction is:
$$u_\alpha (\stackrel{}{x})=\left(\begin{array}{c}g_{njl}\\ i\stackrel{}{\sigma }\widehat{\stackrel{}{r}}f_{njl}\end{array}\right)Y_{jl}^m=\left(\begin{array}{c}g_{njl}Y_{jl}^m\\ if_{njl}Y_{jl^{}}^m\end{array}\right),$$
(12)
where $`g`$ and $`f`$ are real functions, $`n`$ is the radial quantum number, and $`Y_{jl}^m(\widehat{\stackrel{}{r}})`$ are the vector spherical harmonics. The second equality of Eq.(12) follows from $`\stackrel{}{\sigma }\widehat{\stackrel{}{r}}Y_{jl}^m=Y_{jl^{}}^m`$ with $`l^{}=2jl`$. For computational convenience, we will use exactly the same form for $`v_\beta (x)`$. Since for the antiquark solution the lower component is the large component, for $`v_\beta (x)`$ $`l`$ is actually the orbital quantum number of the small component, and $`|E_{j=l+1/2}|>|E_{j=l1/2}|`$. Thus for the antiquarks the sequence is inversed.
Denoting the initial and final quark states as $`u_i`$ and $`u_f`$ respectively, the contribution of the diagrams of Fig. 1 is:
$`\delta s`$ $`=`$ $`{\displaystyle \frac{1}{F_\pi ^2}}{\displaystyle }d^3xd^4x_1d^4x_2\mathrm{\Delta }(x_2,x_1)\overline{u}_f(x_2)S(r_2)\gamma ^5\lambda ^i\times `$ (14)
$`D(x_2,x)\gamma ^3\gamma ^5D(x,x_1)S(r_1)\gamma ^5\lambda ^iu_i(x_1).`$
Here we use $`\delta s`$ to indicate that it is only the contribution from a single quark state. Inserting the explicit expressions for the propagators, we get
$`\delta s`$ $`=`$ $`{\displaystyle \frac{1}{F_\pi ^2}}{\displaystyle }d^4x_1d^4x_2\overline{u}_f(x_2)S(r_2)\gamma ^5\lambda ^i\times `$ (20)
$`[\theta (t_2t)\theta (tt_1){\displaystyle \underset{\alpha \alpha ^{}}{}}u_\alpha (x_2)\mathrm{\Delta }_{\alpha \alpha ^{}}\overline{u}_\alpha ^{}(x_1)+`$
$`\theta (t_1t)\theta (tt_2){\displaystyle \underset{\beta \beta ^{}}{}}v_\beta (x_2)\mathrm{\Delta }_{\beta \beta ^{}}\overline{v}_\beta ^{}(x_1)`$
$`\theta (t_2t)\theta (t_1t){\displaystyle \underset{\alpha \beta ^{}}{}}u_\alpha (x_2)\mathrm{\Delta }_{\alpha \beta ^{}}\overline{v}_\beta ^{}(x_1)`$
$`\theta (tt_2)\theta (tt_1){\displaystyle \underset{\beta \alpha ^{}}{}}v_\beta (x_2)\mathrm{\Delta }_{\beta \alpha ^{}}\overline{u}_\alpha ^{}(x_1)]\times `$
$`S(r_1)\gamma ^5\lambda ^iu_i(x_1){\displaystyle \frac{i}{(2\pi )^4}}{\displaystyle d^4q\frac{\delta _{ij}e^{iq(x_1x_2)}}{q^2m_i^2+iϵ}}.`$
where $`\mathrm{\Delta }_{\alpha \alpha ^{}}=d^3x\overline{u}_\alpha \gamma ^3\gamma ^5u_\alpha ^{}`$, and similarly for $`\mathrm{\Delta }_{\beta \beta ^{}}`$ etc. The four time-ordered terms in Eq.(20) correspond to the time-ordered diagrams of Fig. 4.
We omit here the details for calculating $`\delta s`$ of Eq. (20). The integrals of Eq. (20) can be reduced analytically to radial integrations at the vertex points ($`r_1`$ and $`r_2`$) and of loop momentum $`|\stackrel{}{q}|`$. The the remaining integrations are carried out numerically. $`\mathrm{\Delta }s`$ for the whole nucleon is just $`\delta s`$ times a spin-isospin factor which can be straightforwardly calculated to be 2.
In Table I we list our model parameters. Since $`F_\pi =93`$MeV and $`m_K=495`$Mev are fixed by experiment, our model contains four free parameters: the two quark masses $`m_{u,d}`$, $`m_s`$ and the two strength constants of the scalar and vector potential denoted by $`c`$ and $`\alpha `$. The parameter $`\alpha `$ is fixed by the long-wavelength, transverse fluctuations of the QCD based static-source flux-tube picture . It was obtained to be $`0.26`$ in and $`0.30`$ in , while a much larger value of about $`0.52`$ was used by the Cornell group . Recent lattice calculation got a value around $`0.32`$ in the quenched approximation, and suggested that relaxing the quenched approximation may lead to $`\alpha 0.40`$. Quark masses and confinement strength are rather uncertain quantities. To study the variation of $`\mathrm{\Delta }s`$ over all the parameters, we choose in our calculation four different sets of parameters, including both current and constituent quark masses.
We study very different parameters because they often vary significantly from one model to another. For this model one possible choice of parameters to produce the correct nucleon mass, $`g_A`$, etc. is given in Ref. .
Table II gives the numerical results of $`\mathrm{\Delta }s`$ for the first two sets of parameters. The intermediate quark/antiquark states are summed over up to a radial quantum number of $`n=8`$ and total angular momentum $`j=17/2`$. We also list the intermediate results with the summation including states up to $`n=6`$ and $`j=11/2`$. The contributions from the four time-ordered diagrams in Fig. 4 are given separately. We note significant contributions from Fig. 4C and Fig. 4D, in which a quark-antiquark pair is created or annihilated by the axial vector current; these processes are usually referred to as the “$`Z`$” diagrams. On the other hand the diagram of Fig. 4B gives a fairly large positive contribution, therefore if the “$`Z`$” diagrams are neglected we would incorrectly conclude that $`\mathrm{\Delta }s`$ in the nucleon is positive.
From Table II one would conclude that a stronger confinement also gives a larger $`\mathrm{\Delta }s`$. This is due to the coupling of the meson field to the quark field which is proportional to the effective quark mass $`S(r)=cr+m`$. However, to compare with the energy scale in the lattice QCD calculation of $`\mathrm{\Delta }s`$, we should sum the excited states up to a given energy instead of given radial and orbital quantum number. The resummed $`\mathrm{\Delta }s`$ according to energy are given in Fig. 5. Since the quark states are discrete, we get plateaus in Fig. 5 at the energies where no new states emerge.
Since the strange axial current is a non-conserved composite operator, it has divergent matrix element (as is seen in Fig. 5), and therefore must be renormalized. Analogous to the lattice renormalization, we cut the quark intermediate states at an energy of $`1.7`$GeV, which is roughly the inverse of the lattice spacing in the lattice calculation of $`\mathrm{\Delta }s`$ ($`a^1=1.74`$GeV in ). The “renormalized” results are given in Table III. (In principle, we can also do renormalization by imposing a cutoff on the meson momentum, such as using the Pauli-Villars regulator $`(q^2\mathrm{\Lambda }^2)^1`$. But then the quark intermediate states have to be summed up to convergence. In practice this is not workable. An illustration with $`\mathrm{\Lambda }=1.7`$GeV is given in Fig. 5.)
We note a very interesting phenomena in Fig. 5: the result for $`\mathrm{\Delta }s`$ summed up to a given energy is rather robust against the variation of all the parameters. The insensitivity is especially impressive compared to the huge variation of $`m`$ and $`c`$. Table III shows that the “$`Z`$” diagram’s contribution is still significant.
The insensitivity of $`\mathrm{\Delta }s`$ on the parameter sets can be attributed to the fact that the increase of $`m`$ and $`c`$ (see Eq. (3)) enhances the quark-meson coupling and moves up the single quark state energy. Thus the contribution from a single quark state increases due to the stronger coupling but less states are accessible to be summed over up to a given energy. Similarly, the increase of $`\alpha `$ suppresses the contribution of a single quark state since the lower Dirac components of the quark wavefunctions are increasing. But it also reduces the quark state energy, so we have more states to sum over.
The main results of this paper can be summarized as follows: 1) Strange quark polarization is a very clean and robust prediction of the chiral potential model. Up to second order the only contribution arises from the diagram of Fig. 1. $`\mathrm{\Delta }s`$ depends only weakly on the model parameters, and our calculation shows further that the variation of these parameters does not influence $`\mathrm{\Delta }s`$ too much, provided we sum over the intermediate quark state up to a given energy. 2) The contribution from the intermediate excited quark states are important. It is not enough to restrict the intermediate state to the ground or the first few states. 3) Among the time-ordered diagrams, the quark-antiquark pair creation and annihilation “$`Z`$” diagrams are significant. It is the “$`Z`$” diagrams (Figs. 4C and 4D) that introduce a negative value for $`\mathrm{\Delta }s`$ in the nucleon, while the intermediate negative-energy states (Fig. 4B) gives a fairly large positive contribution. The importance of the pair creation and annihilation contribution to $`\mathrm{\Delta }q`$ has also been noticed by some of us previously in a valence and sea quark mixing model . 4) Our numerical result is quite consistent with experiments ($`\mathrm{\Delta }s`$($`Q^2=3`$GeV<sup>2</sup>)$`=0.10\pm 0.01\pm \mathrm{}`$, where the second $`\pm `$ sign represents further sources of error, principally the low $`x`$ extrapolation ) and lattice QCD calculations ($`\mathrm{\Delta }s=0.12(1)`$ ,$`0.109(30)`$ ), and is also consistent with a schematic calculation in the context of chiral quark model by Cheng and Li . To the best of our knowledge this is first time that $`\mathrm{\Delta }s`$ is consistently calculated up to the one-loop level in a quark model.
This work is supported by the CNSF (19675018), CSED, CSSTC, the DFG (FA67/25-1), and the DAAD. |
warning/0002/astro-ph0002402.html | ar5iv | text | # Intrinsic constraints on very high energy emission in gamma-ray loud blazars
## 1. Introduction
Clearly, the high energy $`\gamma `$-ray emission is an important piece in the blazar puzzle because the $`\gamma `$-ray observations of blazars provide a new probe of dense radiation field released through accretion onto a supermassive black hole in the central engine (Bregman 1990). The Energetic Gamma Ray Experiment Telescope (EGRET) which works in the 0.1–10GeV energy domain has now detected and identified 66 extragalactic sources in 3th catalog (Mukherjee et al 1999). All these objects are blazar-type AGNs whose relativistic jets are assumed to be close to the line of sight to the observer. It seems unambiguous that the intense gamma-ray emission is related with highly relativistic jet.
It has been generally accepted that the luminous gamma-ray emission is radiated from inverse Compton, but the problem of seed photons remains open for debate. The following arguments have been proposed: (1) synchrotron photons in jet (inhomogeneous model of synchrotron self Compton) (Maraschi, Ghisellini & Celotti 1992); (2) optical and ultraviolet photons directly from the accretion disk (Dermer & Schlikeiser 1993); (3) diffusive photons in broad line region (BLR) (Sikora, Begelman & Rees 1994, Blandford & Levinson 1995); (4) the reflected synchrotron photons by electron mirror in broad line region, namely, the reflected synchrotron inverse Compton (RSC) (Ghisellini & Madau 1996). These mechanisms may operate in different kinds of objects, however there is not yet a consensus on how these mechanisms work. Also it is not clear where the $`\gamma `$-ray emission is taking place largely because of uncertainties of soft radiation field in the central engine.
On the other hand, VHE observations (Kerrick et al 1995, Chadwick et al 1999, Roberts et al 1999, Aharonian et al 1999) are making attempts to explore the radiation mechanism because they may provide some restrictive constraints (Begelman, Rees & Sikora 1994, Mastichiadis & Kirk 1997, Tavecchio, Maraschi & Ghisellini 1998, Coppi & Aharonian 1999, Harwit, Protheroe and Biermann 1999). Based on the simple version of SSC model, Stecker, de Jager & Salamon (1996) predicted a large number of low redshift X-ray selected BL Lacs as TeV candidates, taking into account that the presence of intergalactic infrared radiation field including cosmic background leads to strong absorption of TeV photons from cosmological emitters (Stecker & de Jager 1998). It is suggested to form an extended pair halo in cosmological distance due to the external absorption (Aharonian, Coppi, & Voelk 1994). However, so far only three X-ray selected BL Lacs have been found to be TeV emitters by Whipple telescope ($`E>300`$GeV), in addition, photons higher than 0.3TeV in the X-ray-selected PKS 2155-304 with redshift $`z=0.116`$ has been detected photons 0.3TeV by Durham Mrk 6 telescope (Chadwick et al 1999). The recent measurements of intergalactic infrared field is quite different from the previous observations (Madau et al 1998, Steidel 1998). Although this external absorption is definitely important, the critical redshift $`z_c`$ beyond which cosmological back ground radiation and intergalactic infrared fields will absorb VHE photons remains uncertain. Especially the recent VHE observations show that Mrk 501 emits 25 TeV photons (Aharonian et al 1999). Evidently this suggests that the external absorption can not efficiently attenuate the VHE photons from reaching us across distances of 100 Mpc. It is highly desired to accurately probe the star formation rate in order to determine the critical redshift $`z_c`$.
Thus it seems significant to study the intrinsic mechanism for the deficiency of TeV photons from $`\gamma `$-ray loud AGNs disregarding the absorption by intergalactic infrared radiation field. A larger Lorentz factor of the jet implies higher density of the external photons in the blob, if the reflection of clouds in broad line region works, and therefore stronger absorption of high energy $`\gamma `$-rays (Celotti, Fabian & Rees 1998). Here we argue based on the hypothesis of Ghisellini & Madau (1996) that the energy density of reflected synchrotron photon is high enough for pair production via interaction of gamma-ray photons by inverse Compton scattering with reflected synchrotron photons if the bulk velocity is high enough. Further we apply the present constraint to the representative individual objects, Mrk 421 and 3C 279.
## 2. Constraints on VHE
Ghisellini & Madau (1996) have calculated the energy density of reflected synchrotron (Rsy) emission, and compared with the other reflected components. They draw a conclusion that the energy density of Rsy component dominates over 10 times of that of reflected component of accretion disk radiation. In this section we make an attempt to use the observables quantities to express the intrinsic constraints on very high energy emission. The overall $`\nu F_\nu `$ spectrum of blazars shows that there are two power peaks: the first is low energy one between IR/soft X-ray band, and the second is high energy one peaking in the MeV/GeV range (von Montigny et al 1995, Sambruna, Maraschi & Urry 1996, Comastri et al 1997, Kubo et al 1998). This characteristic can be explained by the simple context of one-zone homogeneous SSC or EC model. The low energy peak denoted $`\nu _\mathrm{s}`$ is caused by synchrotron radiation of relativistic electrons, and the second peak denoted $`\nu _\mathrm{c}`$, or $`\nu _{\mathrm{rsc}}`$ results from the Compton scattering off the synchrotron or reflected synchrotron photons by the same population of electrons, respectively. We take the two peaks and their corresponding fluxes as four observable quantities.
From the RSC model the magnetic field $`B`$ can be approximately expressed by the observational quantities. The observed frequency of synchrotron photon is $`\nu _\mathrm{s}=𝒟\nu _0\gamma _\mathrm{b}^2B`$ ($`\nu _0=2.8\times 10^6`$), and the frequency of reflected synchrotron Compton photons reads $`\nu _{\mathrm{rsc}}=𝒟(2\mathrm{\Gamma })^2\gamma _\mathrm{b}^4\nu _0B`$, and we can get the estimation of magnetic field $`B`$
$$B=\frac{(2\mathrm{\Gamma }\nu _\mathrm{s})^2}{𝒟\nu _0\nu _{\mathrm{rsc}}}\frac{𝒟\nu _\mathrm{s}^2}{\nu _0\nu _{\mathrm{rsc}}},$$
(1)
while in pure SSC model the magnetic field is approximately as $`B=\nu _\mathrm{s}^2/(𝒟\nu _0\nu _\mathrm{c})`$, where $`\nu _\mathrm{c}`$ is the frequency of photons emitted by SSC. The Doppler factor $`𝒟=1/\mathrm{\Gamma }[1\mu (1\mathrm{\Gamma }^2)^{1/2}]`$, where $`\mu =\mathrm{cos}\theta `$ is the cosine of the orientated angle of jet relative to the observer. Equation (1) is similar to the model of Sikora, Begelman & Rees (1994) (also see Sambruna, Maraschi & Urry 1996). Comparing with the above two formula, we learn that RSC model needs stronger magnetic field than SSC model does whereas the energy of relativistic electrons is lower in RSC model than in SSC model. The reflected synchrotron Compton (RSC) mainly depends on two parameters: the reflection albedo, namely, the Thomson scattering optical depth ($`\tau _{_{\mathrm{BLR}}}`$), and the Lorentz factor $`\mathrm{\Gamma }`$ of the relativistic jet.
### 2.1. Reflected Synchrotron Compton Emission
In the case of power-law distribution of electrons, $`N=N_0\gamma ^\alpha `$, ($`\gamma _{\mathrm{min}}\gamma \gamma _{\mathrm{max}}`$), where $`N`$ is the number density of relativistic electrons, and $`\gamma `$ is the Lorentz factor of electron, the synchrotron emission coefficiency is approximately given by $`\epsilon _\nu =c_5(\alpha )N_0B^{\frac{1+\alpha }{2}}(\nu /2c_1)^{\frac{1\alpha }{2}}`$. Here $`c_1=6.27\times 10^{18}`$, and $`c_5(\alpha )`$ is tabulated in Pacholczyk (1970) within the frequency range $`\nu _1\nu \nu _2`$, where $`\nu _{1,2}=\nu _0B(\gamma _{\mathrm{min}}^2,\gamma _{\mathrm{max}}^2)`$. The average energy density per frequency $`u_{\mathrm{syn},\nu ^{}}^{}`$ in a region with dimension $`s`$ in the jet comoving frame can be obtained
$$u_{\mathrm{syn},\nu ^{}}^{}=4\pi c^1c_5(\alpha )N_TB^{\frac{1+\alpha }{2}}\left(\frac{\nu ^{}}{2c_1}\right)^{\frac{1\alpha }{2}},$$
(2)
where $`N_T=N_0s`$. The number density of synchrotron photons can be obtained by
$$n_ϵ^{}^{^{}}=n_0(\alpha )N_TB^{\frac{1+\alpha }{2}}ϵ_{}^{}{}_{}{}^{\frac{1+\alpha }{2}},$$
(3)
and $`n_0(\alpha )`$ reads
$$n_0(\alpha )=\frac{4\pi c_5(\alpha )}{hc}\left(\frac{2hc_1}{m_ec^2}\right)^{\frac{\alpha 1}{2}},$$
here $`h`$ is Planck constant, and $`ϵ^{}=h\nu ^{}/m_ec^2`$. We have employed relationship $`n_ϵ^{}=n_\nu ^{}d\nu ^{}/dϵ^{}`$ to derive equation (3). The mean energy density is expressed by
$$u_{\mathrm{syn}}^{}\frac{8\pi c_5(\alpha )}{(3\alpha )c(2c_1)^{\frac{1\alpha }{2}}}N_TB^{\frac{1+\alpha }{2}}\nu _2^{^{\frac{3\alpha }{2}}},$$
(4)
for $`\alpha <3`$. Defining $`l_{\mathrm{ssc}}`$ as
$$l_{\mathrm{ssc}}=\frac{L_\mathrm{s}}{L_{\mathrm{ssc}}}=\frac{u_\mathrm{B}^{}}{u_{\mathrm{syn}}^{}},$$
(5)
we have
$$N_T=\frac{(3\alpha )c(2c_1)^{\frac{1\alpha }{2}}}{64\pi ^2c_5(\alpha )l_{\mathrm{ssc}}}B^{\frac{3\alpha }{2}}\nu _{\mathrm{s}}^{}{}_{}{}^{\frac{\alpha 3}{2}},$$
(6)
where $`\nu _\mathrm{s}^{}`$ denotes $`\nu _2^{}`$.
Since the opening angle of jet ($`\pi /\mathrm{\Gamma }^2`$) is much less than $`2\pi `$, it is then reasonable to assume that the BLR reflection approximates to plane mirror with thickness $`\mathrm{\Delta }R_{\mathrm{BLR}}`$ and electron number density $`n_e`$. The distance distribution of reflected synchrotron photons has been discussed by Ghisellini & Madau (1996). The angular distribution has not been issued. Since the thickness of mirror is zero, the energy density of reflected synchrotron emission sharply increases when blob is close to the mirror. In fact if we drop the assumption of zero-thickness of mirror, this characteristic will disappear. We will deal with this sophisticate model in future. Because the reflected synchrotron emission is isotropic in observer’s frame, the blob receives the reflected photon beamed within a solid angle $`\pi /\mathrm{\Gamma }^2`$. The subsequent section will pay attention to this effects. Neglecting the angle-dependent distribution of reflected photon field, we approximate the Doppler factor $`𝒟2\mathrm{\Gamma }`$ ($`\theta 0`$). For simplicity, we assume the mirror (reflecting clouds in broad line region) has Thomson scattering optical depth $`\tau _{_{\mathrm{BLR}}}=\sigma __\mathrm{T}n_e\mathrm{\Delta }R_{_{\mathrm{BLR}}}\sigma __\mathrm{T}n_eR_{_{\mathrm{BLR}}}`$. The received photon density $`n^{\prime \prime }(ϵ^{\prime \prime })`$ in the jet comoving frame can be then approximately written as
$`n^{\prime \prime }(ϵ^{\prime \prime })`$ $``$ $`(2\mathrm{\Gamma })^2\tau _{_{\mathrm{BLR}}}n_0(\alpha )N_TB^{\frac{1+\alpha }{2}}ϵ_{}^{}{}_{}{}^{\frac{1+\alpha }{2}}`$ (7)
$`=`$ $`(2\mathrm{\Gamma })^{3+\alpha }\tau _{_{\mathrm{BLR}}}n_0(\alpha )N_TB^{\frac{1+\alpha }{2}}ϵ_{}^{\prime \prime }{}_{}{}^{\frac{1+\alpha }{2}},`$
where we use $`ϵ^{\prime \prime }=(2\mathrm{\Gamma })^2ϵ^{}`$ (there are two Doppler shifts due to mirror effects). The convenient form of cross section of photon-photon collision is (Coppi & Blandford 1990)
$$\sigma _{\gamma \gamma }=\frac{\sigma __\mathrm{T}}{5ϵ}\delta \left(ϵ_0\frac{1}{ϵ}\right),$$
(8)
where $`\delta `$ is the usual $`\delta `$-function. This approximation is only valid for case of isotropic radiation field. The Rsy component is seen by the blob within the solid angle $`\mathrm{\Delta }\mathrm{\Omega }2\pi (1\mathrm{cos}\mathrm{\Gamma }^1)\pi /\mathrm{\Gamma }^2`$. Although the cross section of photon-photon interaction holds, the interacting possibility among photons reduces by a factor of $`\mathrm{\Delta }\mathrm{\Omega }/4\pi =1/(2\mathrm{\Gamma })^2`$ due to the beaming effects, which effectively reduces the opacity. Thus the pair production optical depth for photon with energy $`ϵ_\gamma `$ reads
$`\tau _{\gamma \gamma }^{\mathrm{rsc}}(ϵ_\gamma )`$ $`=`$ $`{\displaystyle \frac{0.2\sigma __\mathrm{T}s}{(2\mathrm{\Gamma })^2}}{\displaystyle ϵ_\gamma ^1n^{\prime \prime }(ϵ^{\prime \prime })\delta (ϵ^{\prime \prime }ϵ_\gamma ^1)𝑑ϵ^{\prime \prime }}`$ (9)
$`=`$ $`(2\mathrm{\Gamma })^{1+\alpha }\tau _{_{\mathrm{BLR}}}\tau _{\gamma \gamma }^0(ϵ_\gamma ),`$
here $`\tau _{\gamma \gamma }^0(ϵ_\gamma )`$ is
$$\tau _{\gamma \gamma }^0(ϵ_\gamma )=0.2\sigma __\mathrm{T}sn_0(\alpha )N_TB^{\frac{\alpha +1}{2}}ϵ_\gamma ^{\frac{\alpha 1}{2}}.$$
(10)
We will show the validity of the above approximation in the next subsection. Supposing that RSC operates efficiently in $`\gamma `$-ray loud blazars, we can get $`l_{\mathrm{rsc}}`$ from the observations
$$l_{\mathrm{rsc}}=\frac{L_\mathrm{s}}{L_{\mathrm{rsc}}}=\frac{u_\mathrm{B}^{}}{(2\mathrm{\Gamma })^2\tau _{_{\mathrm{BLR}}}u_{\mathrm{syn}}^{}}=\frac{l_{\mathrm{ssc}}}{(2\mathrm{\Gamma })^2\tau _{_{\mathrm{BLR}}}}.$$
(11)
From equation (11) we have
$$\tau _{_{\mathrm{BLR}}}=\frac{l_{\mathrm{ssc}}}{(2\mathrm{\Gamma })^2l_{\mathrm{rsc}}},$$
(12)
and
$$\tau _{\gamma \gamma }^{\mathrm{rsc}}(ϵ_\gamma )=(2\mathrm{\Gamma })^{\alpha 1}\left(\frac{l_{\mathrm{ssc}}}{l_{\mathrm{rsc}}}\right)\tau _{\gamma \gamma }^0(ϵ_\gamma ).$$
(13)
From equation (11) we know that the observed $`l_{\mathrm{rsc}}`$ represents the reflection ratio and Doppler factor of jet motion as long as the Compton catastrophe does not occur. If we set $`l_{\mathrm{rsc}}l_{\mathrm{ssc}}`$, we get $`\tau _{_{\mathrm{BLR}}}\mathrm{\Gamma }^2=0.01`$ for $`\mathrm{\Gamma }=10`$. This value is the lowest one in the model of Sikora, Begelman & Rees (1994) who suggest $`\tau _{_{\mathrm{BLR}}}=0.10.01`$. In fact we can roughly adopt $`\tau _{_{\mathrm{BLR}}}`$ as the covering factor which is usually taken to be 0.1 in fitting the broad emission line by photoionization model.
Inserting $`B`$ and $`N_T`$ \[eqs(1) and (6)\] into $`\tau _{\gamma \gamma }^0`$ (Eq. 10), and letting $`ϵ_\gamma =ϵ_{\mathrm{obs}}/𝒟`$ and $`s=c𝒟\mathrm{\Delta }t_{\mathrm{obs}}`$, we have the pair production optical depth for $`ϵ_{\mathrm{obs}}`$ in the observer’s frame
$`\tau _{\gamma \gamma }^{\mathrm{rsc}}(ϵ_{\mathrm{obs}})`$ $`=`$ $`K_\alpha {\displaystyle \frac{\nu _\mathrm{s}^{\frac{5+\alpha }{2}}}{\nu _{\mathrm{rsc}}^2}}(2\mathrm{\Gamma })^{3+\alpha }𝒟^{1\alpha }ϵ_{\mathrm{obs}}^{\frac{\alpha 1}{2}}l_{\mathrm{rsc}}^1\mathrm{\Delta }t_{\mathrm{obs}}`$ (14)
$``$ $`K_\alpha {\displaystyle \frac{\nu _\mathrm{s}^{\frac{5+\alpha }{2}}}{\nu _{\mathrm{rsc}}^2}}𝒟^4ϵ_{\mathrm{obs}}^{\frac{\alpha 1}{2}}l_{\mathrm{rsc}}^{}{}_{}{}^{1}\mathrm{\Delta }t_{\mathrm{obs}},`$
where $`K_\alpha =\frac{(3\alpha )\sigma __\mathrm{T}c}{80\pi h\nu _0^2}\left(\frac{h}{m_ec^2}\right)^{\frac{\alpha 1}{2}}`$ ($`K_\alpha =7.9\times 10^{18}`$ for $`\alpha =2.4`$). There are five observational parameters: $`\nu _\mathrm{s}`$, $`\nu _{\mathrm{rsc}}`$, $`\alpha `$, $`\mathrm{\Delta }t_{\mathrm{obs}}`$ and $`l_{\mathrm{rsc}}`$; and the unknown Doppler factor $`𝒟`$. For the typical value of parameters, $`\alpha =2.4`$, $`\nu _\mathrm{s}=4.0\times 10^{14}`$Hz, and $`\nu _{\mathrm{rsc}}=1.0\times 10^{25}`$Hz, $`\mathrm{\Delta }t_{\mathrm{obs}}=1`$ day, and $`𝒟=10`$, we have
$`\tau _{\gamma \gamma }^{\mathrm{rsc}}(ϵ_{\mathrm{obs}})`$ $`=`$ $`1.9l_{\mathrm{rsc}}^1𝒟_{10}^4\left({\displaystyle \frac{ϵ_{\mathrm{obs}}}{\mathrm{TeV}}}\right)^{0.7}\left({\displaystyle \frac{\mathrm{\Delta }t_{\mathrm{obs}}}{\mathrm{day}}}\right)`$ (15)
$`\left({\displaystyle \frac{\nu _\mathrm{s}}{4.0\times 10^{14}\mathrm{Hz}}}\right)^{3.7}\left({\displaystyle \frac{\nu _{\mathrm{rsc}}}{10^{25}\mathrm{Hz}}}\right)^2,`$
where $`𝒟_{10}=𝒟/10`$. Figure 1 shows the opacity due to pair production of photons with very high energy encountering with the reflected synchrotron photons. The equation (15) tells us the constraints on VHE from jet: (1) smaller $`l_{\mathrm{rsc}}`$, i.e. stronger reflection, will leads to the absorption of TeV photons. This parameter represents the energy density reflected by the BLR cloud including the bulk relativistic motion. From this estimation we know that TeV photon will be absorbed by the reflected synchrotron photons provided that $`l_{\mathrm{rsc}}<1.9`$. (2) $`\tau _{\gamma \gamma }`$ is sensitive to $`\nu _\mathrm{s}`$ and $`\nu _{\mathrm{rsc}}`$. (3) $`\tau _{\gamma \gamma }`$ is proportional to $`𝒟^4`$, in contrast to the usual down-limit (see Mattox et al 1993, and Dondi & Ghisellini 1995), providing the upper limit Doppler factor of bulk motion from $`\tau _{\gamma \gamma }1`$,
$`𝒟`$ $``$ $`8.5l_{\mathrm{rsc}}^{1/4}\left({\displaystyle \frac{ϵ_{\mathrm{obs}}}{\mathrm{TeV}}}\right)^{0.175}\left({\displaystyle \frac{\mathrm{\Delta }t_{\mathrm{obs}}}{\mathrm{day}}}\right)^{0.25}`$ (16)
$`\left({\displaystyle \frac{\nu _\mathrm{s}}{4.0\times 10^{14}\mathrm{Hz}}}\right)^{0.925}\left({\displaystyle \frac{\nu _{\mathrm{rsc}}}{10^{25}\mathrm{Hz}}}\right)^{0.5}.`$
This is a new constraint, which is expressed by the observational quantities. It lends us a simple and efficient way to select TeV candidates from known blazars in term of their known characteristics.
### 2.2. Angular Distribution of Reflected Photons
The received photons reflected by BLR in comoving frame is anisotropic, therefore, the pair opacity should be carefully treated. We have made important approximations that the BLR is thought to be a plane mirror and treated the photon-photon interaction in an approximate way. Now let us show the validity of this approximation. We adopt the geometry shown in Fig 1c of Ghisellin & Madau (1996). They show that the energy density of reflected synchrotron photon strongly depends on the location of emitting blob. We should admit that the aximal symmetry holds in the reflected synchrotron emission. We approximate the Thomson scattering event by isotropic scattering with cross section $`\sigma __\mathrm{T}`$ and neglect recoil, which is a very good approximation when $`ϵ_s1`$. The angular distribution of reflected synchrotron emission is given by $`n_{\mathrm{ph}}(ϵ_s,\mu ,r_0)=n_0f(\mu ,r_0)ϵ_s^q`$ ($`n_0`$ is a constant), the function $`f(\mu ,r_0)`$ determines the angular distribution of reflected synchrotron photons (Ghisellini & Madau 1996)
$$f(\mu ,r_0)=𝒟^2(\mu )g(\mu ,r_0),$$
(17)
where $`g(\mu ,r_0)=\left[\left(1r_0^2+r_0^2\mu ^2\right)^{1/2}r_0\mu \right]^2`$, and $`r_0=R_\gamma /R_{_{\mathrm{BLR}}}(0,1)`$, where $`R_\gamma `$ is the distance of blob to the center. Figure 2 shows the angular distribution in blob comoving frame. It can be seen that the geometry effect of reflecting mirror isotropizes the radiation at some degrees, but the beaming effect still dominates. It is still a good approximation that the radiation is beamed with a cone of solid angle $`\pi /\mathrm{\Gamma }^2`$. Thus the pair opacity can be written as (Gould & Schréder 1967)
$$\tau _{\gamma \gamma }(ϵ_\gamma )=2\pi R_{_{\mathrm{BLR}}}_0^1𝑑r_0_1^1𝑑\mu (1\mu )_{ϵ_c}n_{\mathrm{ph}}\sigma _{\gamma \gamma }𝑑ϵ_s,$$
(18)
where $`ϵ_c=2/(1\mu )ϵ_\gamma `$ and the photon-photon cross section $`\sigma _{\gamma \gamma }`$ reads
$$\sigma _{\gamma \gamma }=\frac{3\sigma __\mathrm{T}}{16}(1\beta ^2)\left[(3\beta ^4)\mathrm{ln}\left(\frac{1+\beta }{1\beta }\right)2\beta (2\beta ^2)\right],$$
(19)
where $`\beta `$ is the speed of the electron and positron in the center of momentum frame $`\beta =\left[12/ϵ_\gamma ϵ_s(1\mu )\right]^{1/2}`$. Performing the integral we have
$$\tau _{\gamma \gamma }=n_0\sigma __\mathrm{T}R_{_{\mathrm{BLR}}}ϵ_\gamma ^{q1}A(q),$$
(20)
where $`A(q)`$ reads
$$A(q)=2^{3q}\pi A_0(q)_1^1𝑑\mu (1\mu )^qA_1(\mu ),$$
(21)
with
$$A_0(q)=_0^1𝑑\beta (1\beta ^2)^{q2}\beta \sigma _{\gamma \gamma }(\beta ),$$
(22)
and $`A_1(\mu )=_0^1f(\mu ,r_0)𝑑r_0`$ is the integral of $`f(\mu ,r_0)`$ over the entire broad line region, which can be evaluated as
$`A_1`$ $`=`$ $`\sqrt{1\mu ^2}`$ (23)
$`\left\{\mathrm{cos}\theta _0\mathrm{ln}\left[{\displaystyle \frac{\mathrm{tan}\frac{\varphi _1}{2}}{\mathrm{tan}\frac{\varphi _2}{2}}}\right]+\mathrm{sin}\theta _0\mathrm{ln}\left[{\displaystyle \frac{\mathrm{sin}\varphi _1}{\mathrm{sin}\varphi _2}}\right]\right\},`$
with $`\varphi _1=\pi /2\theta _0`$, $`\varphi _2=\mathrm{arccos}\sqrt{1\mu ^2}\theta _0`$, and $`\theta _0=\mathrm{arcsin}\sqrt{1\mu ^2}`$. The function $`A(q)`$ is plotted in Fig 3. Since the beamed radiation field reduces the effective cross section of photon-photon interaction by a factor $`1/(2\mathrm{\Gamma })^2`$, it would be convenient to check our approximation by the quantity $`(2\mathrm{\Gamma })^2A(q)`$. We can easily find that it is close to 0.3$``$0.4 when $`q1.7`$, suggesting our approximation is accurate enough. It should be pointed out that the present treatments of reflected synchrotron radiation can be conveniently extended to the inclusion of the radiation from the secondary electrons if we further study the pair cascade in the future.
### 2.3. The Dimension of External Absorption
The last two subsections are devoted to the internal absorption of TeV photons, the developments of pair cascade due to the present mechanism will be treated in a preparing paper (Wang, Zhou & Cheng 2000). However it would be useful to compare the dimensions and radiation of the pair cloud due to the internal absorption and the pair halo suggested by Aharonian, Coppi, & Veolk (1994), who argue the formation of pair halo due to the interaction of TeV photons from AGNs with infrared photons of cosmological background radiation. This external absorption produces pairs, which are quickly isotropized by an ambient random magnetic field, forming a extended halo of pairs with typical dimension of ($`R>1`$Mpc). Without specific mechanism we know that the time scale of halo formation is of about $`10^6`$ yr. Usually this absorption is regarded as the main mechanism of deficiency of TeV emission from EGRET-loud blazars (Stecker & de Jager 1998). Let us simply estimate the scale of pair halo before it is isotropized by the ambient magnetic field. Assuming the intergalactic magnetic field $`B=10^9`$ Gauss, then the mean free path of pair electrons in halo reads
$$\lambda _\mathrm{e}1.0\left(\frac{E}{1.0\mathrm{TeV}}\right)^{0.5}\left(\frac{B}{10^9\mathrm{G}}\right)^{0.5}\mathrm{Kpc},$$
(24)
The initial halo is of such a dimension, which is much larger than that of intrinsic absorption case. Aharonian, Coppi, & Volk (1994) have suggested some signatures of such an extended halo, especially for the light curves in high energy bands (Coppi & Aharonian 1999). Anyway this is much larger than that of the present internal pair cloud. Thus it is easier to distinguish the two cases.
## 3. Applications
We have set a new constraint on the very high energy emission in term of observable quantities. As the applications of the present model, we would like to address some properties of very high energy from blazars.
### 3.1. Broadband Continuum and Mirror
The broadband continuum of blazars show attractive features which indicate the different processes powering the objects. The ratio $`L_\gamma /L_{\mathrm{op}}`$ of $`\gamma `$-ray luminosity to optical in flat spectrum radio quasars (FSRQs) is quite different from that in BL Lacs (Dondi & Ghisellini 1995). Comastri et al (1997) confirmed this result in a more larger samples and found this mean ratio is roughly of unity in BL Lacs and $`L_\gamma /L_{\mathrm{op}}30`$, namely $`l_{\mathrm{rsc}}0.03`$ in FSRQs. Ghisellini et al (1993), using the classical limit of SSC model, show that there is a systematical difference in Doppler factors $`𝒟`$ between BL Lacs and core-dominated quasars, $`\mathrm{log}𝒟=0.12`$ for BL Lacs and $`\mathrm{log}𝒟=0.74`$ for core-dominated quasars. These differences have been confirmed by Güijosa & Daly (1996) who assume that the particles and magnetic field are in equipartition. This difference would lead to more prominent difference of reflected synchrotron photon energy density, suggesting a different mechanism in these objects. The two systematically different features in $`l_{\mathrm{rsc}}`$ and Doppler factor $`𝒟`$ strongly suggest that the different mechanism of $`\gamma `$-ray radiation may operate in these objects. From eq.(15) it is believed that the deficiency of TeV emission in radio-loud quasars may be intrinsic due to the present mechanism.
### 3.2. On Mrk 421 and 3C 279
Whipple observatory had ever searched for TeV gamma-ray emission for 15 EGRET-AGNs with low redshift, but only Mrk 421 has positive signal(Kerrick et al 1995). Even at present stage only three X-ray selected BL Lacs have been reported as TeV emitters, Mrk 421, Mrk 501 and 1ES2344+514 (Cataness et al 1997), and PKS 2155-304 is a potential TeV emittor (Chadwick et al 1999). We can apply the present model to the two representative sources: Mrk 421 and 3C 279 for specific illustration.
Mrk 421: This is an X-ray selected BL Lac object, and has been detected GeV $`\gamma `$-ray emission by EGRET (Lin et al 1992), and the first TeV emission by Whipple (Punch et al 1992). It has been extensively and frequently observed by telescopes from radio to TeV bands(Kerrick et al 1995, Macomb et al 1995, Takahashi et al 1996, Krennrich et al 1999). TeV observations of Mrk 421 by Whipple show that the TeV photon did not flare much more dramatically than the X-rays, suggesting that the enhanced high-energy electrons were scattering off a part of the synchrotron spectral energy distribution that remained constant (Takahashi et al 1996). Roughly speaking this object satisfies the energy equipartition for the two power peaks(Zdziarski & Krolik 1993, Macomb et al 1995), suggesting $`l_{\mathrm{ssc}}1`$ (Macomb et al 1995), and pure SSC model agrees with the observations (Krennrich et al 1999), suggesting $`l_{\mathrm{rsc}}1`$. This indicates the RSC process is not important. The synchrotron component peaks in luminosity at UV to soft X-ray energies and continues into KeV X-rays(Maraschi, Ghisellini & Celotti 1994). The gamma-ray emission extends from 50 MeV to an astounding TeV. Data combined over several periods (Lin et al 1992) reveal a hard GeV spectrum ($`\alpha _{\mathrm{GeV}}0.7`$) by EGRET and a steeper one at TeV energies ($`\alpha _{\mathrm{TeV}}1.30`$) from the Whipple observatory (Schubnell et al 1994), implying a spectral break. The multiwavelegenth spectrum shows that $`\nu _\mathrm{s}=3\times 10^{16}`$Hz, $`\nu _{\mathrm{rsc}}=6.5\times 10^{25}`$Hz, and $`\alpha =2.0`$ (Macomb et al 1995, Kubo et al 1998). The shortest timescale of $`\gamma `$-ray variability is about $`\mathrm{\Delta }t_{\mathrm{obs}}20`$ minutes (Gaidos et al 1996). If we take $`𝒟=10`$ and $`l_{\mathrm{rsc}}=1`$, we have $`\tau _{\gamma \gamma }^{\mathrm{rsc}}(ϵ_{\mathrm{obs}})=5.0\times 10^5`$ for TeV photons. This means $`l_{\mathrm{rsc}}1`$, which is consistent with pure SSC model. From this we can estimate the scattering medium $`\tau _{_{\mathrm{BLR}}}=2.0\times 10^4`$ \[$`\tau _{\gamma \gamma }^{\mathrm{rsc}}(ϵ_{\mathrm{obs}})=1`$\] when we take $`l_{\mathrm{rsc}}=5.\times 10^5`$ and $`2\mathrm{\Gamma }=10`$, suggesting the mirror effects can be ruled out in this object. This result agrees to the absence of any evident emission lines in Mrk 421. Interestingly, Celotti, Fabian & Rees (1998) have suggested from rapid TeV variability of Mrk 421 that its accretion rate is lower than $`10^210^3`$ Eddington rate. They thus propose that advection-dominated accretion flow or ion pressure supported tori (Ichimaru 1977, Rees et al 1982, Narayan & Yi 1994) may power the luminosity.
3C 279: This is a typical FSRQ. The first Whipple observation shows negative signal in TeV (Kerrick et al 1995), however the multiwavelength simultaneous observations by EGRET, ASCA, RXTE, ROSAT, IUE in 1996 January-February show an intensive flare with very flat spectrum in EGRET band (Wehrle et al 1998), showing $`l_{\mathrm{rsc}}10^1`$ at the high state of $`\gamma `$-ray emission. It has been argued that RSC may explain the 1996 gamma-ray flare (Wehrle et al 1998). The observed flare show that the synchrotron emission peaks at $`\nu _\mathrm{s}=5.0\times 10^{12}`$Hz, $`\nu _{\mathrm{rsc}}=1.0\times 10^{23}`$Hz, $`l_{\mathrm{rsc}}=0.1`$, and $`\mathrm{\Delta }t_{\mathrm{obs}}=`$8hr (Wehrle et al 1998). If we take $`𝒟=10`$, then we have $`\tau _{\gamma \gamma }^{\mathrm{rsc}}(ϵ_{\mathrm{obs}})=3.2`$. Therefore it is expected that no TeV emission occurs in this object. However it is interesting to note that this is due to the intrinsic mechanism. We hope that there will be some effects due to the presence of pair production in the VHE flare (Wang, Zhou & Cheng 2000). The Q1633+382 (Mattox et al 1993) is quite similar to 3C 279, but it shows much smaller $`l_{\mathrm{rsc}}<10^2`$. The strong reflected synchrotron photons as seed photons may appear in this object, however its redshift ($`z=0.181`$) is too large to detect VHE photons due to the absorptions of back ground photons. Here we suggest that the deficiency of VHE emission may be intrinsic. It is expected to make simultaneous observations at other bands to test its light curves in order to reach a decision.
## 4. Conclusions and discussions
The present paper focuses our attention on the effects of BLR mirror on the attenuation of $`\gamma `$-ray in blazars. The mirror effect mainly depends on two parameters: Lorentz factor of the bulk motion ($`\mathrm{\Gamma }`$) and the Thomson scattering depth ($`\tau _{_{\mathrm{BLR}}}`$) of broad line region. Based on the calculations, we would like to draw the conclusions:
1) The parameters $`l_{\mathrm{rsc}}`$ and Doppler factor $`𝒟`$ in FSRQs are systematically greater than that in BL Lacs. This will cause the more stronger “intrinsic” absorption of VHE photons in FSRQs than that in BL Lacs. It is predicted that there is general absence of very high energy emission in FSRQs, owning to the attenuation of VHE photons by the BLR reflection of synchrotron emission.
2) The mirror model provides a new constraint on relativistic bulk motion. That intrinsic absorption of TeV photons may operate in some objects, especially in FSRQs. This constraint is cause by the motion of blob itself.
Although the origin of $`\gamma `$-ray emission in blazars still remains open, VHE observations sets strong constraints on blazar’s radiation mechanism. These constraints are: (1) brightness temperature exceeding the Kellermann-Pauliny-Toth (Begelman et al 1994), (2) multiwavelength light curves based on the homogeneous model (Mastichiadis & Kirk 1997), (3) high energy variations in X-ray and $`\gamma `$-ray including interaction with background IR radiation (Coppi & Aharonian 1999). These constraints are mainly based on SSC model. The deficiency of VHE photons from high redshift $`\gamma `$-ray loud blazars may be explained by the interaction of the cosmic background radiation fields with the VHE photons. However the possible alternative mechanism may be due to the intrinsic attenuation by the reflected synchrotron photons. Three BL Lac objects have been found to show $`H\alpha `$ and $`H\beta `$ emission, indicating the existence of broad line region in these objects (Vermeulen et al 1995, Corbett et al 1996), even Mrk 421 has been detected weak luminosity of a broad emission line (Morgani, Ulrich & Tadhunter 1992). The increasing evidence of the presence of broad emission lines in BL Lacs lends the possibility that the reflected external photons might be the main source of seed photon in this kind of blazars. Distinguishing the two different mechanisms might be traced by the following-up observations in other wavebands because a pair cascade process may be developed, forming a pair halos in the external absorption (Aharonian, Coppi & Voelk 1994). Such an extended halo due to external absorption may be of very long variable timescale at least $`10^3`$ yr (corresponding to one mean free path) \[see equation (24)\]. However if the intrinsic absorption works in the central engine, the time-dependent synchrotron self-Compton model including pair cascade (Wang, Zhou & Cheng 2000) could predict the interesting spectrum and light curves, which may interpret the variations of PKS 2155-304 (Urry et al 1997). The other radiative properties of such an extended pair halo are needed to be studied in order to distinguish the intrinsic absorption from the external one.
The author is very grateful to the anonymous referee for the physical insight of comments and suggestions, especially on the discussions on the angular distribution of reflected synchrotron photons and its effects on the opacity of pair production. I thank C.-C. Wang and F.-J. Lu for their careful reading of the manuscript and interesting discussions. The simulating discussions with Y.-Y. Zhou, T.-P. Li, M. Wu and B.-F. Liu are gratefully acknowledged. This research is supported by ’Hundred Talents Program of CAS’ and Natural Science Foundation of China under Grant No. 19803002. |
warning/0002/gr-qc0002018.html | ar5iv | text | # 1 Introduction
## 1 Introduction
If the classical gravitational collapse of a star ends in a naked singularity, then it is of interest to know how the quantum evaporation of such a star (via particle creation) compares with the quantum evaporation (via Hawking radiation) of another star that ends as a black hole. A few studies have been carried out in recent years, with an attempt to answer this question. These studies can be traced back to the pioneering works of Ford and Parker , and of Hiscock et al. .
Ford and Parker studied the outgoing quantum flux of a massless scalar field, using the geometric optics approximation, in the background geometry of a classical, spherical dust star forming a shell-crossing naked singularity. They found that this flux does not diverge in the approach to the Cauchy horizon. Their method was used by Barve et al. to calculate the quantum flux in the geometry of a spherical self-similar dust cloud which develops a shell-focusing naked singularity - in this model the flux diverges on the Cauchy horizon, in a positive sense.
Hiscock et al. studied the outgoing flux in the self-similar Vaidya model by obtaining the stress-energy tensor in the effective 2-d Vaidya background from the trace anomaly and the conservation equations. This calculation is in the same spirit as the geometric optics approximation in the sense that only the lowest angular momentum modes are accounted for. However, it has the advantage of yielding the components of the stress tensor everywhere in the spacetime and not just asymptotically. Once again they obtained a divergent flux on the Cauchy horizon. Barve et al. repeated their analysis for the self-similar dust collapse model and obtained similar results.
A new facet was added to these studies by Vaz and Witten , who demonstrated how one could calculate the spectrum of the outgoing radiation, in the approach to the Cauchy horizon, for the self-similar dust model. The spectrum was shown to be non-thermal, and this is a characteristic feature distinguishing the evaporation of naked singularities from thermal black hole evaporation. A detailed justification on the use of Bogoliubov transformations in the presence of a Cauchy horizon has been given by us recently .
In the present paper, we apply the geometric optics approximation to the quantization of a massless scalar field in the background of a spherical collapsing self-similar Vaidya star. Such a star results in black hole formation for part of the initial data, and a naked singularity for the remaining initial data. In this approximation, we calculate the flux of the scalar field emitted to infinity, and the spectrum of the emitted radiation. The essential quantity needed for these calculations is the map from an ingoing null ray to an outgoing null ray which passes through the center of the collapsing cloud. This map is found in Section 2. In Section 3 we show that the flux diverges on the Cauchy horizon, and that the spectrum of the radiation is non-thermal.
These results are identical to those that have been found earlier for the model of self-similar dust collapse , . They provide additional evidence for the conjecture that the divergence of the outgoing flux on the Cauchy horizon and a non-thermal radiation spectrum are generic features of naked singularities.
## 2 The Classical Solution
In this Section, we find the mapping from an ingoing null ray to an outgoing null ray, for those initial conditions which result in a naked singularity or black hole, in the self-similar Vaidya model. The metric is given by
$$ds^2=\left(1\frac{2m(v)}{R}\right)dv^22dvdRR^2d\mathrm{\Omega }^2$$
(1)
where $`v`$ is the advanced time coordinate $`(\mathrm{}<v<\mathrm{})`$ and $`R`$ is the area radius $`(0R<\mathrm{})`$. The mass function is zero for $`v<0`$, it is $`m(v)=\mu v`$ for $`0<v<v_0`$, and equal to a constant $`M=\mu v_0`$ for $`vv_0`$. This represents a collapsing ball of null dust bounded in the region $`0<vv_0.`$ The linearity of the mass function ensures that this Vaidya spacetime is self-similar, i.e. it possesses a homothetic Killing vector field. Clearly, spacetime is Minkowski for $`v<0`$, and here the metric may be written as
$$ds^2=dUdvR^2d\mathrm{\Omega }^2$$
(2)
where $`U=v2R`$ is the retarded time coordinate. Outside the Vaidya region, the geometry is Schwarzschild, and in terms of the Eddington-Finkelstein null coordinates, is given by
$$ds^2=\left(1\frac{2M}{R}\right)d\widehat{U}dvR^2d\mathrm{\Omega }^2,$$
(3)
with $`\widehat{U}=v2R^{}`$, and $`R^{}=R+2M\mathrm{ln}\left(R/2M1\right)`$.
We now obtain double null coordinates in the Vaidya region, by first defining $`x=v/R`$ and $`z=\mathrm{ln}R,`$ so that the metric (1) becomes
$$ds^2=\mathrm{exp}(2z)\left[A(x)dx^2+B(x)dz^2+2C(x)dxdz\right].$$
(4)
Here,
$$A(x)=12\mu x,B(x)=\left(12\mu x\right)x^22x,C(x)=\left(12\mu x\right)x1.$$
(5)
Next, we define
$$d\tau =dz+\frac{C(x)}{B(x)}dx,$$
(6)
and
$$d\chi =\left[\frac{C^2(x)A(x)B(x)}{B^2(x)}\right]^{1/2}dx,$$
(7)
so that the metric can be written as
$$ds^2=\mathrm{exp}(2z)B(x)\left[d\tau ^2d\chi ^2\right].$$
(8)
We define the double-null coordinates
$$u=\mathrm{exp}(\tau \chi ),v=\mathrm{exp}(\tau +\chi ),$$
(9)
so that
$$ds^2=\frac{\mathrm{exp}(2z)B(x)}{uv}dudv.$$
(10)
By using the identity $`C^2(x)A(x)B(x)=1`$, the null coordinates may be written as
$$u=r\mathrm{exp}I_+,v=r\mathrm{exp}I_{}$$
(11)
where
$$I_\pm =\frac{C(x)\pm 1}{B(x)}𝑑x.$$
(12)
It is easily shown that $`I_{}=\mathrm{ln}x`$, so that $`v`$ is trivially the same as the advanced time coordinate introduced in the beginning. We have chosen to write the null coordinates in this manner essentially to show that they can be constructed in a symmetric fashion. The integral $`I_+`$ can also be easily carried out, and we get
$$u=R\left(x\alpha _+\right)^{A_+}(x\alpha _{})^A_{},$$
(13)
where
$$\alpha _\pm =\frac{1\pm \left(116\mu \right)^{1/2}}{4\mu }$$
(14)
are the roots of the quadratic equation
$$x^2\frac{x}{2\mu }+\frac{1}{\mu }=0,$$
(15)
and the coefficients $`A_+`$ and $`A_{}`$ are given by
$$A_+=\frac{12\mu \alpha _+}{\alpha _+\alpha _{}},A_{}=\frac{12\mu \alpha _{}}{\alpha _{}\alpha _+}.$$
(16)
We note that $`A_++A_{}=1`$. Also, the ingoing null coordinate used by Hiscock et al. is obtained by taking the logarithm of the coordinate $`u`$ defined in (13).
The flat spacetime limit of the Vaidya metric is obtained by setting $`\mu =0`$. In this limit, $`v=t+R`$, and it is easily shown that in this limit the coordinate $`u`$ reduces to $`tR`$.
The roots $`\alpha _\pm `$ are real for $`\mu 1/16`$ and complex for $`\mu >1/16`$. We note that when the roots are complex, $`u`$ continues to be real, and can be written as
$$u=R|x\alpha _+|^2\mathrm{exp}(2\xi ImA_+)$$
(17)
where $`\xi `$ is the phase of $`\alpha _+`$.
As is known, the collapse of the cloud results in the formation of a curvature singularity, when the area radius of a shell shrinks to zero. The Kretschmann scalar is given by
$$R^{\alpha \beta \gamma \delta }R_{\alpha \beta \gamma \delta }=\frac{48m^2(v)}{R^6}=48\mu ^2x^6v^4.$$
(18)
The central singularity forms when the radius of the innermost shell, given by $`v=0`$, becomes zero. The curvature diverges when the center $`R=0`$ is approached along the line $`v=0`$, with $`x=v/R`$ a finite non-zero constant. It is known that this singularity is naked for $`\mu 1/16`$ and covered for $`\mu >1/16`$. The singularity corresponding to all shells with $`v>0`$ is covered.
The nature of the central singularity, i.e. as to whether it is naked or not, can also be deduced from the construction of the double null coordinates given above. We note from (13) that for $`v<0`$ the center (i.e. $`R=0`$) is given by the line $`u=v`$. Let us assume that the singular point $`v=0,R=0`$ is approached along the line $`x=x_0`$, with $`x_0`$ a finite, non-zero constant. Then, from (17) we see that the ratio $`u/v`$ is given by
$$\frac{u}{v}=d=\frac{|x_0\alpha _+|^2\mathrm{exp}(2\xi ImA_+)}{x_0}$$
(19)
with $`d`$ in general not equal to one. The singular point $`v=0,R=0`$ is thus the intersection of two lines, $`u=v`$ and $`u=dv`$. Equivalently, it is given by $`u=v=0`$.
If the central singularity is naked, there will be at least one outgoing null ray from the singularity. This means that such a ray, for which $`u=0`$, has $`R0`$. From (13) we see that this is possible if and only if $`x=\alpha _+`$ or $`x=\alpha _{}`$, with $`\alpha _\pm `$ real. The existence of a real root is necessary for the central singularity to be naked, implying that $`\mu 1/16`$. The Cauchy horizon, which is the first null ray to leave the singularity, is given by $`x=\alpha _{}`$, where $`\alpha _{}`$ is the smaller of the two roots. The Eddington-Finkelstein null coordinate $`\widehat{U}`$ corresponding to this ray is found by noting that when the ray meets the boundary $`v_0`$ of the star, $`R=v_0/\alpha _{}`$. Using this value of $`R`$ in the definition of $`\widehat{U}`$ we get
$$\widehat{U}=v_0\frac{2v_0}{\alpha _{}}4M\mathrm{ln}|\frac{v_0}{2M\alpha _{}}1|.$$
(20)
This value is not infinite, which shows that the singularity is globally naked.
We can now work out the map $`(v)`$ between an ingoing null ray and the outgoing null ray to which it corresponds, i.e. we can find the function $`\widehat{U}=(v)`$, for both the black hole case and the naked singularity case. In the naked case, we do this for outgoing rays with a value of $`u`$ close to zero. Consider a null ray, incoming from $`^{}`$, given by $`v=`$ constant, with $`v<0`$. After passing through the center, this becomes an outgoing ray $`u`$, with $`u=v`$. For $`u`$ nearly zero, $`x`$ is close to $`\alpha _{}`$, so we write $`x=\alpha _{}+\widehat{x}`$. From (13) we can write
$$u=BR\widehat{x}^A_{},B=(\alpha _{}\alpha _+)^{A_+}.$$
(21)
By noting that
$$R=\frac{v_0}{\alpha _{}+\widehat{x}}\frac{v_0}{\alpha _{}}\left(1\frac{\widehat{x}}{\alpha _{}}\right)$$
(22)
and by using this relation in (21) we get
$$\widehat{x}=\left(\frac{\alpha _{}}{Bv_0}\right)^{1/A_{}}u^{1/A_{}}.$$
(23)
Using these expressions for $`R`$ and $`\widehat{x}`$ in the definition of $`\widehat{U}`$ we can conclude that
$$\widehat{U}=\widehat{U}^0Qv^{1/A_{}}=(v)$$
(24)
where $`\widehat{U}^0`$ is the Cauchy horizon, and $`Q`$ is an (irrelevant) constant. This is the desired map in the naked singularity case. The inverse function, $`v=𝒢(\widehat{U})`$ is given by
$$v=𝒢(\widehat{U})=\left(\frac{\widehat{U}^0\widehat{U}}{Q}\right)^A_{}.$$
(25)
In order to find the function $`(v)`$ for the case in which the collapse ends in a black hole (i.e. $`\mu >1/16`$), we examine rays which arrive at $`^+`$ just before the event horizon. The event horizon is the outgoing ray which reaches the boundary $`v_0`$ when the boundary has an area radius $`R=2M`$, which implies that $`x=v_0/2M`$ for such a ray. By writing $`x=\widehat{x}+v_0/2M`$, we easily note that to leading order $`\widehat{U}=4M\mathrm{ln}\widehat{x}`$ for such a ray. Further, since $`\widehat{x}`$ is small, $`u`$ is of the form $`u=a+b\widehat{x}`$, and using $`u=v`$ for a ray that goes through the center, we can conclude that
$$\widehat{U}=4M\mathrm{ln}\left(\frac{va}{b}\right)=(v)$$
(26)
which is the map in the black hole case. The inverse function, $`v=𝒢(\widehat{U})`$ is now given by
$$v=𝒢(\widehat{U})=a+be^{\widehat{U}/4M}.$$
(27)
## 3 The Quantum Flux
We consider next the quantization of a massless scalar field $`\varphi (x)`$ in the background geometry of the collapsing Vaidya cloud. We will calculate the flux of the energy radiated to $`^+`$, and the spectrum of the emitted radiation, in the geometric optics approximation. The flux radiated to infinity is given by the off-diagonal component of the stress-energy tensor of the massless scalar field, as
$$P={}_{M}{}^{}0|T_T^R|0_{M}^{}R^2\mathrm{sin}\theta d\theta d\varphi =\frac{1}{24\pi }[\frac{^{\prime \prime \prime }}{(^{})^3}\frac{3}{2}\left(\frac{^{\prime \prime }}{^{}_{}{}^{}2}\right)^2]$$
(28)
where the function $`(v)`$ has been calculated above. $`|0_M`$ is the Minkowski vacuum. For the naked singularity case we substitute the expression (24) for $`(v)`$ and find the radiated power to be
$$P(v)=\frac{1}{48\pi Q^2}\left[\frac{A_{}^21}{v^{2/A_{}}}\right].$$
(29)
We note that $`A_{}`$ is greater than one. Hence this expression diverges, in a positive sense, in the approach to the Cauchy horizon, i.e. as $`v0`$.
In the black hole case we use the $`(v)`$ from (26) in the expression for the radiated power, to get
$$P(v)=\frac{1}{48\pi M^2}$$
(30)
as expected.
In order to calculate the spectrum we recall that the number distribution of Minkowski particles observed on $`^+`$ is simply given by the Bogoliubov coefficient
$$\beta (\omega ^{},\omega )=_{\mathrm{}}^{\mathrm{}}\frac{d\widehat{U}}{4\pi \sqrt{\omega \omega ^{}}}e^{i\omega \widehat{U}}e^{i\omega ^{}𝒢(\widehat{U})}$$
(31)
(where the integral is performed over all of $`^+`$) as
$$_M0|N(\omega )|0_M=_0^{\mathrm{}}d\omega ^{}|\beta (\omega ^{},\omega )|^2$$
(32)
where $`|0_M`$ is the Minkowski vacuum.
Substituting for $`𝒢(\widehat{U})`$ from (27) for the black hole case gives
$$|\beta (\omega ^{},\omega )|^2=\frac{2M}{\pi \omega ^{}}\frac{1}{e^{8\pi M\omega }1}$$
(33)
which is the expected Hawking spectrum.
On the other hand, for the naked singularity case, we take the $`𝒢(\widehat{U})`$ from (25) to get
$$\beta (\omega ^{},\omega )=\frac{1}{2\pi }\sqrt{\frac{\omega }{\omega ^{}}}_{\mathrm{}}^{\widehat{U}^o}𝑑\widehat{U}e^{i\omega \widehat{U}}e^{i\omega ^{}Q^A_{}(\widehat{U}^o\widehat{U})^A_{}}.$$
(34)
Changing variables to $`z=(\widehat{U}^o\widehat{U})`$, one has
$$\beta (\omega ^{},\omega )=\frac{1}{2\pi }\sqrt{\frac{\omega }{\omega ^{}}}e^{i\omega \widehat{U}^o}_0^{\mathrm{}}𝑑ze^{i\omega z}e^{i\omega ^{}(z/Q)^A_{}}$$
(35)
which gives
$$|\beta (\omega ^{},\omega )|^2=\frac{1}{4\pi ^2\omega \omega ^{}}|\underset{k=0}{\overset{\mathrm{}}{}}\frac{(iQ^A_{}\omega ^{}\omega ^A_{}e^{i\pi A_{}/2})^k}{k!}\mathrm{\Gamma }(kA_{}+1)|^2.$$
(36)
This is a non-thermal spectrum, similar to the spectrum in the case of self-similar dust collapse .
## 4 Conclusion
The divergence of the radiated power on the Cauchy horizon has been shown by us to be independent of the assumption of self-similarity. The non-thermality of the spectrum is due to the fact that not all of $`^+`$ is actually probed by infalling waves from $`^{}`$. We have, of course, implicitly assumed that all of $`^+`$ exists, i.e., that spacetime may be analytically continued beyond the Cauchy horizon. This assumption is reasonable in that one does not expect a local collapse to terminate the universe. However, only a complete theory of quantum gravity can provide the answer to the question of the existence of a complete $`^+`$, as this depends on the final fate of the collapse. Therefore it is interesting to consider what might happen if the continuation beyond the Cauchy horizon is not physically acceptable. If the Cauchy horizon is to be regarded as the end point of spacetime, then $`\widehat{U}`$ is not a good asymptotic coordinate. However, one can transform to an asymptotically flat retarded time and it can be shown that the radiation is thermal, but with a temperature that is different from the Hawking temperature. This is analogous to the marginally naked singularity treated by Hiscock et al. . |
warning/0002/astro-ph0002254.html | ar5iv | text | # 1 Introduction
## 1 Introduction
The principle of equivalence of inertial and gravitational masses underlines the theory of gravity. In the Einstein theory of gravity , this leads to that the free gravitational field is nonlocalized. Under the presence of the matter, the gravity is described by the Einstein equations
$$G_{ik}=T_{ik}$$
(1)
where $`G_{ik}`$ is the Einstein tenzor, $`T_{ik}`$ is the tenzor of momentum-energy of the matter. Free gravitational field defined by the absence of the matter $`T_{ik}=0`$ is described by the equations
$$R_{ik}=0$$
(2)
where $`R_{ik}`$ is the Ricci tenzor. The localized field must be described by the tenzor of momentum-energy. Einstein characterized the momentum-energy of the gravitational field by the pseudo-tenzor defined as
$$t^{ik}=H_{,lm}^{ilkm}G^{ik}$$
(3)
where $`H_{,lm}^{ilkm}`$ is the linearized part of $`G_{ik}`$. Thus in the Einstein theory gravitational field is nonlocalized.
The natural way proposed by Lorentz and Levi-Civita is to take $`G_{ik}`$ as the momentum-energy of the gravitational field. However in this case $`G_{ik}`$ is equal to zero for the free gravitational field $`G_{ik}=R_{ik}=0`$. Such a situation may be interpreted as that the gravitational interaction occurs without gravitational field. Then the problem arises as to how to introduce gravitational radiation. The possible resolution of the problem is to introduce some material field as a radiation.
## 2 Theory
Consider gravitational interaction within the framework of the Newton theory and the quantum field theory. The Lagrangian of the Newton gravity is given by
$$L=G\frac{m^2}{r},$$
(4)
with the mass m being the gravitational charge. While expressing the Newton constant $`G`$ via the charge $`g=(\mathrm{}c)^{1/2}`$ and via the Planck mass $`m_{Pl}=(\mathrm{}c/G)^{1/2}`$, the Lagrangian (4) can be rewritten in the form
$$L=G\frac{m^2}{r}=\frac{g^2}{m_{Pl}^2}\frac{m^2}{r}.$$
(5)
The Lagrangian of the Newton gravity in the form (5) describes gravitational interaction by means of the charge $`g`$. In this way gravity may be implemented into the quantum field theory.
Rewrite the Lagrangian (5) in the form of the effective Lagrangian of interaction of the spinor fields
$$L=\frac{g^2}{m_{Pl}^2}J_\mu (x)J^\mu (x).$$
(6)
The term $`1/m_{Pl}^2`$ in the Lagrangian (6) reads that gravitational interaction takes place at the Planck scale. At the same time gravity is characterized by the infinite radius of interaction. To resolve the problem consider the scheme of gravitational interaction which includes both the short-range interaction and the long-range interaction
$$L=L_{short}+L_{long}.$$
(7)
Let us introduce the Planck neutrino $`\nu _{Pl}`$. Let the Planck neutrino is the massless particle of the spin 1/2. Suppose that the Planck neutrino interacts with the other fields at the Planck scale
$$\psi \nu _{Pl}$$
(8)
where $`\psi `$ denotes all the fields of the spin 1/2. This interaction is of short-range and is governed by the Lagrangian (6)
$$L_{short}=\frac{g^2}{m_{Pl}^2}J_\mu (x)J^\mu (x)$$
(9)
where the current $`J_\mu `$ transforms the field $`\psi `$ into the field $`\nu _{Pl}`$. Let the interaction of the Planck neutrinos $`\nu _{Pl}\nu _{Pl}`$ is of long-range and is governed by the Lagrangian identically equal to zero
$$L_{long}0.$$
(10)
The considered scheme allows one to decribe both the classical gravity and the decay of the field $`\psi `$ into the Planck neutrino. In this scheme gravitational radiation can be identified with the Planck neutrino.
Within the framework of the standard quantum field theory, the above scheme of gravitational interaction should include two intermediate fields $`\psi \nu _{Pl}`$ and $`\nu _{Pl}\nu _{Pl}`$. Since the Lagrangian of the interaction $`\nu _{Pl}\nu _{Pl}`$ is identically equal to zero, the energy of the field $`\nu _{Pl}\nu _{Pl}`$ is identically equal to zero. The field $`\psi \nu _{Pl}`$ is defined by the Planck mass. In the theory of gravity there is the limit of ability to measure the length equal to the Planck length $`\mathrm{\Delta }l2(\mathrm{}G/c^3)^{1/2}=2l_{Pl}`$. From this it follows that there is no possibility to measure the field $`\psi \nu _{Pl}`$ in the physical experiment. Thus both intermediate fields $`\psi \nu _{Pl}`$ and $`\nu _{Pl}\nu _{Pl}`$ cannot be measured. This means that the intermediate fields $`\psi \nu _{Pl}`$ and $`\nu _{Pl}\nu _{Pl}`$ do not exist. We arrive at the conclusion that gravitational interaction occurs without intermediate fields.
## 3 The lifetime of proton relative to the decay into Planck neutrino
In view of eq. (8), the decay of proton into Planck neutrino occurs at the Planck scale
$$p\nu _{Pl}.$$
(11)
The lifetime of proton relative to the decay into Planck neutrino is defined by the Lagrangian (6)
$$t_p=t_{Pl}\left(\frac{m_{Pl}}{2m_p}\right)^5$$
(12)
where the factor 2 takes into account the transition from the massive particle to the massless one. This lifetime corresponds to the rest frame. Consider the lifetime of proton in the moving frame with the Lorentz factor
$$\gamma =\left(1\frac{v^2}{c^2}\right)^{1/2}.$$
(13)
In the moving frame the rest mass and time are multiplied by the Lorentz factor
$$m^{}=\gamma m$$
(14)
$$t^{}=\gamma t.$$
(15)
The Planck mass $`m_{Pl}=(\mathrm{}c/G)^{1/2}`$ and the Planck time $`t_{Pl}=(\mathrm{}G/c^5)^{1/2}`$ are built from three fundamental constants $`\mathrm{}`$, $`c`$ and $`G`$. According to the special relativity , the speed of light is fixed in all the inertial frames. Extend the special relativity principle and suppose that the three constants $`\mathrm{}`$, $`c`$ and $`G`$ are fixed in all the inertial frames
$$\mathrm{}^{}=\mathrm{}c^{}=cG^{}=G.$$
(16)
Hence the Planck mass and time are fixed in all the inertial frames. Then the lifetime of proton in the moving frame is given by
$$t_p^{}=t_{Pl}\left(\frac{m_{Pl}}{2\gamma m_p}\right)^5.$$
(17)
For comparison consider the decay of muon which is governed by the Lagrangian of electroweak interaction
$$L=\frac{g^2}{m_W^2}J_\mu (x)J^\mu (x)$$
(18)
where $`m_W`$ is the mass of W-boson. In the rest frame the lifetime of muon is given by
$$t_\mu =t_W\left(\frac{m_W}{m_\mu }\right)^5.$$
(19)
In the moving frame the lifetime of muon is given by
$$t_\mu ^{}=\gamma t_W\left(\frac{m_W}{m_\mu }\right)^5.$$
(20)
Thus unlike the usual situation when the lifetime of the particle, e. g. muon, grows with the Lorentz factor as $`\gamma `$, the lifetime of proton relative to the decay into Planck neutrino decreases with the Lorentz factor as $`\gamma ^5`$. State once again that such a behaviour is due to that the Planck mass built from three fundamental constants $`\mathrm{}`$, $`c`$ and $`G`$ is fixed in all the inertial frames.
## 4 Extra high energy cosmic rays spectrum in view of the decay of proton into Planck neutrino
In view of eq. (17), the lifetime of proton relative to the decay into Planck neutrino decreases with the increase of the kinetic energy of proton. Then the decay of proton can be observed for the extra high energy protons. In particular the decay of proton can be observed as a cut-off in the energy spectrum of extra high energy cosmic rays (EHECRs).
The EHECRs spectrum above $`10^{10}\mathrm{eV}`$ can be divided into three regions: two ”knees” and one ”ankle” . The first ”knee” appears around $`3\times 10^{15}\mathrm{eV}`$ where the spectral power law index changes from $`2.7`$ to $`3.0`$. The second ”knee” is somewhere between $`10^{17}\mathrm{eV}`$ and $`10^{18}\mathrm{eV}`$ where the spectral slope steepens from $`3.0`$ to around $`3.3`$. The ”ankle” is seen in the region of $`3\times 10^{18}\mathrm{eV}`$ above which the spectral slope flattens out to about $`2.7`$.
Consider the EHECRs spectrum in view of the decay of proton into Planck neutrino. Let the earth be the rest frame. For protons arrived at the earth, the travel time meets the condition
$$tt_p.$$
(21)
From this the time required for proton travel from the source to the earth defines the limiting energy of proton
$$E_{lim}=\frac{m_{Pl}}{2}\left(\frac{t_{Pl}}{t}\right)^{1/5}.$$
(22)
Within the time $`t`$, protons with the energies $`E>E_{lim}`$ decay and do not give contribution in the EHECRs spectrum. Thus the energy $`E_{lim}`$ defines a cut-off in the EHE proton spectrum. Planck neutrinos appeared due to the decay of the EHE protons may give a contribution in the EHECRs spectrum. If the contribution of Planck neutrinos in the EHECRs spectrum is less compared with the contribution of protons one can observe the cut-off at the energy $`E_{lim}`$ in the EHECRs spectrum.
Determine the range of the limiting energies of proton depending on the range of distances to the EHECRs sources. Take the maximum and minimum distances to the source as the size of the universe and the thickness of our galactic disc respectively. For the lifetime of the universe $`\tau _0=14\pm 2\mathrm{Gyr}`$ , the limiting energy is equal to $`E_1=3.9\times 10^{15}\mathrm{eV}`$. This corresponds to the first ”knee” in the EHECRs spectrum. For the thickness of our galactic disc $`300\mathrm{pc}`$, the limiting energy is equal to $`E_2=5.5\times 10^{17}\mathrm{eV}`$. This corresponds to the second ”knee” in the EHECRs spectrum. Thus the range of the limiting energies of proton due to the decay of proton into Planck neutrino lies between the first ”knee” $`E3\times 10^{15}\mathrm{eV}`$ and the second ”knee” $`E10^{17}10^{18}\mathrm{eV}`$.
From the above consideration it follows that the decrease of the spectral power law index from $`2.7`$ to $`3.0`$ at the first ”knee” $`E3\times 10^{15}\mathrm{eV}`$ and from $`3.0`$ to around $`3.3`$ at the second ”knee” $`E10^{17}10^{18}\mathrm{eV}`$ can be explained as a result of the decay of proton into Planck neutrino. From this it seems natural that, below the ”ankle” $`E<3\times 10^{18}\mathrm{eV}`$, the EHECRs events are mainly caused by the protons. Above the ”ankle” $`E>3\times 10^{18}\mathrm{eV}`$, the EHECRs events are caused by the particles other than protons.
If Planck neutrinos take part in the strong interactions, they must give some contribution in the EHECRs events. To explain the observed EHECRs spectrum it is necessary to assume that the contribution of Planck neutrinos in the EHECRs spectrum is less compared with the contribution of protons. Suppose that proton decays into 5 Planck neutrinos. Then the energy of the Planck neutrino is $`1/5`$ of the energy of the decayed proton. For the spectral power law index equal to $`2.7`$, the ratio of the proton flux to the Planck neutrino flux is given by $`J_p/J_\nu =5^{1.7}=15.4`$.
From the above consideration it is natural to identify EHE particles with the energies $`E>3\times 10^{18}\mathrm{eV}`$ with the Planck neutrinos. Continue the curve with the spectral power law index $`2.7`$ from the ”ankle” $`E3\times 10^{18}\mathrm{eV}`$ to the first ”knee” $`E3\times 10^{15}\mathrm{eV}`$ and compare the continued curve with the observational curve. Comparison gives the ratio of the proton flux to the Planck neutrino flux $`J_p/J_\nu 15`$. |
warning/0002/quant-ph0002002.html | ar5iv | text | # Upper Bound on the region of Separable States near the Maximally Mixed State
## I Introduction
A key distinguishing feature of quantum physics from classical physics is the prediction of a new kind of correlation between physical quantities, called entanglement. Quantum entanglement has been often been referred to as the inseparability of composite quantum systems. Such an entangled composite system is said to be inseparable because it cannot be prepared by manipulating each subsystem separately, using only measurements and operations local to one subsystem at a time. If a composite quantum mechanical state is specified by some density matrix, how can we tell if the system is entangled?
Much work has been done studying the particular case of two subsystems, with each in a two dimensional Hilbert space (qubits). There is a good understanding of the entanglement for such systems and in fact a criterion, the partial transposition condition of Peres, indicates whether the subsystems are entangled. This, however, is a necessary and sufficient condition only when there are two subsystems, one with Hilbert space dimension $`2`$, and the other of dimension $`2`$ or $`3`$, as was shown by the Horodeckis. For more complex systems, it only determines whether the state contains distillable entanglement, however there are also some states with bound entanglement, which are lumped together with the separable states by this criterion.
Lewenstein et. al. have used the Peres condition to consider two subsystems, but with each subsystem now in a $`N>2`$ dimensional Hilbert space. Życzkowski et. al. have, among other results, shown that all the mixed states in a sufficiently small neighbourhood of the maximally mixed state are separable. They also gave a bound on the size of this neighbourhood for small composite systems. Vidal and Tarrach gave a lower bound on the size of this neighbourhood for arbitrary composite states, of any number of subsystems. Schack and Caves gave bounds for composite systems composed of many qubits ($`N=2`$).
It has been pointed out by Braunstein et. al. that Maximally-entangled states of the GHZ type, with noise added, are connected to recent proposals for NMR quantum computing. They are also relevant to fundamental tests of quantum mechanics using Bell inequalities. The separability of such noisy generalised singlet states has been considered by various authors recently, mostly for the case of qubits (subsystems with Hilbert space dimension two). Schack and Caves, using the approach of Braunstein et. al. have obtained an exact boundary condition for the Werner states (two qubits), while Caves and Milburn extended the approach to two q-trits (Hilbert space dimension three). The Horodeckis gave exact bounds for the case of two subsystems.
In this paper we extend the approach of Peres to consider such noisy generalised singlet states in the general case of $`D`$ subsystems with each subsystem in a $`N`$ dimensional Hilbert space. A parameter $`ϵ`$ specifies the amount of the pure GHZ-type maximally-entangled states present compared with the maximally-mixed noise state.We then ask and answer the following question: what is the maximum value of $`ϵ`$ for this state to be entangled? Another fundamental question is also considered. To create as entangled a state as possible, is it simply better to increase the dimension of the subsystems or is it better to increase the number of subsystems?
## II Generalised Werner states
The states considered here, consist of a mixture of the maximally-mixed noise states, and the GHZ-type entangled states, where the number $`D`$ of subsystems over which entanglement occurs, and the Hilbert space dimension $`N`$ of each subsystem (equal for all $`D`$ subsystems) can take on any values $`N>1`$, $`D>1`$. The relative proportion of GHZ-type states is controlled by the parameter $`\epsilon `$, where $`\epsilon =1`$ corresponds to a pure GHZ-type state, while $`\epsilon =0`$ corresponds to the maximally mixed state. Explicitly, the state is given by the density operator
$`\widehat{\rho }`$ $`=`$ $`(1\epsilon )\widehat{\rho }_n+\epsilon \widehat{\rho }_e,`$ (2)
$`\widehat{\rho }_n`$ $`=`$ $`{\displaystyle \frac{1}{N^D}}{\displaystyle \underset{i_1,i_2,\mathrm{},i_D=1}{\overset{N}{}}}|i_1i_2\mathrm{}i_Di_1i_2\mathrm{}i_D|,`$ (3)
$`\widehat{\rho }_e`$ $`=`$ $`{\displaystyle \frac{1}{N}}{\displaystyle \underset{n,m=1}{\overset{N}{}}}|nn\mathrm{}nmm\mathrm{}m|,`$ (4)
where
$`|i_1i_2\mathrm{}i_D=|i_1|i_2\mathrm{}|i_D.`$ (5)
Here $`|i_k`$ represents one of $`N`$ complete orthogonal basis states for the $`k^{\mathrm{th}}`$ subsystem, while $`\widehat{\rho }_n`$ is the density operator for the maximally mixed state, and $`\widehat{\rho }_e`$ for the GHZ-like state.
The Werner state, consisting of a singlet state and some noise, is the simplest $`D=N=2`$ case, hence we call these generalised Werner states. The state $`\widehat{\rho }`$ can be viewed as a generalised singlet state $`\widehat{\rho }_e`$ after it has emerged from a depolarising channel.
## III Separability
The separability of particular cases of the state (II), and of more general states, has been considered previously by a number of authors . Firstly, it was shown that for the two-qubit ($`N=2,D=2`$) case, (The Werner state) $`\widehat{\rho }`$ is separable for $`\epsilon 1/3`$, and entangled otherwise, whereas for three qubits ($`N=2,D=3`$), Schack and Caves found the state to be separable for $`\epsilon 1/5`$. Soon after, Horodecki and Horodecki found the exact result for arbitrary numbers of qubits:
$$\epsilon _{\text{separable}}\frac{1}{1+N}.$$
(6)
Schack and Caves also found bounds on the size of the separable neighbourhood around the maximally mixed state for totally general states of many qubits ($`N=2`$). For $`D=4`$ and $`D=5`$ these are $`\epsilon 1/33`$ and $`\epsilon 1/243`$ respectively, and for higher $`D`$ are given by
$$\epsilon _{\text{separable}}\{\begin{array}{cc}1/(1+2^D+2^{2D2})\hfill & \mathrm{if}\mathrm{D}\mathrm{even}\hfill \\ 1/(12^D+2^{2D2})\hfill & \mathrm{if}\mathrm{D}\mathrm{odd}\hfill \end{array}$$
(7)
What about the more general case when $`N`$ and $`D`$ are arbitrary? For what values of $`\epsilon `$ is the state given by (II) separable? Vidal and Tarrach gave a maximum bound for the random robustness $`R`$ of arbitrary multi-component states. For the states considered here, the critical value of $`\epsilon `$ at which the states change from being entangled to separable is $`\epsilon _c=1/(1+R(\widehat{\rho }_e||\widehat{\rho }_n))`$ (in the notation of ref.). So according to that bound,
$$\epsilon _{\text{entangled}}>\frac{1}{(1+N/2)^{D1}}$$
(8)
We will prove in the next section (section IV) that the states given by (II), are always entangled if
$`\epsilon _{\text{entangled}}>{\displaystyle \frac{1}{N^{D1}+1}}.`$ (9)
That this bound is strong for qubits ($`N=2`$) is shown in appendix (A), and an explicit decomposition into product states is given for the case of separable $`\widehat{\rho }`$.
## IV Outline of the proof of (6)
Peres has shown that a necessary condition for a state consisting of two subsystems to be separable is that the partial transpose of the density matrix over one of the subsystems, and the partial transpose over the other subsystem, have positive eigenvalues. However, this is a necessary and sufficient condition only when one of the subsystems has Hilbert space dimension $`2`$ or $`3`$, and the other dimension $`2`$, as was shown by the Horodeckis. That paper went on to give a necessary and sufficient condition for a state to be separable. Nevertheless the Peres condition is just what is needed for an upper bound on $`\epsilon `$ for separable states. i.e. any states which break the condition are entangled, although some which satisfy it may also be, but do not have to be, entangled. It is worth noting that any states which satisfy the Peres condition but are entangled are said to contain only “bound” entanglement, as it cannot be used for teleportation, nor distilled by the process of entanglement distillation.
Firstly, note that the Peres condition is easily extended to more than two entangled subsystems. If there are $`D`$ subsystems, one simply chooses some group of $`M<D`$ original subsystems to be called half-system number $`1`$, and the remaining subsystems to be called half-system number $`2`$. If for any such group of subsystems, an eigenvalue of the partial transpose of $`\widehat{\rho }`$ over half-system number $`1`$ (say) is negative, then $`\widehat{\rho }`$ is entangled. Thus to use the Peres condition to full advantage, one must consider all such groups of subsystems.
The state (II) is convenient in this respect, because it is unchanged under relabeling of the subsystems (evident by inspection). Thus the eigenvalues of $`\rho _\alpha ^T`$, the partial transpose of $`\widehat{\rho }`$ over the set $`\alpha `$ of subsystems, need only be looked at for $`D/2`$ (rounded down) sets of subsystems to extract the maximum benefit from the Peres condition. In particular, a choice of sets of subsystems can be $`\alpha _M=\{1,2,\mathrm{},M\}`$ where $`M=1,2,\mathrm{},D/2`$, and $`\alpha _M`$ contains the labels of the subsystems to be considered as members of the $`M`$th half-system.
Firstly let’s consider $`M=1`$, i.e. the first subsystem’s entanglement with the remaining $`D1`$ of them. The partially transposed density matrix is
$`\rho _1^T`$ $`=`$ $`(1\epsilon )\rho _n^T+\epsilon \rho _e^T,`$ (11)
$`\rho _n^T`$ $`=`$ $`{\displaystyle \frac{1}{N^D}}{\displaystyle \underset{i_1,i_2,\mathrm{},i_D=1}{\overset{N}{}}}|i_1i_2\mathrm{}i_Di_1i_2\mathrm{}i_D|,`$ (12)
$`\rho _e^T`$ $`=`$ $`{\displaystyle \frac{1}{N}}{\displaystyle \underset{n,m=1}{\overset{N}{}}}|mn\mathrm{}nnm\mathrm{}m|.`$ (13)
Since all of the elements of $`\widehat{\rho }`$, (hence $`\rho _1^T`$) are finite, the eigenvalues of $`\rho _1^T`$ must also be finite. Now, exploiting the general continuity property of eigenvalues of $`\widehat{\rho }_1^T`$, we conclude that at the value of $`\epsilon =\epsilon _o`$ above which the Peres condition indicates the state is entangled, one or more eigenvalues of $`\rho _1^T`$ must be zero, since they are all positive for $`\epsilon <\epsilon _o`$, and at least one is negative for $`\epsilon >\epsilon _o`$. i.e. for some nonzero eigenvector
$$|\psi =\underset{j_1,j_2,\mathrm{},j_D=1}{\overset{N}{}}\psi _{j_1j_2\mathrm{}j_D}|j_1j_2\mathrm{}j_D0$$
(14)
we must have $`\rho _1^T(\epsilon _o)|\psi =0`$. Expanded, this gives
$`{\displaystyle \frac{1\epsilon _o}{N^D}}{\displaystyle \underset{i_1,i_2,\mathrm{},i_D=1}{\overset{N}{}}}`$ $`\psi _{i_1i_2\mathrm{}i_D}`$ $`|i_1i_2\mathrm{}i_D`$ (15)
$`+{\displaystyle \frac{\epsilon _o}{N}}{\displaystyle \underset{n,m=1}{\overset{N}{}}}`$ $`\psi _{nm\mathrm{}m}`$ $`|mn\mathrm{}n=0.`$ (16)
Equation (15) can explicitly be written out as $`N^D`$ equations
$`(1\epsilon _o)`$ $`\psi _{i_1i_2\mathrm{}i_D}`$ (17)
$`+`$ $`\epsilon _oN^{D1}\delta _{i_2i_3}\delta _{i_2i_4}\mathrm{}\delta _{i_2i_D}\psi _{i_2i_1\mathrm{}i_1}=0,`$ (18)
where $`\delta _{ab}=1`$ if $`a=b`$, $`0`$ otherwise. Now if one or more of the $`i_a:a=3,\mathrm{},D`$ does not equal $`i_2`$, then that equation is satisfied only if $`\epsilon =1`$ or $`\psi _{i_1i_2\mathrm{}i_D}=0`$. The first case is not of interest here, as $`\epsilon =1`$ corresponds to our maximally entangled GHZ like states, so we choose $`\psi _{i_1,i_2\mathrm{}i_D}=0`$.
The rest of the equations where $`i_2=i_3=\mathrm{}=i_D`$, separate into $`N(N1)/2`$ coupled sets of two equations of the identical form
$`(1\epsilon _o)\psi _{ab\mathrm{}b}+\epsilon _oN^{D1}\psi _{ba\mathrm{}a}`$ $`=`$ $`0,`$ (19)
$`(1\epsilon _o)\psi _{ba\mathrm{}a}+\epsilon _oN^{D1}\psi _{ab\mathrm{}b}`$ $`=`$ $`0.`$ (20)
These have solutions if $`\psi _{ab\mathrm{}b}=\psi _{ba\mathrm{}a}=0`$, but this would imply $`|\varphi =0`$, which was specifically excluded in (14). Otherwise, these coupled two equations are only satisfied if
$`\epsilon _o={\displaystyle \frac{1}{N^{D1}+1}}.`$ (21)
So for $`|\psi 0`$, at least one such coupled set of two equations leads to the expression (21). This is the only candidate for the point where the Peres condition becomes satisfied.
Now it can be easily seen that in the total-noise case $`\epsilon =0`$, $`\rho _1^T=\widehat{\rho }`$ and all the eigenvalues of $`\rho _1^T`$ are $`1/(N^D)`$. In the no-noise case ($`\epsilon =1`$), proceeding in similar fashion to before, the eigenvalues $`\lambda `$ of $`\rho _1^T`$ must satisfy
$`{\displaystyle \frac{1}{N}}{\displaystyle \underset{n,m=1}{\overset{N}{}}}`$ $`\psi _{nm\mathrm{}m}`$ $`|mn\mathrm{}n`$ (22)
$`=`$ $`\lambda `$ $`{\displaystyle \underset{i_1,i_2,\mathrm{},i_D=1}{\overset{N}{}}}\psi _{i_1i_2\mathrm{}i_D}|i_1i_2\mathrm{}i_D.`$ (23)
This gives $`\psi _{i_1i_2\mathrm{}i_D}=0`$ if for some $`a=3,4,\mathrm{},D`$, $`i_2i_a`$, and leads to sets of two coupled equations of the form
$`\psi _{ab\mathrm{}b}`$ $`=`$ $`\lambda N\psi _{ba\mathrm{}a},`$ (24)
$`\psi _{ba\mathrm{}a}`$ $`=`$ $`\lambda N\psi _{ab\mathrm{}b}.`$ (25)
These have the solutions $`\psi _{ab\mathrm{}b}=\psi _{ba\mathrm{}a}=0`$ or
$`\lambda =\pm {\displaystyle \frac{1}{N^2}}.`$ (26)
As before, for nonzero eigenvectors, the first cannot be true. So, finally, for the no-noise state, at least one eigenvalue of $`\rho _1^T`$ must be negative, and equal to
$`\lambda ={\displaystyle \frac{1}{N^2}}.`$ (27)
Thus, finally, since $`\epsilon _o`$ is the only value of $`\epsilon `$ where an eigenvalue of $`\rho _1^T`$ is zero, all eigenvalues are positive for $`\epsilon =0`$, and an eigenvalue is negative for $`\epsilon =1`$, there must be at least one negative eigenvalue for $`\epsilon >\epsilon _o`$ given by (21).
When one proceeds in the same fashion when $`M=2,3,\mathrm{},D/2`$, one always gets coupled sets of two equations identical in form to (19), as can be seen by inspection, so nothing new is found. Thus the result of section (II) is indicated.
## V Comparison to known bounds
For the qubit case ($`N=2`$), as seen in figure (1), equation (9) gives an exact bound on the values of $`\epsilon `$ that divide separable from entangled states of the form (II). One sees that the upper bound (8) derived from the work of Vidal and Tarrach comes very close to the exact value for the qubit case. This exact bound is also greater than those lower bounds previously found by Schack and Caves, as expected.
For the two-subsystem ($`D=2`$) case, the bound (9) agrees with the exact one found by the Horodeckis.
For other values of $`D`$ and $`N`$, the results on random robustness by Vidal and Tarrach, lead to an upper bound which is considerably weaker than the upper bound (9) given by the partial transposition condition. The upper bound found here, actually gives a stronger bound on the random robustness of entanglement of states given by (II).
$$R(\widehat{\rho }_e||\widehat{\rho }_n)N^{D1}$$
(28)
It is interesting to note that in the border cases when $`N=2`$ or $`D=2`$, the bound (9) is in fact an exact bound. This is despite the fact that the Peres condition does not necessarily give a strong bound for such states. This leads one to the tentative conjecture that for noisy GHZ-type states of the form (II), the Peres condition may in general give a strong upper bound.
Looking at figure (2), one sees that as the Hilbert space dimension of the subsystems increases, the upper bound on $`\epsilon `$ rapidly decreases, indicating that the entanglement becomes stronger.
As the bound (9) is completely general in $`N`$ and $`D`$, it does provide some answers to the question of what raises entanglement more: creating more entangled subsystems, or increasing their dimension? Since the bound on $`\epsilon `$ decreases exponentially with $`D`$, but only polynomially with $`N`$, one concludes that increasing the number of subsystems is a more effective way of increasing the entanglement.
## VI Conclusion
The results presented here based on the Peres condition provide a lower bound on the parameters $`\epsilon `$ in (II) above which the generalised Werner states are always entangled. Furthermore, it gives an exact bound on this parameter for the case of many qubits. An explicit simple expression is derived that depends on $`D`$, the number of subsystems over which entanglement occurs, and $`N`$ the Hilbert space dimension of each subsystem. Apart from the few cases ($`D=2`$;$`N=2`$ and $`D=3`$) where this bound was known exactly previously, the new bounds are stronger than previously known ones.
This work also sheds light on the question of whether to increase quantum entanglement in a system, is it better to create more entangled subsystems, or to increase the dimension of the existing subsystem. As the bound on $`\epsilon `$ decreases exponentially with $`N`$, but only polynomially with $`D`$, increasing the number of subsystems is a much more effective way of increasing entanglement.
To conclude, the Peres partial transposition condition has provided a good upper bound for determining the separability of a generalised $`N`$, $`D`$ Werner state. However for other systems it is known that this transposition condition fails to give strong results, thus when this condition is useful, and when not, remains an interesting question.
###### Acknowledgements.
We are grateful to Gerard Milburn for discussion about entanglement and separability. W.J.M acknowledges the support of the Australian Research Council.
## A Proof that (6) is exact for qubits
We wish to show that the bound (9) is exact for qubits ($`N=2`$), and to find the product states which combine to give the state when it is separable.
Starting off similarly to Schack and Caves, the $`D`$-subsystem generalisation of the Werner state ((II) with $`N=2`$) can be written in terms of Pauli matrices:
$$\widehat{\rho }(\epsilon )=\frac{1}{2^D}\left\{(1\epsilon )I^D+\frac{\epsilon }{2}\widehat{E}\right\},$$
(A2)
where
$`\widehat{E}=`$ $`(I+\sigma _3)^D+(I\sigma _3)^D`$ (A3)
$`+(\sigma _1+i\sigma _2)^D+(\sigma _1i\sigma _2)^D,`$ (A4)
and $`\sigma _i`$ are the Pauli matrices, $`I`$ is the two-dimensional identity matrix, and for conciseness, the following notation is used:
$$(A)^D=(A)\mathrm{}(A)\text{D}\text{ times}.$$
(A5)
To show that (9) is actually a strong bound, is suffices to find an expansion of $`\rho `$ in terms of a positive sum of direct tensor products of density matrices for this value of
$$\epsilon =\epsilon _c=\frac{1}{2^{D1}+1}.$$
(A6)
By analogy with the results of Schack and Caves , it has been guessed that an expansion of $`\rho (\epsilon _c)`$ in terms of the density matrices
$$P_{\pm i}=\frac{1}{2}(I\pm \sigma _i)\text{ (}\text{i = 1,2,3}\text{)},$$
(A7)
is given by the following:
$$\widehat{\rho }_g=\frac{\epsilon _c}{2}\left(P_3^D+P_3^D\right)+\frac{\epsilon _c}{2^{D1}}P_{x_1}P_{x_2}\mathrm{}P_{x_D}.$$
(A8)
Here the sum is over all permutations of $`D`$ indices $`x_1,x_2,\mathrm{},x_D`$ satisfying the following conditions:
1. $`x_i\{1,1,2,2\}`$.
2. The number of $`x_i\{2,2\}`$ is even (or zero).
3. If the number of $`x_i\{2,2\}`$ is a multiple of four (or is zero), then the number of $`x_i\{1,2\}`$ is even (or zero).
4. If the number of $`x_i\{2,2\}`$ is not a multiple of four, then the number of $`x_i\{1,2\}`$ is odd.
The proof of this proceeds by starting with the above guess, and showing that it is a correct one. As it turns out, the hard work is in actually writing down the guess mathematically, so let us begin with this.
To begin with, note that
$`𝒯_0`$ $`=`$ $`(P_1+P_2+P_1+P_2)^D`$ (A9)
$`=`$ $`{\displaystyle P_{x_1}P_{x_2}\mathrm{}P_{x_D}},`$ (A10)
where the sum is over *all* permutations of $`D`$ indices $`x_1,x_2,\mathrm{},x_N\{1,2,1,2\}`$. Also note that
$$𝒮_0=(P_1+P_2P_1P_2)^D$$
(A11)
will give a sum over all these permutations, except that all terms which have an odd number of indices in $`\{1,2\}`$ will be subtracted rather than added like in $`𝒯_0`$. Thus one can see that
$$_0^e=\frac{1}{2}[𝒯_0+𝒮_0]$$
(A13)
will give a sum over permutations of $`D`$ indices, like $`𝒯_0`$, except that only terms where the number of indices in $`\{1,2\}`$ is even will be included. Similarly
$$_0^o=\frac{1}{2}[𝒯_0𝒮_0]$$
(A14)
will give a sum over only those terms in which the number of indices in $`\{1,2\}`$ is odd.
Now consider some more similar expressions.
$$𝒯_1=(P_1P_2+P_1P_2)^D.$$
(A15)
$`𝒯_1`$ will give a sum over all index permutations, except that all terms which have an odd number of indices in $`\{2,2\}`$ will be subtracted rather than added like for $`𝒯_0`$.
$$𝒯_2=(P_1+iP_2+P_1+iP_2)^D.$$
(A16)
$`𝒯_2`$ will give a similar sum over all permutations, but terms in which the number of indices in $`\{2,2\}`$ is a multiple of four (or is zero) will be added, terms in which this number is even, but not a multiple of four, will be subtracted, terms in which this number is one more than a multiple of four will be added and multiplied by $`i`$, and terms for which this number is one less than a multiple of four will be subtracted and multiplied by $`i`$.
$$𝒯_3=(P_1iP_2+P_1iP_2)^D.$$
(A17)
$`𝒯_3`$ is the complex conjugate of $`𝒯_2`$. It can be seen (after a little thought) that $`\frac{1}{4}[𝒯_0+𝒯_1+𝒯_2+𝒯_3]`$ will give a sum over only those terms in which the number of indices in $`\{2,2\}`$ is a multiple of four (or is zero). Similarly, $`\frac{1}{4}[𝒯_0+𝒯_1𝒯_2𝒯_3]`$ will give a sum over only those terms in which the number of indices in $`\{2,2\}`$ is even, but not a multiple of four.
Analogously to equation (A11) define
$`𝒮_1=(P_1P_2P_1+P_2)^D,`$ (A19)
$`𝒮_2=(P_1+iP_2P_1iP_2)^D,`$ (A20)
$`𝒮_3=(P_1iP_2P_1+iP_2)^D,`$ (A21)
and one can define expressions for $`_i^e`$ and $`_i^o`$ for $`i=0,1,2,3`$ analogously to equations (A11). So following the same reasoning as previously, the sum of all terms in which the number of indices in $`\{2,2\}`$ is a multiple of four, and the number of indices in $`\{1,2\}`$ is even is given by
$$\frac{1}{4}[_0^e+_1^e+_2^e+_3^e].$$
(A22)
And thus finally, the sum in the second term of the guess $`\rho _g`$ ( equation (A8) ) can be written
$`{\displaystyle \frac{1}{4}}[_0^e+_1^e+_2^e+_3^e]+{\displaystyle \frac{1}{4}}[_0^o+_1^o_2^o_3^o]`$ (A23)
$`={\displaystyle \frac{1}{4}}[𝒯_0+𝒯_1+𝒮_2+𝒮_3].`$ (A24)
So the guess that has been made (i.e. equation (A8)) can be rewritten
$$\rho _g=\frac{\epsilon _c}{2}\left(P_3^N+P_3^N\right)+\frac{\epsilon _c}{2^{N+1}}\left[𝒯_0+𝒯_1+𝒮_2+𝒮_3\right],$$
(A25)
which is explicitly (via the expressions (A9), (A11), (A15), (A)) a positive sum of direct tensor products of density matrices, and thus is separable. The only question that remains is whether $`\rho _g=\rho (\epsilon _c)`$ ?
This is the easy part. It is seen using the expression (A7) that
$`𝒯_0`$ $`=`$ $`(2I)^N=2^NI^N,`$ (A27)
$`𝒯_1`$ $`=`$ $`0,`$ (A28)
$`𝒮_2`$ $`=`$ $`(\sigma _1+i\sigma _2)^N,`$ (A29)
$`𝒮_3`$ $`=`$ $`(\sigma _1i\sigma _2)^N,`$ (A30)
$`P_3^N`$ $`=`$ $`2^N(I+\sigma _3)^N,`$ (A31)
$`P_3^N`$ $`=`$ $`2^N(I\sigma _3)^N,`$ (A32)
so
$$\rho _g=\frac{\epsilon _c}{2}I^N+\frac{\epsilon _c}{2^{N+1}}\widehat{E},$$
(A33)
which only seemingly differs from equation (A) by the first term, but using the expression for $`\epsilon _c`$ (A6), one finds that these first terms are equal also.
$$\frac{1\epsilon _c}{2^N}=\frac{\epsilon _c}{2},$$
(A34)
so $`\rho _g=\rho `$, the guess was correct, and thus the bound $`\epsilon _c`$ is strong.
$`𝒬𝒟`$ |
warning/0002/quant-ph0002005.html | ar5iv | text | # Discrete Phase Measurements and the Bell inequality
## I Introduction
The Bell inequalities and their tests of quantum mechanics have generally been seen only as a fundamental test of quantum mechanics until recently. Now however their relevance has become much more important with the proposal for quantum computing and quantum encryption. Fundamental to these schemes is quantum entanglement which is at the core of quantum mechanics. The Bell inequalities provide a mechanism now by which strongly entangled two mode systems can be distinguished from all classical systems. Current tests of the Bell inequalities have suffered from a number of loopholes , including the fair sampling assumption and low detection efficiencies. To date, no definitive experiment test has been achieved although specific experiments have been performed which close one or more of the loopholes. This has lead to proposals from a number of authors including ourselves to consider high efficiency detection. Work by Smithery et. al. has indicated how high efficiency photon detectors can be used to detect large photon numbers but with an uncertainty in the actual photon number. Reid et. al. indicated how high efficiency photon detectors, with large particle numbers, can test macroscopic quantum mechanics.
In a recent paper by Gilchrist et. al. they showed how the predictions of quantum mechanics are in disagreement with those of local hidden variable theories for a situation involving continuous quadrature phase amplitude ( position and momentum ) measurements. More explicitly they showed that the quantum predictions for the probability of obtaining results $`x`$ and $`p`$ for position and momentum (or various linear combinations) cannot be predicted by any local hidden variable theory. The test could be achieved by binning the continuous position and momentum information into two and using the binary results in the strong Clauser Horne Bell inequality test. There predicted violation was small (less than 2%). Munro considered various different strong Bell inequalities and indicated how larger violations may be achieved.
Also published at a similar time by us we proposed how the novel use of phase measurements on a simple correlated photon number triplet could be used to test the GHZ correlations via the Mermin $`F`$ inequality. More explicitly we showed how simple correlated photon number triplets which ideally could be produced in nondegenerate parametric oscillation, where we have signal, idler and pump modes, in conjunction with discrete binary phase measurements could be used to provide a definitive test of quantum mechanics versus all local realism. In fact a maximal violation of the mermin $`F`$ inequality was indicated. However as we mentioned in this paper discrete phase measurements in optical systems have yet to be experimentally realized in the ultra high detector efficiency limit. As an approximation to this binary phase measurement a homodyne quadrature phase amplitude measurements was considered and also found to violate the inequality. The potential violation is significantly reduced however because information in the quadrature phase amplitude measurements must be discarded to achieve the binary result.
In this paper we will be restricting our attention to systems involving correlated photon number states. To obtain the maximum entanglement information we are proposing the use of canonical phase measurements. Initially we consider a discrete binary phase measurement but then generalise to phase measurements involving larger (greater than two) numbers of phase results. To achieve a Bell inequality test in such cases requires binning of the results from the phase measurements. A critical part of this paper is to determine the effect of this binning process. Until recently only two limits have been considered, binary discrete phase measurements and continuous quadrature phase amplitude measurements. We analyse a number of cases between these extreme limits. While our phase measurements may not be easily experimentally achievable valuable insight into the binning of data for the Bell inequalities is obtained. While we will not explicitly consider correlated spin systems, the results indicated by the binning of phases could be applied to these spin systems. For instance with binning a correlated $`3/2`$ system could be used to test the Bell inequality. Historically only correlated $`1/2`$ systems have been used. This opening the possibility for new novel tests of quantum mechanics with higher spin particles.
## II Entanglement
It is important to begin by explaining the reason for considering phase measurements especially when we are restricting our attention to correlated photon number state systems. As has long been known quantum entangled states shows correlations that cannot be explained in terms of the correlations between local classical properties of the subsystems. In this paper we are describing a pure entangled state of two modes in which the correlations are in photon number, that is the nature of correlation can be succinctly stated by saying that there are equal number of photons in each mode. Now as there are many different ways to realize this, the total state is the sum over amplitudes for all possible ways in which this correlation can be realized. This kind of sum over amplitudes for correlations is characteristic of an entangled state.
How do we best see the quantum nature of the correlation? Obviously it is not enough to measure photon number as this would not distinguish a mixed state, with equal photon numbers in each mode, from the equivalent entangled pure state. In some sense we need to measure an observable which carries as little information as possible about photon number in order to see the interference between all the possible ways in which the correlation in photon number can be realized. We conjecture that the best choice is the observable canonically conjugate to photon number; the phase. Before we consider the quantum mechanical phase (Section IV) we will in the next section specify more precisely the nature of the entangled states we are considering and realisable systems that generate them.
## III Entangled Photon Pair States
Tests of quantum mechanics generally require an entanglement between particles. A test of the Bell inequality requires an entanglement between two subsystems. As mentioned previously we are examining a two mode state where there are an equal number of photons in each mode. This correlated photon number pair state is given by
$`|\mathrm{\Psi }={\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}c_n|n|n`$ (1)
where at present we leave $`c_n`$ arbitrary but specify that due to normalisation
$`{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}\left|c_n\right|^2=1.`$ (2)
For convenience in the calculations that follow, we will place the condition on $`c_n`$ that it must be real. Eqn (1) actually describes a number of idealised but real systems. The most well known and extensively examined case is the nondegenerate parametric amplifier specified by an ideal Hamiltonian of the form
$`H=\mathrm{}\chi ϵ\left(ab+a^{}b^{}\right).`$ (3)
where $`ϵ`$ is field amplitude of a nondepleting classical pump and $`\chi `$ is proportional to the susceptibility of the medium. $`a,b`$ are the boson operators for the two spatially separated orthogonal signal and idler modes systems. After a time $`\tau `$, the state of the system is given by (1) with $`c_n`$ specified by
$`c_n={\displaystyle \frac{\mathrm{tanh}^n\left[\chi ϵ\tau \right]}{\mathrm{cosh}\left[\chi ϵ\tau \right]}}`$ (4)
For small $`\chi ϵ\tau `$ in Eqn (4) a state of the form
$`|\mathrm{\Psi }=c_0|0|0+c_1|1|1`$ (5)
can be generated as a reasonably good approximation.
Another source of highly correlated photon number states exists in the steady state by nondegenerate parametric oscillation as modeled by the following Hamiltonian,
$`H=i\mathrm{}ϵ\left(aba^{}b^{}\right)+ab\mathrm{\Gamma }^{}+a^{}b^{}\mathrm{\Gamma }.`$ (6)
Here we have neglected the effect of linear single photon loss. We assume that the coupled signal-idler loss dominates over linear single-photon loss. In Eqn. (6) $`ϵ`$ represents a coherent driving source which generates signal-idler pairs, while $`\mathrm{\Gamma }`$ represents the reservoir systems which gives rise to the coupled signal-idler loss. $`a,b`$ again are the boson operators for the orthogonal signal and idler modes. In the limit of very large parametric nonlinearity and high Q cavities, a state of the form
$$|\text{circle}=𝒩_0^{2\pi }|re^{i\varsigma }_a|re^{i\varsigma }_b𝑑\varsigma $$
(7)
can be generated. Here $`𝒩`$ is the normalisation coefficient while $`|\mathrm{}_a`$ and $`|\mathrm{}_b`$ are coherent states of amplitude $`r`$ in the spatially separated modes $`\widehat{a}`$ and $`\widehat{b}`$. This state can be written in the form of eqn (1) with the $`c_n`$’s now specified by
$`c_n={\displaystyle \frac{r^{2n}}{n!I_0\left(2r^2\right)}}`$ (8)
where $`I_0`$ is the zeroth order modified Bessel function.
Given that we now have exactly specified the nature of the our correlation we return to consider the observable canonically conjugate to photon number; the phase.
## IV Phase states
Above we have proposed that that phase states may be the best way to observe entanglement in photon number. There have been a number of phase states proposed over recent years. A canonical phase state $`|\theta `$ can be generally be defined as
$`|\theta ={\displaystyle \frac{1}{\sqrt{2\pi }}}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}\mathrm{exp}\left[in\theta \right]|n`$ (9)
where $`|n`$ represent the fock states and $`\theta `$ specifies the phase. These phase states are unnormalized as $`\theta |\theta >1`$. In fact they are unnormalisable. A normalizable phase state was proposed by Pegg and Barnett who considered the phase state definition
$`|\theta =\underset{s\mathrm{}}{lim}{\displaystyle \frac{1}{\sqrt{s+1}}}{\displaystyle \underset{n=0}{\overset{s}{}}}\mathrm{exp}\left[in\theta \right]|n.`$ (10)
Here $`|n`$ are the $`s+1`$ number states that span the $`(s+1)`$ dimensional state space. The limit in Eqn (10) is absolutely necessary in order to normalize the states. The parameter $`\theta `$ can take on any real value, although distinct states will occur in a $`2\pi `$ range.
Pegg and Barnett showed that a set of $`s+1`$ orthonormal phase states, with values of $`\theta `$ differing by $`2\pi /(s+1)`$, can be generated from
$`|\theta _\mu =\mathrm{exp}\left[i\widehat{N}\mu 2\pi /(s+1)\right]|\theta _0,\mu =0,\mathrm{},s`$ (11)
where $`|\theta _0`$ is the reference (or zero) phase state. Here $`\widehat{N}`$ is the number operator and $`\mu `$ the particular discrete phase we are interested in. $`\mu `$ generally varies in integer steps from $`0`$ to $`s`$. Our values for $`\theta _\mu `$ are given by
$`\theta _\mu =\theta _0+{\displaystyle \frac{2\mu \pi }{(s+1)}}`$ (12)
which are spread evenly over the range $`\theta _0\theta _\mu \theta _0+2\pi `$, where $`\theta _0`$ is the initial (or reference) phase.
## V Probability distributions
Now given the form (1) and the definition of discrete phase states in section (IV) we now proceed to calculate the joint probability of finding our correlated photon number system with phase $`\theta _{\mu _1}`$ for in the first subsystem/mode and $`\theta _{\mu _2}`$ for the second subsystem/mode. We specify that there is an initial phase angle $`\theta _{0,i}`$ for each of the subsystems $`i`$ (i=1,2). The joint probability $`P_{\mu _1,\mu _2}(\theta _{0,1},\theta _{0,2})`$ is given by
$`P_{\mu _1,\mu _2}(\theta _{0,1},\theta _{0,2})`$ $`=`$ $`{\displaystyle \underset{n=0}{\overset{s}{}}}{\displaystyle \frac{\left|c_n\right|^2}{\left(s+1\right)^2}}+{\displaystyle \underset{n>n^{}=0}{\overset{s}{}}}{\displaystyle \frac{2c_nc_n^{}}{\left(s+1\right)^2}}`$ (13)
$`\times `$ $`\mathrm{cos}\left[\mathrm{\Delta }n\left({\displaystyle \frac{2\pi \left[\mu _1+\mu _2\right]}{s+1}}+\psi _0\right)\right]`$ (14)
Here $`\psi _0`$ is the sum of the two initial phase angle ($`\psi _0=\theta _{0,1}+\theta _{0,2}`$) and $`\mathrm{\Delta }n=nn^{}`$. Here each $`\mu _i`$ can vary from $`0`$ to $`s`$ in integer steps.
It is useful to calculate the marginal probability $`P_{\mu _i}\left(\theta _{0,i}\right)`$ (where $`i`$ equals either 1 or 2) of finding the correlated photon number system with phase $`\theta _{\mu _i}`$ for the $`\mathrm{i}^{\mathrm{th}}`$ subsystem (i=1,2), while having no information about the remaining subsystem (that is the remaining subsystem may be in any phase state). Here the marginal probability $`P_{\mu _i}\left(\theta _{0,i}\right)`$ is given by
$`P_{\mu _i}\left(\theta _{0,i}\right)={\displaystyle \frac{1}{s+1}}{\displaystyle \underset{n=0}{\overset{s}{}}}\left|c_n\right|^2`$ (15)
It is interesting to note that $`P_{\mu _i}\left(\theta _{0,i}\right)`$ is actually independent of all angular dependence, both from $`\mu _i`$ and $`\theta _{0,i}`$, the initial or reference choice of phase for that subsystem. In fact the marginal probability is uniformly distributed over all the possible results. In the limit that $`s`$ becomes very large $`P_{\mu _i}\left(\theta _{0,i}\right)0`$.
## VI Binary Choice
We generally require large $`s`$ to describe an arbitrary phase for a general system. In fact to specify the phase as precisely as possible we require $`s\mathrm{}`$. However in the case of the measurement schemes required for testing the various quantum paradox such as the Bell inequalities, all that is required and actually necessary is a binary result . Hence it would be logical to have a discrete phase measurement with say $`s=1`$, that is only two phase states. Such a scheme has been analyzed by Munro and Milburn for the GHZ state.
It may not always be possible to have only two phase states. If more phase states are present (for example $`s=3`$), a binary result is still required for these particular quantum paradoxes (and especially the Bell inequalities we are interested in, although some of the inequalities such as the Information theoretic Bell inequality or the Mermin higher spin inequality allow for other than a binary result). This could be achieved by dividing or binning the phase states into two discrete distinct sets. We could for instance label $`P_,(\theta _{0,1},\theta _{0,2})`$ as the probability of finding both particles in a $``$ state (where the $``$ state is one of the two possible binary results, the other being $``$). $`P_,(\theta _{0,1},\theta _{0,2})`$ would correspond to the probability of finding both modes in the $``$ state. How exactly this binary choice is achieved will be discussed on a per case basis in the next few sections of this paper. We could for instance specify that an $``$ result corresponds results containing the phase results $`\mu =0\mathrm{}s/2`$. Such a process however discards information about the system and hence we would expect a lessening of our potential Bell inequality violation.
## VII The Bell inequalities
A number of Bell inequalities exist, and the particular one to be considered in this article are the Clauser Horne and the spin Bell inequality. A detailed derivation of the various inequalities will not be given, the reader is referred to references . In Fig (1) we depict a very idealized setup for general Bell inequality experiment.
To formulate the Bell inequalities it is necessary to postulate the existence of a local hidden variable theory. We can write the probability $`P_,(\theta _{0,1},\theta _{0,2})`$ for obtaining a result $``$ and $``$ respectively upon the simultaneous measurements the phase at $`A`$ and the phase at $`B`$ in terms of the hidden variables $`\lambda `$ as
$`P_,(\theta _{0,1},\theta _{0,2})={\displaystyle \rho (\lambda )p_{}^A(\theta _{1,0},\lambda )p_{}^B(\theta _{2,0},\lambda )𝑑\lambda }`$ (16)
The $`\rho (\lambda )`$ is the probability distribution for the hidden variable state denoted by $`\lambda `$. $`p_{}^A(\theta _{0,2},\lambda )`$ is the probability of obtaining a result $``$ upon measurement at $`A`$ of the phase, given the hidden variable state $`\lambda `$. Similarly $`p_{}^B(\theta _{0,2},\lambda )`$ is the probability of obtaining a result $``$ upon measurement at $`B`$ of the phase, also given the hidden variable state $`\lambda `$. Our locality assumption is due to the fact that a measurement at $`A`$ cannot be influenced by the choice of parameter $`\theta _{2,0}`$ at the location $`B`$ (and vice versa).
A number of Bell inequalities exist and in this article we will consider only two cases. The first is the strong Clauser-Horne Bell inequality that can then be written in the form of
$`|𝐁_{\mathrm{CH}}|1`$ (17)
where
$`B_{\mathrm{CH}}={\displaystyle \frac{P_{}(\theta _{1,0},\theta _{2,0})P_{}(\theta _{1,0},\theta _{2,0}^{})+P_{}(\theta _{1,0}^{},\theta _{2,0})+P_{}(\theta _{1,0}^{},\theta _{2,0}^{})}{P_{}\left(\theta _{1,0}^{}\right)+P_1\left(\theta _{2,0}\right)}}`$ (18)
Here $`P_{}`$ is the probability that a “$``$” results occurs at each analyzer $`A,B`$ given $`\theta _{1,0}`$,$`\theta _{2,0}`$. Similarly $`P_{}`$ is the probability that a “$``$” occurs at a detector while having no information about the second. For many of the actual experimental considerations an angle factorization occurs so that $`P_{}(\theta _{1,0},\theta _{2,0})`$ depends only on $`\theta _{1,0}+\theta _{2,0}`$. Also $`P_{}\left(\theta _{i,0}\right)`$ is independent of both $`\theta _{1,0},\theta _{2,0}`$. In this case $`B_{\mathrm{CH}}`$ can be simplified to
$`B_{\mathrm{CH}}={\displaystyle \frac{2P_{}\left(\psi \right)+P_{}\left(\psi \right)P_{}\left(3\psi \right)}{2P_{}}}`$ (19)
where $`\psi =\theta _{1,0}+\theta _{2,0}=\theta _{1,0}^{}\theta _{2,0}^{}=\theta _{1,0}+\theta _{2,0}^{}`$ and $`3\psi =\theta _{1,0}^{}+\theta _{2,0}`$. In some cases it can be shown that $`P_{}\left(\psi \right)=P_{}\left(\psi \right)`$ and hence this expression further simplifies to
$`B_{\mathrm{CH}}={\displaystyle \frac{3P_{}\left(\psi \right)P_{}\left(3\psi \right)}{2P_{}}}`$ (20)
This is the commonly used form for $`𝐁_{\mathrm{CH}}`$. A violation of this strong Clauser Horne Bell inequality is possible if $`|𝐁_{\mathrm{CH}}|>1`$.
The second form of the Bell inequality (sometimes referred to as the spin or original Bell inequality) is given by
$`|B_\mathrm{S}|2`$ (21)
where
$`B_\mathrm{S}=|E(\theta _{1,0},\theta _{2,0})`$ $``$ $`E(\theta _{1,0},\theta _{2,0}^{})`$ (22)
$`+`$ $`E(\theta _{1,0}^{},\theta _{2,0})+E(\theta _{1,0}^{},\theta _{2,0}^{})|2`$ (23)
Here the correlation function $`E(\theta _{1,0},\theta _{2,0})`$ is given by
$`E(\theta _{1,0},\theta _{2,0})`$ $`=`$ $`P_{}(\theta _{1,0},\theta _{2,0})+P_{}(\theta _{1,0},\theta _{2,0})`$ (24)
$``$ $`P_{}(\theta _{1,0},\theta _{2,0})P_{}(\theta _{1,0},\theta _{2,0})`$ (25)
As has been discussed above $`P_{}`$ is probability that a “$``$” results occurs at each analyzer $`A,B`$ given $`\theta _{1,0},\theta _{2,0}`$. $`P_{}`$ is probability that a “$``$” results occurs at each analyzer $`A,B`$, while $`P_{}`$ ($`P_{}`$) is probability that a “$``$” (“$``$”) results occurs at the analyzer $`A`$ and a “$``$” (“$``$”) at $`B`$. With the angle factorization given above, the inequality (21) can be rewritten as
$`B_\mathrm{S}=2E\left(\psi \right)+E\left(\psi \right)E\left(3\psi \right)`$ (26)
When $`E\left(\psi \right)=E\left(\psi \right)`$ this expression further simplifies to
$`B_\mathrm{S}=3E\left(\psi \right)E\left(3\psi \right)`$ (27)
A violation of this inequality is possible if $`|𝐁_\mathrm{S}|>2`$.
These are the two Bell inequalities that will be tested with our ideal correlated photon pairs and phase states.
## VIII Binary Phase Measurements
As has been mentioned above a logical choice for our discrete phase measurements would be to set $`s=1`$, that is a binary result. For an initial state consisting only of a pair of correlated photon number states (given by (5)) we have
$`P_{\mu _1\mu _2}\left(\psi _0\right)={\displaystyle \frac{1}{4}}+{\displaystyle \frac{1}{2}}c_0c_1\mathrm{cos}\left[\left(\mu _1+\mu _2\right)\pi \psi _0\right]`$ (28)
We now specify that the result $`\mu _i=0`$ correspond to a $``$ measurement (one of our binary choices) while $`\mu _i=1`$ correspond to a $``$ measurement. Hence the joint probability of obtaining an $``$ (or $``$) result is
$`P_{}\left(\psi _0\right)=P_{}\left(\psi _0\right)={\displaystyle \frac{1}{4}}+{\displaystyle \frac{1}{2}}c_0c_1\mathrm{cos}\left[\psi _0\right]`$ (29)
while the probability of obtaining an $``$ ($``$) result is
$`P_{}\left(\psi _0\right)=P_{}\left(\psi _0\right)={\displaystyle \frac{1}{4}}{\displaystyle \frac{1}{2}}c_0c_1\mathrm{cos}\left[\psi _0\right]`$ (30)
The marginal probability for this case is simply given by $`P_{}=1/2`$. It is then easy to calculate the expressions for $`B_{\mathrm{CH}}`$ and $`B_\mathrm{S}`$ and in fact we find
$`B_{\mathrm{CH}}`$ $`=`$ $`{\displaystyle \frac{1}{2}}+{\displaystyle \frac{1}{2}}c_0c_1\left[3\mathrm{cos}\left(\psi _0\right)\mathrm{cos}\left(3\psi _0\right)\right]`$ (31)
$`B_\mathrm{S}`$ $`=`$ $`2c_0c_1\left[3\mathrm{cos}\left(\psi _0\right)\mathrm{cos}\left(3\psi _0\right)\right]`$ (32)
Setting $`\psi _0=\pi /4`$ to maximise the values of $`B_{\mathrm{CH}}`$ and $`B_\mathrm{S}`$ we have
$`B_{\mathrm{CH}}`$ $`=`$ $`{\displaystyle \frac{1}{2}}+\sqrt{2}c_0c_1`$ (33)
$`B_\mathrm{S}`$ $`=`$ $`4\sqrt{2}c_0c_1`$ (34)
Our initial condition on $`c_0,c_1`$ for normalization requires that $`c_0^2+c_1^2=1`$, and hence the maximum value of the product $`c_0c_1`$ is one half. Hence
$`B_{\mathrm{CH},\mathrm{max}}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(1+\sqrt{2}\right)`$ (35)
$`B_{\mathrm{S},\mathrm{max}}`$ $`=`$ $`2\sqrt{2}`$ (36)
which is leads to a violation of both inequalities. In fact it is the same size of violation as is obtained when single photon detector schemes are considered.
Above we have considered an initial state consisting only of a linear combination of two correlated photon states. An ideal parametric amplifier given by (1) with the $`c_n`$ coefficient specified by (4) actually has a infinity (or at least very large) number of correlated photon number pair states. It is nearly impossible via parametric amplification to achieve the simplified state considered previously with our choice of $`c_0,c_1`$. The paramp only produces the state given by (5) (with $`c_0c_1`$) in the weakly pumped case. As the pump power increases higher order terms become significant. Generally in considering the typical photon detection Bell inequality schemes, the effect of the higher order terms of the form $`|n|n`$ (with $`n>1`$ has been to dramatically decrease (and eventually destroy) the violation. Hence what is the effect of the higher order pairs in our discrete phase measurement scheme? How rapidly will the higher order terms destroy our violation. In fact a quite surprising result occurs. It can be easily shown that the Bell inequality, with an arbitrary number of higher order pair correlated number state terms included, is
$`B_{\mathrm{CH}}`$ $`=`$ $`{\displaystyle \frac{1}{2}}+{\displaystyle \frac{1}{2}}{\displaystyle \frac{c_0c_1}{c_0^2+c_1^2}}\left[3\mathrm{cos}\left(\psi _0\right)\mathrm{cos}\left(3\psi _0\right)\right]`$ (37)
$`B_\mathrm{S}`$ $`=`$ $`2c_0c_1\left[3\mathrm{cos}\left(\psi _0\right)\mathrm{cos}\left(3\psi _0\right)\right]`$ (38)
Here we notice that the expression for $`B_\mathrm{S}`$ is unchanged for what we had indicated above and hence we will not analyse it further in this case. Focussing our attention on the expression for $`B_{\mathrm{CH}}`$ and optimising for angle we find
$`B_{\mathrm{CH}}`$ $`=`$ $`{\displaystyle \frac{1}{2}}+\sqrt{2}{\displaystyle \frac{c_0c_1}{c_0^2+c_1^2}}`$ (39)
The maximum is as previously provided $`c_0=c_1`$. We do not however has such a signient condition on $`c_0`$ which was needed when we have only two correlated pair states. Thus the addition of more states actually leads to a lessening of the condition for a violation. With the coefficient $`c_n`$ given for the parametric amplifier given by (4), large values of $`\chi ϵ\tau `$ give $`c_0c_1`$ and hence a maximal violation of the Bell inequality. The reason for this effect can be easily understood when it is noted that a binary phase measurement only involves the number states $`|0`$ and $`|1`$.
## IX Tests with larger number of phase states
In our approach consider so far, we have examined a discrete (binary, $`s=1`$) phase measurement. Let us now consider the case where more phase states are present ($`s>1`$), an inefficient binary phase measurement. We have termed this an inefficient binary phase measurement because to formulate the Bell inequality it is generally necessary to bin the phase results into a binary category. Such a process discards information and hence we would not expect it to be as ideal as the binary phase measurement.
There are a number of ways that we can divide or segment our phase space measurements into two categories. Probably the simplest is to divide the phase space into two equal parts. Our phase label $`\mu `$ has values from $`0`$ to $`s`$. We could define the $``$ bin to be contain phase results from $`\mu =0`$ to $`\mathrm{Int}[s/2]`$<sup>*</sup><sup>*</sup>*By $`\mathrm{Int}[s/2]`$ we mean the integer value of s/2, eg. $`\mathrm{Int}[3/2]=1`$ for s odd ($`s/21`$ for s even) and the $``$ bin to contain phase results from $`\mathrm{Int}[s/2]+1`$ for s odd ( $`s/2+1`$ for s even) to $`s`$. For example if $`s=3`$ we say the $``$ bin corresponds to the sum of the $`\mu =0,1`$ phase results while the $``$ bin corresponds to the sum of $`\mu =2,3`$ results. If $`s`$ is even, then we have an odd number of phase states and we must decide which bin to put the extra phase result into. In the cases we are considering this extra phase state for $`s`$ even is assigned to the $``$ bin. For example with $`s=2`$, we have the following values possible values for $`\mu `$ (0,1,2). We say that the $``$ bin corresponds to the sum of the phase results from $`\mu =0`$, while the $``$ bin contains phase results for $`\mu =1,2`$. We note here that the extra phase state $`\mu =1`$ has been put in the $``$ bin. This choice is arbitrary.
Another logical choice of binning could be to put only a single phase state into the $``$ bin, say $`\mu =0`$ while assigning all other phase results to the $``$ bin. For example if we consider $`s=3`$, we say the $``$ state corresponds to the $`\mu =0`$ phase results while the $``$ state corresponds to the sum of $`\mu =1,2,3`$ results. There are also a number of other choices for how this binning process may be achieved but we now focus our attention initially on the equal binning scheme.
### A The equal binning scheme
Let us consider the equal binning schemes with $`s`$ odd. As we have discussed above, for the equal binning scheme we define the $``$ bin to be contain phase results from $`\mu =0`$ to $`\mathrm{Int}[s/2]`$ and the $``$ bin to contain phase results from $`\mathrm{Int}[s/2]+1`$ to $`s`$. For the Clauser Horne Bell inequality we need to determine two quantities, the probability $`P_{}`$ of obtaining an $``$ result and $`P_{}`$, the probability of knowing that an $``$ result occurred for one of the particles while having zero information about the seconds result. The probability $`P_{}`$ for obtaining an $``$ result is simply
$`P_{}(\psi _0)`$ $`=`$ $`{\displaystyle \underset{\mu _1,\mu _2=0}{\overset{\mathrm{Int}[s/2]}{}}}P_{\mu _1\mu _2}(\psi _0)`$ (40)
where $`P_{\mu _1,\mu _2}(\psi _0)`$ is given by Eqn (13). Similarly the marginal probability $`P_{}`$ is given by
$`P_{}={\displaystyle \underset{\mu _1=0}{\overset{\mathrm{Int}[s/2]}{}}}{\displaystyle \underset{\mu _2=0}{\overset{s}{}}}P_{\mu _1,\mu _2}(\psi _0)={\displaystyle \frac{1}{2}}{\displaystyle \underset{n=0}{\overset{s}{}}}\left|c_n\right|^2`$ (41)
For this particular case we have specified that the first particle is in the $``$ bin (which is why the sum over $`\mu _1`$ ranges from $`0`$ to $`\mathrm{Int}[s/2]`$) while having zero information about which bin the second particle is in (this is why the sum over $`\mu _2`$ ranges from $`0`$ to $`s`$). Substituting these probability expressions into the expression for $`B_{\mathrm{CH}}`$ given by (19) we have
$`B_{\mathrm{CH}}={\displaystyle \frac{1}{2}}`$ $`+`$ $`{\displaystyle \frac{1}{s+1}}{\displaystyle \underset{\mu _1,\mu _2=0}{\overset{\mathrm{Int}[s/2]}{}}}{\displaystyle \underset{n>n^{}=0}{\overset{s}{}}}{\displaystyle \frac{c_nc_n^{}}{\sqrt{_{m=0}^s\left|c_m\right|^2}}}`$ (42)
$`\times `$ $`\{2\mathrm{cos}\left[\mathrm{\Delta }n({\displaystyle \frac{2\pi (\mu _1+\mu _2)}{s+1}}\psi _0)\right]`$ (43)
$`+`$ $`\mathrm{cos}\left[\mathrm{\Delta }n\left({\displaystyle \frac{2\pi (\mu _1+\mu _2)}{s+1}}+\psi _0\right)\right]`$ (44)
$``$ $`\mathrm{cos}\left[\mathrm{\Delta }n({\displaystyle \frac{2\pi (\mu _1+\mu _2)}{s+1}}3\psi _0)\right]\}`$ (45)
For an equal superposition of correlated number states, that is $`c_n`$ is given by
$`c_n=1/\sqrt{s+1},`$ (46)
we are now able to calculate $`B_{\mathrm{CH}}`$. The angles $`\psi _0`$ is chosen so as to optimise $`B_{\mathrm{CH}}`$. We need however to be careful with what values $`\psi _0`$ can range over. In fact $`\psi _0`$ is restricted to the range $`[0,2\pi /3]`$ as our expression for $`B_{\mathrm{CH}}`$ also involves $`3\psi _0`$. The angle must be chosen uniquely.
It is found for the equal binning case that only $`s=1`$ violates the inequality which is the result we obtained previously. No violation is possible for higher $`s`$. Hence we consider an alternate binning scheme in the next subsection.
### B A single phase result in the $``$ bin
An alternate choice for dividing the discrete phase space in two would be to choose a single phase state result to be in our $``$ bin. For simplicity we choose the $`\mu =0`$ phase state to be in our $``$ bin. In this case it is easy to show
$`P_,\left(\psi _0\right)`$ $`=`$ $`{\displaystyle \underset{n=0}{\overset{s}{}}}{\displaystyle \frac{\left|c_n\right|^2}{\left(s+1\right)^2}}+{\displaystyle \frac{2}{\left(s+1\right)^2}}{\displaystyle \underset{n>n^{}=0}{\overset{s}{}}}c_nc_n^{}\mathrm{cos}\left[\mathrm{\Delta }n\psi _0\right]`$ (47)
and
$`P_{}={\displaystyle \frac{1}{s+1}}{\displaystyle \underset{n=0}{\overset{s}{}}}\left|c_n\right|^2.`$ (48)
The expression for $`B_{\mathrm{CH}}`$ is then given by
$`B_{\mathrm{CH}}={\displaystyle \frac{1}{s+1}}+{\displaystyle \frac{2}{s+1}}{\displaystyle \underset{n>n^{}=0}{\overset{s}{}}}{\displaystyle \frac{c_nc_n^{}}{\sqrt{_{m=0}^s\left|c_m\right|^2}}}\left\{3\mathrm{cos}\left[\mathrm{\Delta }n\psi _0\right]\mathrm{cos}\left[3\mathrm{\Delta }n\psi _0\right]\right\}`$ (49)
Results for various odd $`s`$ are plotted in Fig (2). A violation of the Bell inequality is possible for $`|B_{\mathrm{CH}}|>1`$. We notice that as $`s`$ increases the violation does decrease but seems to be preserved for quite large $`s`$ (at least to $`s=201`$).
For the results displayed on Fig (2) the $`c_n`$’s are given by (46). The angle $`\psi _0`$ bounded in the range $`[0,2\pi /3]`$ was chosen to maximise the value of $`B_{\mathrm{CH}}`$. It is interesting to note in this specific binning case that as $`s`$ increases, the width of the $``$ binning window decreases and hence the probability of finding it in that bin decreases. This is seen in Eqn (48). Here we observe that the $`P_{}`$, the probability of detecting one subsystem in the $``$ bin while having zero information about the second, decreases as $`P_{}=\frac{1}{s+1}`$. For large $`s`$ this probability becomes very small.
The state considered above basically being an equal weighed superposition of the correlated pairs from $`|0|0\mathrm{}|s|s`$ is extremely ideal. Let us briefly consider the effect to the potential violation if the $`c_n`$ coefficients are those for the ideal parametric amplifier. The $`c_n`$ coefficients were specified by eqn (4). For simplicity in the figure below we define a parameter $`\lambda =\mathrm{cosh}\left[\chi ϵ\tau \right]`$ that varies from $`0`$ to $`1`$. $`\lambda =0`$ corresponds to only the state $`|0|0`$ being present while $`\lambda =1`$ corresponds to an equally weighed superposition. In Fig (3), we plot $`B_{\mathrm{CH}}`$ versus $`\psi _0`$ and $`\lambda `$ for two values of $`s`$, $`s=3`$ and $`s=7`$.
The violation of the Bell inequality is seen in the Fig (3) above as an Island. As $`s`$ increases we notice that the size of this island goes smaller. For even larger $`s`$ a violation is possible but the overall size of the Island of violation decreases significantly. $`\lambda 1`$ is the optimal choice of this parameter to maximise the violation but the angular dependence of $`\psi _0`$ does change with $`s`$.
Other binning choices are also available. We will not discuss these but some do lead to violations of the Bell inequality as seen above.
The last issue to be discussed in this section involves entanglement. By binning of the data we gain some insight on entanglement in two subsystem where each subsystem now has a binary state. What our results show is that for some choices of binning there is enough entanglement left in the total system to violate the strong Bell inequality. A question we pose but leave unanswered is whether a binned system could violate a Bell inequality but the original system not violate the corresponding multistate test of quantum mechanics.
## X The spin analogy
In this paper we have considered correlated photon number pairs and discrete phase measurements. It is also possible to consider correlated spin systems. For instance we could write a correlated $`s=1`$ system as
$`|\mathrm{\Psi }={\displaystyle \frac{1}{\sqrt{3}}}\left[|1|1+|0|0+|1|1\right]`$ (50)
Using similar binning schemes of spin measurements, it is possible to observe similar results to the phase measurements. Care does need to be taken in binning. There are optimal choices as was seen above.
## XI Conclusion
Historically the Bell inequalities have restricted themselves to two physical subsystem where each subsystem has a binary state. Recent work by a number of authors have considered test with two subsystems, but where each subsystem has a larger number of states, sometimes infinite. In some of this work quadrature phase measurements were used and the continuous results binned. In this paper we have considered correlated photon number pairs of the form $`c_n|n|n`$ with a specific type of ideal discrete phase measurement. First we showed how binary phase measurements could produce a maximal violation of the Bell inequality. Then we considered cases where more than two phase results were possible. This was the primarily purpose of the paper.
We showed how by binning the phase measurement results into two categories $``$ and $``$, a test of the Clauser Horne Bell inequality is possible (similar results do occur for the original Bell inequality). As was expected, the potential violation does decrease as the as the number of possible results from the phase measurements increases (as $`s`$ increases). This is to be expected as information in the phase measurement results must be discarded to achieve the binary nature required for the test. We examined two specific binning cases with $`s`$ phase states. The first was to split the $`s`$ phase results into two equal sets. In this case a violation of the Bell inequality was only possible for $`s=1`$. Higher $`s`$ did not violate the inequality. The second case considered was where we took a single phase result ($`\mu =0`$) into one bin, with all the remaining possible results in the other bin. Such a binning scheme achieved a good potential violation of the inequality for quite high $`s`$ for a very idealised correlated state (in fact $`\frac{1}{s+1}_{n=0}^s|n|n`$). The violation however did decrease as $`s`$ increases. When more realistic but still idealised states were considered (those from parametric down conversion), the effect of binning could truly be seen. A large violation was possible, but the parameter region over which a violation occurred decreased rapidly as $`s`$ increased. However for moderate values of $`s`$, this region is still quite large.
Finally while our results here have been applied to correlated photon number pair systems, they can be equally applied to higher spin systems. This opens the possibility for new novel tests of quantum mechanics. Insight can also be potentially about entanglement in this binned systems. What our results show for different types of binning is whether the system is entangled enough to be about to violate a strong Bell inequality.
###### Acknowledgements.
WJM would like to acknowledge the support of the Australian Research Council. |
warning/0002/math0002219.html | ar5iv | text | # Trees and Branches in Banach Spaces
## 1. Introduction
A recurrent theme in Banach space theory takes the following form. One has some property $`(P)`$ and one assumes that in a given separable infinite dimensional Banach space $`X`$, every normalized weakly null sequence (or perhaps every normalized block basis of a given basis for $`X`$) admits a subsequence with $`(P)`$. One then tries to deduce that $`X`$ has some other property $`(Q)`$. In this paper we consider a stronger hypothesis on $`X`$. Namely that every countably infinitely branching tree of $`\omega `$-levels of some type (e.g., the successors of every node are a normalized weakly null sequence or perhaps a block basis of some FDD) admits a branch with $`(P)`$. As we show this is sometimes the proper hypothesis to conclude that $`X`$ has $`(Q)`$.
An example of this type is given in Theorem 4.1 where the following is proved: If $`X`$ is reflexive and there exists $`1<p<\mathrm{}`$ and $`C<\mathrm{}`$ so that every normalized weakly null tree in $`X`$ admits a branch $`C`$-equivalent to the unit vector basis of $`\mathrm{}_p`$ then for all $`\epsilon >0`$ there exists a finite codimensional subspace of $`X`$ which $`C^2+\epsilon `$-embeds into some space $`(F_i)_p`$, an $`\mathrm{}_p`$-sum of finite dimensional spaces. Hence this characterizes when a reflexive space embeds into such a sum.
The motivation for working with branches of trees in place of subsequences comes from the notion of asymptotic structure (\[MT\], \[MMT\]), the recent paper of N.J. Kalton \[K\] and \[KOS\]. In its simplest version suppose $`X`$ has an FDD $`(E_i)`$ and let $`k`$. Then the $`k^{th}`$-asymptotic space of $`X`$ with respect to $`(E_i)`$ may be described as the smallest closed set $`C_k`$ of normalized bases of length $`k`$ with the property that every countably infinitely branching tree of $`k`$ levels in $`S_X`$ whose nodes are all block bases of $`(E_n)`$ must admit for every $`\epsilon >0`$ a branch $`1+\epsilon `$-equivalent to some member of $`C_k`$.
Moreover given $`\epsilon _n0`$ one can then block $`(E_n)`$ into an FDD $`(F_n)`$ with the property that for all $`k`$ any normalized skipped block basis $`(x_i)_1^k`$ of $`(F_n)_{n=k}^{\mathrm{}}`$ is $`1+\epsilon _m`$-equivalent to a member of $`C_k`$ \[KOS\]. We cannot achieve this in the infinite setting, $`k=\omega `$. There is in general no unique infinite asymptotic structure, $`C_\omega `$. However if $`C`$ is big enough so that every such $`\omega `$-level tree has a branch in $`C`$ then one can produce for $`\epsilon >0`$ a blocking $`(F_n)`$ of $`(E_n)`$ so that all normalized skipped block bases of $`(F_n)`$ starting after $`F_1`$ are in $`\overline{C_\epsilon }`$, the pointwise closure (in the product topology of the discrete topology on $`S_X`$) of $`\frac{\epsilon }{2^n}`$-perturbations of elements of $`C`$. This is done in section 3. (We note that an in between ordinal notion of asymptotic structure for $`\alpha <\omega _1`$ has been considered in \[W\], using the generalized Schreier sets $`S_\alpha `$.)
Actually we need to study more general forms of asymptotic structure than that w.r.t. an FDD. We consider the version where one uses arbitrary finite codimensional subspaces rather than just the tail subspaces of a given FDD. While this version is coordinate free we show in section 3 that one may embed $`X`$ into a space with an FDD in such a way that the two notions coincide. Section 2 contains our preliminary work and terminology. In section 5 we apply our results to the more general notion of V.D. Milman’s \[Mi\] spectra of a function. We are indebted to W.B. Johnson for showing us the proof of Lemma 3.1.
## 2. Games in a Banach space $`X`$
Assume that $`X`$ is a separable Banach space of infinite dimension. The set of all subspaces of $`X`$ having finite codimension is denoted by $`\mathrm{cof}(X)`$. $`S_X^\omega `$ and $`S_X^k`$, $`k`$, denote the set of all infinite sequences in $`S_X`$, the unit sphere of $`X`$, respectively all sequences in $`S_X`$ of length $`k`$.
For a set $`𝒜S_X^\omega `$ or $`𝒜S_X^k`$ we consider the following $`𝒜`$-game between two players, having infinitely many, respectively $`k`$, rounds:
$`\text{Player I chooses }Y_1\mathrm{cof}(X)`$
$`\text{Player II chooses }y_1S_{Y_1}`$
$`\text{Player I chooses }Y_2\mathrm{cof}(X)`$
$`\text{Player II chooses }y_2S_{Y_2}`$
$`\mathrm{}`$
Player I wins if the resulting sequence $`(y_i)`$ is in $`𝒜`$.
Note that by replacing a set $`𝒜S_X^k`$, $`k`$, by $`𝒜\times S_X^\omega `$, we need only consider games with infinitely many steps.
We say that Player I has a winning strategy in the $`𝒜`$-game if the following condition, $`\text{W}_I(𝒜)`$ holds.
$`(\text{W}_I(𝒜))`$
$$\{\begin{array}{cc}& \text{There is a family of finite codimensional subspaces of }X\hfill \\ \multicolumn{2}{c}{}\\ & \left(Y_{(x_1,x_2,\mathrm{}x_{\mathrm{}})}\right)_{(x_1,x_2,\mathrm{}x_{\mathrm{}})_{j=0}^{\mathrm{}}S_X^j},S_X^0=\{\mathrm{}\},\hfill \\ \multicolumn{2}{c}{}\\ & \text{indexed over all finite sequences in }S_X\text{, so that:}\hfill \\ & \text{If }(x_n)_n\text{ satisfies the following recursive condition:}\hfill \\ \multicolumn{2}{c}{}\\ & (1)x_1S_Y_{\mathrm{}},\text{ and, for }n2,x_nS_{Y_{(x_1,\mathrm{}x_{n1})}},\hfill \\ \multicolumn{2}{c}{}\\ & \text{then }(x_n)𝒜\text{.}\hfill \end{array}$$
The following Proposition can be deduced immediately from the definition of $`(\text{W}_I(𝒜))`$.
###### Proposition 2.1.
The set of all $`AS_X^^\omega `$ for which Player I has a winning strategy is closed with respect to taking finite intersections.
Similarly, we say that Player II has a winning strategy if
$`(\text{W}_{II}(𝒜))`$
$$\{\begin{array}{cc}& \text{There is a family in }S_X\hfill \\ \multicolumn{2}{c}{}\\ & \left(x_{(Y_1,Y_2,\mathrm{}Y_{\mathrm{}})}\right)_{(Y_1,Y_2,\mathrm{}Y_{\mathrm{}})_{j=1}^{\mathrm{}}\mathrm{cof}^j(X)},\hfill \\ \multicolumn{2}{c}{}\\ & \text{indexed over all finite sequences in }\mathrm{cof}(X)\text{ (of length at least 1) so that}\hfill \\ \multicolumn{2}{c}{}\\ & (2)x_{(Y_1,Y_2,\mathrm{}Y_{\mathrm{}})}S_Y_{\mathrm{}}\text{ if }\mathrm{}\text{ and }Y_1,\mathrm{},Y_{\mathrm{}}\mathrm{cof}(X),\text{ and}\hfill \\ \multicolumn{2}{c}{}\\ & (3)\text{for every sequence }(Y_i)_i\mathrm{cof}(X),\left(x_{(Y_1,Y_2,\mathrm{}Y_i)}\right)_{i=1}^{\mathrm{}}𝒜.\hfill \end{array}$$
###### Remark.
Informally $`(\text{W}_I(𝒜))`$ means the following:
$$Y_1\mathrm{cof}(X)y_1S_{Y_1}Y_2\mathrm{cof}(X)y_2S_{Y_2}\mathrm{}\text{ so that }(y_i)𝒜.$$
Since this is an infinite phrase (unless we considered a game of finitely many draws), it has to be defined in a more formal way as it was done in $`(\text{W}_I(𝒜))`$.
It is not true in general that an $`𝒜`$-game is determined, i.e., that either Player I or Player II has a winning strategy. Note that this would mean that if the above infinite phrase is false then we can formally negate it.
From a result of D. A. Martin \[Ma\] it follows that if $`𝒜`$ is a Borel set with respect to the product topology of the discrete topology in $`S_X`$ then the $`𝒜`$-game is determined. We actually will only need a special case of this theorem which is much easier (see \[GS\] or section 1 of \[Ma\]).
###### Proposition 2.2.
For every $`𝒜S_X^\omega `$ $`(\text{W}_I(𝒜))`$ and $`(\text{W}_{II}(𝒜))`$ are mutually exclusive and if $`𝒜`$ is closed with respect to the product of the discrete topology, then it follows that the failure of $`(\text{W}_I(𝒜))`$ implies $`(\text{W}_{II}(𝒜))`$.
We furthermore note that both statements remain true if we change the game to a game in which Player I has to choose his spaces among some given subset $`\mathrm{\Gamma }\mathrm{cof}(X)`$ and/or Player II has to choose his vectors among a subset $`DS_X`$ or can choose his vector in some neighborhood of $`S_{Y_n}`$, with $`Y_n`$ being the $`n`$-th choice of Player I.
For a more detailed description of these variations of the $`𝒜`$-game we refer to Proposition 2.3, where we discuss the existence of winning strategies. In that Proposition we will show that we can reduce the game into a game in which Player I, assuming he has a winning strategy, can determine a countable collection of finite codimensional spaces before the game starts, then make his choices among this countable collection and still win the game.
We need the following notion of trees and some terminology.
###### Definition.
$`[]^{<\omega }`$ denotes the set of nonempty finite subsets of $``$ and $`[]^k`$ denotes the nonempty subsets of $``$ of cardinality at most $`k`$. These are regarded as countably branching trees of infinite length, respectively, of length $`k`$, under the order $`AB`$ if $`A`$ is an initial segment of $`B`$. A countably branching tree of infinite length in $`S_X`$ is a family $`(x_A)_{A[]^{<\omega }}`$ in $`S_X`$, where the order is that induced by $`[]^{<\omega }`$.
Similarly a countably branching tree of length $`k`$ in $`S_X`$ is a family $`(x_A)_{A[]^k}`$ in $`S_X`$.
Since these are the only kinds of trees we will consider we will simply refer to them as trees of infinite or finite length in $`S_X`$.
If $`(x_A)_{A[]^{<\omega }}`$ or $`(x_A)_{A[]^k}`$ is a tree and $`A[]^{<\omega }\{\mathrm{}\}`$, or $`A[]^{k1}\{\mathrm{}\}`$ respectively, we call the sequence $`(x_{A\{n\}})_{n>\mathrm{max}A}`$ the $`A`$-node of that tree.
If $`(n_i)`$ is an increasing sequence in $``$ of infinite length, respectively of length $`k`$, we call the sequence $`(x_{\{n_1,\mathrm{}n_i\}})_{i=1}^{\mathrm{}}`$, respectively $`(x_{\{n_1,\mathrm{}n_i\}})_{i=1}^k`$, a branch of the tree.
Assume that $`(x_A)_{A[]^{<\omega }}`$ or $`(x_A)_{A[]^k}`$ is a tree of infinite length or length $`k`$, respectively, and $`[]^{<\omega }`$, or $`[]^k`$ has the following property:
1. $``$ is hereditary, i.e., if $`A`$, and $`\mathrm{}B`$ is an initial segment of $`A`$ then $`B`$.
2. Assume that $`A\{\mathrm{}\}`$, and that $`\mathrm{card}(A)<k`$, if we consider the case of a tree of length $`k`$. Then there are infinitely many direct successors of $`A`$ in $``$, i.e., the set $`\{n:A\{n\}\}`$ is infinite.
Then we call the family $`(x_A)_A`$ a subtree of $`(x_A)`$. Note that in that case we can relabel the family $`(x_A)_A`$ as a tree $`(y_A)_{A[]^{<\omega }}`$ or $`(y_A)_{A[]^k}`$, respectively, so that every node and every branch of $`(x_A)_A`$ is node or branch, respectively, of $`(y_A)`$ and vice versa.
If $`(Y_n)`$ is a decreasing sequence of finite codimensional subspaces of $`X`$, we call a tree $`(x_A)`$ (indexed over $`[]^{<\omega }`$ or $`[]^k`$) a $`(Y_n)`$-block-tree if for every $`A[]^{<\omega }`$, respectively every $`A[]^k`$, $`x_AS_{Y_{\mathrm{max}A}}`$.
Let $`\delta _i(0,1]`$, for $`i`$, $`\delta _i0`$. We call a tree $`(x_A)_{A[]^{<\omega }}`$ of infinite length in $`S_X`$ a $`(\delta _i)`$\- approximation of a $`(Y_n)`$-block tree, if
$$\mathrm{dist}(x_A,S_{Y_{\mathrm{max}A}})<\delta _{\mathrm{card}A},\text{ whenever }A[]^{<\omega }$$
If $`𝒯`$ is a topology on $`X`$ (for example the weak topology), we call a tree $`𝒯`$-null if every node is a $`𝒯`$-null sequence.
###### Remark.
For a sequence $`(x_n)X`$ we can define a tree $`(x_A)_{A[]^{<\omega }}`$, by setting $`x_A:=x_{\mathrm{max}A}`$, for $`A[]^{<\omega }`$. Note that then the set of all subsequences of $`(x_n)`$ coincides with the set of all branches of $`(x_A)_{A[]^{<\omega }}`$.
We will be interested in conditions of the following form and relate them to the existence of winning strategies of the above discussed games.
1. Assume that all trees all of whose nodes have a certain property (A) (for example being weakly null), have a branch with a certain property (B) (for example being equivalent to the unit vector basis of $`\mathrm{}_p`$).
From the above, such a condition is a strengthening of the following assumption:
1. All normalized sequences having property (A) have a subsequence with property (B).
Continuing with our notation, if $`𝒜S_X^\omega `$ and $`\epsilon >0`$, we let
$$𝒜_\epsilon =\{(x_i)S_X:(y_i)𝒜,x_iy_i<\epsilon /2^i\text{ for all }i\}$$
and let $`\overline{𝒜_\epsilon }`$ be the closure of $`𝒜_\epsilon `$ with respect to the product of the discrete topology. We note that for $`\epsilon ,\delta >0`$
(4)
$$\overline{(\overline{𝒜_\epsilon })_\delta }\overline{𝒜_{\epsilon +\delta }}.$$
If $`Y\mathrm{cof}(X)`$ and $`\delta >0`$ then
$$(S_Y)_\delta =\{xS_X:xy<\delta \text{ for some }yS_Y\}.$$
Let $`\epsilon >0`$, $`\mathrm{\Gamma }\mathrm{cof}(X)`$ and $`DS_X`$. We define what it means to say Player I has a winning strategy for $`𝒜S_X^\omega `$ given that Player I can only choose $`Y\mathrm{\Gamma }`$ or that II can only choose elements of $`D`$.
$`(\text{W}_I(𝒜,\mathrm{\Gamma },\epsilon ))`$
$$\{\begin{array}{cc}& \text{There exists a family}\hfill \\ \multicolumn{2}{c}{}\\ & \left(Y_{(x_1,x_2,\mathrm{}x_{\mathrm{}})}\right)_{(x_1,x_2,\mathrm{}x_{\mathrm{}})_{j=0}^{\mathrm{}}S_X^j}\mathrm{\Gamma },\hfill \\ \multicolumn{2}{c}{}\\ & \text{so that for every sequence }(x_n)_n\text{ satisfying the following}\hfill \\ & \text{recursive condition:}\hfill \\ \multicolumn{2}{c}{}\\ & (5)x_1(S_Y_{\mathrm{}})_{\epsilon /2},\text{ and, for }n2,x_n(S_{Y_{(x_1,\mathrm{}x_{n1})}})_{\epsilon /2^n}\hfill \\ \multicolumn{2}{c}{}\\ & \text{one has }(x_n)𝒜\text{.}\hfill \end{array}$$
###### Remark.
It is easy to see by (4) that for any $`\epsilon ,\delta >0`$,
$$(\text{W}_I(\overline{𝒜_\epsilon },\{Y_n\},\epsilon ))(\text{W}_I(\overline{𝒜_{\epsilon +\delta }},\{\stackrel{~}{Y}_n\},\epsilon ))$$
whenever $`\{\stackrel{~}{Y}_n\}\mathrm{cof}(X)`$ is a refinement of $`\{Y_n\}`$, by which we mean that
$$Y\{Y_n\}\delta >0\stackrel{~}{Y}\{\stackrel{~}{Y}_n\}\text{ with }S_{\stackrel{~}{Y}}(S_Y)_\delta .$$
$`(\text{W}_I(𝒜,D,\epsilon ))`$
$$\{\begin{array}{cc}& \text{There is a family}\hfill \\ \multicolumn{2}{c}{}\\ & \left(Y_{(x_1,\mathrm{}x_{\mathrm{}})}^{(\epsilon )}\right)_{(x_1,\mathrm{}x_{\mathrm{}})_{j=0}^{\mathrm{}}D^j}\mathrm{cof}(X),\hfill \\ \multicolumn{2}{c}{}\\ & \text{so that for any sequence }(x_n)\text{, such that }x_nD\text{, and}\hfill \\ \multicolumn{2}{c}{}\\ & x_n(S_{Y_{(x_1,\mathrm{}x_{n1})}^{(\epsilon )}})_{\epsilon /2^n},\text{ }n=1,2,\mathrm{}\text{,}\hfill \\ \multicolumn{2}{c}{}\\ & \text{one has }(x_n)𝒜\text{.}\hfill \end{array}$$
###### Proposition 2.3.
1. If $``$ is a countable collection of subsets of $`S_X^\omega `$, then there is a decreasing sequence $`(Y_n)`$ in $`\mathrm{cof}(X)`$ so that the following are equivalent for each $`𝒜`$
1. $`\epsilon >0(\text{W}_I(\overline{𝒜_\epsilon }))`$.
2. $`\epsilon >0(\text{W}_I(\overline{𝒜_{2\epsilon }},\{Y_n\},\epsilon ))`$.
3. For every $`\epsilon >0`$ every $`(\epsilon /2^n)`$-approximation to a $`(Y_n)`$ block tree of infinite length in $`S_X`$ has a branch in $`\overline{𝒜_\epsilon }`$.
4. For every $`\epsilon >0`$ every $`(Y_n)`$ block tree of infinite length in $`S_X`$ has a branch in $`\overline{𝒜_{2\epsilon }}`$.
2. If $`X`$ has a separable dual, then $`(Y_n)\mathrm{cof}(X)`$ can be chosen so that the equivalences in 1. hold for all subsets $`𝒜S_X^\omega `$. In that case it follows that for any $`𝒜S_X^\omega `$ that (1)(a) is equivalent to
1. For every $`\epsilon >0`$ every weakly null tree of infinite length in $`S_X`$ has a branch in $`\overline{𝒜_\epsilon }.`$
###### Proof of Proposition 2.3:.
Let $`D`$ be a countable dense set in $`S_X`$. Using (4) we note that for any $`𝒜S_X^\omega `$ and any $`\epsilon >0`$ it follows that
(6)
$$(\text{W}_I(\overline{𝒜_\epsilon }))(\text{W}_I(\overline{𝒜_{2\epsilon }},D,\epsilon )).$$
Assuming now that for all $`\epsilon >0`$ the condition $`(\text{W}_I(\overline{𝒜_{2\epsilon }},D,\epsilon ))`$ is satisfied we can choose a countable subset of $`\mathrm{cof}(X)`$,
(7)
$$\mathrm{\Gamma }_𝒜=\{Y_{(x_1,\mathrm{}x_{\mathrm{}})}^{(\epsilon )}:\epsilon >0\text{ rational },x_nD\text{ and }x_n(Y_{(x_1,\mathrm{},x_{n1})}^{(\epsilon )})_{\epsilon /2^n}\text{ for }n\},$$
and observe that
(8) $`\epsilon >0(\text{W}_I(\overline{𝒜_{2\epsilon }},D,\epsilon ))`$ $`\text{there exists a countable }\mathrm{\Gamma }\mathrm{cof}(X)\text{ so that}`$
$`\epsilon >0,(\text{W}_I(\overline{𝒜_{2\epsilon }},\mathrm{\Gamma },D,\epsilon )).`$
where $`(\text{W}_I(\overline{𝒜_{2\epsilon }},\mathrm{\Gamma },D,\epsilon ))`$ is defined just like $`(\text{W}_I(\overline{𝒜_{2\epsilon }},\mathrm{\Gamma },\epsilon ))`$ with the difference that the family $`\left(Y_{(x_1,x_2,\mathrm{}x_{\mathrm{}})}\right)`$ is indexed over $`_{j=0}^{\mathrm{}}D^j`$
Using standard approximation arguments and the fact that $`D`$ is dense in $`S_X`$ we observe for any $`\mathrm{\Gamma }\mathrm{cof}(X)`$ and any $`𝒜S^\omega `$
(9)
$$(\text{W}_I(\overline{𝒜_{2\epsilon }},\mathrm{\Gamma },D,\epsilon ))(\text{W}_I(\overline{𝒜_{3\epsilon }},\mathrm{\Gamma },\epsilon ))(\text{W}_I(\overline{𝒜_{3\epsilon }})).$$
Finally assume that $`\stackrel{~}{\mathrm{\Gamma }}\mathrm{cof}(X)`$ is a refinement of $`\mathrm{\Gamma }\mathrm{cof}(X)`$. Then by (4) it follows for $`\epsilon >0`$ that
(10)
$$(\text{W}_I(\overline{𝒜_\epsilon },\mathrm{\Gamma }))(\text{W}_I(\overline{𝒜_{2\epsilon }},\stackrel{~}{\mathrm{\Gamma }})).$$
Let $``$ be any countable collection of subsets of $`S_X^\omega `$. For $`𝒜`$, if for all $`\epsilon >0`$ (W$`{}_{I}{}^{}(\overline{𝒜_\epsilon })`$) is true let $`\mathrm{\Gamma }_A`$ be as in (7), and, otherwise, we set $`\mathrm{\Gamma }_𝒜=\{X\}`$. Since $`_𝒜\mathrm{\Gamma }_𝒜`$ is countable we can choose a decreasing sequence $`(Y_n)\mathrm{cof}`$ which is a refinement of $`_𝒜\mathrm{\Gamma }_𝒜`$.
From (6)–(10) we deduce that for all $`𝒜`$
$$\epsilon >0\text{W}_I(\overline{𝒜_\epsilon })\epsilon >0\text{W}_I(\overline{𝒜_{2\epsilon }},\{Y_n\},\epsilon ).$$
Now $`\text{W}_I(\overline{𝒜_{2\epsilon }},\{Y_n\},\epsilon )`$ says that Player I has in the $`\overline{𝒜_{2\epsilon }}`$-game a winning strategy, even if he has to choose his finite codimensional subspaces among $`\{Y_n\}`$, and even if Player II “can cheat a little bit” by choosing his vectors in $`(S_{Y_n})_{\epsilon /2^n}`$. From Proposition 2.2 we deduce that this is equivalent to the condition that Player II does not have a winning strategy which means that every $`(\epsilon /2^n)`$ approximation to a $`(Y_n)`$-block-tree has a branch in $`\overline{𝒜_{2\epsilon }}`$.
We therefore have proven the equivalence of (a), (b) and (c). Note also that (c)$``$(d) is trivial and since (d) means that Player II has no winning strategy even if Player I has to choose form the set $`\{Y_n\}`$ it follows that (d) implies (a).
In order to prove the second part of the Proposition we note that in the case that $`X`$ has a separable dual we can find a universal countable refinement, i.e., a countable refinement of the whole set $`\mathrm{cof}(X)`$. Indeed, choose a dense sequence $`(\xi _n^{})`$ in $`S_X^{}`$ and let
$$Y_n=𝒩(\xi _1^{},\xi _2^{},\mathrm{},\mathrm{}\xi _n^{})=\{xX:i\{1,\mathrm{}n\}\xi _i^{}(x)=0\}.$$
Secondly note that in this case every $`(Y_n)`$-block-tree is weakly null, and, conversely, that for $`\delta _i0`$, every weakly null tree $`(x_A)_{A[]^{<\omega }}`$ has a subtree $`(y_A)_{A[]^{<\omega }}`$ which is a $`(\delta _i)`$-approximation of a $`(Y_n)`$-block-tree. ∎
## 3. A fundamental combinatorical result
For the games in $`X`$, introduced in Section 2, we want to discuss how a winning strategy of Player I or Player II can be formulated in terms of a coordinate system on $`X`$.
Recall that a Banach space $`Z`$ has an FDD $`(F_i)`$, where, for $`i`$, $`F_i`$ is a finite dimensional subspace of $`Z`$, if every $`zZ`$ can be written in a unique way as $`z=_{i=1}^{\mathrm{}}z_i`$ with $`z_iF_i`$, for all $`i`$. In this case we write $`Z=_{i=1}^{\mathrm{}}F_i`$ and denote by $`\text{c}_{00}(_{i=1}^{\mathrm{}}F_i)`$ the dense linear subspace of $`Z`$ consisting of all finite linear combinations of vectors $`x_i`$, $`x_iF_i`$. For $`mn`$ we denote by $`P_{_{i=m}^nF_i}`$ the canonical projection form $`Z`$ onto $`_{i=m}^nF_i`$.
Using a result of W. B. Johnson, H. Rosenthal and M. Zippin \[JRZ\] we derive the following Lemma.
###### Lemma 3.1.
Let $`(Y_n)`$ be a decreasing sequence of subspaces of $`X`$, each having finite codimension. Then $`X`$ is isometrically embeddable into a space $`Z`$ having an FDD $`(E_i)`$ so that (we identify $`X`$ with its isometric image in $`Z`$)
1. $`\text{c}_{00}(_{i=1}^{\mathrm{}}E_i)X`$ is dense in $`X`$.
2. For every $`n`$ the finite codimensional subspace $`X_n=_{i=n+1}^{\mathrm{}}E_iX`$ is contained in $`Y_n`$.
3. There is a $`c>0`$, so that for every $`n`$, there is a finite set $`D_nS_{_{i=1}^nE_i^{}}`$ such that whenever $`xX`$
(11)
$$x_{X/Y_n}=\underset{yY_n}{inf}xyc\underset{w^{}D_n}{\mathrm{max}}w^{}(x).$$
From (a) it follows that $`\text{c}_{00}(_{i=n+1}^{\mathrm{}}E_i)X`$ is a dense linear subspace of $`X_n`$.
Moreover if $`X`$ has a separable dual $`(E_i)`$ can be chosen to be shrinking (every normalized block sequence in $`Z`$ with respect to $`(E_i)`$ converges weakly to $`0`$, or, equivalently, $`Z^{}=_{i=1}^{\mathrm{}}E_i^{}`$), and if $`X`$ is reflexive $`Z`$ can also be chosen to be reflexive.
###### Remark.
We will prove that $`X`$ is isomorphic to a space $`\stackrel{~}{X}`$ having above properties. Then we consider on $`\stackrel{~}{X}`$ the norm, $`I()_X`$, where $`I:\stackrel{~}{X}X`$ is an isomorphism, and extend this norm to all of $`Z`$. We might loose monotonicity, or bimonotonicity, and we will not be able to assume that the constant $`c`$ in (c) can be chosen close to the value 1. But for later purposes we are more interested in an isometric embedding.
###### Proof of Lemma 3.1.
We consider the following three cases. If $`X`$ is a reflexive space we can choose according to \[Z\] a reflexive space $`Z`$ with an FDD $`(F_i)`$ which contains $`X`$. If the dual $`X^{}`$ is separable we can use again a result in \[Z\] and choose a space $`Z`$ having a shrinking FDD $`(F_i)`$. In the general case we choose $`Z`$ to be a C$`(K)`$-space containing $`X`$, $`K`$ compact and metric (for example $`K=B_X^{}`$ endowed with the $`w^{}`$-topology) and choose an FDD $`(F_i)`$ for $`Z`$.
We first write $`Y_n`$ as the null space $`𝒩(U_n)`$ of a finite dimensional space $`U_nX^{}`$ . We choose a finite set in $`S_{U_n}`$, which norms all elements of $`X/Y_n`$ up to a factor $`1/2`$ and choose for each element of this set a Hahn-Banach extension to an element in $`Z^{}`$. We denote the set of all extensions by $`D_n`$ and let $`V_n`$ be the finite dimensional subspace of $`Z^{}`$ generated by $`D_n`$. We will produce an FDD $`(E_i)`$ for $`Z`$ so that $`D_n_{i=1}^nE_i^{}`$. Hence (c) will hold.
Now
(12)
$$Y_n=𝒩(V_n)X,\text{ with }V_nZ^{},\text{ and }dim(V_n)<\mathrm{}.$$
Secondly we choose a subspace $`\stackrel{~}{W}_nX`$, $`dim(\stackrel{~}{W}_n)=dim(U_n)<\mathrm{}`$, so that $`X`$ is a complemented sum of $`Y_n`$ and $`\stackrel{~}{W}_n`$, $`X=Y_n\stackrel{~}{W}_n`$. Note that in general we do not have control over the norm of the projection onto $`Y_n`$. Given a dense countable subset $`(\xi _n)`$ in $`S_X`$, we inflate $`\stackrel{~}{W}_i`$ to $`W_i=\text{span}(\stackrel{~}{W}_i\{\xi _1,\mathrm{}\xi _i\})`$. Thus the closure of $`_{i=1}^{\mathrm{}}W_i`$ is $`X`$.
Then we choose as follows a separable subspace $`\stackrel{~}{Z}`$ of $`Z^{}`$ which is 1-complemented in $`Z^{}`$, $`Z`$-norming, and contains all the spaces $`V_n`$, $`n`$. In the case that $`X`$ has a separable dual (thus also $`Z^{}`$ is separable) we simply take $`\stackrel{~}{Z}=Z^{}`$. In the general case we let $`\stackrel{~}{Z}`$ be a separable $`L_1`$-space containing a $`Z`$-norming set, all the spaces $`V_n`$, and all the spaces $`F_n^{}`$ (considered as subspaces of $`Z^{}`$).
For $`n`$ let $`P_n:Z_{i=1}^nF_i`$ be the projection from $`Z`$ onto $`_{i=1}^nF_i`$, and let $`T_n:Z^{}\stackrel{~}{Z}`$ be the adjoint $`P_n^{}`$ if $`X^{}`$ is separable. In the general case we choose $`(T_n)`$ to be a sequence of projections of norm 1 from $`Z^{}`$ onto a finite dimensional subspace of $`\stackrel{~}{Z}`$ with the property $`T_1(Z^{})T_2(Z^{})T_3(Z^{})\mathrm{}`$ so that $`_nT_n(Z^{})`$ is dense in $`\stackrel{~}{Z}`$ (as a separable L<sub>1</sub>-space $`\stackrel{~}{Z}`$ is complemented in $`Z^{}`$ and has an FDD).
We are now in the situation of Lemma 4.2 of \[JRZ\], i.e., the following statements hold:
(13)
$$P_n^{}(Z^{})\stackrel{~}{Z}\text{ and }T_n(Z^{})\stackrel{~}{Z},$$
(14)
$$\underset{n\mathrm{}}{lim}P_n(z)=z,\underset{n\mathrm{}}{lim}T_n(y^{})=y^{}\text{ for all }zZ,y^{}\stackrel{~}{Z},\text{ and }$$
(15)
$$K:=\underset{n}{sup}T_n\underset{n}{sup}P_n<\mathrm{}.$$
We conclude from Lemma 4.2 in \[JRZ\] that:
1. Let $`E`$ and $`F`$ be finite dimensional subspaces of $`X`$ and $`\stackrel{~}{Z}`$ respectively. Then there is a projection $`Q`$ on $`Z`$ with finite dimensional range so that the following three conditions (16), (17) and (18) hold
(16)
$$Q|_E=\mathrm{Id}|_E\text{ and }Q^{}|_F=\mathrm{Id}|_F$$
(17)
$$Q^{}(Z^{})\stackrel{~}{Z}$$
(18)
$$Q4(K+K^2)$$
Using $`()`$ we can proceed as in the proof of Theorem 4.1 in \[JRZ\] to inductively define for each $`n`$ a finite dimensional projection $`(Q_n)`$ on $`Z`$ so that for all $`1i,jn`$
(19)
$$Q_iQ_j=Q_jQ_i=Q_{ij},$$
(20)
$$Q_i(X)\underset{s=1}{\overset{i}{}}W_s,$$
(21)
$$\stackrel{~}{Z}Q_i^{}(Z^{})\underset{s=1}{\overset{i}{}}V_s\text{ (in particular }D_iQ_i^{}(Z^{})\text{), and}$$
(22)
$$Q_i4(K+K^2).$$
Indeed, for $`n=1`$ we apply $`()`$ to $`E=W_1`$ and $`F=V_1`$. If $`Q_1,Q_2,\mathrm{}Q_{n1}`$ are chosen we apply $`()`$ to $`E=[Q_{n1}(Z)W_n]`$ and $`F=\text{span}(Q_{n1}^{}(Z^{})V_n)`$. We deduce (20), (21) and (22), and we observe that for $`i<n`$, $`Q_nQ_i=Q_i`$ and $`Q_n^{}Q_i^{}=Q_i^{}`$. Since for $`zZ`$ and $`z^{}Z^{}`$ the second equality implies that
$$Q_iQ_n(z),z^{}=z,Q_n^{}Q_i^{}(z^{})=z,Q_n^{}(z^{})=Q_i(z),z^{},$$
we also deduce that $`Q_iQ_n=Q_i`$.
Now we let $`E_i=(Q_iQ_{i1})(Z)`$ ($`Q_0=0`$) and deduce from (19) and (22), that $`(E_i)`$ is an FDD of a subspace of $`Z`$ which, by (20) still contains $`X`$. (20) also implies that $`\text{c}_{00}(F_i)X`$ is dense in $`X`$. Putting $`X_n=_{i=n+1}^{\mathrm{}}F_iX`$, we note that for $`xX_n`$ and $`z^{}V_n`$ it follows from (21) that $`z^{},x=Q_n^{}(z^{}),x=z^{},Q_n(x)=0`$, and thus, that $`X_n𝒩(V_n)X=Y_n`$.
We also deduce that for $`n`$, $`\text{c}_{00}(_{i=n+1}^{\mathrm{}}F_i)X`$ is dense in $`X_n`$ using the following Lemma which seems to be folklore. ∎
###### Lemma 3.2.
If $`Y`$ is a linear and dense subspace of $`X`$ and $`\stackrel{~}{X}`$ has finite codimension in $`X`$, then $`\stackrel{~}{X}Y`$ is also dense in $`\stackrel{~}{X}`$.
###### Proof.
Let $`FX`$ be a subspace of dimension $`dim(X/\stackrel{~}{X})`$, admitting a continuous projection $`Q:XF`$, so that $`(\mathrm{Id}Q)(X)=\stackrel{~}{X}`$.
Let $`x\stackrel{~}{X}`$. By assumption we find a sequence $`(y_n)Y`$ converging to $`x`$. Let $`V`$ be the (finite dimensional) vector space generated by $`(Q(y_n))_n`$ and choose a basis of $`V`$ of the form $`\{Q(y_{n_1}),\mathrm{}Q(y_n_{\mathrm{}})\}`$. We represent each vector $`Q(y_n)`$ as
$$Q(y_n)=\underset{i=1}{\overset{\mathrm{}}{}}\lambda _i^{(n)}Q(y_{n_i}),$$
and put $`x_n=y_n_{i=1}^{\mathrm{}}\lambda _i^{(n)}y_{n_i}`$. Note that $`x_nY`$ and that $`Q(x_n)=0`$, for all $`n`$. Furthermore it follows that since $`lim_n\mathrm{}Q(y_n)=0`$ and since $`(Q(y_{n_i}))_{i=1}^{\mathrm{}}`$ is basis of $`V`$, that $`lim_n\mathrm{}\lambda _i^{(n)}=0`$ for all $`1i\mathrm{}`$. Therefore it follows that $`lim_n\mathrm{}x_n=lim_n\mathrm{}y_n=x`$. ∎
We are now ready to state and to prove the main result of this section. If a Banach space $`Z`$ has an FDD $`(E_i)`$, we will call a sequence $`(z_i)`$ in $`Z`$ a block sequence with respect to $`(E_i)`$, if for some $`0=k_0<k_1<k_2\mathrm{}`$ for every $`i`$, $`z_i_{j=1+k_{i1}}^{k_i}E_j`$. We will call a tree $`(z_A)_{A[]^{<\omega }}`$ or $`(z_A)_{A[]^k}`$ in $`S_Z`$ a $`(E_i)`$-block tree if every node is a block sequence with respect to $`(E_i)`$. In a similar way given $`\delta _n0`$ we define trees which are $`(\delta _n)`$ approximations to $`(E_i)`$-block trees.
$`(G_i)`$ is a blocking of $`(E_i)`$ if there exist integers $`0=m_0<m_1<\mathrm{}`$ so that $`G_i=_{j=m_{i1}+1}^{m_i}E_j`$ for all $`i`$. $`(x_n)S_Z`$ is a skipped block w.r.t $`(G_i)`$ if
* for some sequence $`1=k_0<k_1<\mathrm{}<`$ in $``$, $`x_n_{j=k_{n1}+1}^{k_n1}G_j`$ for all $`n`$.
If $`\delta =(\delta _i)`$ with $`\delta _i0`$ and $`(x_n)S_Z`$ we say $`(x_n)`$ is a $`(\delta _i)`$-skipped block w.r.t. $`(G_i)`$ if
($`\delta `$-SB) for some sequence $`1=k_0<k_1<\mathrm{}`$ in $``$,
$$(\mathrm{Id}P_{_{j=k_{n1}+1}^{k_n1}G_j})x_n<\delta _n\text{ for all }n.$$
###### Theorem 3.3.
Let $``$ be a countable collection of subsets of $`S_X^\omega `$. Then there exists an isometric embedding of $`X`$ into a space $`Z`$ having an FDD $`(E_i)`$, so that for $`𝒜`$ the following are equivalent.
1. $`\epsilon >0(\text{W}_I(\overline{𝒜_\epsilon }))`$.
2. For every $`\epsilon >0`$ there is a blocking $`(G_i)`$ of $`(E_i)`$ and a sequence $`\delta _i0`$, so that for every sequence $`(x_n)S_X`$, satisfying $`(\delta `$-SB) w.r.t. $`(G_i)`$, $`(x_n)\overline{𝒜_\epsilon }`$.
3. For every $`\epsilon >0`$ there is a blocking $`(G_i)`$ of $`(E_i)`$, so that for every sequence $`(x_n)S_X`$ $`(SB)`$ w.r.t. $`(G_i)`$, $`(x_n)\overline{𝒜_\epsilon }`$.
If $`X`$ has a separable dual $`(E_i)`$ can be chosen to be shrinking and independent from $``$, and, furthermore, if $`X`$ is reflexive, $`Z`$ can be chosen to be reflexive. In these cases (a) is equivalent to
1. For every $`\epsilon >0`$ every weakly null tree in $`S_X`$ has a branch in $`\overline{𝒜_\epsilon }`$.
###### Remark.
Note that Theorem 3.3 means the following. Assume for all $`\epsilon >0`$ Player I has a winning strategy for the $`\overline{𝒜_\epsilon }`$-game. Then given $`\epsilon >0`$, Player I can embed $`X`$ into a space with an appropriate FDD $`(F_i)`$, and use the following strategy:
* Take $`Y_1=_{i=2}^{\mathrm{}}F_iX`$.
* If Player II has chosen the vector $`x_{n1}`$ in the $`n1`$st round,
* choose $`N`$ so that $`P_{_{i=N}^{\mathrm{}}F_i}(x_{n1})<\delta _n`$ and put
* $`Y_n=_{i=N+1}F_iX`$.
The proof of Theorem 3.3 also gives the following. Suppose $`XZ`$ where $`Z`$ has an FDD $`(E_i)`$ and suppose Player I is only allowed to choose subspaces in $`\mathrm{\Gamma }=\{X_{i=n}^{\mathrm{}}E_i:n\}`$ then a) and b) are equivalent for all $`𝒜`$.
###### Proof of Theorem 3.3.
We first choose a decreasing sequence of finite codimensional spaces $`(Y_n)`$ in $`X`$ so that for each $`𝒜`$ the equivalences (a)$``$(b)$``$(c)$``$(d), and, if $`X^{}`$ is separable, (d)$``$(e), of Proposition 2.3 hold. Then we choose the space with an FDD $`(E_i)`$ as in Lemma 3.1.
We note that trivially (b) of the statement of Theorem 3.3 implies (c). Since the conclusion of Lemma 3.1 implies that every $`(X_n)`$-block tree (recall, $`X_n=_{i=n+1}^{\mathrm{}}E_iX`$) has for given sequence $`\delta _i0`$ a subtree which is a $`(\delta _i)`$-approximation of an $`(E_i)`$-block tree for which some branch is (SB) w.r.t. $`(G_i)`$, condition (c) implies condition (a) (Player II cannot have a winning strategy). If $`X^{}`$ is separable the statement (a)$``$(d) is exactly the statement of the second part of Proposition 2.3.
Thus, we are left with the verification of the implication (a)$``$(b).
Let $`\epsilon >0`$ and $`𝒜`$. We put $`\eta _i=\epsilon /c2^{i+2}`$, where the constant $`c>1`$ comes from the conclusion of Lemma 3.1 (c).
###### Claim.
Every tree $`(x_A)_{A[]^{<\omega }}`$ in $`S_X`$ having the property that
(23)
$$x_AX\left(S_{_{i=\mathrm{max}A+1}^{\mathrm{}}E_i}\right)_{\eta _{\mathrm{card}A}},\text{ whenever }A[]^{<\omega },$$
is an $`(\epsilon /2^n)`$-approximation to a $`(Y_n)`$-block tree, and therefore must have a branch in $`\overline{𝒜_{2\epsilon }}`$ (Proposition 2.3 (a)$``$(c)).
###### Remark.
Note that it is in general not true that if $`xX\left(S_{_{i=m}^nE_i}\right)_\delta `$, then we will be able to aproximate $`x`$ by an element in $`X_{m1}=_{j=m}^{\mathrm{}}E_jX`$ up to some $`r(\delta )`$, which converges to $`0`$ if $`\delta `$ tends to 0, and which only depends on $`\delta `$, but not on $`m`$ and $`n`$. But condition (c) of Lemma 3.1 will ensure that we can at least approximate $`x`$ by an element of $`Y_n`$, up to a fixed multiple of $`\delta `$.
In order to prove the claim it suffices to show
$`()`$
$$\{\begin{array}{cc}& \text{Let }\delta >0\text{ and }xX\left(S_{_{i=n+1}^{\mathrm{}}E_i}\right)_\delta .\hfill \\ \multicolumn{2}{c}{}\\ & \text{Then there is a }yS_{Y_n}\text{ with }xy4\delta c\text{.}\hfill \end{array}$$
In order to verify the claim we can assume without loss of generality that $`\delta <1/2c`$ (otherwise the claim is trivial). Choose $`u_{i=n}^{\mathrm{}}E_i`$ and $`vZ`$, $`v<\delta `$, so that $`x=u+v`$. From Lemma 3.1(c) we deduce (recall that $`D_n_{i=1}^nE_i^{}`$) that
$$x_{X/Y_n}c\underset{w^{}D_n}{\mathrm{max}}w^{}(x)=c\underset{w^{}D_n}{\mathrm{max}}w^{}(v)<c\delta .$$
We can therefore write $`x=\stackrel{~}{y}+d`$, with $`\stackrel{~}{y}Y_n`$ and $`dX`$, satisfying $`d<c\delta `$. Since $`x=1`$, we have $`1c\delta <\stackrel{~}{y}<1+c\delta `$. Letting $`y=\stackrel{~}{y}/\stackrel{~}{y}`$ this implies that $`xy4c\delta `$, and finishes the proof of $`()`$.
We next show that there is an increasing sequence $`N_iN`$ so that if we let $`G_i=_{s=1+N_{i1}}^{N_i}E_s`$ then for every sequence $`(x_k)S_X`$ for which there exist integers $`m_0=1<m_1<\mathrm{}`$ so that
$$\mathrm{dist}(x_k,_{s=1+m_{k1}}^{m_k1}G_s)=\mathrm{dist}(x_k,_{i=1+N_{1+m_{k1}}}^{N_{m_k1}}E_i)<\eta _k,k,$$
then $`(x_k)\overline{𝒜_{4\epsilon }}`$.
Since for all $`xS_X`$ it follows that ($`K`$ depends on the basis constant of $`(E_i)`$)
$$\mathrm{dist}(x,_{i=m+1}^nE_i)K(\mathrm{Id}P_{_{i=m+1}^nE_i})(x)$$
this will finish the proof of b) taking $`\delta _i=\eta _i/K`$.
For $`\overline{N}=(N_i)_{i=1}^{\mathrm{}}[]^\omega `$ (the set of infinite subsequences of $``$) we put ($`N_0=0`$)
(24)
$$F_i^{\overline{N}}=_{j=1+N_{i1}}^{N_i}E_i,i=1,2,\mathrm{},\text{ and }$$
(25)
$$^{\overline{N}}=\{(x_i)_{i=1}^{\mathrm{}}S_X:i\mathrm{dist}(x_i,F_{2i}^{\overline{N}})<\eta _i\}.$$
###### Remark.
For $`\overline{N}[]^\omega `$ and every $`(z_i)_{i=1}^{\mathrm{}}S_X`$, having the property that
$$\mathrm{dist}(S_{_{j=1+m_{i1}}^{m_i1}F_j^{\overline{N}}},z_i)<\eta _i,i=1,2,\mathrm{},$$
for some sequence $`1m_0<m_1<m_1+1<m_2<m_2+1<m_3<\mathrm{}`$ there is a sequence $`\overline{M}[\overline{N}]^\omega `$ so that $`(z_i)^{\overline{M}}`$.
Indeed, let $`\stackrel{~}{z}_iS_{_{j=1+m_{i1}}^{m_i1}F_j^{\overline{N}}}`$, for $`i`$ so that $`\stackrel{~}{z}_iz_i<\eta _i`$ and put $`M_{2i1}=N_{m_{i1}}`$ and $`M_{2i}=N_{m_i1}`$. Then it follows that
$$\stackrel{~}{z}_iS_{_{j=1+m_{i1}}^{m_i1}F_j^{\overline{N}}}=S_{_{s=N_{1+m_{i1}}}^{N_{m_i1}}E_s}=S_{F_{2i}^{\overline{M}}}.$$
Thus, $`(z_i)^{\overline{M}}`$.
###### Completion of the proof of Theorem 3.3.
We put
$$𝒞=\{\overline{N}[]^\omega :^{\overline{N}}\overline{𝒜_{4\epsilon }}\}.$$
It is easy to see that $`𝒞`$ is closed in the pointwise topology on $`[]^\omega `$, since $`\overline{𝒜_{4\epsilon }}`$ is closed with respect to the product of the discrete topology on $`S_X^\omega `$.
By the infinite version of Ramsey’s theorem (cf.\[O\]) we deduce that one of the following two cases occurs.
$`\text{Either there exists an }\overline{N}[]^\omega \text{ so that }[\overline{N}]^\omega 𝒞.`$
$`\text{Or there exists an }\overline{N}[]^\omega \text{ so that }[\overline{N}]^\omega []^\omega 𝒞.`$
If the first alternative occurs we are finished by the above remark. Assuming the second alternative, we will show that there is a tree in $`S_X`$ satisfying (23) without any branch in $`\overline{𝒜_{2\epsilon }}`$. This would be a contradiction and imply that the second alternative cannot occur.
If we assume the second alternative we can pick for each $`\overline{M}[\overline{N}]^\omega `$ a sequence $`(y_i^{\overline{M}})_{i=1}^{\mathrm{}}^{\overline{M}}`$ which is not in $`𝒞`$. Let $`\overline{N}=\{N_1,N_2,\mathrm{}\}`$.
Note that for any $`\overline{M}\{N_3,N_4,\mathrm{}\}`$,
$$y_1^{(N_1,N_2,\overline{M})}S_X\left(S_{_{i=1+N_1}^{N_2}E_i}\right)_{\eta _1}.$$
Here $`(N_1,N_2,\overline{M})`$ is the infinite sequence starting with $`N_1`$ and $`N_2`$ and then consisting of the elements of $`\overline{M}`$).
Using the finite version of Ramsey’s theorem and the compactness of $`S_{_{i=1+N_1}^{N_2}E_i}`$ we can find a vector
$$x_{\{1\}}S_X\left(S_{_{i=1+N_1}^{N_2}E_i}\right)_{\eta _1}$$
and an $`\overline{M}^{(1)}\{N_3,N_4,\mathrm{}\}`$ such that
(26)
$$x_{\{1\}}y_1^{(N_1,N_2,\overline{M})}<2\eta _1\text{ for all }\overline{M}[\overline{M}^{(1)}]^\omega .$$
Doing the same procedure again, we can find an
$$x_{\{2\}}S_X\left(S_{_{i=1+N_1^{(2)}}^{N_2^{(2)}}E_i}\right)_{\eta _1}$$
and an $`\overline{M}^{(2)}[\overline{M}_{(1)}]^\omega `$ so that
$$x_{\{2\}}y_1^{N_1^{(2)},N_2^{(2)},\overline{M}}<2\eta _1\text{ for all }\overline{M}[\overline{M}^{(2)}]^\omega ,$$
where $`N_1^{(2)}`$ and $`N_2^{(2)}`$ are the first two elements of the sequence $`\overline{M}^{(1)}`$. Proceeding this way we construct a sequence $`x_{\{i\}}`$ and a decreasing sequence $`(\overline{M}^{(i)})`$ of infinite subsequences of $`\overline{N}`$ so that
$$x_{\{i\}}S_X\left(S_{_{j=1+N_1^{(i)}}^{N_2^{(i)}}E_j}\right)_{\eta _1},\text{ and }$$
$$x_{\{i\}}y^{(N_1^{(i)},N_2^{(i)},\overline{M})}<\eta _1,\text{ for all }\overline{M}[\overline{M}^{(i)}]^\omega .$$
This sequence will be the first level of a tree and the beginning of the level by level recursive construction of this tree as follows.
Assume that for some $`\mathrm{}`$ and every $`A[]^{\mathrm{}}`$ we have chosen an $`x_AS_X`$, a pair of natural number $`N_1^{(A)}`$, and $`N_2^{(A)}`$, and a sequence $`\overline{M}^{(A)}[\{N\overline{N}:N>N_2^{(A)}\}]^\omega `$ so that the following conditions (27) and (28) are satisfied.
(27) $`\text{If }A[]^<\mathrm{}\{\mathrm{}\}\text{ and }n>m>\mathrm{max}A\text{ then}`$
$`N_1^{(A)}<N_2^{(A)}<N_1^{(A\{m\})}<N_2^{(A\{m\})}<N_1^{(A\{n\})}<N_1^{(A\{n\})}`$
$`[N_1^{(\mathrm{})}=N_2^{(\mathrm{})}=0]`$
$`\overline{M}^{(A)}\overline{M}^{(A\{n\})}`$
(28) $`\text{If }n_1<n_2<\mathrm{}<n_{\mathrm{}}\text{ are in },\text{ we put}`$
$`A_j=\{n_1,n_2,\mathrm{}n_j\}\text{ for }j=1,2,\mathrm{}\mathrm{}.`$
Then:
$`x_{A^{(j)}}S_X\left(S_{_{s=N_1^{(A_j)}+1}^{N_2^{(A_j)}}E_s}\right)_{\eta _j}`$
$`x_{A^{(j)}}y^{(N_1^{\left(A_1\right)},N_2^{\left(A_1\right)},\mathrm{}N_1^{\left(A_j\right)},N_2^{\left(A_j\right)},\overline{M})}<\eta _j`$
$`\text{ whenever }\overline{M}[\overline{M}^{(A_j)}]^\omega `$
Then we can choose for $`A[]^{\mathrm{}}`$ the elements $`x_{A\{1+\mathrm{max}A\}}`$, $`x_{A\{2+\mathrm{max}A\}}`$ etc., and the numbers $`N_1^{(A\{1+\mathrm{max}A\})}`$, $`N_2^{(A\{1+\mathrm{max}A\})}`$, $`N_1^{(A\{2+\mathrm{max}A\})}`$, $`N_2^{(A\{2+\mathrm{max}A\})}`$ etc. and the sets $`\overline{M}^{(A\{1+\mathrm{max}A\})}`$, $`\overline{M}_1^{(A\{2+\mathrm{max}A\})}`$, etc. exactly in the same way we chosed $`x_{\{1\}}`$, $`x_{\{2\}}`$ etc. and the numbers $`N_1^{(1)},N_2^{(1)},N_1^{(2)},N_2^{(2)}`$ etc. for the first level.
The condition (28) implies that for every branch $`(z_n)`$ of the constructed tree there is an $`\overline{M}[]^\omega `$ so that $`z_ny_n^{\overline{M}}|2\eta _n`$, for all $`n`$. Since $`(y_n^{\overline{M}})\overline{𝒜_{4\epsilon }}`$ it follows that (recall that $`\eta _n\epsilon /2^n`$) $`(z_n)\overline{𝒜_{2\epsilon }}`$, which is a contradiction and finishes the proof. ∎
## 4. Subspaces of $`(_{i=1}^{\mathrm{}}F_i)_p`$
The purpose of this section is to use Theorem 3.3 to produce an intrinsic characterization of a necessary and sufficient condition that ensures a given Banach space $`X`$ will embed into an $`\mathrm{}_p`$-sum of finite dimensional spaces.
Let $`1p<\mathrm{}`$ and let $`F_i`$ be a finite dimensional space for $`i`$. The $`\mathrm{}_p`$-sum of $`(F_i)`$, $`(F_i)_p`$, is the space of all sequences $`(x_i)`$, with $`x_iF_i`$, for $`i=1,2\mathrm{}`$, so that
$$(x_i)_p=\left(\underset{i=1}{\overset{\mathrm{}}{}}x_i_{F_i}\right)^{1/p}<\mathrm{}.$$
###### Theorem 4.1.
Assume that $`X`$ is reflexive and that there are $`1<p<\mathrm{}`$, and $`C>1`$ so that every weakly null tree in $`S_X`$ has a branch which is $`C`$-equivalent to the unit vector basis of $`\mathrm{}_p`$.
Then $`X`$ is isomorphic to a subspace of an $`\mathrm{}_p`$-sum of finite dimensional spaces.
More precisely, for any $`\epsilon >0`$ there exists a finite codimensional subspace $`\stackrel{~}{X}`$ of $`X`$, so that $`\stackrel{~}{X}`$ is $`(C^2+\epsilon )`$-isomorphic to a subspace of an $`\mathrm{}_p`$-sum of finite dimensional spaces.
Before we start the proof, some remarks are in order.
###### Remark.
The assumption that $`X`$ is reflexive is necessary. Indeed, James’ space $`J`$ \[Ja1\] is not reflexive but has the property that every weakly null tree in $`S_J`$ has a branch which is 2-equivalent to the unit vector basis of $`\mathrm{}_2`$. Actually every normalized skipped block with respect to the shrinking basis of $`J`$ is 2-isomorphic to the unit vector basis of $`\mathrm{}_2`$. Since every $`\mathrm{}_2`$ sum of finite dimensional spaces must be reflexive, $`J`$ cannot be isomorphic to a subspace of such a space.
In \[KW\] Kalton and Werner showed a special version of above result. They proved the conclusion of Theorem 4.1 (with $`C=1`$) under the condition that $`X`$ does not contain a copy of $`\mathrm{}_1`$ and every weakly-null type is an $`\mathrm{}_p`$ type. This means that for every $`xS_X`$ and every normalized weakly null sequence $`(x_n)S_X`$ for $`t>0`$ one has
(29)
$$\underset{n\mathrm{}}{lim}x+tx_n=(1+t^p)^{1/p}.$$
In \[KW\] it was shown that this condition implies that $`X`$ must be reflexive, and it is easy to see that it also implies the hypothesis of Theorem 4.1 with $`C=1+\epsilon `$ for any $`\epsilon >0`$.
Secondly, let us explain the reason for the $`C^2`$ term rather than $`C`$ in the conclusion of Theorem 4.1. A normalized basis $`(x_i)`$ is $`C`$-equivalent to the unit vector basis of $`\mathrm{}_p`$ if there exist constants $`A,B`$ with $`ABC`$ and
($``$)
$$A^1\left(\underset{i=1}{\overset{\mathrm{}}{}}|a_i|^p\right)^{1/p}\underset{i=1}{\overset{\mathrm{}}{}}a_ix_iB\left(\underset{i=1}{\overset{\mathrm{}}{}}|a_i|\right)^{1/p}$$
for all scalars $`(a_i)`$. If we had the hypothesis that every weakly null tree in $`S_X`$ admitted a branch $`(x_i)`$ with this property then we could obtain the conclusion of Theorem 4.1 with $`C^2`$ replaced by $`C`$. However the constants $`A,B`$ above could vary with each such tree and so we can only use $`()`$ with $`A`$ and $`B`$ replaced by $`C`$. In this case we only get $`C^2`$-embedding into $`\mathrm{}_p`$.
We also note that Kalton \[K\] proved the following analogous theorem for $`c_0`$: Let $`X`$ be a separable Banach space not containing $`\mathrm{}_1`$. If there exists $`C<\mathrm{}`$ so that every weakly null tree in $`S_X`$ has a branch $`C`$-equivalent to the unit vector basis of $`c_0`$ then $`X`$ embeds into $`c_0`$.
W.B. Johnson \[J2\] showed that in the case $`XL_p`$ $`(1<p<\mathrm{})`$, if there exists $`K<\mathrm{}`$ so that every normalized sequence in $`X`$ has a subsequence $`K`$-equivalent to the unit vector basis of $`\mathrm{}_p`$ then $`X`$ embeds into $`\mathrm{}_p`$. The tree hypothesis of Theorem 4.1 cannot in general be weakened to the subsequence condition as the following example shows. (Theorem 4.1 and this example solve some questions raised in \[J2\].)
###### Example 4.2.
Let $`1<p<\mathrm{}`$. There exists a reflexive space $`X`$ with an unconditional basis so that $`X`$ satisfies: for all $`\epsilon >0`$ every normalized weakly null sequence in $`X`$ admits a subsequence $`1+\epsilon `$-equivalent to the unit vector basis of $`\mathrm{}_p`$. Yet $`X`$ is not a subspace of an $`\mathrm{}_p`$-sum of finite dimensional spaces.
###### Proof.
Fix $`1<q<p`$. We define $`X=(X_n)_p`$ where each $`X_n`$ is given as follows. $`X_n`$ will be the completion of $`c_{00}([]^n)`$ under the norm
$$x_n=sup\{(\underset{i=1}{\overset{m}{}}x|_{\beta _i}_q^p)^{1/p}:(\beta _i)_1^m\text{ are disjoint segments in }[]^n\}.$$
By a segment we mean a sequence $`(A_i)_{i=1}^k[]^n`$ with $`A_1=\{n_1,n_2,\mathrm{}n_{\mathrm{}}\}`$, $`A_2=\{n_1,n_2,\mathrm{}n_{\mathrm{}},n_{\mathrm{}+1}\}`$ $`\mathrm{}`$ $`A_k=\{n_1,n_2,\mathrm{}n_{\mathrm{}},n_{\mathrm{}+1}\mathrm{}n_{\mathrm{}+k1}\}`$, for some $`n_1<n_2<\mathrm{}n_{\mathrm{}+k1}`$. Thus a segement can be seen as an interval of a branch (with respect to the usual partial order in $`[]^n`$), while a branch is a maximal segment.
Clearly the node basis $`(e_A^{(n)})_{A[]^n}`$ given by $`e_A(B)=\delta _{(A,B)}`$ is a 1-unconditional basis for $`X_n`$. Furthermore the unit vector basis of $`\mathrm{}_q^n`$ is 1-equivalent to $`(e_{A_i}^{(n)})_1^n`$, if $`(A_i)_1^n`$ is any branch of $`[]^n`$.
Thus no extension of the tree $`(e_A^{(n)})_{A[]^n}`$ to a weakly null tree of infinite length in $`S_X`$ has a branch whose basis distance to the $`\mathrm{}_p`$-unit vector basis is closer than $`\mathrm{dist}_b(\mathrm{}_p^{(n)},\mathrm{}_q^{(n)})=n^{\frac{1}{q}\frac{1}{p}}\mathrm{}`$ for $`n\mathrm{}`$. Since it is clear that in every subspace $`Y`$ of an $`\mathrm{}_p`$ sum of finite dimensional spaces every weakly null tree in $`S_Y`$ must have a branch equivalent (for a fixed constant) to the unit vector basis of $`\mathrm{}_p`$ it follows that $`X`$ cannot be embedded into a subspace of an $`\mathrm{}_p`$-sum of finite dimensional spaces.
Also each $`X_n`$ is isomorphic to $`\mathrm{}_p`$ and thus $`X`$ is reflexive.
It remains to show that if $`(x_j)`$ is a normalized weakly null sequence in $`X`$ and $`\epsilon >0`$ then a subsequence is $`1+\epsilon `$-equivalent to the unit vector basis of $`\mathrm{}_p`$. By a gliding hump argument it suffices to prove this in a fixed $`X_n`$. We proceed by induction on $`n`$.
For $`n=1`$ the result is clear since $`X_1`$ is isometric to $`\mathrm{}_p`$. Assume the result has been proved for $`X_{n1}`$. By passing to a subsequence and perturbing we may assume that $`(x_i)_1^{\mathrm{}}`$ is a normalized block basis of the node basis for $`X_n`$.
Let $`\epsilon _i0`$ rapidly. For $`j`$ let $`P_j`$ be the basis projection of $`X_n`$ onto $`[e_A:A[]^n`$, $`\mathrm{min}A=j]`$. Passing to a subsequence we may assume that $`lim_i\mathrm{}P_jx_i_n=a_i`$ and from the definition of $`_n`$ we have $`(a_i)_{i=1}^{\mathrm{}}B_\mathrm{}_p`$. Choose $`a_00`$ so that $`(a_i)_{i=1}^{\mathrm{}}S_\mathrm{}_p`$.
Passing to a subsequence of $`(x_i)`$ we may assume that there exist integers $`1=N_0<N_1<\mathrm{}`$ so that
1. $`x_i(\{j\})0j[N_i,N_{i+1})`$
2. $`P_jx_i=0`$ for $`jN_{i+1}`$
3. $`{\displaystyle \underset{j[N_i,N_{i+1})}{}}P_jx_i_n=\left({\displaystyle \underset{j[N_i,N_{i+1})}{}}P_jx_i_n^p\right)^{1/p}`$ is within $`\epsilon _i`$ of $`a_0`$.
4. If $`j[N_i,N_{i+1})`$, $`i1`$, then if $`a_j0`$, $`(a_j^1P_jx_{\mathrm{}})_{\mathrm{}>i}`$ is $`1+\epsilon _j`$-equivalent to the unit vector basis of $`\mathrm{}_p`$.
5. If $`j[N_0,N_1)`$ and $`a_j0`$ then $`(a_j^1P_jx_{\mathrm{}})_{\mathrm{}=1}^{\mathrm{}}`$ is $`1+\epsilon _j`$-equivalent to the unit vector basis of $`\mathrm{}_p`$.
6. $`\left({\displaystyle \underset{N_1}{\overset{\mathrm{}}{}}}a_j^p\right)^{1/p}<\epsilon _1`$
7. If $`j[N_0,N_1)`$ and $`a_j=0`$ then $`P_jx_i_n\epsilon _i`$ for all $`i`$.
8. If $`j[N_i,N_{i+1})`$ and $`a_j=0`$ then $`P_jx_{\mathrm{}}_n<\epsilon _{\mathrm{}}`$ for $`\mathrm{}>i`$.
Conditions (iv) and (v) use the induction hypothesis and the fact that for all $`j`$,span$`(\{e_{\{j\}A)}:A[]^{n1}`$, $`\mathrm{min}A>j\})`$ is isometric to $`X_{n1}`$. Our conditions are sufficient to yield (for suitably small $`\epsilon _j`$’s) that $`(x_i)`$ is $`1+\epsilon `$-equivalent to the unit vector basis of $`\mathrm{}_p`$. We omit the standard yet tedious calculations. ∎
For the proof of Theorem 4.1 we need a result which was shown in \[KOS\]. It is based on a trick of W. B. Johnson \[J2\] where part (a) was shown.
###### Lemma 4.3.
(Lemma 5.1 in \[KOS\]) Let $`X`$ be a subspace of a space $`Z`$ having a boundedly complete FDD $`(F_n)`$ and assume $`X`$ is w closed (since $`(F_n)`$ is boundedly complete $`Z`$ is naturally a dual space). Then for all $`\epsilon >0`$ and $`m`$ there exists an $`n>m`$ such that if $`x=_1^{\mathrm{}}x_iB_X`$ with $`x_iF_i`$ for all $`i`$, then there exists $`k(m,n]`$ with
1. $`x_k<\epsilon `$ and
2. $`\mathrm{dist}(_{i=1}^{k1}x_i,X)<\epsilon `$.
###### Corollary 4.4.
Let $`X`$ be a subspace of the reflexive space $`Z`$ and let $`(F_i)`$ be an FDD for $`Z`$. Let $`\delta _i0`$. There exists a blocking $`(G_i)`$ of $`(F_i)`$ given by $`G_i=_{j=N_{i1}+1}^{N_i}F_j`$ for some $`0=N_0<N1<\mathrm{}`$ with the following property. For all $`xS_X`$ there exist $`(x_i)_1^{\mathrm{}}X`$ and $`t_i(N_{i1},N_i]`$ for $`i`$ so that
* $`x={\displaystyle \underset{i=1}{\overset{\mathrm{}}{}}}x_i`$.
* For $`i`$ either $`x_i<\delta _i`$ or $`P_{_{j=t_{i1}+1}^{t_i1}F_j}(x_i)x_i<\delta _ix_i`$
* For $`i`$, $`P_{_{j=t_{i1}+1}^{t_i1}F_j}xx_i<\delta _i`$.
###### Proof.
We choose an appropriate sequence $`\epsilon _i0`$ depending upon $`(\delta _i)`$ and the basis constant $`K`$ of $`(F_i)`$. $`N_1`$ is chosen by the lemma for $`\epsilon =\epsilon _1`$ and $`m=1`$. We choose $`N_2>N_1`$ by the lemma for $`\epsilon =\epsilon _2`$ and $`m=N_1`$ and so on.
If $`xS_X`$ the lemma yields for $`i`$, $`t_i(N_{i1},N_i]`$ with $`P_{F_{t_i}}(x)<\epsilon _i`$ and $`z_iX`$ with $`P_{_{j=1}^{t_i1}F_j}(x)z_i<\epsilon _i`$. We then let $`x_1=z_1`$ and for $`i>1`$, $`x_i=z_iz_{i1}`$. Thus $`_{i=1}^nx_i=z_nx`$ and so a) holds.
To see c) we note the following
$$P_{_{j=t_{i1}+1}^{t_i1}F_j}(x)x_iP_{_{j=1}^{t_i1}F_j}(x)z_i+P_{_{j=1}^{t_{i1}}F_j}(x)z_{i1}<\epsilon _i+2\epsilon _{i1}.$$
Thus
$$P_{_{j=t_{i1}+1}^{t_i1}F_j}(x_i)x_i=(\mathrm{Id}P_{_{j=t_{i1}+1}^{t_i1}F_j})(x_iP_{_{j=t_{i1}+1}^{t_i1}F_j}x)<(2K+1)(\epsilon _i+2\epsilon _{i1})$$
which can be made less than $`\delta _i^2`$. This yields b). ∎
###### Remark.
The proof yields that the conclusion of the corollary remains valid for any further blocking of the $`G_i`$’s (which would redefine the $`N_i`$’s).
###### Proof of Theorem 4.1.
We first show that $`X`$ embeds into $`(G_n)_\mathrm{}_p`$ for some sequence $`(G_n)`$ of finite dimensional spaces. Then to obtain the $`C^2+\epsilon `$ estimate we adapt an averaging argument similar to the one of \[KW\].
Applying Theorem 3.3 to the set
$$𝒜=\{(x_i)S_X^\omega :(x_i)\text{ is }C\text{-equivalent to the unit vector basis of }\mathrm{}_p\}$$
we find a reflexive space $`Z`$ with an FDD $`(F_i)`$ with basis constant $`K`$ which isometrically contains $`X`$ and $`\delta _i0`$ so that whenever $`(x_i)S_X`$ satisfies
(30)
$$P_{_{j=n_{i1}+1}^{n_i1}F_j}(x_i)x_i<\delta _i$$
for some sequence $`1=n_0<n_1<\mathrm{}`$ in $``$ it follows that $`(x_i)`$ is $`2C`$-equivalent to the unit vector basis of $`\mathrm{}_p`$. Let $`G_i=_{j=N_{i1}+1}^{N_i}F_j`$ be the blocking given by Corollary 4.4.
Let $`xS_X`$, $`x=\overline{x}_i`$ with $`\overline{x}_iG_i`$ for all $`i`$. Choose $`(x_i)`$ and $`(t_i)`$ as in Corollary 4.4. It follows from (30) that (for $`\delta _i`$’s sufficiently small) that
$$(3C)^1\left(x_i^p\right)^{1/p}3C$$
and
$$(4C)^1\left(\underset{i}{}P_{_{j=t_{i_1}+1}F_j}^{t_i1}x^p\right)^{1/p}4C.$$
Let $`y_i=P_{_{j=t_{i1}+1}^{t_i1}F_j}x`$.
Since
$$\frac{1}{2(K+1)}\mathrm{max}(y_i,y_{i+1})\delta _i\overline{x}_i(2K+1)y_i+\delta _i$$
it follows that $`X`$ embeds isomorphically into $`(G_i)_\mathrm{}_pW`$.
We now renorm $`W`$ so as to contain $`X`$ isometrically. Thus $`W`$ has $`(G_i)`$ as an FDD and there exists $`\stackrel{~}{C}`$ so that if $`(w_i)`$ is any block basis of a permutation of $`(G_i)`$ then
(31)
$$\stackrel{~}{C}^1(w_i^p)^{1/p}w_i\stackrel{~}{C}(w_i^p)^{1/p}.$$
We repeat the first part of the proof. Let $`\epsilon >0`$. From Theorem 3.3 we may assume that there exist $`\delta _i0`$ so that if $`(x_i)S_X`$ satisfies
(32)
$$P_{_{j=n_{i1}+1}^{n_i1}G_j}(x_i)x_i<\delta _i$$
for some $`1=n_0<n_1<\mathrm{}`$ then $`(x_i)`$ is $`C+\epsilon `$-equivalent to the unit vector basis of $`\mathrm{}_p`$. Moreover we may assume that this is valid for any further blocking of $`(G_j)`$. From now on we will replace $`X`$ by the finite codimensional subspace $`_{i=2}^{\mathrm{}}G_iX`$ and $`W`$ by $`_{i=2}^{\mathrm{}}G_i`$ and replace $`G_i`$ by $`G_{i+1}`$. We will show that this new $`X`$ can be $`C^2+\epsilon `$-embedded into an $`\mathrm{}_p`$ sum of finite dimensional spaces.
Let $`H_i=_{j=N_{i1}+1}^{N_i}G_j`$ be the blocking given by Corollary 4.4. Thus (for appropriately small $`\delta _i`$’s) from (32) and Corollary 4.4 we have that if $`xS_X`$ there exist $`t_i(N_{i1},N_i]`$ so that
(33)
$$(C+2\epsilon )^1\left(\underset{i=1}{\overset{\mathrm{}}{}}\underset{j=t_{m_{i1}+1}}{\overset{t_{m_i}}{}}x_j^p\right)^{1/p}x(C+2\epsilon )\left(\underset{i=1}{\overset{\mathrm{}}{}}\underset{j=t_{m_{i1}+1}}{\overset{t_{m_i}}{}}x_j^p\right)$$
where $`x=x_i`$ is the expansion of $`X`$ w.r.t. the FDD $`(G_j)`$ for $`W`$.
Chose $`M`$ so that
(34)
$$\frac{\stackrel{~}{C}^{2^{1/p}}}{M}\epsilon \text{ and }(C+2\epsilon )^1\frac{\stackrel{~}{C}^2}{M^{1/p}}(C+3\epsilon )^1.$$
For $`i=1,2,\mathrm{},M`$ and $`j=0,1,2,\mathrm{}`$ set $`L(i,j)=_{s=(j1)M+i+1}^{jM+i1}H_sW`$ (using $`H_n=\{0\}`$ if $`n0`$) and let $`Y_i=(_{j=0}^{\mathrm{}}L(i,j))_p`$. Let $`Y=(_{i=1}^MY_i)_p`$. We shall prove that $`X`$ $`C^2+\eta (\epsilon )`$-embeds into $`Y`$ where $`\eta (\epsilon )0`$ as $`\epsilon 0`$ which will complete the proof.
To do this we first define maps $`T_i:XY_i`$ for $`1iM`$. If $`x=x_j`$ is the expansion of $`x`$ w.r.t. $`(H_j)`$ we let
$$T_ix=\underset{s=1}{\overset{\mathrm{}}{}}\left(\underset{u=(s1)M+i+1}{\overset{sM+i1}{}}x_s\right)\left(_{s=1}^{\mathrm{}}L(i,s)\right)_p=Y_i.$$
Let $`1iM`$ and $`xS_X`$, $`x=x_j`$ as above. Write $`x_j=_{u=N_{j1}+1}^{N_j}x(j,u)`$ as the expansion of $`x_jH_j`$ w.r.t. $`(G_i)`$. Let $`(t_i)`$ be given by Corollary 4.4 (w.r.t. $`(G_j)`$). From several applications of the triangle inequality and (31) and (33) we have
$`T_i(x)`$ $`=\left[{\displaystyle \underset{j=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{s=(j1)M+i+1}{\overset{jM+i1}{}}}x(s)^p\right]^{1/p}`$
$`[{\displaystyle \underset{j=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{u=t_{(j1)M+i}}{\overset{N_{(j1)N+i}}{}}}x((j1)M+i,u)+{\displaystyle \underset{s=(j1)M+i+1}{\overset{jM+i1}{}}}x(s)`$
$`+{\displaystyle \underset{u=1+N_{jM+i1}}{\overset{t_{jM+i}}{}}}x(jM+i,u)^p]^{1/p}`$
$`+\left[{\displaystyle \underset{j=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{u=t_{(j1)M+i}}{\overset{N_{(j1)N+i}}{}}}x((j1)M+i,u)+{\displaystyle \underset{u=1+N_{jM+i1}}{\overset{t_{jM+i}}{}}}x(jM+i,u)^p\right]^{1/p}`$
$`(C+2\epsilon )x+\left[{\displaystyle \underset{j=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{u=1+N_{jM+i1}}{\overset{N_{jM+i}}{}}}x(jM+i,u)^p\right]^{1/p}`$
$`(C+2\epsilon )x+\stackrel{~}{C}{\displaystyle \underset{s=i(modM)}{}}x_s.`$
Similarly one has
$$T_ix(C+2\epsilon )^1x\stackrel{~}{C}\underset{s=i(modM)}{}x_s.$$
Finally we define $`T:XY=(_1^MY_i)`$, by $`Tx=\frac{1}{M^{1/p}}_{i=1}^MT_ix`$. Note that
$`Tx`$ $`{\displaystyle \frac{1}{M^{1/p}}}(C+2\epsilon )\left({\displaystyle \underset{i=1}{\overset{M}{}}}x^p\right)^{1/p}+{\displaystyle \frac{\stackrel{~}{C}}{M^{1/p}}}\left({\displaystyle \underset{i=1}{\overset{M}{}}}{\displaystyle \underset{j=i(modM)}{}}x_j\right)`$
$`(C+2\epsilon )x+{\displaystyle \frac{\stackrel{~}{C}^2}{M^{1/p}}}x<(C+3\epsilon )x`$
using (31) and (34).
Similarly one deduces that for $`xX`$ it follows that $`T(x)\frac{1}{C+3\epsilon }x`$. ∎
###### Remark.
The proof of Theorem 4.1 had two steps. In the first we started with an embedding of $`X`$ into a certain reflexive space $`Z`$ with an FDD $`(F_i)`$ and showed that $`(F_i)`$ can be blocked to an FDD $`(G_i)`$ so that $`X`$ is isomorphic to a subspace of $`(G_i)_\mathrm{}_p`$. In that step we could not deduce any bound for the constant of that isomorphism. In the second step we “inflated” $`(G_i)_\mathrm{}_p`$ to the space $`(_{i=1}^M_{ji(modM)}G_j)_\mathrm{}_p`$ and showed that this space contains a finite codimensional subspace which is $`C^2+\epsilon `$-equivalent to $`X`$.
The following example shows that even if the space $`X`$ has a basis to begin with, it is in general not possible to pass to a blocking $`(F_n)`$ of that basis and deduce that for some $`n_0`$ the identity is a $`C^2+\epsilon `$-isomorphism between $`_{n=n_0}^{\mathrm{}}F_n`$ and $`(_{n=n_0}^{\mathrm{}}F_n)_\mathrm{}_p`$.
###### Example 4.5.
Let $`𝒟`$ be the set of all sequences $`(D_n)`$ of pairwise disjoint subsets of $``$, so that for each $`n`$, $`D_n`$ is either a singleton or it is of the form $`D_n=\{k,k+1\}`$ for some $`k`$. We give $`\mathrm{}_2`$ the following equivalent norm $`||||||:`$
$$|x|=\underset{(D_n)𝒟}{sup}\left(\underset{n=1}{\overset{\mathrm{}}{}}\left(\underset{jD_n}{}|x_j|\right)^2\right)^{1/2},$$
whenever $`x=(x_j)\mathrm{}_2`$.
It is easy to see that every normalized skipped block $`(x^{(n)})`$ in $`X=(\mathrm{}_2,||||||)`$ is isometrically eqivalent to the $`\mathrm{}_2`$ unit vector basis. Thus the assumptions of Theorem 4.1 are satisfied for any $`C>1`$. On the other hand for any blocking $`(F_n)`$ of the unit vector basis $`(e_i)`$ of $`X`$ it follows for any $`n`$ and $`N_n=\mathrm{max}\{N|e_NF_n\}`$ that $`e_{N_n+1}F_{n+1}`$ and that the span of $`e_{N_n}`$ and $`e_{N_n+1}`$ is isometric to $`\mathrm{}_1^2`$. Therefore the norm of the identity between $`(_{n=2}^{\mathrm{}}F_n)_\mathrm{}_2`$ and $`(_{n=2}^{\mathrm{}}F_n)_{\mathrm{}_2,||||||}`$ is at least $`\sqrt{2}`$.
The following result shows that the property that every normalized weakly null tree contains a branch which is $`C`$-equivalent to the $`\mathrm{}_p`$ unitvector basis dualizes. It can be seen as the isomorphic version of Theorem 2.6. in \[KW\].
###### Corollary 4.6.
Assume $`X`$ is a reflexive Banach space. For $`1<p<\mathrm{}`$ and $`\frac{1}{p}+\frac{1}{q}=1`$ the following statements are equivalent.
1. There is a $`C1`$ so that every normalized weakly null tree in $`X`$ has a branch which is $`C`$-equivalent to the unit vector basis of $`\mathrm{}_p`$.
2. There is a $`C1`$, a finite codimensional subspace $`\stackrel{~}{X}`$ of $`X`$, a sequence of finite dimensional spaces $`(E_i)_{i=1}^{\mathrm{}}`$, and an operator $`T:\stackrel{~}{X}(_{i=1}^{\mathrm{}}E_i)_\mathrm{}_p`$, so that $`C^1xT(x)Cx`$ for all $`x\stackrel{~}{X}`$.
3. There is a $`C1`$ so that every normalized weakly null tree in $`X^{}`$ has a branch which is $`C`$-equivalent to the unit vector basis of $`\mathrm{}_q`$.
4. There is a $`C1`$, a finite codimensional subspace $`Y`$ of $`X^{}`$, a sequence of finite dimensional spaces $`(E_i)_{i=1}^{\mathrm{}}`$, and an operator $`T:Y(_{i=1}^{\mathrm{}}E_i)_\mathrm{}_q`$, so that $`C^1xT(x)Cx`$ for all $`xY`$.
###### Proof.
The implications (a)$``$(b) and (c)$``$(d) follow from Theorem 4.1 and its proof. If we prove (b)$``$(c) then (d)$``$(a) will follow.
Assume that $`C1`$, $`(E_i)_{i=1}^{\mathrm{}}`$, $`\stackrel{~}{X}X`$ and $`T:\stackrel{~}{X}Z=(_{i=1}^{\mathrm{}}E_i)_\mathrm{}_p`$ are given as in the statement of (b). By passing to the renorming $`||||||`$, $`|x|=T(x)`$, for $`x\stackrel{~}{X}`$ we can assume without loss of generality that $`\stackrel{~}{X}`$ is isometric to a subspace of $`Z`$.
We will show that $`\stackrel{~}{X}^{}`$ satisfies the condition (c). Since $`\stackrel{~}{X}^{}`$ is isomorphic to a subspace of $`X^{}`$ of finite codimension the claim will follow.
Thus let $`E:\stackrel{~}{X}(_{i=1}^{\mathrm{}}E_i)_\mathrm{}_p`$ be an isometric embedding and let $`(x_A^{})_{A[]^{<\omega }}`$ be a normalized weakly null tree in $`\stackrel{~}{X}^{}`$.
We will need the following observation.
###### Claim.
If $`(x_n^{})`$ is a normalized and weakly null sequence in $`\stackrel{~}{X}^{}`$, then there are normalized weakly null sequences $`(z_n^{})`$ and $`(x_n)`$ in $`Z^{}`$ and $`\stackrel{~}{X}`$ respectively so that, $`E^{}(z_n^{})=x_n^{}`$ and $`x_n^{}(x_n)=1`$ for $`n`$.
To see this use the Hahn-Banach theorem to choose a normalized sequence $`(z_n^{})_n`$ in $`Z^{}`$ so that $`E^{}(z_n^{})=x_n^{}`$. The sequence $`(z_n^{})`$ is weakly null. Indeed, otherwise we could choose a $`y^{}Z^{}`$, $`y^{}0`$, a subsequence $`(z_{n_k}^{})`$ and a weakly null sequence $`(y_k^{})`$ in $`Z^{}`$ so that $`z_{n_k}^{}=y^{}+y_k^{}`$ for all $`k`$. Thus, $`x_{n_k}^{}=E^{}(y^{})+E^{}(y_k^{})`$, which implies that $`E^{}(y^{})=0`$ and therefore that $`E^{}(y_k^{})=x_{n_k}^{}`$. Since $`lim\; sup_k\mathrm{}y_k^{}=lim\; sup_k\mathrm{}(z_{n_k}^{}^qy^{}^q)^{(1/q)}<1`$, we get a contradiction.
Then we choose $`(x_n)\stackrel{~}{X}`$ so that $`x_n^{}(x_n)=1`$. By a similar argument we have that $`(x_n)`$ is also weakly null.
Using the claim we can find a normalized weakly null tree $`(z_A^{})_{A[]^{<\omega }}`$ in $`Z^{}`$ and a normalized weakly null tree $`(x_A)_{A[]^{<\omega }}`$ in $`\stackrel{~}{X}`$, so that $`E^{}(z_A^{})=x_A^{}`$ and $`x_A^{}(x_A)=1`$ for $`A[]^{<\omega }`$.
Given an $`\epsilon >0`$ we can choose a branch $`(x_n^{})=(x_{A_n}^{})`$ so that $`(z_{A_n}^{})`$ is $`(1+\epsilon )`$ equivalent to the unit vector basis of $`\mathrm{}_q`$, and $`(x_{A_n})`$ is $`(1+\epsilon )`$ equivalent to the unit vector basis of $`\mathrm{}_p`$. This easily implies that $`(x_{A_n}^{})`$ is $`(1+\epsilon )`$ equivalent to the unit vector basis of $`\mathrm{}_q`$. ∎
###### Remark.
W.B. Johnson and M. Zippin \[JZ\] proved the following. Let $`C_p=(_{i=1}^{\mathrm{}}E_i)_\mathrm{}_p`$ where $`(E_i)`$ is dense, in the Banach-Mazur sense, in the set of all finite dimensional spaces. Then $`X`$ embeds into $`C_p`$ if and only if $`X^{}`$ embeds into $`C_q`$ (where $`\frac{1}{p}+\frac{1}{q}=1`$). Thus Corollary 4.6 could be deduced from \[JZ\] and Theorem 4.1 (and \[JZ\] could be deduced from the corollary and theorem).
Furthermore the proof of Corollary 4.6 yields some quantitative information. If a) holds then b) is true with $`C`$ replaced by $`C+\epsilon `$ for any $`\epsilon >0`$. If b) holds then c) is valid with $`C`$ replaced by $`C^2+\epsilon `$.
## 5. Spectra and asymptotic structures
In \[Mi\] Milman introduced the notion of the spectra of a function defined on $`S_X^n`$. Let $`(M,\rho )`$ be a compact metric space and let $`f:S_X^nM`$ be Lipschitz. $`\sigma (f)`$ is defined to be the set of all $`aM`$ for which the following condition (35) is true
(35) $`\epsilon >0Y_1\mathrm{cof}(X)`$ $`y_1S_XY_2\mathrm{cof}(X)y_2S_X`$
$`\mathrm{}Y_n\mathrm{cof}(X)y_nS_X\text{ so that}`$
$`\rho (f(y_1,y_2,\mathrm{}y_n),a)<\epsilon `$
In terms of the game we introduced in Section 2, $`\sigma (f)`$ is the set of all $`aM`$ so that for any $`\epsilon >0`$ Player II has a winning strategy in the $`𝒜^\epsilon `$-game, where
$$𝒜^\epsilon =\{(y_i)_{i=1}^nS_X^n:\rho (a,f(y_1,\mathrm{}y_n))>\epsilon \}$$
(which means that Player II is able to get $`f(y_1,\mathrm{}y_n)`$ arbitrarily close to $`a`$).
As mentioned in \[Mi\] one can also define the spectrum relative to any filtration $`𝒮\mathrm{cof}(X)`$, meaning that $`𝒮`$ has the property that if $`X,Y𝒮`$ there is a $`Z𝒮`$ for which $`XYZ`$. The spectrum of $`f`$ relative to $`S`$ is the set $`\sigma (f,𝒮)`$ of all $`aM`$ for which
(36) $`\epsilon >0Y_1𝒮y_1S_{Y_1}Y_2𝒮y_2S_{Y_2}\mathrm{}Y_n𝒮y_nS_{Y_n}\text{ with}`$
$`\rho (f(y_1,y_2,\mathrm{}y_n),a)<\epsilon .`$
It is obvious that $`\sigma (f,𝒮)\sigma (f,\stackrel{~}{𝒮})`$ whenever $`\stackrel{~}{𝒮}𝒮`$. In particular it follows that $`\sigma (f)\sigma (f,𝒮)`$ for any filtration $`𝒮`$.
If $`X`$ is a subspace of a space $`Z`$ with FDD $`(E_i)`$ we can consider the filtration $`𝒮=\{X_{i=n}^{\mathrm{}}F_i:nN\}`$ and we write $`\sigma (f,(F_i))=\sigma (f,𝒮)`$.
On one hand the unrelativized spectrum $`\sigma (f)`$ seems to be the right concept to study geometric and structural properties of $`X`$, since it is “coordinate free”. On the other hand spectra with respect to an FDD is combinatorically easier to use and understand.
But from Theorem 3.3 we deduce that $`\sigma (f)`$ is equal to the spectrum with respect to a certain FDD (of some super space).
###### Proposition 5.1.
Let $`f:S_X^nM`$ be Lipschitz. Then
(37)
$$\sigma (f)=\{C:C\text{ is a closed subset of }M\text{ and }(\text{W}_I(f^1(C)))\}.$$
Moreover for any $`\epsilon >0`$, $`(\text{W}_I(f^1(C))_\epsilon )`$.
Furthermore $`X`$ can be embedded into a space $`Z`$ with FDD $`(F_i)`$ so that for every $`\epsilon >0`$ there is a $`\delta >0`$ and an $`M_0`$ with the following property.
Whenever $`M_0<M_1<M_2<\mathrm{}M_n`$ and $`(x_i)_{i=1}^nS_X`$ satisfies
$$d(x_i,S_{_{j=1+M_{i1}}^{M_i1}F_j}X)<\delta \text{ for }i=1,\mathrm{},n$$
then $`\rho (f(x_1,x_2,\mathrm{}x_n),\sigma (f))<\epsilon `$.
In the case that $`X^{}`$ is separable, $`\sigma (f)`$ is the minimal closed subset of $`M`$ so that for any $`\epsilon >0`$ any weakly null tree in $`S_X`$ of length $`n`$ has a branch $`(x_1,\mathrm{}x_n)`$ so that $`\rho (f(x_1,\mathrm{}x_n),\sigma (f))<\epsilon `$.
###### Proof.
Let $`𝒞`$ denote the set of all closed subsets of $`M`$ for which (W$`{}_{I}{}^{}(f^1(C))`$ holds. For $`aM`$ we denote the $`\epsilon `$-neighborhood by $`U_\epsilon (a)`$ and observe the following equivalences
$`a`$ $`\sigma (f)`$
$`\epsilon >0Y_1\mathrm{cof}(X)y_1S_{Y_1}\mathrm{}Y_n\mathrm{cof}(X)y_nS_{Y_n}`$
$`\rho (f(y_1,..y_n),a)>\epsilon `$
$`\epsilon >0(\text{W}_I(f^1(MU_\epsilon (a)))`$
$`C𝒞,aC.`$
Thus $`\sigma (f)=\{C:C𝒞\}`$. If $`\eta >0`$ then $`M(\sigma (f))_\eta `$ is compact and is contained in the open covering $`_{C𝒞}MC`$. Thus there exists a finite $`\stackrel{~}{𝒞}𝒞`$ so that $`M(\sigma (f))_\eta _{C\stackrel{~}{𝒞}}MC`$ and thus $`(\sigma (f))_\eta _{C\stackrel{~}{𝒞}}C`$ which implies by Proposition 2.1 that Player I has a winning strategy for $`f^1((\sigma (f))_\eta )`$. By the uniform continuity of $`f`$, $`\eta `$ can be chosen small enough so that $`f^1((\sigma (f))_\eta )_n`$ contained in a given neighborhood of $`f^1(\sigma (f))`$ which finishes the proof of the first part. The remainder of the proposition follows easily from Theorem 3.3. ∎
A special example of spectra was considered by Milman and Tomczak \[MT\], the asymptotic structure of $`X`$. A finite dimensional space $`E`$ together with a normalized monoton basis $`(e_i)_1^n`$ is called an element of the $`n^{th}`$-asymptotic structure of $`X`$ and we write $`(E,(e_i)_{i=1}^n)\{X\}_n`$ if
(38) $`\epsilon >0Y_1\mathrm{cof}(X)y_1S_X\mathrm{}Y_n\mathrm{cof}(X)y_nS_X`$
$`\mathrm{dist}_b((y_i)_{i=1}^n,(e_i)_{i=1}^n)<1+\epsilon `$
where $`\mathrm{dist}_b`$ denotes the basis distance, i.e., if $`(e_i)_{i=1}^n`$ and $`(f_i)_{i=1}^n`$ are two bases of $`E`$ and $`F`$ respectively then $`\mathrm{dist}_b((e_i)_{i=1}^n,(f_i)_{i=1}^n)`$ is defined to be $`TT^1`$ where $`T:EF`$ is given by $`T(e_i)=f_i`$, for $`i=1,\mathrm{}n`$. Note that the space $`(M_n,\mathrm{log}\mathrm{dist}_b)`$ of all normalized bases of length $`n`$ and basis constant not exceeding a fixed constant is a compact metric space.
Therefore we deduce from Proposition 5.1 and the usual diagonalization argument the following Corollary (cf. \[KOS\]).
###### Corollary 5.2.
$`X`$ can be embedded into a space $`Z`$ with FDD $`(F_i)`$ so that for every $`k`$ it follows that:
Whenever $`k=M_0<M_1<M_2<\mathrm{}M_k`$ and
$$x_iS_{_{j=1+M_{i1}}^{M_i1}F_j}X\text{ for }i=1,2\mathrm{}k,$$
then $`\mathrm{dist}_b((x_i)_{i=1}^k,\{X\}_k)<1+\epsilon `$.
In the case that $`X^{}`$ is separable, $`\{X\}_k`$ is the minimal closed subset of $`M_k`$ so that for any $`\epsilon >0`$ any weakly null tree in $`S_X`$ of length $`n`$ has a branch $`(x_1,\mathrm{}x_k)`$ so that $`\mathrm{dist}_b((x_i)_{i=1}^k,\{X\}_k)<1+\epsilon `$.
An interesting case is when the asymptotic structure of $`X`$ is as small as possible.
###### Theorem 5.3.
Let $`X`$ be a separable reflexive Banach space with $`|\{X\}_2|=1`$. Then there exists $`p(1,\mathrm{})`$ so that $`X`$ embeds into the $`\mathrm{}_p`$-sum of finite dimensional spaces. Moreover for all $`\epsilon >0`$ there exists a finite codimensional subspace $`X_0`$ of $`X`$ which $`1+\epsilon `$-embeds into the $`\mathrm{}_p`$-sum of finite dimensional spaces.
###### Proof.
Since there exists $`1p\mathrm{}`$ so that the unit vector basis of $`\mathrm{}_p^2`$ is in $`\{X\}_2`$ (see \[MMT\]) we have that $`\{X\}_2`$ must be this unit vector basis. In turn this condition (see \[MMT\] or \[KOS\]) implies that $`X`$ contains an isomorph of $`\mathrm{}_p`$ ($`c_0`$ if $`p=\mathrm{}`$) and so $`1<p<\mathrm{}`$.
Let $`XZ`$, a reflexive space with an FDD $`(E_n)`$. The condition on $`\{X\}_2`$ yields that for all $`\epsilon >0`$ there exists $`n`$ so that if $`x_1S_X[E_i]_{i=n}^{\mathrm{}}`$ then there exists $`m`$ so that if $`x_2S_X[E_i]_{i=m}^{\mathrm{}}`$ then $`(x_i)_1^2`$ is $`1+\epsilon `$-equivalent to the unit vector basis of $`\mathrm{}_p^2`$. From this it follows that $`X`$ satisfies the hypothesis of Theorem 4.1 with $`C=1`$ and thus the theorem follows. ∎
The following problem remains open. We say $`X`$ is Asymptotic $`\mathrm{}_p`$ if there exists $`K<\mathrm{}`$ so that for all $`k`$ and all $`(x_i)_1^k\{X\}_k`$, $`(x_i)_1^k`$ is $`K`$-equivalent to the unit vector basis of $`\mathrm{}_p`$. An FDD $`(E_n)`$ for a space $`Z`$ is asymptotic $`\mathrm{}_p`$ if there exists $`K<\mathrm{}`$ so that for all $`k`$ if $`(x_i)_1^k`$ is a block sequence of $`(E_i)_k^{\mathrm{}}`$ in $`S_Z`$, then $`(x_i)_1^k`$ is $`K`$-equivalent to the unit vector basis of $`\mathrm{}_p`$.
###### Problem 5.4.
Let $`X`$ be a reflexive Asymptotic $`\mathrm{}_p`$ space for some $`1<p<\mathrm{}`$. Does $`X`$ embed into a space $`Z`$ with an asymptotic $`\mathrm{}_p`$ FDD? |
warning/0002/astro-ph0002357.html | ar5iv | text | # K-Band Spectroscopy of an Obscured Massive Stellar Cluster in the Antennae Galaxies (NGC 4038/4039) with NIRSPECData presented herein were obtained at the W.M. Keck Observatory, which is operated as a scientific partnership among the California Institute of Technology, the University of California and the National Aeronautics and Space Administration. The Observatory was made possible by the generous financial support of the W.M. Keck Foundation.
## 1. Introduction
The Antennae (NGC 4038/4039) are a pair of disk galaxies in an early stage of merging which contain numerous massive super star clusters (SSCs) along their spiral arms and around their interaction region (Whitmore & Schweizer, 1995; Whitmore et al., 1999). The molecular gas distribution peaks at both nuclei and in the overlap region (Stanford et al., 1990), but the gas is not yet undergoing a global starburst typical of more advanced mergers (Nikola et al., 1998). Star formation in starbursts appears to occur preferentially in SSCs. We chose to observe the Antennae because their proximity permits an unusually detailed view of the first generation of merger-induced SSCs and their influence on the surrounding interstellar medium.
The Infrared Space Observatory (ISO) $`1218\mu `$m image showed that the hot dust distribution is similar to that of the gas, but peaks at an otherwise inconspicuous point on the southern edge of the overlap region (Mirabel et al., 1998). This powerful starburst knot is also a flat-spectrum radio continuum source (Hummel & van der Hulst, 1986) and may be associated with an X-ray source (Fabbiano et al., 1997). We imaged the region around this knot, and discovered a bright compact star cluster coincident with the mid-IR peak. We obtained moderate-resolution (R $`1900`$) K-band spectra of both the obscured cluster and the NGC 4039 nucleus.
## 2. Observations & Data Reduction
NIRSPEC is a new facility infrared ($`0.955.6\mu `$m) spectrometer for the Keck-II telescope, commissioned during April through July, 1999 (McLean et al., 1998). It has a cross-dispersed cryogenic echelle with R $`25,000`$, and a low resolution mode with R $`2000`$. The spectrometer detector is a 1024 $`\times `$ 1024 InSb ALADDIN focal plane array, and the IR slit-viewing camera detector is a 256 $`\times `$ 256 HgCdTe PICNIC array.
We observed the Antennae with NIRSPEC during the June 1999 commissioning run. Slit-viewing camera (SCAM) images at 2 $`\mu `$m reveal that the mid-IR ISO peak is a bright (K $`=14.6`$) compact star cluster located 20<sup>′′</sup>.4 east and 4<sup>′′</sup>.7 north of the K-band nucleus. This cluster is associated with a faint (V $`=23.5`$) red (V$``$I $`=2.9`$) source (# 80 in Whitmore & Schweizer 1995) visible with Space Telescope (Whitmore & Zhang, private communication). We obtained low resolution (R $`1900`$) $`\lambda 2.032.45\mu `$m spectra through a $`0{}_{}{}^{\prime \prime }.57\times 42^{\prime \prime }`$ slit at PA=77$`^{^{}}`$ located on the obscured star cluster and the nucleus of NGC 4039. The total integration time on source was 2100 s.
We dark-subtracted, mean-sky-subtracted, flat-fielded, and corrected two-dimensional spectra for bad pixels and cosmic rays before rectifying the curved order onto a grid in which wavelength and slit position are perpendicular. We then corrected for residual sky emission and divided by a B1.5 standard star spectrum to correct for atmospheric absorption. The object spectra were extracted using a Gaussian weighting function matched to their strong continuua collapsed in wavelength (intrinsic FWHM = 0<sup>′′</sup>.84 for cluster, 0<sup>′′</sup>.99 for nucleus)<sup>1</sup><sup>1</sup>1These widths are greater than those measured from the SCAM images, $`0.{}_{}{}^{\prime \prime }69`$ and $`0{}_{}{}^{\prime \prime }.83`$ (intrinsic), due to the extended line contribution and rectification errors of order $``$ 1 pixel at the chip edges., and then an aperture correction was applied to recover the full flux in the continuua. Thus we neglected more-extended H<sub>2</sub> emission, which has maximum FWHM $`1{}_{}{}^{\prime \prime }.7`$ in the cluster and $`1{}_{}{}^{\prime \prime }.2`$ in the nucleus. We obtained a flux scale by requiring the 2.2 $`\mu `$m star flux to equal that corresponding to its K magnitude. Reduced spectra are shown in Figures 1 and 2.
Figure 1. - NIRSPEC spectrum of the obscured star cluster shows nebular and fluorescent H<sub>2</sub> emission with a continuum rising toward the red. Scaled sky counts are plotted at 0.1 mJy. $`\omega `$-shaped curve represents an atmospheric CO<sub>2</sub> band at 2.05 $`\mu `$m.
## 3. Massive Star Cluster
The cluster spectrum is characterized by strong emission lines<sup>2</sup><sup>2</sup>2A table of measured line fluxes is available electronically from http://astro.berkeley.edu/$``$agilbert/antennae. and a continuum (detected with SNR $`15`$) dominated by the light of hot, blue stars and dust. The nebular emission lines are slightly more extended than the continuum, and the H<sub>2</sub> emission is even more extended. This suggests a picture in which hot stars and dust are embedded in a giant compact H ii region surrounded by clumpy (see §3.2) clouds of obscuring gas and dust whose surfaces are ionized and photodissociated by FUV photons escaping from the star cluster.
For a distance to the Antennae of 19 Mpc (H<sub>0</sub>=75 km s<sup>-1</sup> Mpc<sup>-1</sup>, 1<sup>′′</sup>= 93 pc) (Whitmore et al., 1999), we find that the cluster has M$`{}_{\mathrm{K}}{}^{}=16.8`$. We estimate the screen extinction to the cluster by assuming a range of (V$``$K)$`{}_{0}{}^{}01`$ as expected from Starburst99 models (Leitherer et al., 1999), and that A$`{}_{\mathrm{K}}{}^{}=0.11`$ A<sub>V</sub> (Rieke & Lebofsky, 1985). We find A$`{}_{\mathrm{V}}{}^{}=910`$ mag, which implies M$`{}_{\mathrm{K}}{}^{}(0)=17.9`$, adopting A$`{}_{\mathrm{K}}{}^{}=1.1`$ (which is confirmed by our analysis of the H II recombination lines in §3.1). We can use the intrinsic brightness along with the Lyman continuum flux inferred from the de-reddened Br$`\gamma `$ flux ($`3.1\times 10^{14}`$ erg s<sup>-1</sup> cm<sup>-2</sup>), Q(H$`{}_{}{}^{+})_0=1.0\times 10^{53}`$ photons s<sup>-1</sup>, to constrain the cluster mass and age. Using instantaneous Starburst99 models we find a total mass of $`7\times 10^6`$ M (with $`2600`$ O stars) for a Salpeter IMF extending from 1 to 100 M, and an age of $`4`$ Myr. This age is consistent with the lack of photospheric CO and metal absorption lines from red supergiants and other cool giants, which would begin to contribute significantly to the 2 $`\mu `$m light at an age of $`7`$ Myr (Leitherer et al., 1999). The cluster’s density is then about 115 M pc<sup>-3</sup> for stars of 0.1$``$100 M within a half-light radius of $`32`$ pc. This density is 30 times less than that of the LMC SSC, R136 (within a radius of 1.7 pc, assuming a Salpeter proportion of low-mass stars) (Hunter et al., 1995). Thus the Antennae cluster may be a complex of clusters rather than one massive cluster.
Figure 2. - NIRSPEC spectrum of NGC 4039 nucleus shows extended collisionally excited H<sub>2</sub> emission and a strong stellar continuum marked by photospheric absorption. No Br$`\gamma `$ is present. Scaled sky counts are shown at 1.5 mJy.
### 3.1. Nebular Emission
The cluster spectrum features a variety of nebular lines that reveal information about the conditions in the ionized gas around the cluster, which in turn allows us to constrain the effective temperature of the ionizing stars.
We detected H I Pfund series lines from Pf 19 to Pf 38, and display their fluxes relative to that of Br$`\gamma `$ in Figure 3. The filled symbols give fluxes for the blends Pf 28+H<sub>2</sub> 2$``$1 S(0) and Pf 29+\[Fe III\]. They fall well above the other points, which follow closely the theoretical expectation for intensities relative to Br$`\gamma `$ (solid curve) with no reddening applied, for a gas with n$`{}_{\mathrm{e}}{}^{}=10^4`$ cm<sup>-3</sup> and T$`{}_{\mathrm{e}}{}^{}=10^4`$ K (Hummer & Storey, 1987). Excluding the two known blends, the best-fit foreground screen extinction is A$`{}_{\mathrm{K}}{}^{}=1.1\pm 0.3`$ mag (dashed curve), assuming the extinction law of Landini et al. (1984) and evaluated at 2.2 $`\mu `$m. We consider this an upper limit on A<sub>K</sub> because a close look at the spectrum shows that the points above the dashed line in Figure 3 for Pf 22$``$24 at 2.404, 2.393, and 2.383 $`\mu `$m may also be blended or contaminated by sky emission, implying a lower A<sub>K</sub> and a much better fit to the theory. Hence the majority of the extinction to the cluster is bypassed by observing it in K band.
Figure 3. - Pfund line fluxes relative to Br$`\gamma `$ flux ($`1.05\times 10^{14}`$ erg s<sup>-1</sup>cm<sup>-2</sup>). Solid curve is unextincted theoretical curve for n$`{}_{\mathrm{e}}{}^{}=10^4`$ cm<sup>-3</sup>, T$`{}_{\mathrm{e}}{}^{}=10^4`$ K (Hummer & Storey 1987). Filled symbols represent lines that are known blends, and the dashed curve shows theoretical fluxes with the best-fit extinction A$`{}_{\mathrm{K}}{}^{}=1.1`$ mag.
The lack of a strong Pfund discontinuity at 2.28 $`\mu `$m indicates that nebular free-free and bound-free continuum is diluted by starlight and dust emission (signaled by the rising continuum toward longer $`\lambda `$) in the cluster.
The ratios of \[Fe III\] lines are nebular density diagnostics; Table 1 presents observed ratios and theoretical predictions of Keenan et al. (1992) for emission from a collisionally excited 10<sup>4</sup> K gas, as tabulated by Luhman et al. (1998). The ratios of \[Fe III\] 2.146 $`\mu `$m and \[Fe III\] 2.243 $`\mu `$m to \[Fe III\] 2.218 $`\mu `$m are consistent with n$`{}_{\mathrm{e}}{}^{}=10^{3.5}10^4`$ cm<sup>-3</sup>. The ratio \[Fe III\] 2.348 $`\mu `$m/\[Fe III\] 2.218 $`\mu `$m is 20% higher than its theoretical value, which is roughly constant over all of parameter space (Keenan et al., 1992), but \[Fe III\] 2.348 $`\mu `$m is blended with Pf 29 and subject to measurement errors that are larger than the difference in extinctions in question (see Figure 3). Even the minimum value we infer for this ratio, with $`A_\mathrm{K}=0`$, is significantly greater than the model prediction. High values of \[Fe III\] 2.348 $`\mu `$m/\[Fe III\] 2.218 $`\mu `$m were also found by Luhman et al. (1998) in Orion. This discrepancy may be due to blending with another unknown line, or to theoretical error; ratios from the latest calculations have an average deviation from data of 10% (Keenan et al., 1992).
He I line ratios can in principal be used to infer nebular temperature T<sub>e</sub>, and are fairly insensitive to n<sub>e</sub>. However, of the three lines we detected, two are not suitable for such an analysis: the He I 2.1615+2.1624 $`\mu `$m blend falls on the wing of strong Br $`\gamma `$ so its flux has a large (50%) measurement error, and the strong He I 2.0589 $`\mu `$m line is subject to radiative transfer and density effects.
The He I 2.0589 $`\mu `$m/Br $`\gamma `$ ratio is an indicator of the T<sub>eff</sub> of hot stars in H II regions (Doyon et al., 1992), although it is sensitive to nebular conditions such as the relative volumes and ionization fractions of He<sup>+</sup> and H<sup>+</sup>, geometry, density, dustiness, etc. (Shields, 1993). Doherty et al. (1995) studied H and He excitation in a sample of starburst galaxies and H II regions. For starbursts they found evidence for high-T<sub>eff</sub>, low-n<sub>e</sub> ($`10^2`$ cm<sup>-2</sup>) ionized gas from He I 2.0589 $`\mu `$m/Br $`\gamma `$ ratios of 0.22 to 0.64. This is consistent with giant extended H II regions expected to dominate the emission-line spectra of typical starbursts. The ultra-compact H II regions were characterized by higher ratios (0.8$``$0.9) and higher densities, $``$ 10<sup>4</sup> cm<sup>-3</sup>. The cluster has a flux ratio of 0.70, a value between the two object classes of Doherty et al. (1995). Assuming the line emission is purely nebular, this ratio is consistent with a high-density (10<sup>4</sup> cm<sup>-3</sup>) model of Shields (1993), and implies T$`{}_{\mathrm{eff}}{}^{}39,000`$ K for the assumed model parameters. This temperature is similar to that derived by Kunze et al. (1996), $`44,000`$ K, from mid-IR SWS line observations in a large aperture on the overlap region of the Antennae.
| TABLE 1 | | | | | |
| --- | --- | --- | --- | --- | --- |
| Cluster \[Fe III\] Line Ratios<sup>a</sup><sup>a</sup>footnotemark: | | | | | |
| | Rest | Observed | Model Ratio<sup>c</sup><sup>c</sup>footnotemark: | | |
| Transition | $`\lambda `$($`\mu `$m)<sup>b</sup><sup>b</sup>footnotemark: | Ratio | $`10^3`$ | $`10^4`$ | $`10^5`$ |
| <sup>3</sup>G$`{}_{3}{}^{}_{}^{3}`$H<sub>4</sub> | 2.1457 | 0.14$`\pm `$0.02 | 0.10 | 0.17 | 0.34 |
| <sup>3</sup>G$`{}_{5}{}^{}_{}^{3}`$H<sub>6</sub> | 2.2183 | 1.00 | 1.00 | 1.00 | 1.00 |
| <sup>3</sup>G$`{}_{4}{}^{}_{}^{3}`$H<sub>4</sub> | 2.2427 | 0.28$`\pm `$0.02 | 0.26 | 0.29 | 0.38 |
| <sup>3</sup>G$`{}_{5}{}^{}_{}^{3}`$H<sub>5</sub> | 2.3485 | 0.80$`\pm `$0.03<sup>d</sup><sup>d</sup>footnotemark: | 0.66 | 0.66 | 0.66 |
<sup>a</sup> Ratios are dereddened fluxes relative to \[Fe iii\] 2.2183 $`\mu `$m, for which the dereddened flux was 9.11$`\times `$10<sup>-16</sup> ergs s<sup>-1</sup> cm<sup>-2</sup>.
<sup>b</sup> Sugar & Corliss (1985).
<sup>c</sup> Models for T<sub>e</sub>=10<sup>4</sup>K, values of n<sub>e</sub> in cm<sup>-3</sup> (Keenan et al., 1992).
<sup>d</sup> Flux determined by subtracting Pf 29 contribution obtained for the best-fit Landini extinction curve with A$`{}_{\mathrm{K}}{}^{}=1.1`$ mag.
The cluster has properties more like those of a compact H II region than a diffuse one. It appears to be a young, hot, high-density H II region, one of the first to form in this part of the Antennae interaction region (see Habing & Israel 1979 for a review of compact H II regions).
### 3.2. Molecular Emission
The spectrum shows evidence for almost pure UV fluorescence excited by FUV radiation from the O & B stars; the strong, vibrationally excited 1$``$0, 2$``$1 & 3$``$2 H<sub>2</sub> emission has T$`{}_{\mathrm{vib}}{}^{}6000`$ K and T$`{}_{\mathrm{rot}}{}^{}970`$, 1600, and 1800 K, respectively, and weak higher-v (6$``$4, 8$``$6, 9$``$7) transitions are present as well. The H<sub>2</sub> lines are extended over $`200`$ pc, about twice the extent of the continuum and nebular line emission, so a significant fraction of the FUV (912$``$1108 Å) light escapes from the cluster to heat and photodissociate the local molecular ISM.
We obtained the photodissociation region (PDR) models of Draine & Bertoldi (1996) and compared them with our data by calculating reduced $`\chi _\nu ^2`$. Models with high densities (n$`{}_{\mathrm{H}}{}^{}=10^5`$ cm<sup>-3</sup>), moderately warm temperatures (T $`=500`$ to 1500 K at the cloud surface), and high FUV fields (G$`{}_{0}{}^{}=10^310^5`$ times the mean interstellar field) can reasonably fit the data. Figure 4 shows $`\chi _\nu ^2`$ contours for all models projected onto the n$`{}_{\mathrm{H}}{}^{}`$G<sub>0</sub> plane. The best-fit Draine & Bertoldi model is n2023b, which has n<sub>H</sub> = 10<sup>5</sup> cm<sup>-3</sup>, T = 900 K, and G$`{}_{0}{}^{}=5000`$. We fit 22 H<sub>2</sub> lines, excluding 3$``$2 S(2) 2.287 $`\mu `$m because it appears to be blended with a strong unidentified nebular line at 2.286 $`\mu `$m found in higher-resolution spectra of planetary nebulae (Smith et al., 1981). The weak high-v transitions are all under-predicted by this model, and appear to come from lower-density gas (n$`{}_{\mathrm{H}}{}^{}10^310^4`$ cm<sup>-3</sup>) exposed to a weaker FUV field (G$`{}_{0}{}^{}10^210^3`$).
The ortho/para ratio of excited H<sub>2</sub> determined from the relative column densities in v=1, J=3 and J=2 inferred from 1$``$0 S(1) and S(0) lines is 1.62$`\pm `$0.07. This is consistent with the ground state v=0, J=1 and J=0 H<sub>2</sub> being in LTE with ortho/para ratio of 3 if the FUV absorption lines populating the non-LTE excited states are optically thick (Sternberg & Neufeld, 1999). Indeed, the best-fit PDR models have temperatures that are comparable with T<sub>rot</sub> in the lowest excited states, as well as with the warm gas kinetic temperature in the Galactic PDR M16, T = 930$`\pm `$50 K, measured by Levenson et al. (1999).
If the extent of the H<sub>2</sub> emission indicates that the mean-free path of a FUV photon is $`200`$ pc, then $``$n$`{}_{\mathrm{H}}{}^{}`$ = 3 cm<sup>-3</sup> for a Galactic gas-to-dust ratio, while in the PDR(s) n<sub>H</sub> = 10$`{}_{}{}^{4}10^6`$ cm<sup>-3</sup>. This implies that the molecular gas is extremely clumpy, which is consistent with the range of densities inferred from the detection of anomalously strong v = 8$``$6 H<sub>2</sub> emission.
## 4. NGC 4039 Nucleus
The spectrum of the nucleus of NGC 4039 is marked by strong stellar continuum and bright, extended H<sub>2</sub> emission. Strong photospheric Mg I, Na I, Ca I absorption and CO $`\mathrm{\Delta }v=2`$ bands indicate that the continuum is dominated by old giants. The CO band head is stronger than that of a M2III, suggesting some contribution from red supergiants. The absence of Br$`\gamma `$ emission implies that star formation is currently extinct in the nucleus. Spatially extended, collisionally excited H<sub>2</sub> emission in the nucleus may be excited by SNR shocks from the last generation of nuclear star formation, or by merger-induced cloud collisions. We defer detailed analysis of the nuclear spectrum to a later paper.
## 5. Conclusions
The highest surface brightness mid-IR peak in the ISO map of the Antennae Galaxies is a massive ($`16\times `$10<sup>6</sup> M), obscured (A$`{}_{\mathrm{V}}{}^{}910`$), young (age $`4`$ Myr) star cluster with half-light radius $``$ 32 pc, whose strong FUV flux excites the surrounding molecular ISM on scales of up to 200 pc. The cluster spectrum is dominated by extended fluorescently excited H<sub>2</sub> emission from clumpy PDRs and nebular emission from compact H II regions. In contrast, the nearby nucleus of NGC 4039 has a strong stellar spectrum dominated by cool stars, where the only emission lines are due to shock-excited H<sub>2</sub>. These observations confirm the potential of near-infrared spectroscopy for exploration and discovery with the new generation of large ground-based telescopes. Our ongoing program of NIRSPEC observations promises to reveal a wealth of information on the nature of star formation in star clusters.
Figure 4. - Comparison of H<sub>2</sub> line strengths with PDR models. Contours of $`\chi _\nu ^2`$ for 22 lines projected onto n$`{}_{\mathrm{H}}{}^{}`$G<sub>0</sub> plane peak at n$`{}_{\mathrm{H}}{}^{}10^5`$ cm<sup>-3</sup> and G$`{}_{0}{}^{}5000`$. Model points (+) are for T<sub>0</sub> = 300 $``$ 2000 K. White + marks best-fit PDR model of Draine & Bertoldi, with T<sub>0</sub> = 900 K and $`\chi _\nu ^2=9.3`$. Contours are 50, 25, 20, 15, 12, 10.
We acknowledge the hard work of past and present members of the UCLA NIRSPEC team: M. Angliongto, O. Bendiksen, G. Brims, L. Buchholz, J. Canfield, K. Chin, J. Hare, F. Lacayanga, S. Larson, T. Liu, N. Magnone, G. Skulason, M. Spencer, J. Weiss and W. Wong. We thank Keck Director Chaffee and all the CARA staff involved in the commissioning and integration of NIRSPEC, particularly instrument specialist T. Bida. We especially thank Observing Assistants J. Aycock, G. Puniwai, C. Sorenson, R. Quick and W. Wack for their support. We also thank A. Sternberg for valuable discussions. We are grateful to R. Benjamin for providing us with He i emissivity data. AMG acknowledges support from a NASA GSRP grant. |
warning/0002/quant-ph0002078.html | ar5iv | text | # Quantum Tomography Via Group Theory
## Abstract
Amongst the multitude of state reconstruction techniques, the so–called “quantum tomography” seems to be the most fruitful. In this letter, I will start by developing the mathematical apparatus of quantum tomography and, later, I will explain how it can be applied to various quantum systems.
03.65.Bz
Quantum tomography, like all state reconstruction methods, is concerned with the problem of measuring the density matrix $`\rho `$ of a physical system. The first observations on this topic are dated 1957 , but the greatest advancements are a result of this decade’s work. Also the first experiments were performed during this decade . The word tomography stems from medicine, as a consequence of the resemblance between the main formula of the first quantum tomography method, i.e. homodyne tomography , and the inverse Radon transform used in the CAT . Nowadays, however, quantum tomography does not have much in common with its medical ancestor. Quantum tomography is a general term, referring to any state measurement procedure descending from an equation of the same form of equation (18). It is a versatile and mighty technique, as it can be applied to a great variety of systems and as it includes other methods as special cases.
Homodyne tomography became a particular case of quantum tomography, when group representation theory was employed . Even though the latter led to a breakthrough in quantum tomography, it was recognized not to be the most general approach. There were, in fact, tomographic formulas (formulas like (18)) that could not be ascribed to standard group representation theory. Hence, in order to include these cases in a general mathematical framework, I will introduce some conditions that comprise the definition of group representation. Let, then, $`𝒢`$ be a group, $``$ a separable Hilbert space, $`()`$ the algebra of linear bounded operators defined in all $``$, namely Banach algebra , and $`T`$ a linear mapping from $`𝒢`$ to $`()`$ (from now on, the word linear, when referring to an operator, will always be understood). If we can find an irriducible, unitary ray representation $`D`$ of $`𝒢`$ (a ray representation is such that $`D(g)D(h)=\mathrm{e}^{\mathrm{i}\xi _{gh}}D(gh)`$, with $`\xi `$ ) and six correspondences $`\alpha ,\beta :𝒢\times 𝒢`$ and $`f_{ij}:𝒢𝒢`$, with $`i,j=1,2`$, satisfying the equations
$`D(g)T(h)=\mathrm{e}^{\mathrm{i}\alpha _{gh}}T(f_{11}(g)hf_{12}(g)),`$ (1)
$`D(g)T^{}(h)=\mathrm{e}^{\mathrm{i}\beta _{gh}}T^{}(f_{21}(g)hf_{22}(g)),`$ (2)
for every $`g,h𝒢`$, then we will say that $`\{T(g)\}`$ is a tomographic set. If $`T`$ is an irriducible, unitary ray representation then $`\{T(g)\}`$ is a tomographic set (choose $`D=T,f_{11}(g)=g,\mathrm{}`$). The converse is not true in general, so the requirement of $`T`$ being a representation is more stringent. When $``$ is finite–dimensional, the hypothesis that $`\{T(g)\}`$ is a tomographic set is sufficient to derive (18), however the case of $`\mathrm{dim}()=\mathrm{}`$ needs a further condition to make sure that every expression converges and can be attributed a precise mathematical meaning. More explicitly, $`T`$ needs to fulfill the following inequality :
$`{\displaystyle \underset{g}{}}|u|T(g)|v|^2<\mathrm{}|u,|v.`$ (3)
Condition (3) could be quite a nuisance, since it must be checked for every couple of vectors $`|u,|v`$. Fortunately, (3) is equivalent to
$`|u,|v:{\displaystyle }_g|u|T(g)|v|^2<\mathrm{},`$ (4)
as I will demonstrate now. Let us start by admitting the validity of (4) for two vectors $`|w_1,|w_2`$ (eventually coincident) and consider the set $`V`$ of vectors $`|v`$ of the form
$`|v={\displaystyle \underset{g}{}}v_gD(g)|w_i,\mathrm{with}{\displaystyle \underset{g}{}}|v_g|<\mathrm{},`$ (5)
$`i`$ being $`1`$ or $`2`$. $`V`$ is a linear manifold of $``$ because $`\alpha |v_1+\beta |v_2`$ belongs to $`V`$ if $`|v_1,|v_2V`$ and $`\alpha ,\beta `$, since $`_g|\alpha v_{1}^{}{}_{g}{}^{}+\beta v_{2}^{}{}_{g}{}^{}||\alpha |_g|v_{1}^{}{}_{g}{}^{}|+|\beta |_g|v_{2}^{}{}_{g}{}^{}|<\mathrm{}`$. The application of any $`D(h)`$ to a vector $`|vV`$, yields a vector that still belongs to $`V`$, in fact $`D(h)|v=_gv_g\mathrm{e}^{\mathrm{i}\xi _{hg}}D(hg)|w_i=_gv_{h^1g}\mathrm{e}^{\mathrm{i}\xi _{h(h^1g)}}D(g)|w_i`$, with $`_g|v_{h^1g}\mathrm{e}^{\mathrm{i}\xi _{h(h^1g)}}|=_g|v_g|<\mathrm{}`$. The irriducibility of the operators $`D(g)`$ implies that $`V=`$ or, in other words, that every vector in $``$ can be written in the form (5). Hence for every $`|u,|v`$, with the help of Cauchy’s inequality, we obtain
$`{\displaystyle \underset{g}{}}|u|T(g)|v|^2=`$ (6)
$`={\displaystyle \underset{g}{}}\left|{\displaystyle \underset{g_1,g_2}{}}u_{g_1}^{}v_{g_1}w_1|D^{}(g_1)T^{}(g)D(g_2)|w_2\right|^2`$ (7)
$`\left({\displaystyle \underset{g_1}{}}|u_{g_1}|\right)^2\left({\displaystyle \underset{g_2}{}}|v_{g_2}|\right)^2{\displaystyle \underset{g}{}}|w_1|T(g)|w_2|^2<\mathrm{},`$ (8)
namely (4) $``$ (3) ((3) $``$ (4) is obvious).
To prove equation (18) I will work out another identity first:
Assertion 1
If $`A`$ is a trace–class operator on $``$ and $`\{T(g)\}`$ is a tomographic set and satisfies (4) (or (3)) then
$`\mathrm{Tr}A={\displaystyle \frac{1}{\stackrel{~}{k}}}{\displaystyle \underset{g}{}}T(g)AT^{}(g),`$ (9)
with $`\stackrel{~}{k}_g|\phi |T(g)|\psi |^2`$ indipendent of the choice of the normalized vectors $`|\phi ,|\psi `$.
Proof: Hypothesis (3) implies
$`{\displaystyle \underset{g}{}}|a|T(g)|bc|T^{}(g)|d|`$ (10)
$`\left[{\displaystyle \underset{g}{}}|a|T(g)|b|^2{\displaystyle \underset{g}{}}|d|T(g)|c|^2\right]^{\frac{1}{2}}<\mathrm{},`$ (11)
defining unambiguously $`_ga|T(g)|bc|T^{}(g)|d`$ for all $`|a`$, $`|b`$, $`|c`$, $`|d`$, in accordance with . Since the complete space $``$, as a consequence of Riesz–Fréchet representation theorem , is also weakly complete, we may infer, by virtue of (11), that the sequence of the partial sums $`_g^nT(g)|uv|T^{}(g)`$ is weakly convergent, as $`n\mathrm{}`$, for all $`|u,|v`$. The operator $`I_{uv}_gT(g)|uv|T^{}(g)`$ can then be defined as the weak limit of such sequence, as $`n\mathrm{}`$. Using equation (1) and its adjoint and rearranging the sum we immediately get $`D(h)I_{uv}D^{}(h)=I_{uv}`$, which, due to Shur’s first lemma (for a proof in the infinite–dimensional case refer to ), is equivalent to $`I_{uv}=k_{uv}I`$, where $`k_{uv}`$ depends on $`|u,|v`$ and $`I`$ is the identity in $``$. Analogously, equation (2) entails that $`\stackrel{~}{I}_{uv}_gT^{}(g)|uv|T(g)`$ is a multiple of the identity. Given a normalized $`|w`$, $`k_{ww}`$ may be easily evaluated, in fact $`k_{ww}=w|k_{ww}I|w=w|I_{ww}|w=_g|w|T(g)|w|^2,`$ where the series could be interchanged with the inner product because of the way $`I_{uv}`$ is defined. The constant $`k_{uv}`$ may be expressed in terms of $`k_{ww}`$ in the following manner (let $`|a,|b`$ be generic vectors): $`a|I_{uv}|b=v|\stackrel{~}{I}_{ba}|u=v|uw|\stackrel{~}{I}_{ba}|w=v|ua|I_{ww}|b,`$ which means that $`I_{uv}=k_{w,w}v|uI`$. The choice $`|u=|v=|\psi `$ normalized and the calculation of the mean value of the last equation on the normalized vector $`|\phi `$ produces $`_g|\phi |T(g)|\psi |^2=k_{w,w}`$, proving that $`_g|\phi |T(g)|\psi |^2`$ is indipendent of the vectors $`|\phi ,|\psi `$ (for as long as their norm is 1) and will therefore be indicated simply by $`\stackrel{~}{k}`$. Schmidt decomposition of a trace–class operator $`A`$, i.e. $`A=_ia_i|u_iv_i|,`$ where $`\{|u_i\}`$ e $`\{|v_i\}`$ are orthonormal sequences and $`_ia_i<\mathrm{}`$, $`a_i>0i,`$ helps showing that $`_gT(g)AT^{}(g)`$ is meaningful. It is indeed sufficient to check the absolute convergence of the expression
$`{\displaystyle \underset{g}{}}a|T(g)AT^{}(g)|b=`$ (12)
$`={\displaystyle \underset{g}{}}{\displaystyle \underset{i}{}}a_ia|T(g)|u_iv_i|T^{}(g)|b,`$ (13)
for all $`|a,|b`$, to insure the validity of the definition of $`_gT(g)AT^{}(g)`$ as the weak limit of $`_g^nT(g)AT^{}(g)`$, as $`n\mathrm{}`$. The inequality
$`{\displaystyle \underset{g}{}}|a|T(g)|u_iv_i|T^{}(g)|b|`$ (14)
$`\left[{\displaystyle \underset{g}{}}|a|T(g)|u_i|^2{\displaystyle \underset{g}{}}|b|T(g)|v_i|^2\right]^{\frac{1}{2}}=\stackrel{~}{k},`$ (15)
together with $`a_i>0`$ and $`_ia_i<\mathrm{}`$, guarantees that the sum of the absolute values of the terms in (13) is $`\stackrel{~}{k}_ia_i<\mathrm{}`$. Because of the absolute convergence we can also rearrange the order of the two sums, obtaining the assertion’s thesis:
$`{\displaystyle \frac{1}{\stackrel{~}{k}}}{\displaystyle \underset{g}{}}T(g)AT^{}(g)={\displaystyle \frac{1}{\stackrel{~}{k}}}{\displaystyle \underset{i}{}}{\displaystyle \underset{g}{}}a_iT(g)|u_iv_i|T^{}(g)=`$ (16)
$`={\displaystyle \underset{i}{}}a_iv_i|u_iI=\mathrm{Tr}A.`$ (17)
Now, finally, equation (18):
Assertion 2
The operator identity
$`A={\displaystyle \frac{1}{\stackrel{~}{k}}}{\displaystyle \underset{g}{}}\mathrm{Tr}\left[AT(g)\right]T^{}(g)`$ (18)
holds, when $`A,T`$ and $`\stackrel{~}{k}`$ are defined as in assertion 1.
Proof: Let $`O`$ be an invertible trace–class operator. Using (9) twice, it is straightforward to check that
$`{\displaystyle \underset{g}{}}\text{Tr}[AT^{}(g)]OT(g)={\displaystyle \underset{g}{}}\text{Tr}[T(g)O]T^{}(g)A.`$ (19)
Expanding the trace on the complete orthonormal sequence $`\{|\phi _i\}`$, with the help of equation (9) again, we may write
$`{\displaystyle \frac{1}{\stackrel{~}{k}}}{\displaystyle \underset{g}{}}\text{Tr}[T(g)O]\phi _i|T^{}(g)A|\phi _j=`$ (20)
$`={\displaystyle \underset{h}{}}\phi _h|O\text{Tr}\left[|\phi _h\phi _i|\right]A|\phi _j=\phi _i|OA|\phi _j.`$ (21)
Equations (19) and (21) give $`\frac{1}{\stackrel{~}{k}}_g\text{Tr}[AT^{}(g)]OT(g)=OA,`$ which is equivalent to (18) because of the invertibility of $`O`$.
Note is devoted to a brief comment on some technical features of assertions 1 and 2.
Before we start to examine physical cases, I will cast light on some aspects of (18). It is well–known that Hilbert–Schmidt operators form a Hilbert space, usually denoted as $`\sigma _c()`$, with inner product $`(A,B)_o\mathrm{Tr}[A^{},B]`$, and that the space $`\tau _c()`$ of trace–class operators is contained in $`\sigma _c()`$. If we define $`P(g)=\stackrel{~}{k}^{\frac{1}{2}}T^{}(g)`$, equation (18) can formally be rewritten as
$`A={\displaystyle \underset{g}{}}(P(g),A)_oP(g).`$ (22)
This equality is of simple interpretation: all vectors of $`\tau _c()`$ can be expanded in terms of a closure relation, resorting to elements $`P(g)`$ that are not necessarily in $`\sigma _c()`$ but that belong to a larger set, in the same way that occurs with the expansion of a vector in terms of generalized vectors. A comparison will make the situation even clearer. If we identify $`\sigma _c()`$ with $`L^2()`$ (the space of square–integrable functions on $``$), then $`\tau _c()`$ corresponds to the space $`𝒮()`$ (test functions on $``$ decreasing rapidly at infinity) and $`()`$ to the space of tempered distributions $`𝒮^{}()`$ dual to $`𝒮()`$. This analogy suggests that formula (22) is also valid for $`A\sigma _c()`$ if we define the inner product $`(P(g),A)_o`$ as the limit of $`(P(g),A_n)_o`$, with the sequence of trace–class operators $`\{A_n\}`$ converging to $`A`$. Naturally, not all closure relation arise from a group. It is possible to use conditions similar to (1) and (2) defined on sets that are not groups and still obtain (18). Or even more generally, there are cases of spectral decompositions in $`\sigma _c()`$ that do not satisfy anything like (1) and (2) (for example the eigenvectors of a self–adjoint operator from $`\sigma _c()`$ to $`\sigma _c()`$ give a closure relation in $`\sigma _c()`$). However, from an operative point of view, these more general approaches to formula (18) are not very useful. Groups are indeed simple objects to be dealing with and they produce quite a large number of interesting results. Moreover, closure relations that do not exhibit a group structure usually derive from a group. To clarify this point, we must keep in mind that nothing in (22) guarantees that $`\{P(g)\}`$ is a complete orthonormal system. We know it is complete, because of (22), but we cannot be sure that it is orthonormal with respect to the inner product $`(,)_o`$. And, in fact, generally it is not orthonormal, i.e. the set $`\{P(g)\}`$ is often overcomplete. If an orthonormal system is extracted from the original overcomplete set, the group properties may disappear and we would be left with a spectral decomposition that is not generated by a group, even if it originated in a group. We will encounter an example of this circumstance afterwards.
Equation (9) is a closure relation itself. Choosing $`A=|vv|`$, with $`|v`$ normalized and arbitrary, it states that $`\{P(g)|v\}`$ is a complete set in $``$. In other terms, the complete set $`\{P(g)|v\}`$ in $``$ corresponds to the complete set $`\{P(g)\}`$ in $`\sigma _c()`$. Note that this connection is not a consequence of the specific context ($`𝒢`$ being a group and $`T`$ satisfying (1) and (2)) in which we proved formulas (9) and (18).
So far, I assumed that $`𝒢`$ was discrete. Nonetheless, (9) and (18) apply to other situations. It is useful to recall that every unitary irriducible representation of a compact group is finite–dimensional (meaning that $``$ is finite–dimensional), whereas every unitary irriducible representation of a non–compact group is infinite–dimensional (with the exception of the trivial representation). Hence, for a finite group or a compact Lie group, the mathematical problem simplifies: $`\mathrm{dim}()<\mathrm{}`$, $`\tau _c()=\sigma _c()=()`$ and the convergence of $`\stackrel{~}{k}`$ is always granted. In particular, for a finite group, by tracing both members of (9), we immediately recognize that $`\stackrel{~}{k}=\frac{[𝒢]}{\mathrm{dim}()}`$, with $`[𝒢]`$ indicating the order of $`𝒢`$. For a compact Lie group, formulas (9) and (18) are obtained by substituting $`\frac{1}{[𝒢]}_g`$, appearing in (9) and (18) for a finite group, with $`_𝒢d\mu (g)`$, where $`\mathrm{d}\mu (g)`$ is Haar’s invariant measure for $`𝒢`$. Similarly, the formal substitution $`_g_𝒢d\mu (g)`$ allows to write (9) and (18) for a non–compact Lie group, with a discrete group as the starting point. A warning is necessary in this case, however. For non–compact groups $`_𝒢d\mu (g)`$ is not convergent; moreover, differently from the compact case, right and left invariance may correspond to two different $`\mathrm{d}\mu (g)`$ . Only when they coincide, i.e. only when $`𝒢`$ is unimodular, formulas (9) and (18) are applicable.
Equation (18) (and (9)) deserves its own self–existence as a pure mathematical result, but my initial goal was different. I was concerned with the physical problem of measuring the density matrix, and, thus far, (18) does not give us any clue on how to solve such problem. I will now show how (18) is, in actuality, much nearer to the solution than what may appear. If $``$ is the Hilbert space associated with the physical system under consideration, then the density matrix $`\rho `$ is an element of $`\tau _c()`$ and, consequently, can be written in place of $`A`$ in formula (18). Moreover, if each $`T(g)`$ is self–adjoint or is a function of a self–adjoint operator, then we can evaluate the trace in (18) over its complete set of eigenvectors $`\{|g,t\}`$. This operation yields an expression containing quantities of the form $`g,t|\rho |g,t`$, which can be interpreted as the probability that a measurement of $`T(g)`$ gives the eigenvalue $`t(g)`$ corresponding to $`|g,t`$ (the case of $`t(g)`$ degenerate could be treated analogously). These quantities are, in principle, experimentally accessible and will be indicated with $`p(g,t)`$. Formula (18) then becomes $`\rho =_{g,t}p(g,t)[\frac{1}{\stackrel{~}{k}}t(g)T^{}(g)]`$, where if the group is not discrete, or if $`T(g)`$ has continuous spectrum, sums must be replaced by integrals (obviously there could eventually be both sums and integrals). The observation that the eigenvectors $`|g,t`$ and $`|g^{},t`$ are not necessarily different even if $`gg^{}`$ suggests to divide $`𝒢`$ in classes $`𝒢_i`$, requiring the property that all $`T(g)`$ corresponding to the same class have the same eigenvectors. We would then write $`\rho =_{i,t}p(i,t)_{g𝒢_i}[\frac{1}{\stackrel{~}{k}}t(g)T^{}(g)]_{i,t}p(i,t)K(i,t)`$, where the operator $`K(i,t)`$ is usually called pattern function (again, if $`i`$ is a continuous index then $`_id\mu (i)`$). The formula
$`\rho ={\displaystyle \underset{i,t}{}}p(i,t)K(i,t)`$ (23)
is the essence of most state reconstruction methods. It states that measuring the probabilities $`p(i,t)`$ and calculating $`K(i,t)`$ is all that is needed in order to obtain $`\rho `$. The peculiarity of quantum tomography is that $`K(i,t)`$ does not need to be determined by solving an inverse problem, it is explicitly given by $`\frac{1}{\stackrel{~}{k}}_{g𝒢_i}t(g)T^{}(g)`$.
It took quite long to work out (23), but its generality will show its strength now that we turn to examples.
I will begin with the spin case, as it is the least complex. The (reduced) spin density matrix of one particle with spin $`S`$ (integer or half–integer) is defined in a Hilbert space $`_S`$, with $`\mathrm{dim}(_S)=2S+1`$. The compact group $`SU(2)`$ is particularly suited for the case of arbitrary $`S`$, since there exists an irriducible, unitary representation of $`SU(2)`$ in every finite–dimensional space. If $`SU(2)`$ is parametrized with $`(\vartheta ,\phi ,\psi )`$ belonging to $`[0,\pi ]\times [0,2\pi )\times [0,2\pi ]`$ and if $`\stackrel{}{n}`$ is defined as $`(\mathrm{cos}\phi \mathrm{sin}\vartheta ,\mathrm{sin}\phi \mathrm{sin}\vartheta ,\mathrm{cos}\vartheta )`$, then the operators $`R(\stackrel{}{n},\psi )=\mathrm{e}^{\mathrm{i}\psi \stackrel{}{S}\stackrel{}{n}}`$, where $`\stackrel{}{S}`$ is the spin operator , constitute an irriducible, unitary representation of $`SU(2)`$ and, consequently, a tomographic set in $`_S`$. Formula (23) then becomes
$`\rho ={\displaystyle _\Sigma }d\mathrm{\Omega }_\stackrel{}{n}{\displaystyle \underset{m=S}{\overset{S}{}}}p(\stackrel{}{n},m)K_S(\stackrel{}{n},m),`$ (24)
where $`\mathrm{d}\mathrm{\Omega }_\stackrel{}{n}`$ is the area element of the unitary spherical surface $`\Sigma `$, $`K_S(\stackrel{}{n},m)`$ is the pattern function given by $`(2S+1)_0^{2\pi }d\psi \frac{\mathrm{sin}^2\frac{\psi }{2}}{4\pi ^2}\mathrm{e}^{\mathrm{i}\psi (\stackrel{}{S}\stackrel{}{n}m)}`$, and $`p(\stackrel{}{n},m)`$ is the probability that $`m`$ is the result of a measurement of $`\stackrel{}{S}\stackrel{}{n}`$. The calculation of $`K_S(\stackrel{}{n},m)`$, the experimental apparatus needed to measure $`p(\stackrel{}{n},m)`$ and some numerical simulations can be found in . Because $`\mathrm{dim}(_S)<\mathrm{}`$, it is evident that the tomographic set $`\{R(\stackrel{}{n},\psi )\}`$ is overcomplete. It is then possible to choose a finite number of operators $`R(\stackrel{}{n},\psi )`$ and still obtain a closure relation in $`\sigma _c(_S)`$. As previously mentioned, the operation of extraction of a smaller complete set from the entire set $`\{R(\stackrel{}{n},\psi )\}`$ usually produces a class of operators not corresponding to a group. However, this need not always be the case. The dihedral and tetrahedral subgroups of $`SU(2)`$ can be represented with a finite number of operators $`R(\stackrel{}{n},\psi )`$, producing tomographic formulas for the cases of $`S=\frac{1}{2}`$ and $`S=1`$ (for explicit formulas and some further considerations refer to ). Unfortunately, there is only a finite number of finite subgroups of $`SU(2)`$; furthermore, the tomographic set associated with the tetrahedral group is still overcomplete. Therefore, waiving the group structure is a necessity for the obtainment of a complete orthonormal system in $`\sigma _c(_S)`$ for a generic S. This does not mean that we have to completely give up the use of the operators $`R(\stackrel{}{n},\psi )`$. Knowing in advance that $`\mathrm{dim}(\sigma _c())=(2S+1)^2`$, we can just choose $`(2S+1)^2`$ linearly independent operators $`R(\stackrel{}{n},\psi )`$, that do not need to form or correspond to a group, and apply Gram–Schmidt orthonormalization procedure: $`B_1\frac{R(\stackrel{}{n}_1,\psi _1)}{R(\stackrel{}{n}_1,\psi _1)_o}`$, $`B_2\frac{R(\stackrel{}{n}_2,\psi _2)(B_1,R(\stackrel{}{n}_2,\psi _2))_oB_1}{R(\stackrel{}{n}_2,\psi _2)(B_1,R(\stackrel{}{n}_2,\psi _2))_oB_1_o}`$,$`\mathrm{}`$, with $`O_o\sqrt{(O,O)_o},O\sigma _c()`$. By definition, $`\{B_i\},i=1,2,\mathrm{},(2S+1)^2,`$ is a basis in $`\sigma _c()`$. Nevertheless also $`\{B_i^{}\}`$ is a basis in $`\sigma _c()`$, therefore every $`A`$ in $`\sigma _c()`$ can be decomposed as $`A=_i(B_i^{},A)_oB_i^{}=_i\mathrm{Tr}[B_iA]B_i^{}`$. Because the operators $`B_i`$ are linear combinations of the operators $`R(\stackrel{}{n}_i,\psi _i)`$, with a little algebra we get $`A=_i\mathrm{Tr}\left[R(\stackrel{}{n}_i,\psi _i)A\right]_i`$, where the $`(2S+1)^2`$ operators $`_i`$ are linear combinations of the operators $`R^{}(\stackrel{}{n}_i,\psi _i)`$ as a result of the reorganization of the sum on $`i`$. One may check that $`(_i^{},R(\stackrel{}{n}_j,\psi _j))_o=\delta _{ij}`$ (Kronecker’s delta), which means that $`\{_i\}`$ is the dual basis of $`\{R(\stackrel{}{n}_i,\psi _i)\}`$ in $`\sigma _c()`$. Calculating the last trace on the eigenstates $`|\stackrel{}{n}_i,m`$, associated with the eigenvalue $`\mathrm{e}^{\mathrm{i}\psi _im}`$ of the operators $`R(\stackrel{}{n}_i,\psi _i)`$, with $`A=\rho `$, we attain a finite version of (24):
$`\rho ={\displaystyle \underset{i=1}{\overset{(2S+1)^2}{}}}{\displaystyle \underset{m=S}{\overset{S}{}}}p(\stackrel{}{n}_i,m)𝒦_S(i,m),`$ (25)
with $`𝒦_S(i,m)=\mathrm{e}^{\mathrm{i}\psi _im}_i`$ (the suffix $`S`$ in $`𝒦_S(i,m)`$ is a remainder of the dependence of the operators $`_i`$ on the dimension of $`_S`$). Equations having the same form of (24) or (25) are quite common in the problem of spin state reconstruction (for example ). Incidentally, we might also observe that the orthonormalization of $`(2S+1)^2`$ linearly independent projectors $`|\stackrel{}{n}_i,S\stackrel{}{n}_i,S|`$, $`i=1,2,\mathrm{},(2S+1)^2`$, instead of the operators $`R(\stackrel{}{n}_i,\psi _i)`$, leads to the same results obtained by Amiet and Weigert .
The Hilbert space $`_C`$ associated with a system of $`n`$ spins is given by the tensor product of the $`n`$ single–spin spaces. If we were to write (18) for the elements of $`\sigma _c(_C)`$, we would tempted to choose $`𝒢=SU(2)\times \mathrm{}\times SU(2)`$ ($`n`$ times) and $`D(g)=T(g)=_{i=1}^nR(\stackrel{}{n}_i,\psi _i)`$, with $`g𝒢`$ ($`\times `$ and $``$ denote respectively the direct product of groups and the tensor product of operators). This choice would actually give a valid closure relation in $`\sigma _c(_C)`$ (equation (18)), but the corresponding probabilities in equation (23) would require measurements on single components of the system, which are not always possible. This difficulty, at least for systems of spins $`\frac{1}{2}`$, can be overcome with a different approach, as illustrated in .
As a model for systems associated with infinite–dimensional Hilbert spaces, we can take the space $`_O`$ of one mode of the electromagnetic field. Although the problem is mathematically identical for other systems (for example, $`_O`$ is isomorphic to the space of a spinless, non–relativistic particle in one dimension), quantum optics gives the unique possibility of measuring the equivalent of linear combinations of position and momentum, namely the so–called quadratures $`X_\varphi \frac{1}{2}(a^{}\mathrm{e}^{\mathrm{i}\varphi }+a\mathrm{e}^{\mathrm{i}\varphi })`$, with $`\varphi `$ and $`a`$ and $`a^{}`$ indicating the annihilation and creation operators respectively . This opportunity can be exploited by choosing the non–compact Lie group of translations in the complex plane, with elements $`\alpha `$, as $`𝒢`$. Since the displacement operators $`𝖣(\alpha )=\mathrm{exp}(\alpha a^{}\alpha ^{}a)`$ form an irriducible, unitary ray representation of $`𝒢`$, such that $`_{}\mathrm{d}^2\alpha |0|𝖣(\alpha )|0|^2=\pi `$ ($`\mathrm{d}^2\alpha \mathrm{d}(\mathrm{Re}\alpha )\mathrm{d}(\mathrm{Im}\alpha )`$ is the invariant measure of the unimodular group of translations in the plane and $`|0`$ is the vacuum state), we can set $`D(\alpha )=T(\alpha )=𝖣(\alpha )`$. For the purpose of writing (23), we should, however, express the tomographic set in terms of quadratures. This can be achieved by parameterizing $`𝒢`$ with $`k`$ and $`\varphi [0,\pi )`$, related to $`\alpha `$ by the equation $`\alpha =\frac{\mathrm{i}}{2}k\mathrm{e}^{\mathrm{i}\varphi }`$, since $`T(\varphi ,k)T(\alpha (\varphi ,k))=\mathrm{e}^{\mathrm{i}kX_\varphi }`$. Equation (23) for this case then reads
$`\rho ={\displaystyle _0^\pi }d\varphi {\displaystyle _{\mathrm{}}^+\mathrm{}}dxp(\varphi ,x)𝖪(\varphi ,x),`$ (26)
where $`p(\varphi ,x)`$ is the probability that measuring $`X_\varphi `$ we get $`x`$ and $`𝖪(\varphi ,x)=\frac{1}{\pi }_{\mathrm{}}^+\mathrm{}dk\frac{|k|}{4}\mathrm{e}^{\mathrm{i}k(xX_\varphi )}`$ (note that $`\stackrel{~}{k}=\pi `$). Equation (26) is the fundamental formula of homodyne tomography, which I will not be discussing here, because the literature on it is already abundant ( and references therein).
Homodyne tomography is not the only technique that allows to reconstruct the state $`\rho `$ of one mode of the electromagnetic field. K. Banaszek and K. Wódkiewicz showed how $`\rho `$ could be determined by measuring, for every $`\alpha `$ and $`n`$, the probability $`p(\alpha ,n)`$ that $`n`$ is the number of photons in the state $`𝖣(\alpha )\rho 𝖣^{}(\alpha )`$ . The same result can be recovered with equation (23), if we maintain the choices $`𝒢=`${Translations in the complex plane} and $`D(\alpha )=𝖣(\alpha )`$, but we select the tomographic set $`\{𝖣^{}(\alpha )\mathrm{e}^{\mathrm{i}ya^{}a}𝖣(\alpha )\}`$, where $`y`$ can be any real number that is not a multiple of $`2\pi `$. It is not hard to check that (1) and (2) are satisfied and that $`_{}\mathrm{d}^2\alpha |0|𝖣^{}(\alpha )\mathrm{e}^{\mathrm{i}ya^{}a}𝖣(\alpha )|0|^2=\frac{\pi }{2(1\mathrm{cos}y)}`$. Then equation (23) is simply
$`\rho ={\displaystyle _{}}\mathrm{d}^2\alpha {\displaystyle \underset{n=0}{\overset{+\mathrm{}}{}}}p(\alpha ,n)\mathrm{K}_y(\alpha ,n),`$ (27)
with $`\mathrm{K}_y(\alpha ,n)=\frac{2(1\mathrm{cos}y)}{\pi }𝖣^{}(\alpha )\mathrm{e}^{\mathrm{i}y(na^{}a)}𝖣(\alpha )`$.
Differently from (24), neither (26) nor (27) can be simplified by extracting a complete subset from the entire tomographic set, since both $`P_1(\alpha )\pi ^{\frac{1}{2}}𝖣^{}(\alpha )`$ and $`P_2(\alpha )(\frac{\pi }{2(1\mathrm{cos}y)})^{\frac{1}{2}}𝖣^{}(\alpha )\mathrm{e}^{\mathrm{i}ya^{}a}𝖣(\alpha )`$ form orthonormal systems: $`(P_1(\alpha ),P_1(\alpha ^{}))_o=(P_2(\alpha ),P_2(\alpha ^{}))_o=\delta (\mathrm{Re}\alpha \mathrm{Re}\alpha ^{})\delta (\mathrm{Im}\alpha \mathrm{Im}\alpha ^{})`$, with $`\delta `$ indicating Dirac’s delta.
Some examples have shown how different state reconstruction problems can be treated as particular cases of a general method, which can be summarized in (18) and in the consequent (23). Although the theory developed in this letter is quiete comprehensive, there is still space for further generalizations. One possibility is to assume that the mappings $`D`$ and $`T`$ are defined on two distinct groups, changing (1) and (2) appropriately. Apparently, this generalization would produce other interesting tomographic formulas. For example, the state $`\rho `$ of one optical mode could be determined by measuring only the presence (or the absence) of photons in the state $`𝖣(\alpha )\rho 𝖣^{}(\alpha )`$, for every $`\alpha `$.
My gratitude goes to Prof. G. M. D’Ariano and L. Maccone of the Quantum Optics Group of Pavia. This work, without their collaboration, would not have existed. |
warning/0002/cond-mat0002327.html | ar5iv | text | # Spin Glass Ordering in Diluted Magnetic Semiconductors: a Monte Carlo Study
## Abstract
We study the temperature-dilution phase diagram of a site-diluted Heisenberg anti-ferromagnet on a fcc lattice, with and without the Dzyaloshinskii-Moriya anisotropic term, fixed to realistic microscopic parameters for $`IIB_{1x}Mn_xTe`$ ($`IIB`$$`=`$$`Cd`$, $`Hg`$, $`Zn`$). We show that the dipolar Dzyaloshinskii-Moriya anisotropy induces a finite-temperature phase transition to a spin glass phase, at dilutions larger than 80$`\%`$. The resulting probability distribution of the order parameter $`P(q)`$ is similar to the one found in the cubic lattice Edwards-Anderson Ising model. The critical exponents undergo large finite size corrections, but tend to values similar to the ones of the Edwards-Anderson-Ising model.
Although most theoretical investigations on spin glass (SG) systems have focused on models based on Ising spins, systems that have been investigated in the recent period, like the diluted magnetic semiconductors $`IIBMnTe`$ series, $`IIB=Cd,Hg,Zn`$, i.e. $`Cd_{1x}Mn_xTe`$ , $`Hg_{1x}Mn_xTe`$ and $`Zn_{1x}Mn_xTe`$ (and somehow even typical experimental samples like $`CuMn`$, $`AgMn`$, $`Eu_xSr_{1x}S`$), are closer in nature to continuous Heisenberg spins. Early computer simulations suggested that in three spatial dimensions neither systems with local interactions and Heisenberg or XY spins, nor systems with long-range RKKY interactions undergo a finite temperature SG phase transition. This fact could be potentially annoying from a phenomenological point of view, but it becomes acceptable after noticing that a small anisotropic interaction, neglected in the above calculations, could induce a finite temperature SG phase transition.
Still the situation is not crystal clear: the work of claimed that the most important anisotropic coupling in $`IIBMnTe`$ materials , the Dyalozhinskii-Moriya (DM) interaction, is not able to induce a SG phase. On the contrary a theoretical analysis of experimental data on the $`IIBMnTe`$ series was used to suggest the presence of SG ordering in $`3D`$ Heisenberg spin glasses (for finite temperature and no anisotropies): recent numerical simulations support the existence of a chiral phase transition in such systems. The role of the anisotropy was reconsidered in , where the Heisenberg spin Edwards-Anderson (HEA) model was considered with the addition of a random pseudo-dipolar interaction: clear signatures of a finite temperature SG phase were found. This result is not consistent with the one of (that was considered as being based on a realistic modelization of $`IIBMnTe`$, even if the direction of the DM vectors, see (1), was chosen at random, while they should be periodic along the lattice ). We remind at last that a diluted Ising anti-ferromagnet on a fcc lattice has been studied in , where the signature of a SG phase transition for low enough densities has been detected.
We have taken here the point of view of trying to be as realistic as possible, analyzing a model as close as possible to the experimental samples. We show that non-random DM terms (selecting a realistic value for the anisotropy) are able to induce a SG phase transition in a $`3D`$ Heisenberg spin glass on a fcc lattice. We analyze and discuss in detail the values of critical exponents: the experimental results for the $`IIB_{1x}Mn_xTe`$ materials are in good agreement with the most accurate calculations for the Edwards-Anderson model with Ising spins (IEA) on the cubic lattice ($`\nu =1.8\pm 0.2)`$, $`\eta =0.26\pm 0.04`$), but the numerical simulations of and of , yielded $`\nu 1.0`$. We will show that the numerical calculation of the critical exponents on the accessible lattice sizes suffer from serious finite-size corrections, and that a systematic analysis of the numerical data establishes clear trend towards values of $`\nu `$ larger than the ones found in , and close to the experimental values.
The site-diluted anti-ferromagnetic (AFM) Heisenberg model on the fcc lattice, with and without DM anisotropy, is a model for the $`IIB_{1x}Mn_xTe`$ series, where the $`Mn`$ atoms form an fcc lattice with localized (Heisenberg) spins interacting through short-range (super-exchange) AFM terms, while the magnetically inert $`IIB`$ atoms randomly replaces the $`Mn`$ over the lattice. An AFM interaction on the fcc lattice is frustrated, and gives rise to some interesting order-disorder phenomena , both with Heisenberg and Ising spins . The dilution disorder deletes some of the sites on the system, thus providing the random combination of frustrated and unfrustrated plaquettes, that is believed to be essential for SG ordering. The Hamiltonian of the system is
$$H=J\underset{𝒙,𝒚^{}}{}\left[𝑺_𝒙𝑺_𝒚+\frac{D}{J}𝑹_{𝒙𝒚}(𝑺_𝒙𝑺_𝒚)\right],$$
(1)
where the fields $`𝑺=(S_1,S_2,S_3),`$ $`S_1`$, $`S_2`$ and $`S_3`$ real with $`𝑺^2=1`$, represent the spin of the $`Mn`$ atoms, and $`J>0`$. A lattice site is randomly occupied by a spin with probability $`p`$. The sum labeled by $`𝒙,𝒚^{}`$ runs over the pairs of occupied nearest neighboring sites of the lattice. The unit-length vectors $`𝑹_{𝒙𝒚}`$ specify the DM anisotropy, and they verify $`𝑹_{𝒙𝒚}=𝑹_{𝒚𝒙}`$. Following we set $`𝑹_{(\frac{1}{\sqrt{2}},\frac{1}{\sqrt{2}},0)}=(\frac{1}{\sqrt{2}},\frac{1}{\sqrt{2}},0)`$, while the other five independent vectors are obtained using the three-fold rotation symmetries of the lattice. In Ref. the ratio $`D/J`$ has been estimated to be $`0.054`$ for $`Zn_{.77}Mn_{.33}Te`$ and $`Cd_{.77}Mn_{.33}Te`$, and is very mildly dependent on the composition of the sample: we have fixed it to $`0.06`$ for simplicity. As the local magnetic field acting on spin $`𝑺_𝒙`$ is (see (1))
$$h_𝒙=J\underset{𝒚,𝒙𝒚=\sqrt{2}}{}𝑺_𝒚+\frac{D}{J}(𝑺_𝒚𝑹_{𝒙,𝒚}),$$
(2)
it is easy to implement a heat-bath algorithm and the over-relaxed micro-canonical algorithm of . We have found that the combination of these two updates tremendously reduces the thermalization effort. In the production runs we have performed a full-lattice heat-bath sweep followed by $`19`$ over-relaxed updates, that will be referred in the following as an elementary Monte-Carlo step (EMCS). In order to define the observables, it is useful to consider a replica (i.e. a thermally independent system with the same set of occupied sites, that we denote $`\stackrel{~}{𝑺}(𝒙)`$). The measured observables can be most easily described in terms of the following three basic fields: the tensor field ($`\tau _{\alpha ,\beta }(𝒙)=S_𝒙^\alpha S_𝒙^\beta \frac{1}{3}\delta ^{\alpha ,\beta }`$, if the lattice site $`𝒙`$ is occupied, and zero otherwise), the tensorial overlap ($`O_{\alpha ,\beta }(𝒙)=S_𝒙^\alpha \stackrel{~}{S}_𝒙^\beta `$), and the scalar overlap field ($`q(𝒙)=𝑺_𝒙\stackrel{~}{𝑺}_𝒙`$).
The rationale for studying the tensor field $`\tau (𝒙)`$ is that previous studies of AFM diluted systems on the fcc lattice , showed that only for moderate dilutions the system ceases to develop AFM ordering. The tensorial magnetization is an ideal order parameter to check this possibility, since it will be non-vanishing for any conceivable type of AFM or helicoidal ordering. Moreover, it would also work on more sophisticated, yet trivial situations like the found in $`D=0`$, $`p=1`$ . The tensorial overlap is most adequate to study the isotropic ($`D=0`$) case, since in this situation the Hamiltonian posses a $`O(3)`$ global symmetry: when the anisotropy is switched on the symmetry reduces to $`Z_2`$, and the use of the scalar overlap becomes natural. For all three fields, one can define straightforwardly the corresponding susceptibility, Binder parameter and a finite lattice correlation length .
The model (1) without impurities ($`p=1.0,D=0.06J`$) undergoes a phase transition at $`T_\mathrm{c}(p=1)0.60J`$ from a paramagnetic phase to an AFM phase, as shown by the behavior of the tensorial magnetization. For larger dilutions a lower temperature value needs to be reached in order to exit from the paramagnetic phase: the critical line, $`T_\mathrm{c}(p)`$, will eventually reach zero temperature at the percolation threshold for the magnetic ions ($`p_\mathrm{c}0.2`$). The first question is for which dilution the system forgets its global AFM ordering. In order to answer this question we have performed slow annealings in $`60`$ samples (and its corresponding replicas) at dilutions $`p=1.0,0.9`$ and $`0.8`$, in lattices $`L=8`$ and $`16`$; at $`p=0.7`$ and $`p=0.6`$ (that will be shown to be in the SG compositional range), we have annealed $`700`$ samples. The results for the Binder cumulant of the tensor and scalar overlap fields at $`p=0.7`$ and $`p=0.9`$ are displayed in figure 1. At $`p=0.9`$ for both observables we find a low temperature, AFM ordered phase, since the tensorial magnetization is non-vanishing. There is a strong dip close to the phase transition point, which probably is very plausibly of first order. On the contrary for $`p=0.7`$ it is clear that the tensorial magnetization is no longer an appropriate order parameter. For $`p=0.8`$ (not shown in the plot), our results indicate a cross-over regime between the two situations. Therefore, the low temperature phase turns from AFM to SG at $`0.7<p_\mathrm{c}<0.8`$, similarly to what happens in the Ising case . Also for $`D=0,p<0.8`$ we do not find an AFM ordered phase.
In order to quantify how strong the effect of the anisotropy is, we compare the system at $`p=0.7`$ with $`D=0`$ and $`D=0.06J`$. In figure 2 we show the correlation length of the tensorial overlap in units of the lattice size. This operator should be zero in the paramagnetic phase, diverging in the SG phase, and at the critical point reaches a finite universal value. In the isotropic case (see zoom in upper part of figure 2) the crossing point of the $`L=12`$ and $`16`$ lattices is not clearly resolved, and their respective crossings with the $`L=8`$ curve shifts to lower temperature with growing lattice size. Moreover, the data (not shown in the plot) for the Binder Cumulant of the tensorial overlap rapidly grow at the crossing temperature of the correlation length, but then saturate at a value which decreases with the lattice size, without a crossing, similarly to the results shown in , where it has been shown that the chiral-glass phase appears precisely at the temperatures at which the Binder cumulant of the tensorial overlap grows.
On the other hand, with a $`6\%`$ DM anisotropy (see the lower frame of figure 2), we find a neat crossing of the correlation length (the tensorial overlap Binder cumulant has a marked dip, in contrast with the scalar overlap shown in the lower frame of figure 1). As the phase transition for the anisotropic system occurs at a temperature $`80\%`$ higher than the one close to the crossings of the $`D=0`$ case, that according to signal a real chiral-glass phase transition, the natural conclusion is that the DM anisotropy is not a smooth perturbation that reveals a hidden chiral-glass ordering.
In order to characterize more precisely the SG phase we have studied the distribution of the scalar overlap at $`p=0.6,D=0.06J`$ and $`T=J/4.50.78T_\mathrm{c}`$. At this temperature we have estimated the mean thermalization time in the $`L=16`$ lattice, by considering a logarithmic plot of the mean overlap susceptibility of $`64`$ samples, as a function of MC-time, starting from a random configuration, and we have found it to be of order $`250`$ EMCS. After that we have performed a run with $`800`$ samples, with $`L=8,12,16`$, performing $`8000`$ EMCS on each sample, and taking a measure every $`4`$ EMCS. We display $`P(q)`$ in figure 3: is central part is remarkably stable for growing lattice size. Therefore, on the lattice sizes that we are able to thermalize, the pattern we obtain is completely analogous to the one found for the Ising EA model in $`3D`$ .
Due to the global $`Z_2`$ symmetry of the Hamiltonian (1), one would expect it to belong to the same universality class of the IEA model in $`3D`$, which seems even more plausible from our measures of the $`P(q)`$. To further investigate this relation we have measured the critical exponents, in the dilution range where we definitively find SG ordering, namely $`p=0.7`$ and $`0.6`$. Since we have at our disposal only a narrow range of lattice sizes, it is important to use a finite size scaling analysis that allows to study the scaling corrections. We have used the quotient method of , that has been particularly useful in the study of scaling corrections in disordered systems . We measure an operator $`O`$, diverging in the infinite size limit at criticality as $`|TT_\mathrm{c}|^{x_O}`$, on two finite lattices of sides $`L`$ and $`sL`$, and we select the temperature value where the two correlation lengths in units of the lattice size coincide (see the crossing of figure 2). For the quotient of these two measures we have
$$\frac{O(sL,T)}{O(L,T)}|_{\frac{\xi (sL,T)}{\xi (L,T)}=s}=s^{x_O/\nu }\left(1+𝒪(L^\omega )\right),$$
(3)
where $`\omega >0`$ is related to the first irrelevant operator in the Renormalization Group sense. The main advantage of the relation (3) is that the large statistical correlation between the measurements of $`O`$ and $`\xi `$ allows to measure the quotient with sufficient accuracy as to uncover the scaling corrections. In our $`Z_2`$ symmetric case we have of course used the scalar overlap correlation length. Our results are displayed in table I, and they do show the presence of significant scaling corrections ($`\nu `$ is computed using the temperature derivative of $`\xi _q`$, $`\eta `$ from the susceptibility of the scalar overlap). Since we only have few lattice sizes it is meaningless to try an infinite size extrapolation as the one of . We should still mention that the critical exponents obtained by a simple log-log fit (for instance, with data measured at the maximum of the specific-heat) is roughly equivalent to an average of the transient exponents displayed in our table. Therefore the value of $`\nu 1`$ found with Ising spins or in the HEA model with pseudo-dipolar anisotropy , is most probably a preasymptotic value. In fact, the best available results for the IEA model in $`3D`$, $`\nu =1.8(2),\eta =0.26(4)`$ are plausible infinite volume extrapolations for our results. However, in order to definitively elucidate this point, it would be helpful to use the extrapolation method of , that allows to work in the paramagnetic region with a significantly smaller thermalization effort.
We have shown for the first time that dipolar DM anisotropic local interactions are able to induce a SG phase transition for Heisenberg spins in three dimensions. This result has been obtained with the very small realistic value of the anisotropy coupling constant in the $`IIBMnTe`$ series. Given the dramatic effect of this small perturbation term in the Hamiltonian, we suggest that the chiral-glass mechanism proposed in is overwhelmed by the neglected DM term. We have studied the temperature-dilution phase diagram of the Hamiltonian (1) in a large dilution range. We have found that the low-temperature phase changes from AFM to SG order between $`p=0.8`$ and $`p=0.7`$. We have used a combination of micro-canonical and heat-bath Monte Carlo update, that have allowed us to thermalize a $`L=16`$ lattice at $`T=0.78T_\mathrm{c}`$, in the SG phase. We have measured the distribution of the overlap, finding results analogous to the ones of Ising EA $`3D`$ model on similar lattice sizes. We have given an estimate the critical-exponents on the SG dilution range. Our results suffer from severe finite size corrections, but it is plausible to deduce that critical exponents are converging to the Ising EA results, as the experimental results for the $`IIBMnTe`$ suggest . Further open questions need clarification: a better measure of the critical exponents using the method of , the precise characterization of the AFM phase, and a detailed study of the order of the paramagnetic anti-ferromagnetic phase transition at low dilution.
We acknowledge interesting discussions with D. Brogioli, B. Coluzzi, A. Geddo-Lehmann, G. Parisi, F. Ricci-Tersenghi, J. J. Ruiz-Lorenzo. V.M.M. is a M.E.C. fellow and has been partially supported by CICyT(AEN97-1708 and AEN99-1693). The simulations have been performed using the Pentium clusters of the Università di Cagliari (Kalix2) and the Universidad de Zaragoza (RTNN collaboration). |
warning/0002/math0002155.html | ar5iv | text | # A Willmore functional for compact surfaces of complex projective plane
## 1 Introduction
Amongst the global conformal invariants for compact surfaces in a Riemannian manifold $`(M,,)`$, perhaps the better known is the Willmore functional. For an immersion $`\varphi :\mathrm{\Sigma }M`$ from a compact surface $`\mathrm{\Sigma }`$, the Willmore functional is defined by
$$W(\varphi )=_\mathrm{\Sigma }\left(|H|^2+\overline{K}\right)𝑑A,$$
where $`H`$ denotes the mean curvature vector of the immersion $`\varphi `$, $`\overline{K}`$ the sectional curvature of $`M`$ restricted to $`\mathrm{\Sigma }`$ and $`dA`$ the canonical measure of the induced metric.
This functional has been extensively studied when $`M`$ is the Euclidean space $`\text{}^n`$ (or the sphere $`\text{𝕊}^n`$ or the hyperbolic space $`\text{}\text{}^n`$, because $`W`$ is invariant under conformal transformations of ambient space). The type of problems which have been studied for this ambient space $`\text{}^n`$ are of a different nature. On the one hand are the papers written with the object of obtaining lower bounds for the functional $`W`$. Given that $`W`$ has $`4\pi `$ as an absolute minimum and is only reached by the round spheres of $`\text{}^n`$, Willmore proposed to study this functional on tori and conjectured that $`2\pi ^2`$ is the minimum value for these surfaces and it is attained only by the Clifford torus. This problem, which still has not been resolved, can be considered as the starting point of a series of important papers in which minimization problems of $`W`$ are studied. Amongst those which we mention are \[K1\], \[K2\], \[K3\], \[LY\], \[M\], \[MoR\], \[R\] and \[S\].
On the other hand, other authors are interested in the study of critical surfaces for $`W`$ (the known Willmore surfaces) whose Euler-Lagrange equations were obtained (only when $`M`$ has constant curvature) by Weiner in \[W1\]. Two papers of interest are \[BB\] and \[P\], where the authors construct Willmore tori in $`\text{𝕊}^3`$ from minimal surfaces of $`\text{}\text{}^3`$ and elastic curves of $`\text{𝕊}^2`$ respectively. Two other papers of interest are \[B\] and \[Mo\] where Willmore spheres in $`\text{𝕊}^3`$ and $`\text{𝕊}^4`$ respectively are classified, describing the relation of these Willmore spheres and certain class of minimal surfaces in $`\text{}^3`$ and $`\text{}^4`$.
In this paper the authors studied the Willmore functional for compact surfaces of the complex projective plane $`\text{}\text{}^2`$. If $`\varphi :\mathrm{\Sigma }\text{}\text{}^2`$ is an immersion of an orientable compact surface $`\mathrm{\Sigma }`$ in the complex projective plane of constant holomorphic sectional curvature $`4`$, then the Willmore functional is given by
$$W(\varphi )=_\mathrm{\Sigma }\left(|H|^2+1+3C^2\right)𝑑A,$$
where $`C`$ is the Kähler function on $`\mathrm{\Sigma }`$ defined by $`\varphi ^{}(\mathrm{\Omega })=CdA`$, $`\mathrm{\Omega }`$ being the Kähler two form on $`\text{}\text{}^2`$ and $`dA`$ the volume two form on $`\mathrm{\Sigma }`$.
This functional can be written (see section 2) as $`W=\frac{1}{2}(W^++W^{})`$, where $`W^+`$ and $`W^{}`$ are also conformal invariant functionals (see Proposition 1), defined by
$$W^+(\varphi )=_\mathrm{\Sigma }\left(|H|^2+6C^2\right)𝑑A,W^{}(\varphi )=_\mathrm{\Sigma }\left(|H|^2+2\right)𝑑A.$$
These functionals $`W^\pm `$ are closely related with the Penrose twistor bundles $`𝒫^\pm `$ over $`\text{}\text{}^2`$ (see Proposition 3), because twistor holomorphic surfaces, i.e., surfaces of $`\text{}\text{}^2`$ whose twistor liftings are holomorphic (which we study in depth in Theorem 1) are critical surfaces for these functionals (in fact they are minimizers for $`W^\pm `$). As $`\text{}\text{}^2`$ with its canonical orientation is a self-dual Riemannian manifold but not an anti-self-dual, the twistor bundle $`𝒫^{}`$ is a complex manifold and the twistor bundle $`𝒫^+`$ is an almost-complex manifold but non-complex (see \[AHS\]). This fact allows us to easily construct (see Proposition 6) twistor holomorphic compact surfaces with negative spin (and therefore Willmore surfaces for $`W^{}`$), but it is really complicated to get non-trivial examples of twistor holomorphic surfaces with positive spin.
The Euler-Lagrange equations for the functionals $`W^\pm `$ (Proposition 5) say that the minimal surfaces of $`\text{}\text{}^2`$ are critical for the functional $`W^{}`$. It is also interesting to remark that the functional $`W^{}`$ restricted to minimal surfaces is twice the area functional. Due to all these considerations about $`W^\pm `$, the authors think that the $`W^{}`$ Willmore functional is the natural one to be studied for surfaces in $`\text{}\text{}^2`$. In this way, in section 4, we study lower bounds for $`W^{}`$ obtaining results which can be sumarized as follows
> Let $`\varphi :\mathrm{\Sigma }\text{}\text{}^2`$ be an immersion of a compact surface $`\mathrm{\Sigma }`$. Then
>
> $`W^{}(\varphi )2\pi \mu `$, being $`\mu `$ the maximum multiplicity of $`\varphi `$, and the equality holds if and only if $`\varphi (\mathrm{\Sigma })`$ is a complex projective line. In particular $`2\pi `$ is the minimum value for $`W^{}`$.
>
> If $`\varphi `$ is Lagrangian, then $`W^{}(\varphi )4\pi \mu `$, and the equality holds if and only if $`\varphi (\mathrm{\Sigma })`$ is either a real projective line or a Whitney sphere. In this case $`4\pi `$ is the minimum value for $`W^{}`$.
Finally, we wish to mention that the proofs of the above results (Theorems $`2`$ and $`3`$ in the paper) can be easily extended for surfaces in the complex projective space $`\text{}\text{}^n`$, so that Theorems $`2`$ and $`3`$ and Corollaries $`1`$ and $`2`$ are also true when we change $`\text{}\text{}^2`$ by $`\text{}\text{}^n`$. We also wish to mention that almost all the results in the paper can be generalized for surfaces of the complex hyperbolic plane.
## 2 The Willmore functional for surfaces in four manifolds
Let $`(M,,)`$ be an n-dimensional Riemannian manifold and $`\varphi :\mathrm{\Sigma }M`$ an immersion of a compact surface $`\mathrm{\Sigma }`$. The Willmore functional $`W(\varphi )`$ is defined by
$$W(\varphi )=_\mathrm{\Sigma }\left(|H|^2+\overline{K}\right)𝑑A,$$
where $`H`$ is the mean curvature of $`\varphi `$, $`dA`$ the area two form on $`\mathrm{\Sigma }`$ and $`\overline{K}=\overline{R}(e_1,e_2,e_2,e_1)`$, being $`\overline{R}`$ the curvature of $`,`$ and $`\{e_1,e_2\}`$ an orthonormal basis in $`\mathrm{\Sigma }`$. This functional is invariant under conformal changes of the metric $`,`$.
From now on we suppose that $`M`$ is an oriented four–manifold, and that $`\mathrm{\Sigma }`$ is an oriented surface. If $`T^{}\mathrm{\Sigma }`$ is the normal bundle of $`\varphi `$, then we have the orthogonal decomposition
$$\varphi ^{}TM=T\mathrm{\Sigma }T^{}\mathrm{\Sigma }.$$
Let $`\overline{}`$ be the connection on $`\varphi ^{}TM`$ induced by the Levi-Civita connection of $`TM`$ and let $`\overline{}=+^{}`$ be the corresponding decomposition.
If $`\{e_1,e_2,e_3,e_4\}`$ is an oriented orthonormal local reference on $`\varphi ^{}TM`$ such that $`\{e_1,e_2\}`$ is an oriented reference on $`T\mathrm{\Sigma }`$, then we define the normal curvature $`K^{}`$ of the immersion $`\varphi `$ by
$$K^{}=R^{}(e_1,e_2,e_3,e_4),$$
where $`R^{}`$ is the curvature tensor of the normal connection $`^{}`$. Also we will denote by $`\overline{K}^{}`$ the function on $`\mathrm{\Sigma }`$ given by
$$\overline{K}^{}=\overline{R}(e_1,e_2,e_3,e_4).$$
When $`\mathrm{\Sigma }`$ is compact, the Euler characteristics of $`T\mathrm{\Sigma }`$ and $`T^{}\mathrm{\Sigma }`$ are given respectively by
$$\chi =\frac{1}{2\pi }_\mathrm{\Sigma }K𝑑A\mathrm{and}\chi ^{}=\frac{1}{2\pi }_\mathrm{\Sigma }K^{}𝑑A.$$
Now the Willmore functional can be decomposed as
$$W(\varphi )=\frac{1}{2}\left(W^+(\varphi )+W^{}(\varphi )\right),$$
where
$`W^+(\varphi )={\displaystyle _\mathrm{\Sigma }}\left(|H|^2+\overline{K}\overline{K}^{}\right)𝑑A,`$
$`W^{}(\varphi )={\displaystyle _\mathrm{\Sigma }}\left(|H|^2+\overline{K}+\overline{K}^{}\right)𝑑A.`$
###### Proposition 1
The functionals $`W^+`$ and $`W^{}`$ are invariant under conformal changes of the metric $`,`$ on $`M`$.
Proof: Let $`,_{}=e^{2u},`$ a metric on $`M`$ conformal to $`,`$, being $`u:M\text{}`$ a smooth function. Then it is very well-known that the second fundamental forms $`\sigma `$ and $`\sigma _{}`$ of $`\varphi `$ with respect to $`,`$ and $`,_{}`$ are related by
$$\sigma _{}(u,v)=\sigma (u,v)u,v(\overline{}u)^{},$$
where means normal component. From here, it is an easy exercise to check that
$`\left(|H_{}|_{}^2+\overline{K}_{}K_{}\right)dA_{}=\left(|H|^2+\overline{K}K\right)dA,`$
$`\left(\overline{K}_{}^{}K_{}^{}\right)dA_{}=\left(\overline{K}^{}K^{}\right)dA`$
where $``$ means the corresponding object for the metric $`,_{}`$. The above formulae prove the Proposition taking into account the Gauss-Bonnet theorem and that the normal bundles of $`\varphi `$ (with respect to both metrics) are isomorphic. q.e.d.
Now we are going to relate these functionals with the twistor bundles. Given a point $`xM`$, let $`𝒫_x^\pm `$ be the set of almost Hermitian structures $`J_x^\pm `$ over $`T_xM`$ such that if $`\mathrm{\Omega }^\pm (u,v)=J_x^\pm u,v`$, then $`\pm \mathrm{\Omega }\mathrm{\Omega }`$ is the orientation induced on $`T_xM`$ from $`M`$. Then $`𝒫^\pm =_{xM}𝒫_x^\pm `$ are $`\text{}\text{}^1`$-fiber bundles over $`M`$, called the twistor bundles of $`M`$. If $`\pi ^\pm :P^\pm M`$ are the projections, then the vertical distributions $`V^\pm =\mathrm{ker}\pi _{}^\pm `$ inherit from the standard complex structure of the complex projective line $`\text{}\text{}^1`$ an almost complex structure $`J^{v\pm }`$. The Levi-Civita connection of $`M`$ induces a decomposition of the tangent bundles of $`𝒫^\pm `$
$$T𝒫^\pm =H^\pm V^\pm $$
with $`H^\pm TM`$, via $`\pi _{}^\pm `$ and there are also on the horizontal distributions $`H^\pm `$ almost complex structures $`J^{h\pm }`$ defined by
$$J_{J_x^\pm }^{h\pm }=J_x^\pm .$$
So $`𝒥^\pm =J^{h\pm }+J^{v\pm }`$ define almost complex structures on $`𝒫^\pm `$ which only depend of the oriented conformal structure of $`M`$. A central result due to Atiyah, Hitchin and Singer \[AHS\] is that $`(𝒫^+,𝒥^+)`$ is a complex manifold if and only if $`(M,,)`$ is anti-self-dual, and $`(𝒫^{},𝒥^{})`$ is a complex manifold if and only if $`(M,,)`$ is self-dual.
Let $`\varphi :\mathrm{\Sigma }M`$ an immersion of an oriented surface $`\mathrm{\Sigma }`$ and $`\{e_1,e_2,e_3,e_4\}`$ an orthonormal local reference on $`\varphi ^{}TM`$ such that $`\{e_1,e_2\}`$ is an oriented reference on $`T\mathrm{\Sigma }`$. We are going to define two almost complex structures $`J^\pm `$ on $`\varphi ^{}TM`$ by
$$J^\pm (e_1)=e_2,J^\pm (e_3)=\pm e_4.$$
We remark that the $`J^\pm `$ on $`\mathrm{\Sigma }`$ is the complex structure on $`\mathrm{\Sigma }`$ compatible with the given orientation, and that on $`T^{}\mathrm{\Sigma }`$, $`^{}J^\pm =0`$. So the Koszul-Malgrange theorem \[KM\] says that $`J^\pm `$ give to $`T^{}\mathrm{\Sigma }`$ two unique structures of holomorphic line bundles over $`\mathrm{\Sigma }`$ such that a normal section $`\xi `$ is holomorphic if and only if
$$_{J^\pm v}^{}\xi =J^\pm _v\xi $$
for any $`vT\mathrm{\Sigma }`$.
We define the twistor liftings $`\stackrel{~}{\varphi }^\pm :\mathrm{\Sigma }𝒫^\pm `$ by
$$\stackrel{~}{\varphi }^\pm (p)=J_{\varphi (p)}^\pm .$$
Although it is not explicity stated, the following result was proved in \[F\].
###### Proposition 2
Let $`\varphi :\mathrm{\Sigma }M`$ be an immersion from an oriented surface into an oriented four-dimensional Riemannian manifold $`M`$. Then the following assertions are equivalent:
The twistor liftings $`\stackrel{~}{\varphi }^\pm :(\mathrm{\Sigma },J^\pm )(𝒫^\pm ,𝒥^\pm )`$ of $`\varphi `$ are holomorphic.
The second fundamental form $`\sigma `$ of $`\varphi `$ satisfies
$$\sigma (u,v)=\sigma (J^\pm u,J^\pm v)J^\pm \sigma (J^\pm u,v)J^\pm \sigma (u,J^\pm v)$$
for any vectors $`u,vT\mathrm{\Sigma }`$.
The almost complex structures $`J^\pm `$ define on $`\varphi ^{}TM`$ structures of holomorphic bundles.
An immersion $`\varphi :\mathrm{\Sigma }M`$ satisfying one of the three equivalent conditions given in Proposition 2 will be called twistor holomorphic with positive or negative spin. As consequence of Proposition 2, the twistor holomorphicity of $`\varphi `$ does not depend on the chosen orientation on the surface $`\mathrm{\Sigma }`$. Hence we can talk about twistor holomorphic immersions from an orientable surface into an oriented four-dimensional Riemannian manifold. Twistor holomorphic surfaces with positive or negative spin which are also minimal are calles superminimal surfaces with positive or negative spin, (\[B\],\[F\],\[G\]). The surfaces which are simultaneously twistor holomorphic with positive and negative spin are the umbilical ones.
We define bilinear forms $`\sigma ^\pm `$ on $`\mathrm{\Sigma }`$ valuated on $`T^{}\mathrm{\Sigma }`$ by
$$\sigma ^\pm (u,v)=\sigma (u,v)\sigma (J^\pm u,J^\pm v)+J^\pm \sigma (J^\pm u,v)+J^\pm \sigma (u,J^\pm v)$$
for vectors $`u,vT\mathrm{\Sigma }`$. Then it is easy to check that
$`|H|^2+\overline{K}\overline{K}^{}=KK^{}+{\displaystyle \frac{1}{16}}|\sigma ^+|^2,`$
$`|H|^2+\overline{K}+\overline{K}^{}=K+K^{}+{\displaystyle \frac{1}{16}}|\sigma ^{}|^2.`$
Now Proposition 2 gives the following known result (\[F\]), which relates the above functionals $`W^\pm `$ with the twistor theory.
###### Proposition 3
Let $`\varphi :\mathrm{\Sigma }M`$ be an immersion from an orientable compact surface into an oriented four-dimensional Riemanmnian manifold $`M`$. Then
$`W^+(\varphi )2\pi (\chi \chi ^{})`$,
$`W^{}(\varphi )2\pi (\chi +\chi ^{})`$.
Moreover the equality in (i) holds if and only if $`\varphi `$ is twistor holomorphic with positive spin, and the equality in (ii) holds if and only if $`\varphi `$ is twistor holomorphic with negative spin.
In \[ChT\], $`\chi \chi ^{}`$ was called the adjunction number of $`\mathrm{\Sigma }`$ in $`M`$.
Perhaps the first case to study the functionals $`W^\pm `$ was when $`(M,,)`$ is the 4-dimensional Euclidean space, or equivalently (remember that $`W^\pm `$ are invariant under conformal transformation of the ambient space) when $`(M,,)`$ is a sphere $`\text{𝕊}^4`$ with its standard metric of constant curvature one. In this case, if $`\varphi :\mathrm{\Sigma }\text{𝕊}^4`$ is an immersion of an orientable compact surface, then it is easy to check that
$$W^+(\varphi )=W^{}(\varphi )=W(\varphi )=_\mathrm{\Sigma }\left(|H|^2+1\right)𝑑A,$$
and so these functionals are the classical Willmore functional $`W`$.
Moreover, as $`\text{𝕊}^4`$ is self-dual and anti-self-dual, the twistor spaces $`(𝒫^+,𝒥^+)`$ and $`(𝒫^{},𝒥^{})`$ are complex manifolds, and it is well-known that they are biholomorphic to $`\text{}\text{}^3`$ with its standard complex structure. Also the twistor projections $`\pi ^\pm `$ are related by $`\pi ^{}=A\pi ^+`$, where $`A`$ is the antipodal map on $`\text{𝕊}^4`$, and hence the twistor holomorphic surfaces with negative spin of $`\text{𝕊}^4`$ are the images by the antipodal map of the twistor holomorphic surfaces with positive spin. In this case it is remarkable to refer to the papers \[LY\],\[MoR\] where lower bounds for $`W`$ are studied, and also to the paper \[Mo\] where the critical surfaces of genus zero of the functional $`W`$ are classified.
## 3 Case of complex projective plane
In $`\text{}^3`$ we consider the Hermitian product
$$(z,w)=\underset{i=1}{\overset{3}{}}z_i\overline{w}_i,$$
for any $`z,w\text{}^3`$, where $`\overline{z}`$ stands for the conjugate of $`z`$. Then, $`\mathrm{}(,)`$ is the Euclidean metric and $`\mathrm{}(,)`$ the Kähler two–form on $`\text{}^3`$.
Let $`\text{}\text{}^2`$ be the complex projective plane with its canonical Fubini-Study metric $`,`$ of constant holomorphic sectional curvature $`4`$. Then
$$\text{}\text{}^2=\{\mathrm{\Pi }(z)=[z]/z=(z_1,z_2,z_3)\text{}^3\{0\}\}$$
where $`\mathrm{\Pi }:\text{}^3\{0\}\text{}\text{}^2`$ is the standard projection. The metric $`\mathrm{}(,)`$ becomes $`\mathrm{\Pi }`$ in a Riemannian submersion. The complex structure of $`\text{}^3`$ induces via $`\mathrm{\Pi }`$ the canonical complex structure $`J`$ on $`\text{}\text{}^2`$. The Kähler two form $`\mathrm{\Omega }`$ in $`\text{}\text{}^2`$ is defined by $`\mathrm{\Omega }(u,v)=Ju,v`$. We will consider $`\text{}\text{}^2`$ with the orientation $`\mathrm{\Omega }\mathrm{\Omega }`$.
If $`\varphi :\mathrm{\Sigma }\text{}\text{}^2`$ is an immersion of an oriented surface $`\mathrm{\Sigma }`$, the Kähler function $`C`$ on $`\mathrm{\Sigma }`$ is defined by
$$\varphi ^{}\mathrm{\Omega }=CdA$$
where $`dA`$ is the volume form on $`\mathrm{\Sigma }`$. We remark that the sign of $`C`$ depends on the orientation in $`\mathrm{\Sigma }`$. So $`C^2`$ does not depend of the chosen orientation in $`\mathrm{\Sigma }`$ and $`C^2`$ is defined even for non-orientable surfaces. It is clear that the Kähler function satisfies $`1C1`$. Surfaces with $`C=1`$, $`C=1`$ and $`C=0`$ are called respectively holomorphic, anti-holomorphic and Lagrangian. A complex surface will be synonym of either a holomorphic or anti-holomorphic surface.
In addition, if $`\mathrm{\Sigma }`$ is compact, then the topological degree $`d`$ of the map $`\varphi `$ is given by
$$d=\frac{1}{\pi }_\mathrm{\Sigma }C𝑑A.$$
Also it is interesting to remark that the relation between $`J`$ and the almost complex structures $`J^\pm `$ on $`\varphi ^{}T\text{}\text{}^2`$ defined in section 2 is
(1)
$$(Jv)^{}=CJ^+v=CJ^{}v,(J\xi )^{}=CJ^+\xi =CJ^{}\xi ,$$
for any $`vT\mathrm{\Sigma }`$ and $`\xi T^{}\mathrm{\Sigma }`$, where $``$ and $``$ stand for tangent and normal components.
Before studying the functionals $`W^\pm `$ in this case, we are going to point out a strong property of the function $`C`$, which is not true when the codimension of the surface is bigger than two and that we have seen proved only in particular cases.
###### Proposition 4
Let $`\mathrm{\Sigma }`$ be a compact orientable surface of $`\text{}\text{}^2`$ with constant Kähler function $`C`$. Then $`\mathrm{\Sigma }`$ is either a complex or a Lagrangian surface.
Proof: Suppose that $`\mathrm{\Sigma }`$ is not a complex surface, i.e. the constant $`C`$ satisfied $`1<C<1`$. Then, using (1), we can choose an oriented orthonormal local reference $`\{e_1,e_2,e_3,e_4\}`$ such that $`\{e_1,e_2\}`$ is an oriented reference on $`T\mathrm{\Sigma }`$ and
(2)
$$Je_1=Ce_2+\sqrt{1C^2}e_4,Je_2=Ce_1+\sqrt{1C^2}e_3.$$
As $`C=Je_1,e_2`$ and $`C`$ is constant, then derivating $`C`$ with respect to a tangent vector $`v`$ we have
$$\sigma (v,e_1),Je_2=\sigma (v,e_2),Je_1.$$
Using (2) in this formula we get
$$A_{e_3}e_1=A_{e_4}e_2.$$
Now it is straighforward to check, using the above information that $`6C^2K+K^{}=0`$. Integrating this equation, we finally obtain
(3)
$$6C^2\text{ Area }(\mathrm{\Sigma })=2\pi (\chi \chi ^{}).$$
On the other hand, as $`1<C<1`$, $`F:T\mathrm{\Sigma }T^{}\mathrm{\Sigma }`$ defined by $`F(v)=(Jv)^{}`$ defines an isomorphism of vector bundles, which implies that $`\chi =\chi ^{}`$. Using this in (3) we obtain that $`C=0`$, which finishes the proof. q.e.d.
By using (1) and the well-known expression of the curvature of the Fubini–Study metric, it is straightforward to check that these functionals $`W,W^\pm `$ are given in this case by
$`W(\varphi )={\displaystyle _\mathrm{\Sigma }}\left(|H|^2+1+3C^2\right)𝑑A`$
$`W^+(\varphi )={\displaystyle _\mathrm{\Sigma }}\left(|H|^2+6C^2\right)𝑑A,W^{}(\varphi )={\displaystyle _\mathrm{\Sigma }}\left(|H|^2+2\right)𝑑A.`$
We remark that these functionals are defined for any compact surface not necessarily orientable.
As it was showed in section 2, to study the functionals $`W^\pm `$ it is necessary to understand the surfaces wich are twistor holomorphic. From the definition and using (1), it is easy to check that complex surfaces and minimal Lagrangian surfaces of $`\text{}\text{}^2`$ are twistor holomorphic surfaces with positive spin. In fact (as Gauduchon pointed out in \[G\]) they exactly are the superminimal surfaces with positive spin of $`\text{}\text{}^2`$. On the other hand, a complex or Lagrangian twistor holomorphic surface with negative spin, must be an umbilical surface and then from \[KZ\] it must be totally geodesic. In the next result we study some important properties of these surfaces. Before to stablish it, we need to point out a result which was proved in \[EGT\].
###### Lemma 1 (EGT)
Let $`(\mathrm{\Sigma },,)`$ be an oriented compact Riemannian surface, and $`h:\mathrm{\Sigma }\text{}`$ a function of absolute value type, i.e., a smooth function satisfying $`h=|t|f`$ with $`t`$ a holomorphic function and $`f`$ a smooth positive function. Then
$$_\mathrm{\Sigma }\mathrm{\Delta }\mathrm{log}hdA=2\pi N(h),$$
where $`\mathrm{\Delta }`$ is the Laplacian operator of $`\mathrm{\Sigma }`$ and $`N(h)`$ is the sum of all orders for all zeroes of $`h`$.
###### Theorem 1
Let $`\varphi :\mathrm{\Sigma }\text{}\text{}^2`$ be a twistor holomorphic immersion of an oriented surface $`\mathrm{\Sigma }`$. If $`\varphi `$ has positive spin, then
Either $`\varphi `$ is a complex immersion or the complex points of $`\varphi `$ are isolated. Moreover the functions $`\sqrt{1C}`$ and $`\sqrt{1+C}`$ are of absolute value type.
If $`\mathrm{\Sigma }`$ is compact and non-complex, the degree $`d`$ of $`\varphi `$ is given by
$$3d=N_+N_{},$$
where $`N_+=N(\sqrt{1C})`$ and $`N_{}=N(\sqrt{1+C})`$.
Under the conditions of $`ii)`$,
$$W^+(\varphi )=_\mathrm{\Sigma }\left(|H|^2+6C^2\right)𝑑A=2\pi (N_++N_{}).$$
As consequence, if $`\varphi `$ is totally real (i.e. $`\varphi `$ has not complex points), then $`\varphi `$ is a minimal Lagrangian surface.
If $`\varphi `$ has negative spin, then
The mean curvature $`H`$ is a holomorphic vector field on $`T^{}\mathrm{\Sigma }`$ with respect to the holomorphic structure associated to $`J^{}`$. Hence either $`\varphi `$ is superminimal with negative spin or $`H`$ has only isolated zeroes.
If $`\mathrm{\Sigma }`$ is compact and not superminimal, then
$$W^{}(\varphi )=_\mathrm{\Sigma }\left(|H|^2+2\right)𝑑A=2\pi (\chi +N(H)),$$
where $`N(H)`$ is the number of zeroes of $`H`$.
Moreover, if $`\varphi :\mathrm{\Sigma }\text{}\text{}^2`$ is an immersion of a sphere whose mean curvature $`H`$ is a holomorphic vector field on $`T^{}\mathrm{\Sigma }`$ with the holomorphic structure associated to $`J^{}`$, then $`\varphi `$ is either a complex immersion or a twistor holomorphic immersion with negative spin.
Proof: To prove the result we need to use complex coordinates on $`\mathrm{\Sigma }`$. Let $`z=x+iy`$ be a local isothermal parameter on $`\mathrm{\Sigma }`$ compatible with the given orientation. We will denote
$$_z=\frac{1}{2}(_xi_y),\overline{}_{\overline{z}}=\frac{1}{2}(_x+i_y),$$
the Cauchy-Riemann operators. Then
$$|_z|^2=_z,_{\overline{z}}>0,_z,_z=0,$$
where $`,`$ also denote the -linear extension of the metric $`,`$ to the complexified bundles. Then
(4)
$$__z_z=\mathrm{log}|_z|^2_z,\sigma (_z,_{\overline{z}})=|_z|^2H.$$
If $`\{e_3,e_4\}`$ is an orthonormal local reference on $`T^{}\mathrm{\Sigma }`$ such that $`\{_x,_y,e_3,e_4\}`$ is the orientation on $`\varphi ^{}T\text{}\text{}^2`$, we define
$$\xi =\frac{e_3ie_4}{\sqrt{2}}.$$
Then one can check (translating to complex notation the above arguments) that: a) $`\varphi `$ is twistor holomorphic with positive spin if and only if $`\sigma (_z,_z),\xi =0`$; b) $`\varphi `$ is twistor holomorphic with negative spin if and only if $`\sigma (_z,_z),\overline{\xi }=0`$; and c) $`J_z=iC_z+J_z,\xi \overline{\xi }`$.
Suppose that $`\varphi `$ is twistor holomorphic with positive spin. From a) we have that
(5)
$$\sigma (z,z)=\sigma (z,z),\overline{\xi }\xi .$$
Then if $`F=J_z,\xi `$, using (4) and (5) it is clear that
$$F=\left(\mathrm{log}|_z|^2+__z^{}\xi ,\overline{\xi }\right)F.$$
So either $`F`$ vanishes identically or $`F`$ has only isolated zeros. But from c) we obtain that $`|F|^2=|_z|^2(1C^2)`$, and hence $`\varphi `$ is either a complex immersion or the complex points of $`\varphi `$ are isolated. This proves the first part of i).
On the other hand, the Kähler function $`C`$ is written by
$$C=\frac{iJ_{\overline{z}},_z}{|_z|^2}.$$
So, using (4) and (5) again, it is easy to see that
(6)
$$C=iJH,_z=iFH,\overline{\xi }.$$
So derivating (6), using again (4) and (5) and the Codazzi equation, one can see that
(7)
$$\overline{}C=|_z|^2(|H|^2+3C(1C^2)).$$
Now, from (6) and (7), the gradient $`C`$ and the Laplacian $`\mathrm{\Delta }C`$ of the fuction $`C`$ satisfied
$$|C|^2=(1C^2)|H|^2,\mathrm{\Delta }C=2C(|H|^2+3(1C^2)).$$
If $`\varphi `$ is not a complex immersion and from i) the complex points are isolated, then, outside the complex points, one can obtain easily from the above formulae
$$\mathrm{\Delta }\mathrm{log}\sqrt{1C}=\frac{|H|^2}{2}3C(1+C),\mathrm{\Delta }\mathrm{log}\sqrt{1+C}=\frac{|H|^2}{2}+3C(1C).$$
Now, ii) and iii) follow from Lemma 1 by proving that $`\sqrt{1C}`$ and $`\sqrt{1+C}`$ are of absolute value type. This last assertion follows using a similar reasoning as in \[EGT\], Theorem A,(i), and its proof will be omited.
Suppose now that $`\varphi `$ is twistor holomorphic with negative spin. From $`b)`$ we have that
(8)
$$\sigma (z,z)=(z,z),\xi \overline{\xi }.$$
In complex notation, the holomorphicity of $`H`$ with respect to the holomorphic structure associated to $`J^{}`$ means that
(9)
$$_{_{\overline{z}}}^{}(HiJ^{}H)=0.$$
Derivating the second equation of (4), using (4), (8) and the Codazzi equation, it is not difficult to obtain that
$$_{_{\overline{z}}}^{}H=B\xi ,$$
for a certain complex function $`B`$. From here it is clear that $`H`$ satisfied equation (9) and hence $`H`$ is a holomorphic vector field. The remaining assertions in $`i)`$ and $`ii)`$ are easy consequences of this fact.
Finally, let $`\varphi :\mathrm{\Sigma }\text{}\text{}^2`$ be an immersion of a sphere whose mean curvature vector $`H`$ is a holomorphic vector field with the holomorphic structure associated to $`J^{}`$. We consider the very well-known cubic differential form $`\mathrm{\Theta }`$ on $`\mathrm{\Sigma }`$ (see \[ES\],\[EGT\]) defined by
$$\mathrm{\Theta }=\sigma (_z,_z),J_zdz^3.$$
Now using similar arguments to used in the proof of the positive case, it is not difficult to check that the holomorphicity of $`H`$ implies that $`\mathrm{\Theta }`$ is holomorphic, and then as $`\mathrm{\Sigma }`$ is a sphere, $`\mathrm{\Theta }`$ vanishes identically. So
(10)
$$0=\sigma (_z,_z),J_z=\sigma (_z,_z),\overline{\xi }F.$$
As $`H`$ is holomorphic, either $`H0`$ or $`H`$ has only isolated zeroes. If $`H0`$, then (see \[EGT\]) either $`\varphi `$ is a complex immersion or the complex points of $`\varphi `$ are isolated. So from (10), if $`\varphi `$ is not a complex immersion, $`\varphi `$ is twistor holomorphic with positive spin and as consequence superminimal with positive spin. If $`H`$ has only isolated zeroes, the set of points where $`F`$ vanishes has empty interior, and so from (10) $`\varphi `$ is twistor holomorphic with positive spin. This finishes the proof. q.e.d.
Now we are going to get the Euler-Lagrange equations for the functionals $`W^+`$ and $`W^{}`$.
###### Proposition 5
Let $`\varphi :\mathrm{\Sigma }\text{}\text{}^2`$ be an immersion of an compact surface $`\mathrm{\Sigma }`$.
$`\varphi `$ is a critical point of the functional $`W^+`$ if and only if the mean curvature vector $`H`$ of $`\varphi `$ satisfies
$$\mathrm{\Delta }^{}H+(5+9C^22|H|^2)H+\stackrel{~}{A}(H)+12(JJ^+C)^{}=0.$$
$`\varphi `$ is a critical point of the functional $`W^{}`$ if and only if the mean curvature vector $`H`$ of $`\varphi `$ satisfies
$$\mathrm{\Delta }^{}H+(13C^22|H|^2)H+\stackrel{~}{A}(H)=0.$$
In both cases $`\mathrm{\Delta }^{}`$ and $`\stackrel{~}{A}`$ are defined by
$$\mathrm{\Delta }^{}=\underset{i=1}{\overset{2}{}}\{_{e_i}^{}_{e_i}^{}_{_{e_i}e_i}^{}\},\stackrel{~}{A}H=\underset{e=1}{\overset{2}{}}\sigma (A_He_i,e_i),$$
being $`\{e_1,e_2\}`$ an orthonormal reference tangent to $`\mathrm{\Sigma }`$.
###### Remark 1
We note that minimal surfaces of $`\text{}\text{}^2`$ are critical points of the Willmore functional $`W^{}`$. However the only minimal surfaces critical for $`W^+`$ are the superminimal with positive spin, i.e., the complex and minimal Lagrangian surfaces \[G\]. In fact, from i) a minimal surface is critical for $`W^+`$ if and only if $`JJ^+C`$ is tangent to the surface. So $`\mathrm{\Sigma }_0=\{p\mathrm{\Sigma }/(C)(p)0\}`$ is an open subset of $`\mathrm{\Sigma }`$ where $`\varphi `$ is a complex immersion, which is imposible by the very definition of $`\mathrm{\Sigma }_0`$. So $`\mathrm{\Sigma }_0=ø`$, and then Proposition 4 proves the assertion.
On the other hand, notice that twistor holomorphic immersions with positive or negative spin are a kind of critical surfaces for $`W^+`$ or $`W^{}`$ respectively because they minimize the corresponding Willmore functionals.
Proof of Proposition 5: Following the computations got by Weiner in \[W1\], Theorem 2.1, it is easy to see that the first derivative of the functional $`W^{}`$ is given by
(11) $`\delta W^{}(\varphi )={\displaystyle _\mathrm{\Sigma }}\mathrm{\Delta }H+\stackrel{~}{A}H+(13C^22|H|^2),\delta \varphi 𝑑A,`$
where $`\delta \varphi `$ stands for the variation vector field, which can be taken normal to the surface $`\mathrm{\Sigma }`$. Now in order to compute the first derivative of the functional $`W^+`$, we start studying the functional
$$F(\varphi )=_\mathrm{\Sigma }C^2𝑑A.$$
Using the well-known fact that
(12)
$$\delta (dA)=2H,\delta \varphi dA,$$
it follows that the first derivative for the functional $`F`$ is given by
$$\delta F(\varphi )=_\mathrm{\Sigma }\left(2C\delta (C)2C^2H,\delta \varphi \right)𝑑A.$$
In order to compute $`\delta (C)`$, we recall the definition of $`C`$:
$$\varphi ^{}\mathrm{\Omega }=CdA.$$
Then, taking derivatives and using (12)
(13) $`\delta (\varphi ^{}\mathrm{\Omega })=\delta (C)dA2CH,\delta \varphi dA.`$
Now following standard arguments it is easy to check that if $`\{e_1,e_2\}`$ is an oriented orthonormal reference on $`T\mathrm{\Sigma }`$,
$$(\delta (\varphi ^{}\mathrm{\Omega }))(e_1,e_2)=\overline{}_{e_2}\delta \varphi ,e_1\overline{}_{e_1}\delta \varphi ,e_2=\text{div }J^+(J\delta \varphi )^{},$$
where div stands for the divergence operator on $`\mathrm{\Sigma }`$. So, using this equation in (13) we get
$$\delta (C)dA=\left(\text{div }J^+(J\delta \varphi )^{}+2CH,\delta \varphi \right)dA.$$
So the first variation of $`F(\varphi )`$ is
$$\delta F(\varphi )=_\mathrm{\Sigma }\left(2C^2H,\delta \varphi 2C\text{ div }J^+(J\delta \varphi )^{}\right)𝑑A.$$
On the other hand, from the divergence Theorem
$`{\displaystyle _\mathrm{\Sigma }}C\text{ div }(J^+(J\delta \varphi )^{}dA={\displaystyle _\mathrm{\Sigma }}C,J^+(J\delta \varphi )^{}dA`$
$`={\displaystyle _\mathrm{\Sigma }}(JJ^+C)^{},\delta \varphi 𝑑A.`$
Using this we finally get
$$\delta F(\varphi )=_\mathrm{\Sigma }2C^2H+(J(J^+C)^{},\delta \varphi dA.$$
Now using this formula, (11) and (12) we obtain the first variation of $`W^+`$, and the Proposition follows. q.e.d.
To finish this section, we are going to study twistor holomorphic surfaces from the view point of the twistor spaces. If we consider $`\text{}\text{}^2`$ with the orientation $`\mathrm{\Omega }\mathrm{\Omega }`$, then $`\text{}\text{}^2`$ is a self-dual Riemannian manifold but not an anti-self-dual Riemannian manifold. So the twistor bundle $`(𝒫^{},𝒥^{})`$ is a complex manifold and $`(𝒫^+,𝒥^+)`$ is not a complex manifold. In fact it is well-known (see \[ES\]) that $`(𝒫^+,𝒥^+)`$ can be differentiably identified with $`P(T^{2,0}\text{}\text{}^2\text{})`$. Also, the complex manifold $`(𝒫^{},𝒥^{})`$ can be endowed with a Riemannian metric which becomes it in a Einstein-Kähler manifold. Under this identification, $`𝒫^{}`$ is the following complex hypersurface of $`(\text{}\text{}^2\times \text{}\text{}^2,,,,JJ)`$
$$𝒫^{}\{([z],[w])\text{}\text{}^2\times \text{}\text{}^2/z^t\overline{w}=0\},$$
and the twistor projection $`\pi ^{}`$ is nothing but
$$\pi ^{}(([z],[w]))=[\overline{z}\overline{w}],$$
for any $`[z],[w]\text{}\text{}^2`$. Also the two natural projections $`\pi _i:𝒫^{}\text{}\text{}^2`$ with $`i=1,2`$ are holomorphic and antiholomorphic maps respectively.
The non-compact Lie group $`PGL(3,\text{})`$ of the complex transformations of $`\text{}\text{}^2`$ acts over $`𝒫^{}`$ by
$$[A]([z],[w])=([Az],[A^1w]),$$
where $`[A]`$ is the class of a matrix $`AGL(3,\text{})`$ and $`A^{}`$ stands for the transpose conjugate of $`A`$. When $`AU(3,\text{})`$, then
$$\pi ^{}([A]([z],[w]))=[A^t](\pi ^{}([z],[w])).$$
###### Remark 2
Using this twistor space, one can reformulate, in a not very complicated way, the second part of Theorem 1 as follows. Given a non-complex immersion $`\varphi :\mathrm{\Sigma }\text{}\text{}^2`$ of an oriented surface $`\mathrm{\Sigma }`$, then $`H`$ is a holomorphic vector field in $`T^{}\mathrm{\Sigma }`$ with respect to the holomorphic structure associated to $`J^{}`$ if and only if the twistor lifting $`\stackrel{~}{\varphi }`$ of $`\varphi `$ is a harmonic map. So the second part of Theorem 1 is a consequence of the fact that every holomorphic map from an oriented surface into a Kähler manifold is harmonic.
The study (from this view point) of twistor holomorphic surfaces with positive spin is complicated because it is equivalent to study holomorphic curves in the non-complex manifold $`𝒫^+`$. Hovewer, as we will point out now, this view point will allow to understand very well the twistor holomorphic surfaces with negative spin. In fact, if $`\varphi :\mathrm{\Sigma }\text{}\text{}^2`$ is a twistor holomorphic immersion with negative spin, and $`\stackrel{~}{\varphi }=(\stackrel{~}{\varphi }_1,\stackrel{~}{\varphi }_2)`$ its twistor lifting, then $`\stackrel{~}{\varphi }_1:\mathrm{\Sigma }\text{}\text{}^2`$ is a holomorphic curve and $`\stackrel{~}{\varphi }_2:\mathrm{\Sigma }\text{}\text{}^2`$ is an anti-holomorphic curve. In this context, $`\varphi `$ is superminimal with negative spin if and only if $`\stackrel{~}{\varphi }_2`$ is the dual curve of $`\stackrel{~}{\varphi }_1`$. Hence, the Lie group $`PGL(3,\text{})`$ acts on the twistor holomorphic surfaces with negative spin in the following way:
$$[A]\varphi =\pi ^{}([A](\stackrel{~}{\varphi }_1,\stackrel{~}{\varphi }_2))=\pi ^{}([A]\stackrel{~}{\varphi }_1,[A^1]\stackrel{~}{\varphi }_2),$$
for any $`AGL(3,\text{})`$. This action sends superminimal surfaces with negative spin into themselves, and when $`AU(3,\text{})`$, this action is the standard one. We will say that a twistor holomorphic surface with negative spin is a twistor deformation of another one if it is its image under the above action.
If $`\mathrm{\Sigma }`$ is also compact, we will denote by $`d_i`$ the degree of $`\stackrel{~}{\varphi }_i`$, $`i=1,2`$. If $`d_i=0`$ for some $`i`$, then $`\stackrel{~}{\varphi }_i`$ is a point, and then (see definition of $`𝒫^{}`$) $`\stackrel{~}{\varphi }_j`$ with $`ji`$ is a complex projective line. Hence $`d_j=1`$ and $`\varphi (\mathrm{\Sigma })`$ is a complex projective line.
If $`\varphi `$ is also a conformal immersion, then it is not difficult to check that $`\stackrel{~}{\varphi }_i`$ are conformal immersions too and if $`,`$ denotes the metric on $`\mathrm{\Sigma }`$ induced by $`\varphi `$ and $`,_i`$ the metrics on $`\mathrm{\Sigma }`$ induced by $`\stackrel{~}{\varphi }_i`$, $`i=1,2`$, then
(14)
$$,_1=\left(\frac{|H|^2+2(1C)}{4}\right),,,_2=\left(\frac{|H|^2+2(1+C)}{4}\right),.$$
From here it follows the following result.
###### Proposition 6
Let $`\varphi :\mathrm{\Sigma }\text{}\text{}^2`$ be a twistor holomorphic immersion with negative spin of an orientable compact surface $`\mathrm{\Sigma }`$. Then
$`W^{}(\varphi )={\displaystyle _\mathrm{\Sigma }}\left(|H|^2+2\right)𝑑A=2\pi (d_1+d_2)`$,
The degree $`d`$ of $`\varphi `$ is given by $`d=d_2d_1`$.
$`W^{}(\varphi )2\pi `$, and the equality holds if and only if $`\varphi (\mathrm{\Sigma })`$ is a complex projective line $`\text{}\text{}^1`$.
In Proposition 6 above we have found that complex projective lines are twistor holomorphic suarfaces with negative spin attaining the minimum value for $`W^{}`$. Now, we are going to describe other examples of twistor holomorphic surfaces with negative spin of $`\text{}\text{}^2`$ whose $`W^{}`$ are also small.
Twistor holomorphic compact surfaces with negative spin and $`W^{}=4\pi `$.
In this case, $`(d_1,d_2)`$ can take the value $`(1,1)`$. So, from Theorem 1 and Proposition 6, the surface must be a sphere with $`\chi ^{}=0`$ and $`d=0`$. From the remark made before Proposition 6, $`\stackrel{~}{\varphi }_2`$ cannot be the dual curve of $`\stackrel{~}{\varphi }_1`$, hence $`\varphi `$ cannot be superminimal and then from Theorem 1, the holomorphic field $`H`$ has not zeroes. Also, as $`\stackrel{~}{\varphi }_i`$ are unramified, (13) says that $`\varphi `$ is unramified too.
Now, up to an holomorphic transformation of $`\text{}\text{}^2`$, $`\stackrel{~}{\varphi }_1:\text{}\{\mathrm{}\}\text{}\text{}^2`$ can be taken as
$$\stackrel{~}{\varphi }_1(z)=\mathrm{\Pi }(1,z,0).$$
Now, since $`\stackrel{~}{\varphi }=(\stackrel{~}{\varphi }_1,\stackrel{~}{\varphi }_2)`$ lies in $`𝒫^{}`$, easy computations say that $`\stackrel{~}{\varphi }_2:\text{}\{\mathrm{}\}\text{}\text{}^2`$ is given by
$$\stackrel{~}{\varphi }_2(z)=\mathrm{\Pi }(\overline{z},1,\overline{P(z)}),\text{with}P(z)=a+bz,(a,b)\text{}^2\{0\}.$$
So our twistor holomorphic surface, (see definition of $`\pi ^{}`$), is a twistor deformation of $`\varphi _{a,b}:\text{}\{\mathrm{}\}\text{}\text{}^2`$, with $`(a,b)\text{}^2\{0\}`$, where
$$\varphi _{a,b}(z)=\mathrm{\Pi }(a\overline{z}b|z|^2,a+bz,1+|z|^2).$$
It is interesting to remark that $`\varphi _{a,b}`$ are embeddings.
Twistor holomorphic compact surfaces with negative spin and $`W^{}=6\pi `$.
In this case, $`(d_1,d_2)`$ can take the value $`(1,2),(2,1)`$. But up to anti-holomorphic isometries of $`\text{}\text{}^2`$ it is sufficient to study the case $`(d_1,d_2)=(1,2)`$. In this case, the surface must be also a sphere but with $`\chi ^{}=1`$ and $`d=1`$. As $`\stackrel{~}{\varphi }_2`$ cannot be the dual curve of $`\stackrel{~}{\varphi }_1`$, $`\varphi `$ cannot be superminimal and then from Theorem 1, the holomorphic field $`H`$ has only one zero. Also, as $`\stackrel{~}{\varphi }_i`$ are unramified, (13) says that $`\varphi `$ is unramified too.
Now, as $`PGL(3,\text{})`$ acts transitively on the conics of $`\text{}\text{}^2`$, $`\stackrel{~}{\varphi }_2:\text{}\{\mathrm{}\}\text{}\text{}^2`$ can be taken as
$$\stackrel{~}{\varphi }_2(z)=\mathrm{\Pi }(1,\overline{z},\overline{z}^2).$$
So, since $`\stackrel{~}{\varphi }=(\stackrel{~}{\varphi }_1,\stackrel{~}{\varphi }_2)`$ lies in $`𝒫^{}`$, $`\stackrel{~}{\varphi }_1:\text{}\{\mathrm{}\}\text{}\text{}^2`$ is given by
$$\stackrel{~}{\varphi }_1(z)=\mathrm{\Pi }(az,a+bz,b)\text{with}[(a,b)]\text{}\text{}^1.$$
So our twistor holomorphic surface is a twistor deformation of $`\psi _{[(a,b)]}:\text{}\{\mathrm{}\}\text{}\text{}^2`$, with $`[(a,b)]\text{}\text{}^1`$, where
$$\psi _{[(a,b)]}(z)=\mathrm{\Pi }(\overline{b}z(|z|^2+1)\overline{a}z^2,(\overline{b}+\overline{a}|z|^2z),\overline{a}(|z|^2+1)\overline{b}\overline{z}).$$
Twistor holomorphic compact surfaces with negative spin and $`W^{}=8\pi `$.
The next examples are particulary interesting because they will be characterized in the next section. They will be twistor holomorphic spheres with negative spin, $`W^{}=8\pi `$ and Kähler function $`C=0`$, i.e. Lagrangian surfaces. In particular their degrees will be zero, and then their twistor liftings will be a pair of conics. In fact these examples are called Whitney spheres and in \[CU2\] they were characterized as the only twistor holomorphic Lagrangian surfaces with negative spin. Up to isometries of $`\text{}\text{}^2`$ the are only a $`1`$-parameter family of surfaces which can be defined as follow.
For each $`t[0,\mathrm{}[`$, we define $`\varphi _t:\text{𝕊}^2\text{}\text{}^2`$ by
$$\varphi _t(x,y,z)=\mathrm{\Pi }(x,y,z\mathrm{cosh}t+i\mathrm{sinh}t),$$
for any $`(x,y,z)\text{𝕊}^2=\{(x,y,z)\text{}^3/x^2+y^2+z^2=1\}`$.
We remark that $`\varphi _0`$ is the totally geodesic immersion of $`\text{𝕊}^2`$, which is a covering of the totally geodesic embedding of $`\text{}\text{}^2`$. For $`t>0`$, $`\varphi _t`$ is an embedding except at the poles of $`\text{𝕊}^2`$ where it has a double point. Amongst them only $`\varphi _0`$ is a minimal surface.
## 4 Lower bounds for the functional $`W^{}`$
In this section we start obtaining a lower bound for the functional $`W^{}`$.
###### Theorem 2
Let $`\varphi :\mathrm{\Sigma }\text{}\text{}^2`$ be an immersion of a compact surface $`\mathrm{\Sigma }`$ and $`\mu `$ the maximum multiplicity of $`\varphi `$. Then:
$$_\mathrm{\Sigma }\left(|H|^2+3+C^2\right)𝑑A4\pi \mu ,$$
and the equality holds if and only if $`\varphi (\mathrm{\Sigma })`$ is a complex projective line $`\text{}\text{}^1`$. As a consequence,
$$W^{}(\varphi )=_\mathrm{\Sigma }\left(|H|^2+2\right)𝑑A2\pi \mu ,$$
and the equality holds if and only if $`\varphi (\mathrm{\Sigma })`$ is a complex projective line $`\text{}\text{}^1`$.
###### Corollary 1
The area of a compact minimal surface $`\mathrm{\Sigma }`$ immersed in $`\text{}\text{}^2`$ with maximum multiplicity $`\mu `$ satisfies
$$\text{Area}\text{ }(\mathrm{\Sigma })\pi \mu ,$$
and the equality holds if and only if $`\mathrm{\Sigma }`$ is a complex projective line of $`\text{}\text{}^2`$.
Proof of Theorem 2: As the maximum multiplicity of $`\varphi `$ is $`\mu `$, let $`\{p_1,\mathrm{},p_\mu \}`$ be points of $`\mathrm{\Sigma }`$ such that $`\varphi (p_i)=[a]\text{}\text{}^2`$ for any $`i=1,\mathrm{},\mu `$. We define a function $`f:\text{}\text{}^2\text{}`$ by
$$f([z])=\frac{|(z,a)|^2}{|z|^2|a|^2},$$
for any $`[z]\text{}\text{}^2`$. Then $`0f1`$, and $`f([z])=0`$ if and only if $`[z]`$ is in the cut locus $`\text{}\text{}_{[a]}^1`$ of the point $`[a]`$. Also, $`f([z])=1`$ if and only if $`[z]=[a]`$. So $`\mathrm{log}(1f)`$ is a well defined function on $`\text{}\text{}^2\{[a]\}`$.
From now on (in order to simplify the notation) we will consider that $`|a|=1`$, and we will restrict $`\mathrm{\Pi }`$ to the unit sphere $`\text{𝕊}^5\text{}^3`$. So, the function $`f`$ will be nothing but
$$f([z])=|(z,a)|^2,$$
for any $`z\text{𝕊}^5`$.
First we compute the gradient of $`f`$. If $`v`$ is any tangent vector to $`\text{}\text{}^2`$ at $`[z]`$, then
$$v(f)=2\mathrm{}(v^{},(z,a)a),$$
being $`v^{}`$ the horizontal lift to $`T_z\text{𝕊}^5`$ of $`v`$. So,
$$(\overline{}f)_{[z]}=2(d\mathrm{\Pi })_z((z,a)a|(z,a)|^2z),$$
for any $`[z]\text{}\text{}^2`$. It is intereting to remark that $`|\overline{}f|^2=4f(1f)`$.
Now, using that $`\mathrm{\Pi }:\text{𝕊}^5\text{}\text{}^2`$ is a Riemannian submersion, it is easy to check that the Hessian of $`f`$ is given by
(15) $`(\overline{}^2f)(u,v)=2fu,v+2\mathrm{}((u^{},a)(a,v^{})),`$
for any vectors $`u,vT_{[z]}\text{}\text{}^2`$, being $`u^{},v^{}`$ the horizontal lifts to $`T_z\text{𝕊}^5`$ of $`u,v`$. In that follows it will be interesting to take into account the following formula, which can be easily check
(16)
$$f\mathrm{}((u^{},a)(a,v^{}))=\overline{}f,u\overline{}f,v+\overline{}f,Ju\overline{}f,Jv$$
We can define on $`\mathrm{\Sigma }\{p_1,\mathrm{},p_\mu \}`$ the function $`\mathrm{log}(1h)`$, where $`h=f(\varphi )`$.
By decomposing
$$\overline{}f\varphi =h+\xi $$
in its tangencial and normal components, it is easy to see that
$`|h|^2=4h(1h)|\xi |^2,`$
(17)
$`(^2h)(u,v)=(\overline{}^2f)(\varphi _{}u,\varphi _{}v)+\sigma (u,v),\xi ,`$
for any vectors $`u,v`$ tangent to $`\mathrm{\Sigma }`$. From (17) and (15) we obtain that
$$\mathrm{\Delta }\mathrm{log}(1h)=\frac{2}{1h}\underset{i=1}{\overset{2}{}}|(e_i^{},a)|^2\frac{2}{1h}H,\xi +\frac{|\xi |^2}{(1h)^2},$$
where $`\{e_1,e_2\}`$ is an orthonormal reference on $`\mathrm{\Sigma }`$. As
$$\left|H\frac{\xi }{1h}\right|^2=|H|^2+\frac{|\xi |^2}{(1h)^2}\frac{2}{1h}H,\xi ,$$
we obtain that
(18)
$$\mathrm{\Delta }\mathrm{log}(1h)|H|^2\frac{2}{1h}\underset{i=1}{\overset{2}{}}|(e_i^{},a)|^2,$$
and the equality holds if and only if $`H=\xi /(1h)`$.
Now, from (16), we obtain
(19)
$$4h\underset{i=1}{\overset{2}{}}|(e_i^{},a)|^2=(1+C^2)|h|^2+(1C^2)|\xi |^2+2(J\xi )^{},Jh.$$
On the other hand
$`0|\sqrt{2}(J\xi )^{}{\displaystyle \frac{1}{\sqrt{2}}}(Jh)^{}|^2`$
$`=2C^2|\xi |^2{\displaystyle \frac{1}{2}}(1C^2)|h|^2+2(J\xi )^{},Jh.`$
Using this inequality in (19) we get
$`4h{\displaystyle \underset{i=1}{\overset{2}{}}}|(e_i^{},a)|^2{\displaystyle \frac{3+C^2}{2}}|h|^2+(1+C^2)|\xi |^2`$
$`={\displaystyle \frac{3+C^2}{2}}|\overline{}f\varphi |^2{\displaystyle \frac{1C^2}{2}}|\xi |^2{\displaystyle \frac{3+C^2}{2}}4h(1h),`$
and the equality holds if and only if $`(1C^2)\xi =0`$ and $`(J\xi )^{}`$ and $`(Jh)^{}`$ are colinear.
So finally, from (18) we get that
(20)
$$\mathrm{\Delta }\mathrm{log}(1h)|H|^23C^2,$$
and the equality holds if and only $`\varphi (\mathrm{\Sigma })`$ is a complex projective line $`\text{}\text{}^1`$. In fact, the equality holds if and only if $`H=\xi /(1h)`$, $`(1C^2)\xi =0`$ and $`(J\xi )^{}`$ and $`(Jh)^{}`$ are colinear. Let
$$\mathrm{\Sigma }_0=\{p\mathrm{\Sigma }/C^2(p)=1\}.$$
If the interior of $`\mathrm{\Sigma }_0`$ is empty, then $`\xi =0`$ on the whole $`\mathrm{\Sigma }`$. If not the interior of $`\mathrm{\Sigma }_0`$ is a complex curve and then it is a minimal surface. So, as $`\xi =(1h)H`$, we have $`\xi =0`$ in the interior of $`\mathrm{\Sigma }_0`$. So in any case $`\xi 0`$ on the whole surface and in particular the surface is minimal. The third condition says that $`Jh`$ is a tangent vector, and so outside the zeroes of $`h`$, $`\varphi `$ is a complex curve. On this set, and because $`\overline{}f\varphi =h`$, (16) says that $`\sigma (v,h)=0`$ for any vector $`v`$ tangent to $`\mathrm{\Sigma }`$, and so $`\varphi `$ is totally geodesic. As $`|h|^2=4h(1h)`$,
$$\{p\mathrm{\Sigma }/(h)(p)=0\}=\{p_1,\mathrm{},p_\mu \}\varphi ^1\left(\text{}\text{}_{[a]}^1\right).$$
So $`\varphi `$ is a complex totally geodesic surface on the whole $`\mathrm{\Sigma }`$ which finishes our claim.
Let $`B_{[a]}(\epsilon )`$ be the geodesic ball in $`\text{}\text{}^2`$ centered at the point $`[a]`$ with radius $`\mathrm{arccos}\sqrt{1\epsilon ^2}`$, that is, the set
$$B_{[a]}(\epsilon )=\{p\text{}\text{}^2|\mathrm{\hspace{0.17em}1}f(p)\epsilon ^2\},$$
with $`\epsilon `$ too small in order to $`B_\epsilon =\varphi ^1(B_{[a]}(\epsilon ))`$ will be the disjoint union of neighbourhoods $`B_i`$, $`i=1,\mathrm{},\mu `$ around $`p_i`$ in $`\mathrm{\Sigma }`$. Then the divergence theorem on the manifold $`\mathrm{\Sigma }B`$ says that
$$_{\mathrm{\Sigma }B}\mathrm{\Delta }\mathrm{log}(1h)𝑑A=\underset{i=1}{\overset{\mu }{}}_{B_i}\frac{h,\nu _i}{1h}𝑑s,$$
where $`\nu _i`$ is the unit conormal of $`B_i`$ pointing to the interior of $`B_i`$. Since the function $`h`$ attains its maximum value 1 at each $`p_i`$ and $`h`$ is constant along each $`B_i`$, we have that
$$\nu _i=\frac{h}{|h|}_{|B_i}.$$
So, combining these equalities with the integral equality above, we have
$$_{\mathrm{\Sigma }B}\mathrm{\Delta }\mathrm{log}(1h)𝑑A=\underset{i=1}{\overset{\mu }{}}_{B_i}\frac{|h|}{1h}𝑑s=\underset{i=1}{\overset{\mu }{}}\frac{1}{\epsilon ^2}_{B_i}|h|𝑑s.$$
As $`\epsilon `$ tends to zero, $`|h|`$ along $`B_i`$ approaches to $`|f|=2\epsilon \sqrt{1\epsilon ^2}`$ and the length of $`B_i`$ approaches to
$$2\pi \text{ radius }B_i=2\pi \mathrm{arccos}\sqrt{1\epsilon ^2}.$$
Then, we obtain that
$$_\mathrm{\Sigma }\mathrm{\Delta }\mathrm{log}(1h)𝑑A=4\pi \mu .$$
This equality and (20) prove the inequality we were looking for. q.e.d.
In the next result we improve the lower bound obtained in Theorem 2 for $`W^{}`$ in the family of Lagrangian surfaces of $`\text{}\text{}^2`$.
###### Theorem 3
Let $`\varphi :\mathrm{\Sigma }\text{}\text{}^2`$ be a Lagrangian immersion of a compact surface $`\mathrm{\Sigma }`$ and $`\mu `$ the maximum multiplicity of $`\varphi `$. Then
$$W^{}(\varphi )=_\mathrm{\Sigma }\left(|H|^2+2\right)𝑑A4\pi \mu ,$$
and the equality holds if and only if either $`\varphi `$ is totally geodesic and $`\varphi (\mathrm{\Sigma })`$ is a real projective plane with $`W^{}(\varphi )=4\pi `$ or $`\varphi `$ is a Whitney sphere with $`W^{}(\varphi )=8\pi `$.
###### Corollary 2
The area of a minimal Lagrangian compact surface $`\mathrm{\Sigma }`$ immersed in $`\text{}\text{}^2`$ with maximum multiplicity $`\mu `$ satisfies
$$\text{Area}\text{ }(\mathrm{\Sigma })2\pi \mu ,$$
and the equality holds if and only if $`\mathrm{\Sigma }`$ is totally geodesic.
###### Remark 3
We consider the following holomorphic an anti-holomorphic immersions $`\stackrel{~}{\varphi }_i:\text{}\{\mathrm{}\}\text{}\text{}^2`$, $`i=1,2`$ given by
$$\stackrel{~}{\varphi }_1(z)=\mathrm{\Pi }(1,z,0),\stackrel{~}{\varphi }_2(z)=\mathrm{\Pi }(\overline{z}(1+\overline{z}),(1+\overline{z}),\overline{z}).$$
Then the corresponding twistor holomorphic surface with negative spin $`\varphi :\text{}\{\mathrm{}\}\text{}\text{}^2`$ is
$$\varphi (z)=\mathrm{\Pi }(|z|^2,z,(1+z)(1+|z|^2)).$$
It is easy to check that $`\varphi `$ is regular and that is embedded except at $`z=0,\mathrm{}`$ where $`\varphi `$ has a doble point. So $`\mu =2`$. As the degrees of $`\stackrel{~}{\varphi }_1`$ and $`\stackrel{~}{\varphi }_2`$ are $`1`$ and $`2`$, then $`W^{}(\varphi )=6\pi `$. This example shows that even in the family of non-complex compact surfaces of $`\text{}\text{}^2`$, Theorem 3 is not true.
Proof of Theorem 3: Using a similar reasoning like in the proof of Theorem 1, we get the following integral formula
$$_\mathrm{\Sigma }\left(|H|^2+\frac{2}{1h}\underset{i=1}{\overset{2}{}}|(e_i^{},a)|^2\right)𝑑A4\pi \mu ,$$
and the equality holds if and only if $`H=\xi /(1h)`$. In this case, using that $`\varphi `$ is a Lagrangian immersion, i.e., $`C=0`$, in (19) we have
$$\underset{i=1}{\overset{2}{}}|(e_i^{},a)|^2=1h.$$
So we finally obtain
$$_\mathrm{\Sigma }\left(|H|^2+2\right)𝑑A4\pi \mu ,$$
and the equality holds if and only if $`H=\xi /(1h)`$. Now we are going to classify Lagrangian surfaces of $`\text{}\text{}^2`$ whose mean curvature is given in the above way.
From now on we will work on the dense open subset of $`\mathrm{\Sigma }`$ defined by
$$\mathrm{\Sigma }_0=\{p\mathrm{\Sigma }/\overline{}f(\varphi (p))0.\}$$
On $`\mathrm{\Sigma }_0`$ the function $`h`$ satisfies $`0<h<1`$. First, from (17), (15) and using elementary properties of Lagrangian surfaces it follows
(21)
$$_vJ\xi ,w=\sigma (v,w),Jh2\mathrm{}((v^{},a)(a,(Jw)^{})).$$
So derivating $`JH=J\xi /(1h)`$ and using (16) and (21) it follows
$`_vJH,w={\displaystyle \frac{1}{1h}}\sigma (v,w),Jh+{\displaystyle \frac{2h}{h(1h)^2}}h,vJ\xi ,w`$
$`{\displaystyle \frac{2}{h(1h)}}h,wJ\xi ,v,`$
for any $`v,w`$ tangent to $`\mathrm{\Sigma }`$. As $`JH`$ is a closed vector field on $`\mathrm{\Sigma }`$, the first term is symmetric. Also, as the second fundamental form is also symmetric, we obtain that the others terms are symmetric too, and so
(22) $`dh\alpha =0,`$
where $`\alpha `$ is the 1-form on $`\mathrm{\Sigma }`$ given by
$$\alpha (v=v,J\xi .$$
Now we are going to prove that there exists a vector field $`X`$ tangent to $`\mathrm{\Sigma }_0`$ and functions $`a`$ and $`b`$ on $`\mathrm{\Sigma }_0`$ with $`a^2+b^2=h`$ such that
(23)
$$\overline{}f\varphi =aX+bJX.$$
So in particular this vector field $`X`$ verifies $`|X|^2=4(1h)`$.
In fact, let $`A=\{p\mathrm{\Sigma }_0/dh_p=0\}`$ and $`B=\{p\mathrm{\Sigma }_0/\alpha _p=0\}`$. If $`A=\mathrm{\Sigma }_0`$ or $`B=\mathrm{\Sigma }_0`$ the claim is trivial. Otherwise, $`A`$ and $`B`$ are proper closed subsets of $`\mathrm{\Sigma }_0`$. Now, $`dh\alpha =0`$ says that on $`\mathrm{\Sigma }_0/A`$ we can writte $`\alpha =\lambda dh`$ for certain smooth function $`\lambda `$. Taking $`X=\frac{\sqrt{1+\lambda ^2}}{\sqrt{h}}h`$, then on $`\mathrm{\Sigma }_0/A`$ we have $`\overline{}f\varphi =aX+bJX`$ for certain smooth functions $`a`$ and $`b`$ on $`\mathrm{\Sigma }_0/A`$ satisfying $`a^2+b^2=h`$. Making a similar reasoning with $`B`$ we writte on $`\mathrm{\Sigma }_0/B`$, $`\overline{}f\varphi =a^{}X^{}+b^{}JX^{}`$ with $`a^2+b^2=h`$. It is clear that, on the non-empty subset $`\mathrm{\Sigma }_0/(AB)`$, we can take $`X^{}=X`$, $`a^{}=a`$ and $`b^{}=b`$. So we prove the existence of such $`X`$ satisfying (23).
Derivating (23) with respect to a vector $`v`$, taking tangent and normal components and using (15), (16) and (17) we obtain
$`2hv+2v,XX=a,vX+a_vXbA_{JX}v`$
$`b,vX+aA_{JX}v+b_vX=0.`$
From these equations it is easy to obtain that
$$\sigma (X,X)=2\rho JX,\sigma (X,V)=2bJV,$$
for certain function $`\rho `$, where $`V`$ is any orthogonal vector field to $`X`$. From here
$$2H=\frac{\rho +b}{2(1h)}JX+cJV,$$
for certain function $`c`$. But $`\xi =bJX`$, and then
$$H=\frac{b}{1h}JX.$$
So we get that $`\rho =3b`$ and $`c=0`$. This in particular means that $`|\sigma |^2=3|H|^2`$. The Gauss equation implies that $`|H|^2+2=2K`$, and since our surface is Lagrangian, $`K^{}=K`$. So finally we get that $`|H|^2+2=K+K^{}`$, which means that our surface is twistor holomorphic with negative spin. Now the mean result in \[CU2\] finishes our proof. q.e.d.
###### Remark 4
The totally geodesic surfaces and the Whitney spheres of $`\text{}\text{}^2`$ have the property that their mean curvature vectors are given by
$$H=\frac{(\overline{}f\varphi )^{}}{1h},$$
being $`h=f\varphi `$, $`f([z])=|(z,a)|^2`$ and $`\varphi `$ the immersion. The geometric meaning of this property is the following. Let $`M=\text{}\text{}^2[a]`$, and $`g`$ the metric on $`M`$ conformal to the Fubini-Study metric defined by
$$g=\frac{1}{(1f)^2},.$$
Then it is not difficult to see that $`(M,g)`$ is a complete Riemannian manifold (with one end) and with zero scalar curvature. If $`\varphi :\mathrm{\Sigma }\text{}\text{}^2`$ is an immersion with $`\{p_1,\mathrm{},p_\mu \}=\varphi ^1([a])`$, then the mean curvatures vectors $`\widehat{H}`$ and $`H`$ of $`\varphi `$ with respect to the metric induced by $`g`$ and $`,`$ are related (see proof of Proposition 1) by
$$\frac{\widehat{H}}{1f^2}=H\frac{(\overline{}f\varphi )^{}}{1h}.$$
So the condition $`H=(\overline{}f\varphi )^{}/(1h)`$ means that the surface $`\mathrm{\Sigma }\{p_1,\mathrm{},p_\mu \}`$ is minimal in $`(M,g)`$.
To end this section we would like to remark something about the functional $`W^{}`$ acting on tori. When you consider the Willmore functional on compact surfaces of $`\text{}^4`$, there is a very famous conjecture, due to Willmore, which says that the Willmore functional on tori is bounded below by $`2\pi ^2`$ and the Clifford torus is the only torus which achieves this minimum. Since the Clifford torus is Lagrangian, Minicozzi \[M\], studied this problem in this smaller class of Lagrangian tori.
In our case, we will also called Clifford torus to the following torus $`T`$ embedded in $`\text{}\text{}^2`$ and defined by
$$T=\{\mathrm{\Pi }(z)\text{}\text{}^2/|z_i|^2=\frac{1}{3},i=1,2,3\}.$$
It is easy to check that $`T`$ is a minimal Lagrangian torus with area $`4\pi ^2/3\sqrt{3}`$. So its Willmore functional is $`W^{}(T)=8\pi ^2/3\sqrt{3}`$.
For complex tori of $`\text{}\text{}^2`$, the Willmore functional $`W^{}`$ take the value $`2\pi d`$, where $`d`$ is the degree of the torus. So, as $`d`$ must be non smaller that $`3`$, we obtain that in this family $`W^{}6\pi >8\pi ^2/3\sqrt{3}`$.
For twistor holomorphic tori in $`\text{}\text{}^2`$ with negative spin, Proposition 4 says that the Willmore functional $`W^{}`$ satisfied again $`W^{}=2\pi (d_1+d_2)`$. But again the corresponding holomorphic and antiholomorphic curves have genus one. So $`W^{}6\pi `$.
In the family of tori of $`\text{}\text{}^2`$ with non-zero parallel mean curvature vector, very recently Kenmotsu and Zhou in \[KZ\] have proved that they are Lagrangian and flat, and then, up to isometries they can be parametrized by $`T_{r_1,r_2,r_3}/r_1r_2r_3>0,r_1^2+r_2^2+r_3^2=1`$, where
$$T_{r_1,r_2,r_3}=\{\mathrm{\Pi }(z)\text{}\text{}^2/|z_i|^2=r_i^2,i=1,2,3\}.$$
It is clear that $`W^{}(T_{r_1,r_2,r_3})8\pi ^2r_1r_2r_38\pi ^2/3\sqrt{3}`$, and the equality is only achieved by the Clifford torus.
Also, in \[CU1\], Castro and Urbano classified minimal Lagrangian tori of $`\text{}\text{}^2`$ invariant under a $`1`$-parameter group of holomorphic isometries. This family of tori is described in terms of elliptic functions, and it is not a complicated exercise to check that the Willmore functional $`W^{}`$ on these tori satisfy $`W^{}8\pi ^2/3\sqrt{3}`$, with equality only for the Clifford torus.
These considerations make reasonable the following conjecture:
> The Clifford torus achieves the minimum of the Willmore functional $`W^{}`$ either amongst all tori in $`\text{}\text{}^2`$ or amongst all Lagrangian tori in $`\text{}\text{}^2`$.
| Departamento de Geometría y Topología |
| --- |
| Universidad de Granada |
| 18071 Granada |
| SPAIN |
| e-mails:smontiel@goliat.ugr.es and furbano@goliat.ugr.es | |
warning/0002/hep-ph0002196.html | ar5iv | text | # References
HZPP-0002
Feb. 20, 2000
Thermal Equilibration and Ellipsoidal Expansion
of Rotationally-Symmetrical Longitudinal Flow
in Relativistc Heavy Ion Collisions<sup>1</sup><sup>1</sup>1Work supported in part by the NSFC under project 19775018.
Feng Shengqin, Liu Feng and Liu Lianshou
Institute of Particle Physics, Huazhong Normal University, Wuhan, 430079, China
ABSTRACT
A rotationally-symmetrical ellipsoidal flow model is proposed for the relativistic heavy-ion collisions and compared with the 14.6 A GeV/$`c`$ Si-Al and 10.8 A GeV/$`c`$ Au-Au collision data. The large stopping in the heavier collision system and heavier produced particles is accounted for by using the ellipsoidal flow picture. The central dip in the proton and deuteron rapidity distributions for Si-Al collision are reproduced.
I. Introduction
The experimental finding that colliding nuclei are not transparent but undergo a violent reaction in central collisions represents one of the major motivations for the study of ultra-relativistic heavy ion collisions at the CERN/SPS, BNL/AGS and also at the future BNL/RICH and CERN/LHC. Of central importance is the ability of understanding to what extent the nuclear matter has been compressed and heated.
The study of collective flow in high energy nuclear collisions has attracted increasing attention from both experimental<sup></sup> and theoretical<sup></sup> point of view. The rich physics of longitudinal and transverse flow is due to their sensitivity to the system evolution at early time. The expansion and cooling of the heated and highly compressed matter could lead to considerable collectivity in the final state. Due to the high pressure, particles might be boosted in the transverse and longitudinal directions. The collective expansion of the system created during a heavy-ion collision implies space-momentum correlation in particle distributions at freeze-out.
The experimental data of the rapidity distributions of produced particles in 14.6 A GeV/$`c`$ Si-Al collisions have been ultilized to study the collective expansion using a cylindrically-symmetrical flow model<sup></sup>. The model results fit well with the experimental distribution of pion, but is too narrow in the case of heavier particles proton and deuteron. In particular, the central dip, which can be clearly seen in the distribution of proton, is failed to be reproduced.
More recently, E877 Collaboration<sup></sup> has published their data for 10.8 A GeV/$`c`$ Au-Au collisions, which provide a good chance to compare the stopping power in the collision systems of different sizes. The possible central peak of the rapidity distribution of proton at around mid-rapidity, which was obtained through extrapolating the experimental data to mid-rapidity using RQMD model<sup></sup>, has been taken as an evidence for the increasing of stopping power, but the reliability of this extrapolation is model-dependent.
It has been shown earlier<sup></sup> that the ellipsoidal expansion is a simple way to take the nuclear stopping into account. In the present paper we propose a rotationally-symmetrical ellipsoidal flow model to describe the space-time evolusion in relativistic heavy ion collisions. The large stopping in the heavier collision system and heavier produced particles is described by using this picture. The central dips in the proton and deuteron rapidity distributions for Si-Al collisions are reproduced.
In section II the rotationally-symmetrical ellipsoidal flow model is formulated. The results of the model are given and compared with the experimental data in section III. A short summary and conclusions are given in section IV. In order to avoid the complexity in the production of strange particles and concentrate on the expansion of the system, we will discuss in this paper only normal non-strange particles —— pions, protons and deuterons.
II. Rotationally-symmetrical ellipsoidal flow
Firstly, let us briefly recall the fireball scenario of relativistic heavy ion collisions.
Since the temperature at freeze-out exceeds 100 MeV, the Boltzmann approximation is used. Transformed into rapidity $`y`$ and transverse momentum $`p_t`$ this implies<sup></sup>:
$$E\frac{\mathrm{d}^3n}{\mathrm{d}^3p}E\mathrm{e}^{(\mathrm{E}/\mathrm{T})}=\mathrm{m}_\mathrm{t}\mathrm{cosh}(\mathrm{y})\mathrm{e}^{(\mathrm{m}_\mathrm{t}\mathrm{cosh}(\mathrm{y})/\mathrm{T})}$$
(1)
Here $`m_t=\sqrt{m^2+p_t^2}`$ is the transverse mass, $`m`$ is the mass of the produced particles at freeze-out.
The rapidity is defined as $`y=\mathrm{tanh}^1(p_l/E)`$, where $`p_l`$ is the longitudinal momentum of the produced particle. Substituting into Eq.(1) and integrating over $`m_t`$, we get the rapidity distribution of the isotropic thermal source,
$$\frac{\mathrm{d}n_{\mathrm{iso}}}{\mathrm{d}y}\frac{m^2T}{(2\pi )^2}(1+2\xi _0+2\xi _0^2)\mathrm{e}^{(1/\xi _0)}.$$
(2)
Here $`\xi _0=T/m\mathrm{cosh}(y)`$.
However, the momentum distribution of the measured particles is certainly not isotropic. It is privileged in the direction of the incident nuclei. This is because the produced hadrons still carry their parent’s kinematic information, making the longitudinal direction more populated than the transverse ones. The simplest way<sup></sup> to account for this anisotropy is to add the contribution from a set of fire-balls, sketched schematically in Fig.1 as dashed circles, with centers located uniformly in the rapidity region \[$`y_0,y_0`$\]. The corresponding rapidity distribution is obtained through changing the $`\xi _0`$ in Eq,(2) into $`\xi =T/m\mathrm{cosh}(yy^{})`$ and integrating over $`y^{}`$ from $`y_0`$ to $`y_0`$:
$$\frac{\mathrm{d}n_{\mathrm{cyl}}}{\mathrm{d}y}_{y_0}^{y_0}dy^{}\frac{m^2T}{(2\pi )^2}(1+2\xi +2\xi ^2)\mathrm{e}^{(1/\xi )},$$
(3)
$`\xi =T/m\mathrm{cosh}(yy^{})`$. Equivalently, we can also use the angular variable $`\mathrm{\Theta }`$ defined by $`\mathrm{\Theta }=2\mathrm{tan}^1\mathrm{exp}(y^{})`$, and change the integration variable in Eq.(3) to $`\mathrm{\Theta }`$, cf. the solid circle and lines in Fig.1.
This simple approach fits the rapidity distribution of pions well but failed to reproduce the central dip in heavier produced particles, which is clearly seen in the experimental distribution of protons and has some evidence in the distribution of deuterons.
Note that in this model the longitudinal and transverse expansions of the system are totally independent. This is a crude approximation. A more reasonable picture is an ellipsoidal expansion. For simplicity the rotational symmetry arround the longitudinal direction is still assumed, but the emission angle is now
$$\theta =\mathrm{tan}^1(e\mathrm{tan}\mathrm{\Theta }),$$
(4)
where $`e`$ ($`0e1`$) is the ellipticity, cf. Fig.2.
In this model the ellipticity parameter $`e`$ represents the degree of anisotropy of flow in the transverse and longitudinal direction. The smaller is $`e`$, the more anisotropic is the flow. The nuclear stopping can be taken into account in this way.
Subsituting Eq.(4) together with $`y_\mathrm{e}^{}=\mathrm{ln}\mathrm{tan}(\theta /2)`$ into Eq.(3), the rapidity distribution is obtained:
$`{\displaystyle \frac{\mathrm{d}n_{\mathrm{ellip}}}{\mathrm{d}y}}`$ $`=`$ $`eKm^2T{\displaystyle _{\theta _{\mathrm{min}}}^{\theta _{\mathrm{max}}}}\left(1+{\displaystyle \frac{2T}{m\mathrm{cosh}(yy_\mathrm{e}^{})}}+{\displaystyle \frac{2T^2}{m^2\mathrm{cosh}^2(yy_\mathrm{e}^{})}}\right)`$ (5)
$`\times \mathrm{exp}(m\mathrm{cosh}(yy_\mathrm{e}^{})/T)Q(\theta )\mathrm{d}\theta ,`$
$$y_\mathrm{e}^{}=\mathrm{ln}\mathrm{tan}(\theta /2),Q(\theta )=\frac{1}{\sqrt{e^2+\mathrm{tan}^2\theta }|\mathrm{cos}\theta |\mathrm{sin}\theta }.$$
(6)
Here $`\theta _{\mathrm{min}}=2\mathrm{tan}^1(e^{y_{\mathrm{e0}}^{}})`$, $`\theta _{\mathrm{max}}=2\mathrm{tan}^1(e^{y_{\mathrm{e0}}^{}})`$. $`y_{\mathrm{e0}}^{}`$ is the rapidity limit which confines the rapidity interval of ellipsoidal flow. We treat it together with the ellipticity $`e`$ as two free papameters of the model to fit the amount of flow and stopping required by the data.
III. Comparison with experiments
The rapidity distributions of pion, proton and deuteron for 14.6 A GeV/$`c`$ Si-Al collisions<sup></sup>, are given in Fig.3 ($`a,b`$ and $`c)`$. The dashed, dotted and solid lines correspond to the results from isotropical thermal model, cylindrically-symmetrical flow model and rotationally-symmetrical ellipsoidal flow (RSEF) model respectively. The rapidity limit $`y_{\mathrm{e0}}^{}`$ and the ellipticity $`e`$ used in the calculation are listed in Table I. The rapidity limit $`y_0^{}`$ used in the cylindrically-symmetrical flow model of Ref. is also listed for comparison. The parameter $`T`$ is chosen to be 0.12 GeV following Ref..
Table I The value of model-parameters
| | Si-Al Collisions | | | Au-Au Collisions | |
| --- | --- | --- | --- | --- | --- |
| Parameter | $`\pi `$ | p | d | $`\pi `$ | p |
| $`e`$ | 0.28 | 0.52 | 0.56 | 0.32 | 0.58 |
| $`y_{\mathrm{e0}}^{}`$ | 1.35 | 1.35 | 1.35 | 1.05 | 1.05 |
| $`y_0^{}`$ | 1.15 | 1.15 | 1.15 | | |
It can be seen from the figures that the RSEF model reproduces the central dip of rapidity distribution of heavier particles (proton and deuteron) in coincidence with the experimental findings, while for light particles (pions) there is a plateau instead of dip at central rapidity. Note that the appearance or disappearance of central dip is insensitive to the rapidity limit $`y_{\mathrm{e0}}^{}`$ but depends strongly on the magnitude of the ellipicity $`e`$ and the mass $`m`$ of the produced particles. For the heavier particles (proton and deuteron) a central dip appear for $`e<0.8`$, but for light particles (pions) there is no dip even when $`e`$ is as small as 0.28, cf. Table I and Fig. 3.
It can also be seen from Table I that $`e_\mathrm{d}>e_\mathrm{p}>e_\pi `$. It means that the system is less alongated for proton and deuteron than for pion. This describes nuclear stopping.
On the other hand, the width of the rapidity distributions are mianly controlled by the parameter $`y_{\mathrm{e0}}^{}`$. The value $`y_{\mathrm{e0}}^{}=1.35`$, a little bigger than 1.15 used in the cylindrically-symmetrical flow model of Ref. can account for the wide distriution of heavier particles (protons and deuterons) and at the same time fits the pion-distribution well.
In Fig.4 are shown the rapidity distributions of pions and protons for Au+Au collisions at 10.8 A GeV/$`c`$<sup></sup>. The solid and dashed lines correspond to the results of RSEF model (with parameters listed in Table I) and cylindrically-symmetrical flow model respectively. The latter are obtained also using RSEF with the same rapidity limit $`y_0^{}`$ as the $`y_{\mathrm{e0}}^{}`$ listed in Table I but with ellipicity $`e=1`$. The histogram is the result from the RQMD model.
It can be seen from Fig.4 that in the RSEF model there is a shallow dip (plateau) in the central rapidity of the distribution of proton, instead of a central peak as predicted by the cylindrically-symmetrical flow model. However, the presently available experimental data are restricted to the large rapidity. The peak at central rapidity is the extrapolation of data using RQMD and is model dependent. It is intersting to see whether the prediction of a central dip (plateau) or a central peak will be observed in future experiments.
Comparing the parameter values for Si-Al (smaller colliding nuclei) and Au-Au (larger colliding nuclei) collisions listed in Table I, it can be seen that the rapidity limit $`y_{\mathrm{e0}}^{}`$ is smaller and the elliticity $`e`$ is bigger for the larger colliding nuclei than for the smaller ones. Both of these two show that the hadronic system formed from the larger colliding nuclei is less alongated, i.e. there is stronger nuclear stopping in the collision of larger nuclei.
IV. Summary and Conclusions
In high energy heavy-ion collisions, due to the transparency of the nucleus the participants will not lose the historical vestiges and the produced hadrons will carry some of their parent’s memory of motion, leading to the unequivalence in longitudinal and transverse directions. So it is reasonable to assume that the flow of produced particle is privileged in the longitudinal direction. This picture has been used by lots of models<sup></sup>. Here we should mention two thermal and hydrodynamic models, one is the the boost-invariant longitudinal expansion model postulated by Bjorken <sup></sup> which can explain such an anisotropy already at the level of particle production. This model has been formulated for asymptotically high energies, where the rapidity distribution of produced particles establishes a plateau at midrapidity. The second model is the cylindrical symmetry flow model postulated first by Schnedermann, Sollfrank and Heinz<sup></sup> which account for the anisotropy of longitudinal and transverse direction by adding the contribution from a set of fire-balls with centers located uniformly in the rapidity region \[-$`y_0^{}`$,$`y_0^{}`$\] in the longitudinal direction, sketched schematically in Fig.1 as dasheded circles. In this model the centers of fire-balls distribute uniformly in rapidity, and so it gives the picture that longitudinal and transverse expansion are totally independent. It can account for the wider rapidity distribution when comparing to the prediction of the pure thermal isotropically model but failed to reproduce the central dip in the proton and deuteron rapidity distributions.
In this paper, we propose an ellipsoidal expansion model with rotationally-symmetrical longitudinal flow (rotationally-symmetrical ellipsoidal flow model, RSEF) which realizes that the centers of fire-balls are distributed un-uniformly in a rotationally-symmetrical ellipsoidal shape around the longitudinal direction. The ellipticity parameter $`e`$ can account for the extent of anisotropy of phase space in transverse and longitudinal expansion. The central dip in the proton and deuteron rapidity distributions and the central peak in the pion distribution are well reproduced simultaneously from this model.
It is found that the depth of the central dip of the heavier particle distributions depends strongly on the magnitude of the ellipicity $`e`$. In other words, the anisotropy in transverse and longitudinal directions of the ellipsiod of phase space, which is given by the ellipticity $`e`$, determines also the depth of the central dip for heavier particles.
Through comparing the feature of collision systems of different size, we found that the maximum flow velocities are larger for the lighter collision systems than the heavier ones, which suggests, together with smaller $`e`$, a larger stopping in the larger collision system.
Figure captions
Fig.1 Schematic sketch of the cylindrically-symmetrical flow model.
Fig.2 Schematic sketch of the Rotationally-symmetrical ellipsoidal flow model
Fig.3 Rapidity distributions for central 14.6 A GeV/$`c`$ Si+Al collisions. Open circles — data from Si+Al collisions ; Dashed lines — isotropical thermal model; Dotted lines — cylindrically-symmetrical flow model; solid line — rotationally-symmetrical ellipsoidal flow (RSEF) model. Temperature $`T=0.12GeV`$. Fig.3 ($`a`$), ($`b`$) and ($`c`$) are for the pion, proton and deuteron distributions respectively.
Fig.4 Rapidity distributions for pions and protons in central Au+Au collisions with 10.8 A GeV/$`c`$. Full circles represent measured data, open circles reflected data. The solid line is our calculation using the RSEF model. The histogram shows the results from RQMD calculations and the dotted line is the results from the prediction of cylinderically-symmetrical flow model. The temperature $`T=0.14GeV`$. Fig.4 ($`a`$) and Fig.4 ($`b`$) are for the pion and proton distributions respectively.
Fig. 1
Fig. 2
Fig. 3
Fig. 4 |
warning/0002/math0002112.html | ar5iv | text | # Appendix: Polynomials arising from the tautological ring
1. Statement of results. For positive integers $`g`$ and $`k`$ define
$$P_g(k)=\underset{l=1}{\overset{k}{}}\frac{(k1)!}{(kl)!}\frac{1}{k^l}\underset{m=1}{\overset{l}{}}(1)^{lm}\left(\genfrac{}{}{0pt}{}{l}{m}\right)\frac{m^{2g+l1}}{(2g+l1)!}$$
$`1`$
(the inner sum here is a Stirling number), e.g. for $`k3`$,
$$P_g(1)=\frac{1}{(2g)!},P_g(2)=\frac{2^{2g1}+g}{(2g+1)!},P_g(3)=\frac{2(3^{2g+1}+2^{2g+2}g+6g^2+5g)}{9(2g+2)!}.$$
A property of the function $`P_g`$ which is far from obvious—and is false if the number $`2g1`$ on the right-hand side of (1) is replaced by an even number—is that it is a polynomial in $`k`$ for each fixed $`g`$, the first values being
$$P_1(k)=\frac{1}{2},P_2(k)=\frac{k}{24},P_3(k)=\frac{3k^2k}{1440},P_4(k)=\frac{9k^38k^2+2k}{120960}.$$
This fact was discovered and proved in the preceding article by Faber and Pandharipande by an indirect argument in which the coefficients of the polynomials $`P_g(k)`$ were interpreted as intersection numbers of certain cycles in the moduli space of curves of genus $`g`$. Here we will give a more direct combinatorial proof and will also obtain alternative expressions for the polynomial $`P_g(k)`$ and explicit formulas for its highest and lowest coefficients. The formulas for the coefficients of $`k^{g1}`$, $`k^{g2}`$, $`k^2`$ and $`k^1`$ were quoted in Section 5.2 of .
###### Theorem 1
(i) For each integer $`g1`$, the function $`P_g(k)`$ defined by $`(1)`$ is a polynomial of degree $`g1`$ in $`k`$. (ii) Write $`P_g(k)=_{i=0}^{g1}c_{g,i}k^i`$. Then for fixed $`j0`$ and $`g>j`$ we have
$$c_{g,gj1}=\frac{(g1)!}{2^g(2g1)!}C_j(g),$$
$`2`$
where
$$C_0(g)=1,C_1(g)=\frac{g(g2)}{9},C_2(g)=\frac{g(g3)(5g^29g+1)}{810},$$
and in general $`C_j(g)`$ is a polynomial of degree $`2j`$ with leading coefficient $`{\displaystyle \frac{(1/9)^j}{j!}}`$. (iii) For fixed $`i0`$ and $`g>i+1`$ we have $`c_{g,i}=\underset{j=0}{\overset{i}{}}\gamma _{i,j}(g)\beta _{2gj1},`$ where $`\beta _n={\displaystyle \frac{B_n}{n!}}`$ $`(B_n`$ = $`n`$th Bernoulli number$`)`$ and $`\gamma _{i,j}(g)`$ is a polynomial of degree $`ij`$. In particular $`(`$for $`g>2)`$
$$c_{g,0}=0,c_{g,1}=\frac{1}{2}\beta _{2g2},c_{g,2}=\frac{g}{2}\beta _{2g2},c_{g,3}=\frac{g(g+2)}{6}\beta _{2g2}+\frac{1}{24}\beta _{2g4}.$$
Parts (i) and (ii) of Theorem 1 are equivalent to the following amusing result. Let us define numbers $`A(g,n)`$ ($`g1,n0`$) by
$$\underset{n=0}{\overset{\mathrm{}}{}}A(g,n)x^n=e^x\underset{k=0}{\overset{\mathrm{}}{}}P_g(k)\frac{x^k}{k!}$$
or equivalently by
$$A(g,n)=\underset{k=0}{\overset{n}{}}\frac{(1)^{nk}}{k!(nk)!}P_g(k),P_g(k)=\underset{n=0}{\overset{k}{}}\frac{k!}{(kn)!}A(g,n).$$
$`3`$
###### Theorem 2
The numbers $`A(g,n)`$ vanish for $`ng`$. For $`ng1`$ we have
$$A(g,n)=\frac{(g1)!}{2^g(2g1)!}C_{gn1}^{}(gn1),$$
$`4`$
where
$$C_0^{}(h)=1,C_1^{}(h)=\frac{7h^2+5h}{18},C_2^{}(h)=\frac{245h^4+594h^3+283h^242h}{3240},$$
and in general $`C_r^{}(h)`$ is a polynomial of degree $`2r`$ in $`h`$ with leading coefficient $`{\displaystyle \frac{(7/18)^r}{r!}}`$.
This theorem, as well as more general results concerning the numbers
$$A_\nu (g,n)=\underset{k=1}{\overset{n}{}}\frac{(1)^{nk}}{k!(nk)!}k^\nu P_g(k)(\nu 0),$$
which are related to part (iii) of Theorem 1, will be proved in §3. For instance, we have
$$A_1(g,n)=\frac{(1)^{n1}}{2n!}\beta _{2g2},A_2(g,n)=A_1(g,n)\left(g\underset{k=1}{\overset{n}{}}\frac{1}{k}\right)(n+2g>2).$$
$`5`$
To state the remaining results, and for the proofs, we will need some more notation. As in , we write $`𝒞(x^n,f(x))`$ to denote the coefficient of $`x^n`$ in a power series $`f(x)`$ and $`h_n(\alpha _1,\mathrm{},\alpha _l)=𝒞(x^n,_{i=1}^l(1\alpha _ix)^1)`$ for the full symmetric function of degree $`n`$ in variables $`\alpha _1,\mathrm{},\alpha _l`$. For any integer $`n0`$, we define $`S_n(l)`$ by
$$S_n(l)=𝒞(x^n,\left(\frac{e^x1}{x}\right)^l).$$
$`6`$
For $`l`$ we have the formulas
$$\frac{(n+l)!}{l!}S_n(l)=\frac{1}{l!}\underset{m=0}{\overset{l}{}}(1)^{lm}\left(\genfrac{}{}{0pt}{}{l}{m}\right)m^{n+l}=h_n(1,\mathrm{\hspace{0.17em}2},\mathrm{},l)=𝔖_{n+l}^{(l)},$$
where $`𝔖_{n+l}^{(l)}`$ denotes the Stirling number of the second kind (=number of partitions of a set of $`n+l`$ elements into $`l`$ non-empty subsets). In particular, equation (1) can be written
$$P_g(k)=\underset{l=1}{\overset{k}{}}\frac{(k1)!}{(kl)!}\frac{1}{k^l}S_{2g1}(l).$$
$`7`$
However, $`S_n(l)`$ is a polynomial (of degree $`n`$) in $`l`$, the first values being
$$S_0(l)=\mathrm{\hspace{0.33em}1},S_1(l)=\frac{l}{2},S_2(l)=\frac{3l^2+l}{24},S_3(l)=\frac{l^3+l^2}{48},\mathrm{},$$
so it makes sense for any complex value of $`l`$. For $`l=0`$ we clearly have $`S_n(l)=0`$ for all $`n>0`$. For $`l=1`$ we have $`S_n(l)=\beta _n`$ by definition, where $`\beta _n=B_n/n!`$ as in Theorem 1, and more generally $`S_n(l)`$ for fixed negative $`l`$ is a finite combination of Bernoulli numbers (Lemma 3 below), the first three cases for $`n`$ odd being
$$S_{2g1}(1)=0,S_{2g1}(2)=\beta _{2g2},S_{2g1}(3)=\frac{3}{2}(2g3)\beta _{2g2}(g3).$$
Using these numbers, we can now state a formula for $`P_g(t)`$ as a power series in $`t`$.
###### Theorem 3
Define the function $`S_n(l)`$ by eq. $`(6)`$. Then for each integer $`g1`$ we have
$$P_g(t)=\underset{r=1}{\overset{\mathrm{}}{}}\frac{S_{2g1}(r)t^{r1}}{(1+t)\mathrm{}(r+t)}[[t]].$$
$`8`$
In particular, the power series on the right-hand side of $`(8)`$ is in fact a polynomial in $`t`$.
This theorem gives an alternative definition of the polynomials $`P_g(t)`$, but, as with (1), the polynomial property is not clear from this definition, and is not true if the index $`2g1`$ on the right-hand side of (8) is replaced by an even number.
The next result gives a closed form expression for the generating function of the $`P_g(t)`$ as an integral. This looks less elementary than the preceding results, but has the advantage of making it obvious that $`P_g`$ is a polynomial.
###### Theorem 4
Define a power series $`F(x)`$ by
$`F(x)`$ $`={\displaystyle \frac{\mathrm{sinh}x/2}{x/2}}\mathrm{exp}\left({\displaystyle \frac{x/2}{\mathrm{tanh}x/2}}1\right)=\mathrm{exp}\left({\displaystyle \underset{n=2}{\overset{\mathrm{}}{}}}{\displaystyle \frac{n+1}{n}}\beta _nx^n\right)`$ $`9`$
$`=\mathrm{\hspace{0.33em}1}+{\displaystyle \frac{1}{8}}x^2+{\displaystyle \frac{7}{1152}}x^4+{\displaystyle \frac{61}{414720}}x^6+\mathrm{}.`$
Then the $`P_g(t)`$ are given by the generating function identity
$$\underset{g=1}{\overset{\mathrm{}}{}}P_g(t)x^{2g1}=\frac{1}{2}F(x)^t_0^xF(y)^t𝑑y.$$
$`10`$
The polynomiality of the functions $`P_g(t)`$ follows immediately because we can rewrite the generating series identity (10) in the form
$$P_g(t)=\underset{n=0}{\overset{g1}{}}\frac{p_{g1n}(t)p_n(t)}{2(2n+1)},$$
where $`p_n(t)`$ denotes the coefficient of $`x^{2n}`$ in $`F(x)^t`$, which is clearly a polynomial in $`t`$ of degree $`n`$. Equation (10) is also equivalent to the following recursion for the polynomials $`P_g`$.
###### Theorem 5
The polynomials $`P_g(t)`$ can be given recursively by the formulas
$$P_1(t)=\frac{1}{2},P_g(t)=\frac{t}{2g1}\underset{n=1}{\overset{g1}{}}(2n+1)\beta _{2n}P_{gn}(t)(g2).$$
$`11`$
The final result describes the coefficients $`c_{g,i}`$ (which are actually the numbers of interest, since it is they, and not the values of the polynomial $`P_g(k)`$, which occur in as intersection numbers) via a generating series with respect to the variable $`g`$ rather than $`i`$. We begin with the well-known fact that the inverse power series of $`x=ye^y`$ is given by $`y=_{k1}k^{k1}x^k/k!`$. A simple generalization of this says that the power series
$$Q_i(y)=(1)^i\underset{k=1}{\overset{\mathrm{}}{}}\frac{k^{k1i}}{k!}(ye^y)^k$$
$`12`$
is in fact a polynomial in $`y`$ for every integer $`i0`$, the first few values being
$$Q_0(y)=y,Q_1(y)=\frac{1}{2}y^2y,Q_2(y)=\frac{1}{6}y^3\frac{3}{4}y^2+y,Q_3(y)=\frac{1}{24}y^4\frac{11}{36}y^3+\frac{7}{8}y^2y.$$
The polynomials $`Q_i(y)`$ can also be defined and computed using the recursion
$$Q_0(y)=y,Q_{i+1}(y)=_0^y\frac{x1}{x}Q_i(x)𝑑x(i0)$$
$`13`$
or the generating function identity
$$\underset{i=0}{\overset{\mathrm{}}{}}Q_i(y)t^i=\underset{r=1}{\overset{\mathrm{}}{}}\frac{t^{r1}y^r}{(1+t)\mathrm{}(r+t)}.$$
$`14`$
The following theorem provides yet another characterization of these polynomials and a new generating function for the rational numbers $`c_{g,i}`$.
###### Theorem 6
(i) The polynomial $`Q_i`$ is, up to a constant, the unique polynomial with constant term $`0`$ and degree $`i+1`$ satisfying
$$Q_i\left(\frac{x}{1e^x}\right)Q_i\left(\frac{x}{e^x1}\right)=\text{O}\left(x^{2i+1}\right)(x0).$$
$`15`$
(ii) For all integers $`g1`$ and $`i0`$ we have
$$c_{g,i}=𝒞(x^{2g1},Q_i\left(\frac{x}{1e^x}\right)).$$
$`16`$
The proof of this theorem will be given in §5.
2. Polynomials defined by functional equations. We begin by giving two simple (and well-known) lemmas which will be used several times in the sequel.
###### Lemma 1
Let $`r`$ be a non-negative integer and $`z`$ be a variable. Then
$$\frac{1}{z(z1)\mathrm{}(zr)}=\underset{m=0}{\overset{r}{}}\frac{(1)^{rm}}{m!(rm)!}\frac{1}{zm}.$$
Proof. Compare residues on the two sides. $`\mathrm{}`$
###### Lemma 2
Let $`z`$ and $`y`$ be two free variables. Then
$$\underset{r=0}{\overset{\mathrm{}}{}}\frac{y^r}{z(z1)\mathrm{}(zr)}=e^y\underset{m=0}{\overset{\mathrm{}}{}}\frac{y^m}{m!}\frac{1}{zm}.$$
Proof. The equality of the coefficients of $`y^r`$ is Lemma 1. Alternatively, we can prove the identity directly by observing that it holds for $`y=0`$ and that
$`{\displaystyle \frac{}{y}}\left(y^ze^y\text{LHS}\right)`$ $`={\displaystyle \underset{r=0}{\overset{\mathrm{}}{}}}\left({\displaystyle \frac{e^yy^{rz}}{z\mathrm{}(zr)}}{\displaystyle \frac{e^yy^{rz1}}{z\mathrm{}(zr+1)}}\right)`$
$`=e^yy^{z1}={\displaystyle \frac{}{y}}(y^ze^y\text{RHS}).\mathrm{}`$
We now prove several results saying that certain generating functions which are a priori power series are in fact polynomials. We denote by $`(x)_n`$ the ascending Pochhammer symbol $`x(x+1)\mathrm{}(x+n1)`$.
###### Proposition 1
For each $`n0`$, there is a unique polynomial $`B_n(z,y,t)`$ in three variables $`z`$$`y`$ and $`t`$, of degree $`n1`$, satisfying the identity
$$(zt)B_n(z,y,t)yB_n(z1,y,t)=(z)_n\underset{m=0}{\overset{n}{}}\left(\genfrac{}{}{0pt}{}{n}{m}\right)y^m(t)_{nm}.$$
$`17`$
Examples. For $`0n3`$ the polynomials $`B_n`$ are given by
$$B_0=0,B_1=1,B_2=z+y+t+1,B_3=(z+1)(z+2)+(y+t)(z+y+t)+y+3t.$$
Proof. The recursion is equivalent to the functional equation
$$(zt)(z,y,t,u)y(z1,y,t,u)=(1u)^ze^{yu}(1u)^t$$
$`18`$
for the generating function $`(z,y,t,u)=\underset{n=0}{\overset{\mathrm{}}{}}B_n(z,y,t){\displaystyle \frac{u^n}{n!}}`$. The solution of this is
$$(z,y,t,u)=(1u)^t_0(zt,y,u),$$
$`19`$
where $`_0(z,y,u)`$ ($`=(z,y,0,u)`$) satisfies the simpler functional equation
$$z_0(z,y,u)y_0(z1,y,u)=(1u)^ze^{yu}.$$
$`20`$
Write $`_0(z,y,u)`$ as $`_{r0}\beta _r(z,u)y^r`$. Then (20) is equivalent to
$$z\beta _r(z)=\{\begin{array}{cc}(1u)^z1\hfill & \text{if }r=0\text{,}\hfill \\ \beta _{r1}(z1,u)\frac{u^r}{r!}\hfill & \text{if }r>0\text{,}\hfill \end{array}$$
which can be solved by induction on $`r`$ to give the closed formula
$$\beta _r(z,u)=\frac{(1u)^{z+r}}{z(z1)\mathrm{}(zr)}\underset{s=0}{\overset{r}{}}\frac{1}{z(z1)\mathrm{}(zs)}\frac{u^{rs}}{(rs)!}.$$
$`21`$
Using Lemma 1 we can rewrite (21) as
$$\beta _r(z,u)=\underset{m=0}{\overset{r}{}}\frac{(1)^{rm}}{m!(rm)!}\frac{(1u)^{z+r}(1u)^{rm}}{zm}$$
or, going back to the generating function $`_0`$,
$$_0(z,y,u)=e^{y(u1)}\underset{m=0}{\overset{\mathrm{}}{}}\frac{y^m}{m!}\frac{(1u)^{z+m}\mathrm{\hspace{0.17em}1}}{zm}.$$
$`22`$
Substituting this into (19) gives the generating series $`(z,y,t,u)`$ in the form
$$(z,y,t,u)=e^{y(u1)}\underset{m=0}{\overset{\mathrm{}}{}}\frac{y^m}{m!}\frac{(1u)^{z+m}(1u)^t}{zmt}.$$
$`23`$
To see that the coefficients of this with respect to $`u`$ are polynomials, we rewrite (22) as
$`_0(z,y,u)`$ $`=e^{y(u1)}{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{y^m}{m!}}{\displaystyle _0^u}(1v)^{z+m1}𝑑v`$ $`24`$
$`={\displaystyle _0^u}(1v)^{z1}e^{y(uv)}𝑑v`$
$`={\displaystyle \underset{p=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{q=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(z+1)_py^q}{p!q!}}{\displaystyle _0^u}v^p(uv)^q𝑑q`$
$`={\displaystyle \underset{p=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{q=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(z+1)_py^qu^{p+q+1}}{(p+q+1)!}}`$
(the last equality by Euler’s beta integral). Now substituting this into (19) and using the binomial expansion of $`(1u)^t`$ gives the explicit polynomial expression
$$\frac{B_n(z,y,t)=}{p,q,l0p+q+l+1=n\left(\genfrac{}{}{0pt}{}{n}{l}\right)(zt+1)_p(t)_ly^q[z,y,t].\mathrm{}}$$
$`25`$
Of course, we could have simply written down (25) and checked that it satisfies the identity (17); we gave the full derivation for clarity and because some of the formulas found along the way will be needed below. In particular, from (24) and (19) we get the integral representation
$$(z,y,t,u)=(1u)^t_0^u(1v)^{z+t1}e^{y(uv)}𝑑v$$
$`26`$
and from (21) and (19), or (23) and Lemma 2, we get the generating function identity
$$(z,y,t,u)=\underset{r=0}{\overset{\mathrm{}}{}}\frac{(1u)^{z+r}y^r}{(zt)\mathrm{}(ztr)}(1u)^te^{uy}\underset{r=0}{\overset{\mathrm{}}{}}\frac{y^r}{(zt)\mathrm{}(ztr)}.$$
$`27`$
This can also be obtained from (26) by writing $`_0^u=_u^1+_0^1`$ (for $`\mathrm{}(zt)<0`$).
We now consider the specialization of the above functions to the case $`y=t`$.
###### Proposition 2
For each $`n0`$, there is a unique polynomial $`\widehat{B}_n(z,t)`$ in $`z`$ and $`t`$, of degree $`[(n1)/2]`$ in $`t`$, satisfying the identity
$$(zt)\widehat{B}_n(z,t)+t\widehat{B}_n(z1,t)=(z)_n\underset{m=0}{\overset{n}{}}\left(\genfrac{}{}{0pt}{}{n}{m}\right)(t)^m(t)_{nm}.$$
$`28`$
Examples. For $`0n4`$ we have
$$\widehat{B}_0=0,\widehat{B}_1=1,\widehat{B}_2=z+1,\widehat{B}_3=2t+(z+1)_2,\widehat{B}_4=3(z+3)t+(z+1)_3.$$
Proof. Since (28) is just the specialization of (17) to $`y=t`$, its solution is of course given simply by $`\widehat{B}_n(z,t)=B_n(z,t,t)`$; what we have to show is that the degree with respect to $`t`$ drops by a factor of 2 under this specialization. To do this we expand $`(1v)^{z1}`$ in the integral representation (26) by the binomial theorem and change $`v`$ to $`uv`$ to get
$$(z,t,t,u)=\underset{r=0}{\overset{\mathrm{}}{}}\left(\genfrac{}{}{0pt}{}{z+r}{r}\right)u^{r+1}_0^1v^r\left[\frac{1uv}{1u}e^{uvu}\right]^t𝑑v.$$
The expression in square brackets has a power series expansion in $`u`$ beginning $`1+\text{O}(u^2)`$, so the integrand is a power series in $`tu^2`$ and $`u`$. It follows that $`(z,t,t,u)`$ is $`u`$ times a power series in $`tu^2`$ and $`u`$ and hence that the coefficient $`\widehat{B}_n(z,t)`$ of $`u^n`$ has degree $`(n1)/2`$ in $`t`$ for every $`n`$, as claimed. Specifically, from the expansion
$$\frac{1uv}{1u}e^{uvu}=\mathrm{exp}\left(\underset{m=2}{\overset{\mathrm{}}{}}\frac{u^m}{m}(1v^m)\right)$$
we find the closed form
$$\frac{\widehat{B}_n(z,t)=}{r,k_2,k_3,\mathrm{}0r+2k_2+3k_3+\mathrm{}=n1\left(\genfrac{}{}{0pt}{}{z+r}{r}\right)\frac{t^{k_2+k_3+\mathrm{}}}{2^{k_2}k_2!\mathrm{\hspace{0.17em}3}^{k_3}k_3!\mathrm{}}_0^1v^r(1v^2)^{k_2}(1v^3)^{k_3}\mathrm{}𝑑v}$$
from which the coefficients of $`\widehat{B}_n`$ can be computed explicitly. In particular, we see that $`l+2mn1`$ for all monomials $`z^lt^m`$ occurring in $`\widehat{B}_n`$, and that in the case of equality the coefficient of this monomial comes only from the term $`r=l`$, $`k_2=m`$, $`k_3=k_4=\mathrm{}=0`$ in the above sum and equals the beta integral $`_0^1v^l(1v^2)^m𝑑v/2^ml!m!`$. $`\mathrm{}`$
Now comes the second point. The specialization $`y=t`$ had the effect in the above proof of making the linear term in the power series expansion of $`\left(\frac{1uv}{1u}\right)^te^{uy(1v)}`$ vanish, but it also has a second, less obvious effect: if we denote by $`U(x)`$ the power series
$$U(x):=\mathrm{\hspace{0.33em}1}\frac{x}{e^x1}=\frac{x}{2}\frac{x^2}{12}+\frac{x^4}{720}\mathrm{},$$
then we have
$$u=U(x)\frac{e^u}{1u}=\frac{e^x1}{xe^{x/2}}\mathrm{exp}\left(\frac{x}{e^x1}+\frac{x}{2}1\right)=F(x),$$
$`29`$
where $`F(x)`$ is the power series defined in Theorem 4 in §1 and is an even function of $`x`$. This leads immediately to the following definition and proposition:
###### Proposition 3
For each positive integer $`g`$, the function
$$P_g(z,t):=𝒞(x^{2g1},(z,t,t,U(x))$$
$`30`$
is a polynomial of degree $`2g2`$ in $`z`$ and $`g1`$ in $`t`$ and satisfies the identities
$$(zt)P_g(z,t)+tP_g(z1,t)=S_{2g1}(z)$$
$`31`$
and
$$P_g(z,t)=\underset{r=0}{\overset{\mathrm{}}{}}\frac{S_{2g1}(zr)(t)^r}{(zt)\mathrm{}(ztr)}(z)[[t]].$$
$`32`$
Proof. Equation (31) follows by substituting $`y=t`$, $`u=U(x)`$ into the generating series identity (18), since the second term $`e^{tu}(1u)^t`$ on the right is an even power series in $`x`$ by virtue of equation (29), while the coefficient of $`x^{2g1}`$ in the first term $`(1u)^z`$ is $`S_{2g1}(z)`$ by definition. Similarly, equation (32) is obtained by substituting $`y=t`$, $`u=U(x)`$ into (27) and noting that the second term is an even power series in $`x`$.$`\mathrm{}`$
3. Proof of Theorems 2–5. We begin with Theorem 2. From (3) and (7) we have
$$A(g,n)=\underset{1lkn}{}\frac{(1)^{nk}k^{l1}}{(nk)!(kl)!}S_{2g1}(l).$$
For fixed $`l`$ the coefficient of $`S_{2g1}(l)`$ can be rewritten
$`{\displaystyle \underset{k=l}{\overset{n}{}}}{\displaystyle \frac{(1)^{nk}}{(nk)!(kl)!}}k^{l1}`$ $`=𝒞(t^l,{\displaystyle \underset{k=l}{\overset{n}{}}}{\displaystyle \frac{(1)^{nk}}{(nk)!(kl)!}}{\displaystyle \frac{1}{kt}})`$
$`=𝒞(t^l,{\displaystyle \frac{(1)^{nl}}{(nt)(nt1)\mathrm{}(lt)}})`$
(the latter by Lemma 1 with $`r=nl`$, $`z=nt`$), so, replacing $`l`$ by $`r=nl`$,
$$A(g,n)=𝒞(t^n,\underset{r=0}{\overset{n1}{}}\frac{S_{2g1}(nr)(t)^r}{(nt)(nt1)\mathrm{}(nrt)}).$$
$`33`$
The key observation is now that if we replace the summation on the right by one from $`r=0`$ to $`\mathrm{}`$, then its value does not change: the terms $`r=n`$ and $`r=n+1`$ contribute nothing because $`S_{2g1}(0)=S_{2g1}(1)=0`$, and the terms with $`rn+2`$ contribute nothing because the rational function $`1/(nt)(nt1)\mathrm{}(ntr)`$ has only a simple pole at $`t=0`$ and hence its product with $`t^r`$ has no coefficient of $`t^n`$. Hence equation (32) gives
$$A(g,n)=𝒞(t^n,P_g(n,t)).$$
This proves the vanishing of $`A(g,n)`$ for $`ng`$ (since $`P_g(z,t)`$ is a polynomial of degree $`g1`$ in $`t`$ for all $`z`$) and hence also the fact that $`P_g(k)`$ is a polynomial in $`k`$ of degree $`g1`$. The statement (4) about the values of the numbers $`A(g,n)`$ for $`gn`$ fixed can be proved by using the integral representation of the generating function $`(z,t,t,u)`$, but since the argument is similar to the one we give below for equation (2) (to which (4) is in fact equivalent), and since the statement about the form of the $`A(g,n)`$ was included only for amusement, we omit the derivation.
We now turn to $`A_\nu (g,n)`$. The same argument as was used to derive (33) gives
$$A_\nu (g,n)=𝒞(t^{n+\nu },\underset{r=0}{\overset{n1}{}}\frac{S_{2g1}(nr)(t)^r}{(nt)(nt1)\mathrm{}(nrt)})$$
for any $`\nu >0`$, but now changing the sum to one over all $`r0`$ does change the right-hand side, since the terms $`r=n+\mu +1`$ of the sum have non-0 coefficients of $`t^n`$ for $`1\mu \nu `$. Equation (32) therefore now gives
$$A_\nu (g,n)=𝒞(t^{n+\nu },P_g(n,t))𝒞(t^\nu ,\frac{(1)^n}{(nt)\mathrm{}(1t)}\underset{\mu =1}{\overset{\nu }{}}\frac{S_{2g1}(\mu 1)t^\mu }{(1+t)\mathrm{}(\mu +1+t)}).$$
Again the first term vanishes for $`n`$ sufficiently large ($`ng\nu `$), so for small $`\nu `$ we get explicit formulas for $`\nu `$, two examples being given by equation (5). By analyzing these formulas we could deduce the statement in part (iii) of Theorem 1 about the lowest coefficients of $`P_g(k)`$. But it will be easier to work directly with $`P_g(k)`$, using the following result.
###### Proposition 4
For each positive integer $`k`$ the polynomials $`P_g(z,t)`$ defined by $`(30)`$ satisfy the identity
$$P_g(tk,t)=\underset{l=1}{\overset{k}{}}\frac{(k1)!}{(kl)!}t^lS_{2g1}(l+tk).$$
$`34`$
In particular, the function $`P_g(k)`$ defined by $`(1)`$ is equal to the polynomial $`P_g(0,k)`$.
Proof. We prove this by induction on $`k`$: setting $`z=t`$ in (31) gives the case $`k=1`$ of (34), and setting $`z=tk`$ in (31) gives the induction step from $`k`$ to $`k+1`$.$`\mathrm{}`$
The remaining results stated in §1 follow easily from the last statement of Proposition 4. Theorem 3 is obtained immediately by taking $`z=0`$ in equation (32). For Theorem 4, we first use the integral representation (26) to write
$$(0,t,t,u)=\left(\frac{e^u}{1u}\right)^t_0^u\left(\frac{e^v}{1v}\right)^t\frac{dv}{1v}.$$
Now making the substitutions $`u=U(x)`$ and $`v=U(y)`$ and using equation (29) we get
$$(0,t,t,U(x))=F(x)^t_0^xF(y)^t\frac{U^{}(y)}{1U(y)}𝑑y.$$
But
$$\frac{U^{}(y)}{1U(y)}=\frac{e^y}{e^y1}\frac{1}{y}=\frac{1}{2}+\text{(odd power series in }y\text{)},$$
so
$$(0,t,t,U(x))=\frac{1}{2}F(x)^t_0^xF(y)^t𝑑y+\text{(even power series in }x\text{)}.$$
Equation (10) now follows from the equality $`P_g(t)=P_g(0,t)`$ and the definition of $`P_g(z,t)`$. Finally, the recursion (11) is, as already stated in §1, equivalent to equation (10): if we denote by $`𝔓(x,t)`$ the generating function occurring on the left-hand side of (10), then
$$(10)\frac{1}{2}=F(x)^t\frac{}{x}\left(F(x)^t𝔓(t,x)\right)=\frac{𝔓(x,t)}{x}t\frac{F^{}(x)}{F(x)}𝔓(t,x),$$
$`35`$
and this is seen to be equivalent to (11) by substituting $`F^{}(x)/F(x)=_{n1}(2n+1)\beta _{2n}x^{2n1}`$ from (9) and comparing the coefficients of $`x^{2g2}`$ on both sides.
4. Proof of Theorem 1. We now know, from Proposition 4 or Theorem 4 or 5, that $`P_g(k)`$ is a polynomial. It remains to prove the statements made in Theorem 1 about the coefficients $`c_{g,gj1}`$ ($`j`$ fixed) and $`c_{g,i}`$ ($`i`$ fixed). We start with the “top” coefficients $`c_{g,gj1}`$. Writing $`y=vx`$ in (10) we find
$$\underset{g=1}{\overset{\mathrm{}}{}}P_g(t)x^{2g2}=\frac{1}{2}_0^1\mathrm{exp}\left(\underset{r=1}{\overset{\mathrm{}}{}}\lambda _rtx^r(1v^r)\right)𝑑v$$
where $`\lambda _r=𝒞(x^{2r},\mathrm{log}F(x))=(1+1/2r)\beta _{2r}`$. Expanding the integral as in the proof of Proposition 2 and comparing the coefficients of $`x^{2g2}t^{gj1}`$ on both sides, we find
$`c_{g,gj1}`$ $`{\displaystyle \frac{={\displaystyle \frac{1}{2}}{\displaystyle }}{\alpha ,\beta ,\gamma \mathrm{}0}}`$
$`\alpha +2\beta +3\gamma +\mathrm{}=g1`$
$`\beta +2\gamma +\mathrm{}=j{\displaystyle \frac{\lambda _1^\alpha }{\alpha !}}{\displaystyle \frac{\lambda _2^\beta }{\beta !}}{\displaystyle \frac{\lambda _3^\gamma }{\gamma !}}\mathrm{}{\displaystyle _0^1}(1v^2)^\alpha (1v^4)^\beta (1v^6)^\gamma \mathrm{}𝑑v`$
$`={\displaystyle \frac{1}{2}}{\displaystyle \underset{jd2j}{}}{\displaystyle \frac{\lambda _1^{gd1}}{(gd1)!}}{\displaystyle _0^1}(1v^2)^{gj1}H_{j,d}(v^2)𝑑v`$
with
$$\frac{H_{j,d}(x)=}{\alpha ,\beta ,\gamma \mathrm{}0\beta +2\gamma +\mathrm{}=j2\beta +3\gamma +\mathrm{}=d\frac{\lambda _2^\beta }{\beta !}\frac{\lambda _3^\gamma }{\gamma !}\mathrm{}(1+x)^\beta (1+x+x^2)^\gamma \mathrm{}.}$$
This can now be computed by expanding $`H_{j,d}`$ as a polynomial and computing each term $`_0^1(1v^2)^{gj1}v^{2n}𝑑v`$ as a beta integral, and can easily be seen to have the form (2) for some polynomial $`C_j(g)`$. The highest power of $`g`$ occurs for the maximal value $`d=2j`$, corresponding to taking $`\beta =j`$ and $`\gamma =\mathrm{}=0`$. Also, to compute the coefficient of the highest power of $`g`$ we may replace $`H_{j,d}(x)`$ by its constant term $`H_{j,d}(0)`$, since the main contribution to the integral for $`g`$ large comes from $`v`$ near 0, and the asymptotic value of $`_0^1(1v^2)^{gj1}𝑑v`$ is $`C(g)=2^{2g2}(g1)!^2/(2g1)!`$ (independent of $`j`$) by the beta integral formula. It follows that the asymptotic formula for $`c_{g,gj1}`$ is
$$c_{g,gj1}\frac{C(g)}{2}\frac{\lambda _1^{g2j1}}{(g2j1)!}\frac{(2\lambda _2)^j}{j!}\frac{C(g)\lambda _1^{g1}}{2(2g1)!}g^{2j}\frac{(2\lambda _2/\lambda _1^2)^j}{j!},$$
and this agrees with the result stated in Theorem 1 because $`\lambda _1=1/8`$ and $`2\lambda _2/\lambda _1^2=2/9`$. One can also prove equation (2), and obtain explicit recursion relations for the polynomials $`C_j(g)`$, from the recursion relation given in Theorem 5. The details are left to the reader.
For the “bottom” coefficients $`c_{g,i}`$ ($`i`$ fixed) we use the expansion (8) together with the following lemma, which expresses the “negative Stirling numbers” $`S_n(r)`$ for $`r`$ fixed as finite linear combinations of Bernoulli numbers:
###### Lemma 3
For $`nr1`$ we have the identity
$$S_n(r)=\underset{j=0}{\overset{r1}{}}(1)^{r1j}\left(\genfrac{}{}{0pt}{}{nj1}{rj1}\right)S_j(r)\beta _{nj}.$$
Proof. One sees by induction that the powers of the function $`1/(e^x1)`$ are linear combinations of its derivatives. From the formulas
$$\left(\frac{1}{e^x1}\right)^r=\underset{s=1}{\overset{r}{}}S_{rs}(r)\frac{1}{x^s}+\text{O}(1)(x0)$$
and
$$\frac{(1)^{s1}}{(s1)!}\frac{d^{s1}}{dx^{s1}}\left(\frac{1}{e^x1}\right)=\frac{1}{x^s}+(1)^{s1}\underset{l=s}{\overset{\mathrm{}}{}}\left(\genfrac{}{}{0pt}{}{l1}{s1}\right)\beta _lx^{ls}$$
we deduce
$$\left(\frac{1}{e^x1}\right)^r=\underset{s=1}{\overset{r}{}}S_{rs}(r)\left(\frac{1}{x^s}+(1)^{s1}\underset{l=s}{\overset{\mathrm{}}{}}\left(\genfrac{}{}{0pt}{}{l1}{s1}\right)\beta _lx^{ls}\right),$$
and the desired result follows by comparing coefficients of $`x^{nr}`$ on both sides.$`\mathrm{}`$
Part (iii) of Theorem 1 follows immediately from (8) and Lemma 3. Explicitly, we have
$$c_{g,i}=\underset{j=0}{\overset{i}{}}\left(\underset{r=j+1}{\overset{i+1}{}}(1)^{rj}\left(\genfrac{}{}{0pt}{}{2gj2}{rj1}\right)S_j(r)\alpha _{ir+1}(r)\right)\beta _{2gj1},$$
where
$$\alpha _n(r):=𝒞(t^n,\frac{1}{(1+t)\mathrm{}(r+t)})=\frac{(1)^n}{r!}h_n(1,\frac{1}{2},\mathrm{},\frac{1}{r}),$$
$`36`$
and the coefficient of $`\beta _{2gj1}`$ in this formula is a polynomial of degree $`ij`$ in $`g`$. $`\mathrm{}`$
5. The polynomials $`Q_i(y)`$ and the second generating function for the $`c_{g,i}`$. In this section we will discuss the polynomials defined by equations (12)–(14) and prove Theorem 6. We must first check that the power series in (12) is indeed a polynomial of degree $`i+1`$ and that the three definitions are indeed equivalent. For the first statement, note that if $`ni+2`$ then
$$𝒞(y^n,Q_i(y))=\underset{k=1}{\overset{n}{}}\frac{k^{k1i}}{k!}\frac{(k)^{nk}}{(nk)!}=\frac{1}{n!}\underset{k=0}{\overset{n}{}}(1)^{nk}\left(\genfrac{}{}{0pt}{}{n}{k}\right)k^{n1i}=\mathrm{\hspace{0.33em}0}$$
(the $`n`$th difference of a polynomial of degree $`<n`$ vanishes). For the second, note that the system of integral recursions (13) is equivalent to the system of differential recursions
$$Q_0(y)=y,yQ_{i+1}^{}(y)=(y1)Q_i(y)(i0)$$
$`37`$
(no initial values are needed here because the ($`i+1`$)st equation in this system implies that $`Q_{i+1}(0)=0`$, which is the needed initial condition to solve the $`i`$th equation). It is easy to check that the functions satisfied by (12) or by (14) both satisfy the system (37), so they are all equal. We can write out (14) more explicitly as
$$Q_i(y)=\underset{r=1}{\overset{i+1}{}}\alpha _{i+1r}(r)y^r,$$
$`38`$
with $`\alpha _n(r)`$ defined by (36); these numbers obviously satisfy $`\alpha _n(r1)=r\alpha _n(r)+\alpha _{n1}(r)`$, and this is equivalent to the statement that the polynomials given in (38) satisfy (37).
Now set $`Y(x)=x/(1e^x)`$ and $`\stackrel{~}{Q}_i(x)=Q_i(Y(x))`$. Then (37) gives
$$\stackrel{~}{Q}_{i+1}^{}(x)=Y^{}(x)\frac{Y(x)1}{Y(x)}\stackrel{~}{Q}_i(x)=:\gamma (x)Q_i(x).$$
$`39`$
But an easy calculation shows that the function $`\gamma (x)`$ is nothing other than the logarithmic derivative $`F^{}(x)/F(x)`$ of the function defined in (9). In particular it is an odd function of $`x`$, so that from (39) we deduce that also
$$\frac{d}{dx}\left(\stackrel{~}{Q}_{i+1}(x)\stackrel{~}{Q}_{i+1}(x)\right)=\gamma (x)\left(\stackrel{~}{Q}_i(x)\stackrel{~}{Q}_i(x)\right).$$
$`40`$
This equation and the fact that $`\stackrel{~}{Q}_i(x)\stackrel{~}{Q}_i(x)`$ vanishes at $`x=0`$ imply by induction on $`i`$ that $`\stackrel{~}{Q}_i(x)\stackrel{~}{Q}_i(x)`$ vanishes to order $`2i+1`$ at the origin for all $`i0`$, which is the first assertion of Theorem 6. (The uniqueness statement follows immediately from the existence since the polynomials $`Q_0,Q_1,\mathrm{},Q_i`$ form a basis for the space of polynomials of degree $`i+1`$ with no constant term.) Equation (16), which can be written as the generating function identity
$$2\underset{g=1}{\overset{\mathrm{}}{}}P_g(t)x^{2g1}=\underset{i=0}{\overset{\mathrm{}}{}}\left(\stackrel{~}{Q}_i(x)\stackrel{~}{Q}_i(x)\right)t^i,$$
$`41`$
follows at the same time, since the differential equation (40) is equivalent to the differential equation in (35) for the generating series $`P_g(t)x^{2g1}`$ or to the recursion (11) for its coefficients. $`\mathrm{}`$
C. Faber and R. Pandharipande, Logarithmic series and Hodge integrals in the tautological ring, this volume, pp. XX–XX. |
warning/0002/quant-ph0002044.html | ar5iv | text | # Security of classical noise-based cryptography
## I Introduction
Cryptography is used to transmit a message from a sender (referred to as Alice) to a receiver (Bob) without leaking useful information to others. It has been proved that a message can be transmitted securely if it is coded and decoded by a sequence of random bits (key) whose length is equal to that of the message. The problem of secure transmission is then reduced to that of generating a secret key shared by Alice and Bob. Classical cryptosystems rely on computational complexity and may be broken by an effective algorithm or a powerful computer. Quantum key distribution (QKD) protocols, in contrast, provide an unconditionally secure key, the security of which is inherent in the laws of quantum mechanics. This remarkable advantage of QKD protocols has been attracting increasing research interest since the proposal by Bennett and Brassard. Although QKD has been demonstrated in over-20-km-fiber communication channels, its application in practical communication systems is not straightforward. The QKD protocols require single-photon transmission to guarantee the security, and thus are vulnerable to loss and noise inherent in actual transmission channels. Optical amplifiers will not solve this problem, because the noise of the optical amplifiers inevitably destroys the quantum correlation. Single-photon transmission requires the use of complicated and inefficient photon counting techniques instead of conventional analog detection, besides a truly practical single-photon source is not yet available. The QKD protocols are therefore not fully compatible with the current optical fiber communication systems. A secure key distribution protocol compatible with the current systems is desirable. This would be a protocol that uses more than one photon and allows optical amplifiers to be used. Such a protocol would be based on coherent state photons, or classical light.
Maurer has shown perfect cryptographic security can be obtained in a classical noisy channel with the help of a noiseless feedback channel. Yuen and Kim examined the principles underlying the QKD protocol with two non-orthogonal quantum states (B92 protocol.) The security of the B92 protocol relies on two facts: (i) an eavesdropper (Eve) cannot accurately determine the value of each transmitted bit ( i.e., no efficient opaque eavesdropping.) (ii) Eve cannot closely correlate Bob’s measurement results with her own ( i.e., no efficient translucent eavesdropping.) Yuen and Kim pointed out that these two conditions can be satisfied in a classical transmission system, where the detectors of Bob and Eve are under independent additive noise and show a small signal-to-noise-ratio (SNR.) They proposed a classical noise-based protocol for key distribution (referred here to as the YK protocol.) The YK protocol working with classical light, would have advantages in practical implementations. Proving the security in YK protocol is, however, subtler than in QKD protocols. Since many photons are transmitted to carry one-bit information, the conditions specified above will not be satisfied if the SNR of Eve’s detection is sufficiently high. Eve’s SNR can be increased by using low-noise detection equipment, or simply by moving closer to Alice than Bob (because of the fiber loss.) The original analysis of the security of YK protocol assumed the same SNR for Bob and Eve . For practical implementations, it is important to determine the design rules of Eve’s SNR and Bob’s.
In this article we quantitatively examine the security of the YK protocol, and show that the YK protocol is secure, even if Eve’s SNR is 9 dB better than Bob’s. We also show experimental results that demonstrate secure key distribution against translucent attack. Section 2 provides the condition for secure key distribution in terms of the secure key distribution rate. Section 3 describes the experiment on the YK protocol using conventional fiber optics. Section 4 discusses the implementation issues.
## II Theory
We first define the secure key distribution rate. Suppose Alice transmits an equally probable binary string to Bob. Shannon information between Alice and Bob is expressed by
$$I_{AB}=1+e_B\mathrm{log}_2e_B+\left(1e_B\right)\mathrm{log}_2\left(1e_B\right),$$
(1)
where $`e_B`$ denotes the error rate of Bob’s decision. Because of the decision errors, Bob has the information in only $`n_{sift}I_{AB}`$ of the $`n_{sift}`$ sifted bits. Alice and Bob exchange redundant information over the public channel in order to obtain the reconciled key. This procedure is called error correction, the best known practical protocol for which was given by Brassard and Salvail. For successful error correction with Brassard-Salvail’s protocol, the error rate should be less than 0.15. To establish a secret key, Alice and Bob use privacy amplification, random hashing of the reconciled key into a shorter key. If they shorten the reconciled key of length $`n_{rec}`$ by the fraction $`\tau `$ and sacrifice $`n_S`$ bits as a safety parameter, Eves’s Shannon information on the final key of length $`\tau n_{rec}n_S`$ is bounded by
$$I_E\frac{2^{n_S}}{\mathrm{ln}2}.$$
(2)
The fraction $`\tau `$ is given by
$$\tau =1+\left(1/n_{rec}\right)\mathrm{log}P_C,$$
(3)
where the collision probability $`P_C\left(X\right)`$ of $`X`$ is defined as follows: Let $`X`$ be a random variable with an alphabet $`𝒳`$ and distribution $`P_X`$. The collision probability is the probability that $`X`$ takes the same value twice in two independent experiments, that is, $`P_C\left(X\right)=_{x𝒳}P_X\left(x\right)^2`$. The logarithm of the collision probability thus refers to Eve’s information on the key. The collision probability can be expressed by the probability $`p(k)`$ that $`k`$ is the $`i`$-th signal of Bob’s string and the joint probability $`p(k,l)`$ that $`k`$ is the $`i`$-th signal of Bob’s string and $`l`$ is the $`i`$-th signal of Eve’s string. We have the following formula for the fraction $`\tau `$:
$$\tau =1+\mathrm{log}_2\left[\underset{k,l}{}\frac{p(k,l)^2}{p(k)}\right].$$
(4)
According to Bruß and Lütkenhaus, Bob can generate secure bits from his sifted bits at the rate $`R`$ of
$$R=I_{AB}\left(1e_B\right)\tau e_B.$$
(5)
We refer this rate $`R`$ as the secure key distribution rate, and for secure key distribution its value be positive. The actual key generation rate is further reduced by multiplying the generation rate of the sifted key. In the following part of this section, we derive the conditions under which $`R`$ is positive.
In the YK protocol, the bit values (“0” and “1”) are encoded so as to make the probability distribution of the received signal symmetric. Alice sends encoded bits on a weak classical light. Signal $`s_0(t)=S\varphi (t)`$ is transmitted for “0”, and $`s_1(t)=S\varphi (t)`$ is transmitted for “1”, where $`_T\varphi (t)𝑑t=1`$. We here measure the signal value as the voltage on the load resistance $`R_{load}`$ of a photodiode. Mean signal voltage $`S`$ is defined by $`S^2=_Ts_i^2(t)𝑑t`$, and $`S^2/R_{load}`$ represents the signal energy over the duration $`T`$ (signal energy per bit.) The output $`r(t)`$ of the detector contains the noise $`n(t)`$, so $`r(t)=s_i(t)+n(t)`$. If the noise is white Gaussian noise with spectral density $`\sigma ^2`$, the probability distribution of the detected signal $`V`$ is expressed by
$$P(V)=\{\genfrac{}{}{0pt}{}{\left(1/\sqrt{2\pi }\right)\mathrm{exp}\left[\left(VS\right)^2/\left(2\sigma ^2\right)\right]\left(for\mathrm{`}\mathrm{`}0\mathrm{"}\right)}{\left(1/\sqrt{2\pi }\right)\mathrm{exp}\left[\left(V+S\right)^2/\left(2\sigma ^2\right)\right]\left(for\mathrm{`}\mathrm{`}1\mathrm{"}\right)},$$
(6)
where the signal is averaged over the duration $`T`$ as $`V=_Tr(t)𝑑t`$. The SNR $`\beta ^2`$ in this system is define by $`\beta =S/\sigma `$. In a conventional decision scheme the bit values are determined to be “0” if $`V>0`$ and “1” if $`V<0`$. Decision errors will occur at the rate of $`Q\left(\beta \right)`$, where $`Q`$ is the scaled complementary error function defined by
$$Q(x)=\frac{1}{\sqrt{2\pi }}_x^{\mathrm{}}\mathrm{exp}\left(y^2/2\right)𝑑y.$$
(7)
We set a threshold $`V_{th}=mS(m>1)`$ to make a decision:“0” if $`V>`$ $`V_{th}`$ and “1” if $`V<V_{th}`$, but leave inconclusive if $`V_{th}VV_{th}`$. The probability of making a decision is given by the following decision rate:
$$F_+=Q\left(\left(m+1\right)\beta \right)+Q\left(\left(m1\right)\beta \right),$$
(8)
and the error rate is
$$e=\frac{Q\left(\left(m+1\right)\beta \right)}{F_+}.$$
(9)
The sifted key is generated from the raw bit string by the Bob’s decision. The decision rate $`F_+`$ thus refers to the generation rate of the sifted key. As seen in Eqs. (8) and (9), the decision rate $`F_+`$ and the error rate $`e`$ are determined by the values of the SNR and the threshold. As described below, this error rate determines the joint probabilities $`p(k,l)`$ and therefore the secure key distribution rate. The system is thus fully characterized by the SNR and the threshold. The error rate can be reduced by increasing the threshold, but, a high threshold will also reduce the decision rate. Since, as we can see by comparing Fig. 1 and Fig. 2, the decision rate decrease faster than the error rate, the threshold value should not be set too high.
As in the B92 protocol, the inconclusive results play a essential role in guaranteeing the security of the key distribution. A finite threshold value of Bob enables him to make accurate decisions on his sifted key at a cost of the generation rate. Eve, on the other hand, should make a decision with zero threshold in order to obtain conclusive results for all the transmitted bits. If Eve uses a finite threshold in her decision, she will obtain the inconclusive results on the sifted bits. The assumption of independent noise prevents Eve from predicting which bit Bob will obtain a conclusive result. Eve can acquire no information from these inconclusive bits. Since Eve’s error rate $`e_E`$ is less than 1/2, she will obtain more information by making a decision with zero threshold. Therefore, Bob can make more accurate decisions on the sifted key bits than Eve can. That is, Bob has more information than Eve, and can distill secure key bits with Alice.
Now we will examine the conditions for security against eavesdropping. We here consider only two simple kind of eavesdropping, translucent attack and opaque attack. A translucent attack can be made by simply putting a beam splitter in the transmission channel. The translucent attack to the YK protocol, in contrast to those to the QKD protocols, will not change the state of the transmitted light. The probability distribution of Bob’s bits is the same as that of Alice’s, $`p(0)=p(1)=1/2`$, because after error correction Alice and Bob share completely correlated results. The joint probabilities $`p(k,l)`$ are $`p(0,0)=p(1,1)=(1e_E)/2`$ and $`p(0,1)=p(1,0)=e_E/2`$. The fraction $`\tau `$ is calculated from Eq. (4) as
$$\tau =1+\mathrm{log}_2\left(12e_E+2e_E^2\right).$$
(10)
The secure key distribution rate can be estimated by using Eqs. (1), (5), and (10). Figure 3 shows Eve’s required error rate as a function of Bob’s. As Bob’s error rate $`e_B`$ increases, Eve’s error rate should be increased in order to obtain a positive secure key distribution rate . For example, if Bob’s error rate is 0.15, Eve should make errors at a rate greater than 0.27. This implies that SNR of Eve’s system should be less than 0.38 for white Gaussian noise. On the other hand, Bob’s SNR should be better than 0.057 to keep his error rate smaller than 0.15 and his decision rate at 10<sup>-3</sup>. The secure key distribution is therefore possible even if Eve’s SNR is six times (8 dB) as large as Bob’s. The tolerance of the SNR increases as Bob’s error rate decreases, and it reaches 10 dB for $`e_B=0.01`$.
In an opaque attack, Eve receives all the photons in $`\eta n`$ out of the $`n`$ bits sent by Alice. Then Eve sends the $`\eta n`$ bits to Bob according to her decision. Eve never touches the rest of the bits ($`\left(1\eta \right)n`$ bits) and forwards them to Bob. To protect information from opaque attack, Bob should determine his threshold according to the average signal intensity of each bit. If he observes only the average intensity over many bits, Eve can set a finite decision threshold to reduce her error rate and will then obtain conclusive results for $`\gamma \eta n`$ bits ($`\gamma <1.`$) If she sends only the conclusive results with signals $`\gamma ^1`$ times as intense as received, Bob will obtain the same long-time average signal intensity he would if Eve did not intercept the photons. If Bob observes the signal intensity of each bit, Eve must send every bit with the same intensity as she receives it. Eve then should make a decision with zero threshold, otherwise she will lose the information on the inconclusive results. There is a trade-off for Eve on the fraction $`\eta `$: a large $`\eta `$ will increase Eve’s information gain, but will also make her easily detectable from the increase of Bob’s error rate. Bob’s error rate on the unintercepted bits is $`e_B`$, but the error rate on the intercepted bits is $`\left(1e_E\right)e_B+e_E(1e_B)`$. The eavesdropping thus increases Bob’s error rate to
$$e_B^{}=\left(1\eta \right)e_B+\eta \left[\left(1e_E\right)e_B+e_E\left(1e_B\right)\right].$$
(11)
To calculate the secure key distribution rate by using Eqs. (1), (5), and (10), we estimate the joint probabilities $`p(k.l).`$ After the error correction, Bob has $`(1e_B^{})nF_+`$ bits. The probability distribution i symmetric: $`p(0)=p(1)=1/2.`$ Eve obtains $`(1e_E)\left(1e_B\right)\eta nF_+`$ $`+(1/2)(1\eta )(1e_B)nF_+`$ correct results and $`e_Ee_B\eta nF_++(1/2)(1\eta )(1e_B)nF_+`$ incorrect results on Bob’s bits. The joint probabilities are obtained as
$`p(1,1)`$ $`=`$ $`{\displaystyle \frac{\left[(1e_E)\left(1e_B\right)\eta +(1e_B)(1\eta )/2\right]nF_+}{2\left(1e_B^{}\right)nF_+}}`$ (12)
$`=`$ $`{\displaystyle \frac{\left[\left(1e_E\right)\eta +\left(1\eta \right)/2\right]\left(1e_B\right)}{2\left(1e_B^{}\right)}}`$ (13)
$`p(1,0)`$ $`=`$ $`{\displaystyle \frac{e_Ee_B\eta nF_++(1e_B)(1\eta )nF_+/2}{2\left(1e_B^{}\right)nF_+}}`$ (14)
$`=`$ $`{\displaystyle \frac{e_Be_E\eta +(1/2)\left(1e_B\right)\left(1\eta \right)}{2\left(1e_B^{}\right)}}`$ (15)
$`p(0,0)`$ $`=`$ $`p(1,1)`$ (16)
$`p(0,1)`$ $`=`$ $`p(1,0).`$ (17)
Figure 4 shows the minimum required values of Eve’s error rate for secure key distribution $`\left(R>0\right)`$ as a function of Bob’s error rate $`e_B^{}`$. Though Bob can observe only $`e_B^{}`$ values, he can estimate $`e_B`$ from the SNR of his detection system. Eve will be detected if $`e_B^{}e_B`$. The detection is easy if Bob’s error rate is much lower than Eve’s. A high error rate for Bob may hide Eve, but secure key distribution is possible even in this case. Suppose $`e_B=0.1`$ and $`e_B^{}=0.15`$. As shown in Fig. 4, the secure key distribution rate is positive if Eve’s error rate is larger than 0.12. If the system is under white Gaussian noise, this condition on the error rate is satisfied when Eve’s SNR is smaller than 1.35 (1.3 dB.) Since Bob’s SNR should be better than 0.089 (-10.5 dB) to keep the decision rate at 10<sup>-3</sup> and $`e_B=0.1`$, the tolerance in SNR is 11.8 dB. This small SNR for Eve implies that the signal should be sent on a weak light. Increasing the light intensity reduces Eve’s error rate, and makes the secure key distribution impossible.
## III Experiment
In implementing the YK protocol, we should code the bit values in such a way that the probability distribution of the received signals is symmetric . In this experiment we used the unipolar Manchester code. This code represents ”1” as a change from ON to OFF and ”0” as a change from OFF to ON. It can be decoded as follows: divide the incident light into two paths, one of which is set one half of the pulse width longer than the other. Then take a difference of the two light intensities by a balanced detector. The latter half of the pulse slot yields a negative signal for ”1” and a positive signal for ”0”. Binary phase shift keying (BPSK) also yields a symmetric distribution by homodyne detection, and would be more sensitive, but unipolar Manchester code is easier to implement.
Figure 5 shows the experimental setup. Two distributed feedback (DFB) laser diodes (LDs) served as 1.3 $`\mu `$m light sources. A pattern generator provided a signal pulse string to modulate one DFB LD (signal LD) directly. The second pattern generator was synchronized to the first, and provided an 8-ns pulse at the beginning of each pulse string. This pulse modulated the other DFB LD (trigger LD) directly to generate a trigger light pulse. The output of the signal LD was set weaker than that of the trigger LD. The clock frequency in the present experiment was 25 MHz. Only a fixed pattern of 101010$`\mathrm{}`$ was transmitted. The coded signal light then became a square wave with a duty of 50 % and a pulse duration of 20 ns. We sent strings of 30.8 kbits. The outputs of the two LDs were combined and attenuated by an attenuator (ATT1.) To simulate the translucent attack by an eavesdropper, we inserted a 50:50 divider. An attenuator (ATT2) was placed in one arm of the divider to examine the SNR tolerance for the secure key distribution. The signals of both outputs were detected by the receivers. Each receiver consisted of a 50:50 divider, a fiber delay of a half pulse width, and a balanced detector. The balanced detectors made of two commercial InGaAs pin photodiodes loaded by 50 $`\mathrm{\Omega }`$ resisters were operated in analog mode. The catalog data (typical values) for the quantum efficiency and the dark current of the photodiodes at 25 C were 90 % and 5 nA. The photodiodes ware not cooled. The output signals of the receivers were led to amplifiers ($`G=40`$ dB) and then to analog-digital converters.
Figure 6 shows a typical probability distribution of the output signal from the amplifier. It is well represented by the sum of two Gaussians. The intensity of the optical signal was 0.380 $`\mu `$W (-34.2 dBm) at the input port of the receiver, and the SNR of this signal was 1.0 (0 dB.) We averaged the output pulse over the duration (10 ns), and evaluated the decision rate and error rate as a function of the SNR and the threshold. The results are shown in Fig. 1 and Fig. 2. The experimental results agree well with the theory assuming white Gaussian noise. These indicated that white Gaussian noise dominated the present receiver sensitivity, and that the security analysis described in Sec. 2 can be applied to the experiment. The number of the sifted bits became small when the threshold value is high. We had less than 30 bits, if the decision rate is less than 10<sup>-3</sup>. This insufficient sample number caused the error rate fluctuation observed for large $`m`$’s in Fig. 2. For SNR values up to 0 dB, the noise level was almost same as the dark noise level, but for larger SNR, it increased with the signal intensity. Dark current of the photodiodes was negligible compared to the thermal noise. These indicates that, for SNR values up to 0 dB, the sensitivity of the system was dominated by thermal noise, which is constant to the input photon number. As the intensity increased, the thermal noise was exceeded by shot noise, which is proportional to the input photon number. SNR was proportional to the square of the input power for weak signals, and tended to be proportional to the input power as the signal intensity increased.
The error rate shown in Fig. 2 provides a criterion for key distribution secure against opaque attack. A low error rate of 0.038 was obtained for weak signals by setting the threshold at $`m=10`$, where the SNR was -9.25 dB. The decisions were made at the rate of 0.0008, slightly lower than 10<sup>-3</sup>. This error rate was lower than the theoretical value of 0.072 because of the fluctuation described above. Using $`e_B=0`$ line in Fig. 4, we conclude that the key distribution is secure if Eve’s error rate is larger than 0.1, where we use the theoretical value of the error rate (0.072) for Bob. This condition is satisfied if Eve’s SNR is 0 dB, because we obtained the error rate of 0.15 in the experiment. Bob’s advantage in SNR was thus greater than 9.25 dB. This advantage was almost constant for large SNR signals.
Security against the translucent attack was examined as follows. We assigned one receiver that followed ATT2 as Bob, and the other receiver as Eve. ATT1 affected the SNRs of both Bob and Eve, whereas ATT2 determined the ratio of the SNRs. The decisions in Bob were recorded with several values of the threshold $`m`$, while the Eve’s decisions were recorded with the threshold fixed at zero. We measured the error rate $`e_B`$ and decision rate $`F_+`$ of Bob and the error rate $`e_E`$ of Eve. We estimated the joint probabilities $`p(0,0),`$ $`p(0,1)`$, $`p(1,0),`$ and $`p(1,1)`$ from the bit data about which Bob made correct decisions. Finally, we calculated the secure key distribution rate $`R`$ by using Eqs. (1), (4), and (5). Figure 3 shows the secure key distribution rate as a function of the error rates of Eve and Bob. The symbols in Fig. 3 show the secure key distribution rates estimated from the experiment. Experimental results agree well with theoretical results (lines.) The secure key distribution was achieved if error rates of Eve and Bob are in the region above the $`R=0`$ line in Fig. 3. The decrease in Bob’s SNR reduced the range of the signal intensity for secure key distribution. Secure key distribution was impossible when Bob’s SNR was -9 dB smaller than that of Eve. This result also agrees well with the prediction. We obtained the largest actual secure key distribution rate $`F_+R=0.04`$ when the SNRs of both Bob and Eve were unity (0 dB) and Bob’s threshold was set to $`m=2`$. The observed error rates were 0.01 for Bob and 0.15 for Eve. The secure key distribution rate was $`R=0.29`$. Higher secure key distribution rates were obtained by setting larger threshold values, but the reduction in the decision rate decreased the product $`F_+R`$. Alice transmitted signals at 50 Mb/s, so that the key transmission rate in the present experiment was 2 Mb/s. This is a hundred times as fast as the key transmission rate reported in the QKD experiments . The transmission rate was limited only by the electric circuits. The secure key would be transmitted at 400 Mb/s if a 10-Gb/s transmission channel were used.
## IV Discussions
This theoretical analysis has shown that the secure key distribution is possible as long as the ratio of Bob’s SNR to Eve’s is better than -9 dB, and the experimental results presented here confirmed it. A practical cryptosystem should thus be designed to satisfy this condition. Eve may stay much closer to Alice than Bob, and her signal may be larger than Bob’s because of the fiber loss. We estimate a limit of the transmission distance in the following. It would be very difficult to use complicated networks, where the path of a traffic is not fixed. We have to construct a cryptosystem on a simple network or a point-point channel. Suppose, for simplicity, we construct it on a point-to-point channel. SNR is proportional to the square of the light intensity when the system sensitivity is limited by thermal noise. Then Bob’s advantage of 9 dB refers to the fiber length of 22.5 km using a lowest loss fiber (0.2 dB/km) and neglecting connection loss. SNR is proportional to the light intensity in systems limited by shot noise, and the fiber can be as long as 45 km in those systems. This values would be increased assuming the translucent attack, because Eve would tap the channel and receive small part of the signal.
Amplifiers can be used in YK protocol as long as the SNR permits. They will improve the SNR by reducing the effect of the thermal noise, and therefore will be useful when the system is limited by the thermal noise. In the shot noise limit, even an ideal amplifier reduces the SNR by 3 dB. The use of amplifiers is restricted by this degradation in the SNR. At most three amplifiers are thus possible.
The above estimation assumed that Bob and Eve use the same detectors. Bob should reduce system noise as possible to guarantee the security, by cooling the receiver, for example. If he can suppress all the thermal noise, the SNR of his system will be limited by shot noise, the standard quantum limit. The mean photon number transmitted in this system should be reduced to unity, because the SNR of 0 dB refers to mean photon number of unity in the shot noise limited systems. The error rate can be reduced below the standard quantum limit by optimum decision. The improvement will be apparent for small mean photon numbers. Security analysis based on quantum detection theory as well as practical implementation of the optimum decision are open for further study. It would be noteworthy that the security analysis described in the present article will provide a security criteria for the B92 protocols employing dim coherent lights.
The YK protocol provides more efficient key distribution at higher bit rates than do other QKD protocols, but it requires that the signal intensity be controlled to keep Eve’s SNR advantage smaller than 9 dB. This may be a disadvantage compared to the QKD protocols like BB84, where the unconditional security is proved if the photons are generated by a perfect single photon source. However, it has been shown that the Eve’s advantage in SNR will limit the efficiency of the protocol in a lossy channel. The SNR control would be also required in actual BB84 systems.
## V Conclusion
Quantitative analysis of the security of the Yuen-Kim protocol shows that the secure key distribution is possible even if the eavesdropper receives signals with a signal-to-noise-ratio better than that with which the legitimate receiver receives them. It has been shown that the signal-to-noise-ratio of the legitimate receiver may be -9 dB smaller than that of the eavesdropper. The results of an experiment using conventional fiber optics agrees well with the analysis results. These results have demonstrated a practical implementation of a secure key distribution protected by the laws of physics. We think the YK protocol would be a solution for practical cryptography systems. |
warning/0002/astro-ph0002490.html | ar5iv | text | # A New Complete Sample of Sub-millijansky Radio Sources: An Optical and Near-Infrared Study
## 1 INTRODUCTION
The advent of deep radio surveys reaching flux densities well below 1 mJy (Mitchell & Condon 1985; Windhorst et al. 1985; Oort 1987; Windhorst et al. 1993; Hopkins et al. 1998; Richards 2000) revealed a new population of faint sources more numerous than the AGN-powered radio galaxies dominating the strong-source population. This corresponds to a steepening of the differential source counts over non-evolving predictions at levels $`4`$ mJy. The faint counts suggest that significant evolution has occured over the redshift range spanned by the observed population. Photometric and spectroscopic studies (Thuan et al. 1984; Windhorst et al. 1985; Thuan & Condon 1987; Benn et al. 1993) suggest that the faint excess at 1.4 GHz is composed predominately of star-forming galaxies similar to the nearby starburst population dominating the Infrared Astronomical Satellite (IRAS) 60-$`\mu `$m counts (Benn et al. 1993). Indeed, this is supported by the strong correlation between radio (1.4 GHz) and far-infrared (60-$`\mu `$m) flux densities of disk galaxies (Helou, Soifer & Rowan-Robinson 1985), implying a significant proportion of starburst galaxies at faint radio flux densities.
The overall observed source-count distribution from faint ($`\mu `$Jy) to bright ($`S_{1.4}10`$mJy) flux densities cannot be explained by starbursts alone. Evolutionary models of radio source counts need to invoke two separate populations (eg. Danese et al. 1987; Rowan-Robinson et al. 1993; Hopkins et al. 1998). Condon (1989) describes these populations as starbursts and monsters, each powered by different mechanisms: ‘starbursts’ deriving their radio emission from supernova remnants and HII regions and ‘monsters’ from compact nuclear related activity (eg. active galactic nuclei; AGNs). The proportion of AGN is much greater at higher flux densities $`S_{1.4}10`$ mJy (Kron, Koo & Windhorst 1985; Gruppioni et al. 1998), where a majority are associated with classical radio galaxies exhibiting extended (FRI and FRII-type) morphologies (Fanaroff & Riley 1974). The optical counterparts to sources at bright radio flux densities $`S_{1.4}1`$ mJy is comprised mostly of ellipticals while at sub-mJy to $`\mu `$Jy levels, the optical counterparts are identified as blue galaxies exhibiting peculiar (compact, interacting and irregular) morphologies (Kron, Koo & Windhorst 1985; Gruppioni et al. 1998). Studies of faint radio sources, namely their stellar population, how they evolve with redshift, and how they relate to local normal galaxies is progressing rapidly, however, much remains to be learned from the faintest ($`\mu `$Jy) radio populations at redshifts $`1z2`$.
Radio surveys are insensitive to the effects of absorption by dust which is known to bias surveys severely at optical/UV wavelengths. This is particularly important for derivations of the cosmic history of star formation and its relation to hierarchical models of galaxy formation. Optical/UV studies have shown that there is an increase in both the space density of star-forming, morphologically disturbed galaxies (eg. Richards et al. 1999) and also, the global star formation rate with redshift to $`z1`$ (eg. Madau et al. 1996). Similar evidence is emerging from studies of sub-mm sources (Blain et al. 1999) and amongst the faint radio population at $`1z2`$ \- the redshift range probed by the deepest surveys (eg. Cram 1998; Haarsma et al. 1999).
There has been much speculation as to whether global star-formation rates (SFRs) derived from radio observations exceed those determined from optical/UV studies. Cram et al. (1998) note that systematic discrepancies may exist between the various star formation indicators, which are not well understood (see also Schaerer 1999). It is encouraging to see however that Haarsma et al. (2000) find global SFRs derived from 1.4GHz observations of the Hubble Deep Field to exceed optically detetermined values by a factor of a few out to $`z1`$. Indeed, an analysis of Balmer decrements and optical–near-infrared colors in star forming galaxies by Georgakakis et al. (1999) from the Phoenix Deep Survey (Hopkins et al. 1998), finds evidence for visual extinctions from one to a few magnitudes. Currently, about 20% of existing micro-jansky radio samples remain unidentified to $`I=25`$ mag in Hubble Deep Field images (eg. Richards et al. 1999). A majority of these could represent a significant population of dust-enshrouded starbursts and/or AGN at high redshift. These results are in support of efforts to further understand the dust properties of star-forming galaxies.
The primary aim of this paper is to introduce a new complete sample of radio sources selected at 1.4 GHz, uniformly selected over the flux range $`S0.3`$ mJy (5$`\sigma `$) from an area covering $`0.3\mathrm{deg}^2`$. Although much larger area surveys to deeper radio sensitivities have been carried out (eg. Hopkins et al. 1998), the present study reports the results of more sensitive optical observations. Archival near-infrared data for a subset of the sample are also presented. The near-infrared data is from the ongoing 2MASS project, and represents a unique aspect of this study in the identification of radio-selected starbursts. Although we currently lack valuable spectroscopic information, we combine radio–near-infrared–optical flux ratios, radio maps, and optical images to explore the properties of the entire sample. Our deep optical identifications provide the opportunity to asses the importance of dust in star-forming systems via the observed radio–to–optical and near-infrared–to–optical colors. Simple stellar synthesis models that include radio emission and reddening are used to constrain possible amounts of absorption.
This paper is organised as follows. In §2, we discuss the radio observations and data reduction, present the radio catalog, and compare our results with data available from (shallower) all sky radio surveys. The optical photometric observations, their reduction and astrometric calibrations are discussed in §3. Our method for radio-optical identification, the archival near-infrared data, radio-optical image overlays, and our optical/near-infrared catalog are presented in §4. A study of the radio, near-infrared, and optical colors and constraints on synthesis models incorporating dust is presented in §5. An application of our color-color analysis to select high-redshift ultraluminous infrared galaxies from deep radio surveys is discussed in §6. All results and future prospects are summarised in §7.
## 2 RADIO OBSERVATIONS
### 2.1 Strategy
Observations were made with the VLA C-configuration at 1.4 GHz on 1998 December 19. This configuration yields a good compromise for resolution and surface-brigntness sensitivity. The $`5\sigma `$ confusion limit for this configuration is only $`50\mu `$Jy/beam at 1.4 GHz (Mitchell & Condon 1985) since the synthesised beam size (Full Width at Half Power; FWHP) is $`15`$ arcsec. The resulting radio positions have rms errors $`1`$ arcsec, except for extended sources with multiple components, sufficient for making optical identifications. At 1.4 GHz the FWHP of the VLA primary beam is 31 arcmin and approximately corresponds to the diameter of our final imaged field. This relatively large coverage avoids field-to-field variations in source counts induced by high-redshift clustering. Although surveys at higher frequencies (eg. 8 GHz) can reach lower flux densities than at 1.4 GHz, most radio sources have spectral indices $`\alpha 1`$ ($`S\nu ^\alpha `$) so that the population being sampled is similar.
To provide uniform sensitivity over the full area of our field, we observed seven positions arranged in a filled hexagonal pattern with a separation of $`26\mathrm{}`$ between pointing centers. The resulting root mean square (rms) map noise is thus nearly constant (cf. Condon et al. 1998). Our field was centered on RA(2000)=$`00^h13^m12^s.0`$, Dec.(2000)=$`+25\mathrm{°}54\mathrm{}44\mathrm{}`$. This field was chosen for its relatively low foreground galactic-cirrus emission likely to affect optical/near-infrared identifications and also for the absence of bright radio galaxies. The integration time on each pointing was $`1`$ hour. This allowed us to reach an rms noise of $`60\mu `$Jy in regions free from bright contaminating sources (see § 2.3 for more details).
Our observations were made in spectral line mode with four Intermediate Frequencies (IFs), each divided into 7 spectral channels of width 3.125 MHz. The advantage of this mode is to minimize bandwidth smearing (i.e. chromatic aberration) which reduces the point-source sensitivity away from the pointing center and causes appreciable image distortion. Additionally, the spectral line mode is less prone to narrowband interference noise spikes. With continuum mode however, we would have had a little over twice the bandwidth and a factor $`\sqrt{2}`$ lower noise.
### 2.2 Data Reduction
The data were analysed with the NRAO AIPS reduction package. We observed 3C 48 to calibrate the visibility amplitudes, using $`S=16.5`$ and $`S=15.9`$ Jy at 1.365 and 1.435 GHz, respectively. The calibration was applied using the SPLIT task. The calibrated data from each pointing were edited and imaged separately, CLEANed, and restored with 15 arcsec FWHP Gaussian beams. The seven separate images were weighted and combined as described in Condon et al. (1998) to produce a final $`33\mathrm{}\times 33\mathrm{}`$ map with nearly uniform sensitivity and corrected for primary-beam attenuation.
### 2.3 Noise and Source Extractions
The resulting rms noise of our final map after correcting for primary-beam attenuation is not uniform over the entire field, but increases by up to 30% near a strong ($`400`$ mJy) source near the edge of our field. Despite this variation in sensitivity, we were able to divide our $`33\times 33\mathrm{arcmin}^2`$ field into four equal ($`16.5\times 16.5\mathrm{arcmin}^2`$) regions within which the rms noise varies by no more than a few percent. These constant-noise regions simplify the application of an automated source-finding algorithm over a single continuous region (see below). The lowest and highest rms noise amongst these regions was $`60.9\mu `$Jy/beam and $`90.3\mu `$Jy/beam respectively. See Table 1 for the region definitions. The rms noise of each region was estimated from the Gaussian core of the amplitude distribution of the pixel values as produced by the AIPS task IMEAN. In Figure 1, we show the distribution in pixel values of our entire $`33\mathrm{}`$ field. The rms deviation in peak flux density derived from a fit to this histogram is $`69\mu `$Jy/beam. Figure 2 shows a contour map of our entire radio field.
Each constant-noise region in Table 1 was used for the source extractions. Within each region, we searched for radio sources down to a peak flux density $`5`$ times the rms value of the region. The sources were extracted by the AIPS task SAD which uses Gaussian fits to estimate the fluxes, positions, and angular sizes of the selected sources. For faint sources however, unconstrained Gaussian fits may be unreliable (see Condon 1997). For this reason, we adopted the following method for the source extraction: first, we ran the task SAD with a $`3\sigma _{rms}`$ detection threshold to obtain an initial list of candidates, we then derived the peak flux densities of the faint sources (with $`3\sigma <S_{peak}<7\sigma `$) using the MAXFIT task. This task uses a second order interpolation algorithm and is known to be more accurate. Only sources with a MAXFIT peak flux density $`5\sigma _{rms}`$ were retained. For these faint sources, the total flux density was estimated using the task IMEAN, which integrates the (median background subtracted) pixel values in a specific rectangle. The rectangle was chosen to encompass as much of the source as possible. For all other parameters (sizes, positions and position angles) we retained the values obtained from the initial Gaussian fits. Only two sources had irregular morphologies showing multiple components. For these, the total (background subtracted) flux was determined using the task TVSTAT which allows an integration over a specific irregular area defined to encompass as much of the source as possible.
The numbers of sources detected in each constant-noise region are summarised in Table 1. A total of 62 sources were detected with flux densities $`5\sigma _{rms}`$ over an area of $`0.303\mathrm{deg}^2`$. Within poisson uncertainties, this number is consistent with source counts from surveys by Ciliegi et al. (1999) and Hopkins et al. (1998). To our limiting (mean) sensitivity of 0.35mJy and within $`0.303\mathrm{deg}^2`$, they report a source count of typically $`70\pm 8`$. This confirms the reliability of our radio source detections and flux density estimates.
Table 2 shows the full radio catalog which reports (in column order): the source name; RA and Dec.(J2000); errors in RA and Dec.; the peak flux density $`S_P`$; error in $`S_P`$; the total flux density $`S_T`$; error in $`S_T`$; the full width half maximum (FWHM) of the major and minor axes $`\theta _M`$ and $`\theta _m`$ (determined from Gaussian fits); the position angle $`PA`$ of the major axis (in degrees); and the rms errors associated with $`\theta _M`$, $`\theta _m`$, and $`PA`$ respectively. The different components of multiple sources are labelled “C1” and “C2”. In Figure 3, we show the distribution of peak flux densities and the total to peak flux ratio as a function of peak flux for all sources. Sources with ratios $`S_T/S_P<1`$ in Figure 3b are primarily due to uncertainties on measured fluxes as estimated from the two dimensional Gaussian fits. Uncertainties in peak fluxes are typically $`10\%`$, while total flux estimates are more uncertain due to a sensitive dependence on the Gaussian fitting procedure. Errors in total fluxes are typically $`30\%`$. Contour maps of radio sources with available optical data are shown in Figure 8.
### 2.4 Comparison with the NVSS Catalog
The $`33\times 33\mathrm{arcmin}^2`$ region that we observed with the VLA was also covered by the NRAO VLA Sky Survey (Condon et al. 1998, NVSS). The NVSS covers the sky north of $`\delta =40\mathrm{°}`$ at 1.4 GHz with $`45\mathrm{}`$ resolution and limiting flux density of $`2.25`$ mJy ($`5\sigma `$). To this limit, 17 of our sources were found to be in common with the NVSS public catalog. One source in our catalog however is a double component source and is unresolved by the NVSS. A comparison of flux densities derived in this study with those from this catalog is shown in Figure 4. It is evident that our derived flux densities are on average slightly lower than those from the NVSS. This was also reported in a larger comparison study by Ciliegi et al. (1999) using a similar observational set-up and can be explained by the difference in resolutions used in the two surveys: $`15\mathrm{}`$ vs. $`45\mathrm{}`$. High-resolution surveys tend to miss flux from low-surface-brightness emission. Although the effect is only marginal for bright sources, it may become important for attempts to detect faint, low-surface-brightness objects at high redshift.
## 3 OPTICAL PHOTOMETRIC OBSERVATIONS
### 3.1 Observations and Data Reduction
Optical CCD photometry of our radio field was carried out on the 5-m Hale telescope at Palomar<sup>1</sup><sup>1</sup>1Operated by California Institute of Technology, Pasadena, CA 91125 Observatory during the nights of 27-29 August 1997. These fields were initially selected for subsequent deep mid-IR imaging with the Wide Field Infrared Explorer mission (WIRE), but the mission failed to perform to expectations. The Carnegie Observatories Spectroscopic Multislit and Imaging Camera (COSMIC; Kells et al. 1998) mounted at the $`f/3.5`$ prime focus with a TEK 2K CCD was used to image nine $`9.7\times 9.7\mathrm{arcmin}^2`$ fields in the Gunn $`r`$ (6550Å) filter. Each optical field comprised of three optical pointings offset by $`2\mathrm{}`$, each with integrations of 600 sec giving a total 1800 sec per field. This resulted in a limiting magnitude of $`r25`$ mag ($`5\sigma `$). The seeing was typically $`11.4\mathrm{}`$ (FWHM). The optical fields do not cover our entire $`33\times 33\mathrm{arcmin}^2`$ radio field. The nine slightly overlapping optical fields correspond to an areal coverage $`27.5\times 27.5\mathrm{arcmin}^2`$, or about 70% of the radio map.
The CCD data were reduced with standard tasks in the IRAF package. Frames were first bias-subtracted, then flat-fielded using dome flats. Bad pixels and columns were removed by interpolating between adjacent pixels and lastly, the individual dithered frames were median combined to remove cosmic ray hits. Calibration was performed using standard stars from Thuan & Gunn (1976). These were used to correct for atmospheric extinction from varying airmass and provide the instrumental zero point. Photometric uncertainties, estimated using these standards, are no more than $`0.05`$ mag.
Sources in the reduced optical frames were extracted using the DAOPHOT package in IRAF (Stetson 1987). This package has the benefit of performing photometry in crowded fields, which is the case in most regions of our optical fields. It performed the following steps: first, sources were extracted above a given threshold, given as a multiple of the total CCD noise (sky and read noise, $`\sigma _{tot}`$). We adopted threshold values of 4.5-5$`\sigma _{tot}`$, just high enough to avoid large numbers of spurious detections. Second, simple aperture photometry was performed on these identified sources. This required a specification of the aperture radius which is likely to contain most of the light of our target source, and width of a surrounding annulus to estimate and subtract the sky background. We adopted a radius of six arcsec and annulus width of four arcsec. Next, a point spread function (PSF) was determined in each of our nine fields. This involved an interative technique to remove contamination from neighbouring sources in crowded fields near our PSF candidates. Simultaneous PSF fitting on all sources was then performed to identify sources which were previously hidden in crowded regions. Finally, the magnitudes determined from PSF fits were aperture corrected to a common aperture size as used on our standard stars. Aperture corrections were typically 0.22 mag. A final visual inspection removed any spurious detections. A total of $`300`$-390 sources were extracted from each $`9.7\times 9.7\mathrm{arcmin}^2`$ field to a limit of $`r25`$ mag. Previous optical surveys find typically 380-520 sources within this area to this limit, and the variation is primarily due to clustering. Such fluctuations are found to be significant on such relatively small scales (eg. Metcalfe et al. 1991).
### 3.2 Astrometry
The astrometry on the optical images was based on 10-12 APM catlog stars in each field (Maddox et al. 1990). The ccmap and cctran tasks in the IRAF immcoords package were used to compute plate solutions relating pixel positions to astrometric coordinates. Astrometric coordinates for all sources on the frames were then determined. By comparing the positions of several sources common to the APM catalog and our fields, we found that our rms position uncertanties are typically $`0.9\mathrm{}`$.
## 4 OPTICAL/NEAR-INFRARED IDENTIFICATION OF RADIO SOURCES
### 4.1 Method for Optical Identification
We assigned optical identifications and estimate their reliability using a robust likelihood ratio ($`LR`$) analysis. This general method has frequently been used to assess identification probabilities for radio and infrared sources (eg. de Ruiter, Willis & Arp 1977; Prestage & Peacock 1983; Sutherland & Saunders 1992; Lonsdale et al. 1998). The method, which computes the probability that a suggested identification is the ‘true’ optical counterpart, is outlined as follows:
For each optical candidate $`i`$ in the search area of some radio source $`j`$, we calculated the value of the dimensionless difference in radio and optical positions:
$$r_{ij}=\left[\frac{(\alpha _i\alpha _j)^2}{\sigma _{\alpha _i}^2+\sigma _{\alpha _j}^2}+\frac{(\delta _i\delta _j)^2}{\sigma _{\delta _i}^2+\sigma _{\delta _j}^2}\right]^{1/2},$$
(1)
where the $`\alpha `$’s and $`\delta `$’s represent right ascensions and declinations respectively, and $`\sigma `$’s standard deviations. We chose a moderately large search radius of $`10\mathrm{}`$ to allow for the maximal position uncertainties: $`\sigma _{\mathrm{opt}}1.4\mathrm{}`$ and $`\sigma _{\mathrm{rad}}1.5\mathrm{}`$ (assuming $`5\sigma _{eff}`$, where $`\sigma _{eff}=(\sigma _{\mathrm{opt}}^2+\sigma _{\mathrm{rad}}^2)^{1/2}`$). Such a radius is also small enough to avoid large numbers of chance associations.
Given $`r_{ij}`$, we must now distinguish between two mutually exclusive possibilities: (1) the candidate is a confusing background object that happens to lie at distance $`r_{ij}`$ from the radio source or (2) the candidate is the ‘true’ identification that appears at distance $`r_{ij}`$ owing solely to radio and optical position uncertainties. We assume that the radio and optical positions would coincide if these uncertainties were zero. This assumption however is not valid when the centroid of an extended radio source is used, and is further discussed below.
To distinguish between these cases, we compute the likelihood ratio $`LR_{ij}`$, defined as:
$$LR_{ij}=\frac{\mathrm{exp}\left[r_{ij}^2/2\right]}{2\pi N(<m_i)\left[(\sigma _{\alpha _i}^2+\sigma _{\alpha _j}^2)(\sigma _{\delta _i}^2+\sigma _{\delta _j}^2)\right]^{1/2}},$$
(2)
where $`N(<m_i)`$ is the local surface density of objects brighter than candidate $`i`$. The likelihood ratio $`LR_{ij}`$ is simply the ratio of the probability of an identification (the Rayleigh distribution: $`r\mathrm{exp}(r^2/2)`$), to that of a chance association at $`r`$ ($`2\pi N(<m_i)\sigma _\alpha \sigma _\delta `$). $`LR_{ij}`$ therefore represents a ‘relative weight’ for each match $`ij`$, and our aim is to find an optimum cutoff $`LR_c`$ above which a source is taken to be a reliable and likely candidate. Its advantage over alternative methods (purely based on finding the lowest random chance match, eg. Downes et al. 1986) is that it allows for a possible distant candidate to still be the ‘true’ identification even when there is still a high chance of it being a spurious background source.
It is important to note that our form for $`LR_{ij}`$ (eq. 2) slightly differs from that used by earlier studies (eg. Lonsdale et al. 1998) in that it doesn’t contain the multiplicative “$`Q`$” factor in the numerator. This factor represents the apriori probability that a “true” optical counterpart brighter than the flux limit exists amongst the identifications. For our purposes, we will treat $`LR_{ij}`$ as simply a a relative weight measure for each radio-optical match, just for the purposes of assigning an optimal cutoff for reliable identification (see below). We are not concerned with its absolute value, which is required when computing formal probability measures from $`LR`$. For simplicity, we have therefore set $`Q=1`$ in this work.
The optical surface density as a function of magnitude to be used in computing $`LR`$ was determined from the total number of objects visible in our optical frames. The variation in surface density in the vicinity of each radio source caused by possible clustering effects was found to be small: no more than 5% on $`23\mathrm{}`$ scales.
The distribution of $`LR`$ values for all possible radio source-candidate matches is shown by the shaded histogram in Figure 5. Following Lonsdale et al. (1998), we generate a truly random background population with respect to the radio sources by offseting the radio source positions by $`30\mathrm{}`$. $`LR`$ values for each radio source were then re-computed and their distribution is given by the thick-lined histogram in Figure 5. A comparison of the number of associations for (true) radio source positions with the number of associations found for random (offset) positions will enable us to determine a critical value $`LR_c`$ for reliable identification. From these distributions, we compute the reliability as a function of $`LR`$:
$$R(LR_{ij})=1\frac{N_{random}(LR_{ij})}{N_{true}(LR_{ij})},$$
(3)
where $`N_{true}`$ and $`N_{random}`$ are the number of true and random associations respectively. The reliability computed in this way also represents an approximate measure of the identification probability for a candidate with given $`LR`$.
Figure 6 illustrates the reliability as a function of $`LR`$. Above $`\mathrm{log}(LR)0.5`$, the reliabilities are $`85\%`$ because few random associations exceed this value of $`LR`$ (Fig. 5). As a good working measure we therefore assume a cutoff $`\mathrm{log}(LR_c)0.5`$ above which a source is taken to be a likely candidate. The determination of reliabilities via the $`LR`$ method is insensitive to variations in $`N(<m_i)`$ across the field or uncertainties in its derivation, and also the assumption of Gaussian error ellipses in the radio and optical positions. Such uncertainties are “normalised out” when one computes the ratio of random to true number of associations within a search radius when estimating the reliability (eq. 3). Lonsdale et al. (1998) have shown that the absolute value of $`LR`$ itself depends on the characteristics of the source population being identified (eg. stars versus galaxies). Different populations (assuming they could be classified a-priori using some diagnostic) map into different underlying surface densities at the ‘identifying’ wavelength, implying that distributions in $`LR`$ (eq. 2) will also be different. For a robust determination of the reliability in such situations, see Lonsdale et al. (1998).
There are two complications to consider in the above method. The first is when one attempts to identify extended (or resolved) radio sources with this method. For all radio sources, we have used the positions of centroids derived from two-dimensional Gaussian fits in computing the $`LR`$ for optical candidates. For unresolved sources with (Gaussian fitted) sizes $`15\mathrm{}`$ (the synthesised FWHP beam width), the source is likely to have a compact central component and the optical position is expected to lie close to its quoted radio centroid. For an extended (resolved) source however, the radio and optical positions may differ considerably since errors in the radio centroids are only $`2\mathrm{}`$. In such cases, if $`LR<LR_c`$ the identification may still be valid, since its low $`LR`$ value could purely be due to a real large positional offset. The second complication is when a radio source has more than one optical candidate within its search radius with $`LR>LR_c`$. This occured in about 20% of cases and was primarily due to contaminating stars. We assess these ambiguities and increase the robustness of our identifications by visually examining all optical candidates according to the following criteria:
1. If candidates have $`LR<LR_c`$ for a radio source with $`\theta _{min,max}15\mathrm{}`$, then identification is classified as uncertain.
2. If candidates have very low reliability, $`LRLR_c`$ (for unresolved radio sources), or there are no objects in the search radius, then radio source is classified as empty field.
3. If $`LRLR_c`$, i.e. where reliability is moderately “low”, then identification is also uncertain.
4. If more than one optical candidate exists with $`LR>LR_c`$, then only source(s) with extended (galaxy-like) optical profile is taken as the identification. Point sources associated with quasar nuclei are not considered in our identification scheme due to their relatively low surface density compared to galaxies ($`1:4000`$) in sub-millijansky radio samples.
5. For unique, $`LR>LR_c`$ candidates, its optical profile is also checked for confirmation.
### 4.2 Results
Of the 62 radio sources, 43 lie within our optically imaged $`27.5\times 27.5\mathrm{arcmin}^2`$ field. We found optical identifications for 26 to $`r25`$ mag with reliabilities $`R_{id}80\%`$. Four sources have identifications classified as uncertain owing to a moderately low identification reliability of $`R_{id}78\%`$ (and $`\mathrm{log}(LR)0.4`$). Five more are uncertain because they have extended radio structure and large possible positional offsets between optical and radio centroid positions. Eight radio sources lie in “definite” optical empty fields with no candidates brighter than $`r25`$ mag.
Other optical follow-up studies found similar results. Georgakakis et al. (1999) identified $`47\%`$ of sources to $`R=22.5`$ mag from the Phoenix Deep Survey ($`S_{1.4\mathrm{GHz}}>0.2`$ mJy) (Hopkins et al. 1998). Deeper indentifications of sources as faint as $`S_{1.4\mathrm{GHz}}40\mu `$Jy from Hubble Deep Field images revealed a 80% success rate to $`I=25`$ mag (Richards et al. 1999). Ignoring the uncertain identifications in our study (from criteria 1 and 3 above and which are excluded from our analysis), we find that $``$18% of our sources are unidentified to $`r25`$. Accounting for differences in bandpasses and sensitivity, this is broadly consistent with the above studies. Figure 7 shows the distribution of apparent magnitude $`r`$ for all reliable (robust) and uncertain identifications in our sample.
### 4.3 Near-Infrared Data
Near-infrared data in the $`J`$(1.25$`\mu `$m), $`H`$(1.65$`\mu `$m) and $`K_s`$(2.17$`\mu `$m) photometric bandpasses were obtained from the Two Micron All Sky Survey (2MASS) project database. For multi-band detection of point sources, this survey is currently scanning the sky to sensitivities 16.5, 16.0 and 15.5 mag at signal-to-noise ratios $`7`$, $`5`$ and $`7`$ respectively in $`J`$, $`H`$ and $`K_s`$. The data relevant to this study are not yet released in the public catalogs, and was retrieved from the ‘internal working database’ at IPAC<sup>2</sup><sup>2</sup>2The Infrared Processing and Analysis Center, California Institute of Technology.. Photometry in this database was determined using custom PSF-fitting software and algorithms are described in Cutri et al. (2000). Since such data have not been subjected to the rigorous quality assurance as that in the public release catalogs, we have examined individual images for quality and any possible systematic uncertainties in the photometry.
To maximise the possible number of detections, we searched for near-infrared counterparts to each radio source with a conservatively low signal-to-noise ratio threshold of $`5`$ in each band. In cases where a source was detected in only one or two of the three bandpasses: $`J`$, $`H`$ and $`K_s`$, we note its “band-filled” 95% confidence upper-limit in the undetected band. In other words, the 2MASS catalog also reports upper limits to measurements in an undetected band by placing an aperture over the position inferred from detections in other bands. The image pixel scale of the 2MASS detectors is $`2.0\mathrm{}`$ and the positional uncertainties are $`0.5\mathrm{}`$. Owing to the relatively shallow flux limits of the 2MASS survey, the background source surface density is low enough that chance associations with radio positions are very unlikely.
We searched the 2MASS database for near-infrared counterparts to our 43 radio sources that have available optical information from our deep optical survey and examined their images for quality. We found 7 reliable matches with 6 detections in $`J`$, 7 in $`H`$, and 4 in $`K_s`$ at the $`5\sigma `$ level. Upper limits were available for the remainding “band-filled” values.
Results of our optical identification analysis and available near-infrared data are shown in Table 3. In column order, this table reports: the radio source name (see Table 2); RA and Dec.(J2000) of the optical counterpart of the radio source; optical-radio position separation ($`\delta _{radopt}`$) in arcsec; logarithm of the liklihood ratio ($`LR`$); reliability of the optical identification (see eq.); apparent $`r`$-band magnitude; $`J`$, $`H`$ and $`K_s`$ magnitudes with errors or $`2\sigma `$ upper limits; $`rH`$ color; $`rK_s`$ color, and last, the optical morphology if the optical counterpart is visually resolved with size $`5\theta _{FWHM}(PSF)`$.
Optical morphologies were determined from light profile fitting of individual sources and comparison with $`R^{1/4}`$-law (elliptical-like) and exponential (disk-like) profiles. In cases where a disturbed or interacting morphology is apparent, then it is designated to have an irregular (labelled as ‘Irr’) morphology. Sources unresolved in the optical with typically $`5\theta _{FWHM}(PSF)`$ are designated as ‘unknown’ and labelled as ‘?’ in Table 2.
### 4.4 Optical and Radio Map Overlays
In Figure 8 we show the optical image – radio map overlays for the 43 radio sources with available optical information. A visual inspection of the optical images of resolved counterparts shows a diverse morphological mix, consistent with previous studies. About 40% of our optical identifications have elliptical/disk-like morphology, while $`10\%`$ can be identified as exhibiting peculiar (either interacting or disturbed) morphologies. It is important to note that these ‘disturbed’ sources are based on visual inspection alone and their morphology could still be uncertain until future spectroscopy confirms their nature. The elliptical/disk hosts also tend to be associated on average with sources of relatively brighter radio flux density ($`2`$ mJy) than the irregular class. This is consistent with previous radio-optical identification studies which find an increasing fraction of irregular-type galaxies at $`S_{1.4}2`$ mJy (eg. Kron, Koo, & Windhorst 1985, Hammer et al. 1995, Gruppioni et al. 1998) and a significant number of elliptical galaxies hosting the brighter extended radio galaxies and AGN (Condon 1989).
A further observation is the unique radio structure exhibited by our eight optical empty field sources with $`r>25`$ mag. These are represented by maps in Fig. 8 labelled by the letter “E”. All show compact and symmetric (presumably unresolved) structures and could represent either of the following: distant (possibly dusty) AGN where with our radio sensitivity, we could have detected a nominal FR-I galaxy to $`z1.3`$, or, nearby, compact dusty starbursts at $`z0.3`$ as constrained by typical starburst luminosities: $`L_{1.4\mathrm{GHz}}10^{23}\mathrm{W}\mathrm{Hz}^1`$. The second explanation for the nature of the empty fields is more plausible, given that the majority of submillijansky radio sources are associated with star-forming galaxies and less than 5% are usually identified with bright FR-Is at $`z1`$ (Kron, Koo & Windhorst 1985).
## 5 ANALYSIS OF RADIO AND OPTICAL–NEAR-INFRARED COLORS
This section presents an analysis of flux ratios between the radio, near-infrared, and optical bands to explore possible contributions from AGN and starbursts to the observed radio emission as well as the importance of absorption by dust. Because we lack spectroscopic information, our analysis treats the sub-mJy sources as one homogeneous population and uses a simple stellar synthesis model to interpret its properties quantitatively.
### 5.1 A Simple Synthesis Model
We can predict the radio–to–optical(–near-IR) flux ratios and $`rK`$ colors for a range of galaxy types using the stellar population synthesis code of Bruzual & Charlot (1993) (hereafter BC). On its own however, the BC model does not directly predict the amount of radio emission expected from a star-forming galaxy, which we need for the determination of flux ratios involving the radio band. We do this by relating the star-formation rates derived from empirical calibrations involving the UV and radio bands as follows.
The 1.4 GHz radio emission from star-forming systems is believed to be primarily synchrotron emission from cosmic rays accelerated in supernova remnants plus a small ($`10\%`$) thermal contribution from HII regions (Condon & Yin 1990, Condon 1992). Thus to a good approximation, the radio luminosity is taken to be proportional to the formation rate of stars with $`M>5M_{\mathrm{}}`$:
$$SFR(M5M_{\mathrm{}})=\frac{L_{1.4\mathrm{GHz}}}{4\times 10^{28}\mathrm{erg}\mathrm{s}^1\mathrm{Hz}^1}M_{\mathrm{}}\mathrm{yr}^1.$$
(4)
(Condon 1992). These same massive stars will also contribute significantly to the UV continuum emission in the range $`1200`$-$`2500`$Å. In particular, there have been many different calibrations of the SFR from the UV-flux. For a Salpeter initial mass function (IMF) from $`mM_{\mathrm{}}`$ to 100$`M_{\mathrm{}}`$, the calibration of Madau et al. (1998) (which assumes no dust correction) yields
$$SFR(MmM_{\mathrm{}})=Q_m\left(\frac{L_{UV}}{7.14\times 10^{27}\mathrm{erg}\mathrm{s}^1\mathrm{Hz}^1}\right)M_{\mathrm{}}\mathrm{yr}^1.$$
(5)
We have modified the initial relation of Madau et al. (1998) to include the factor $`Q_m`$, which represents the fraction of stellar masses contributing to the observed SFR,
$$Q_m=\frac{_{mM_{\mathrm{}}}^{100M_{\mathrm{}}}M\psi (M)𝑑M}{_{0.1M_{\mathrm{}}}^{100M_{\mathrm{}}}M\psi (M)𝑑M},$$
(6)
where $`\psi (M)M^x`$ is the IMF. For $`m=0.1`$, we have $`Q_m=1`$. Assuming a Salpeter IMF ($`x=2.35`$), we find that for stars with $`M>5M_{\mathrm{}}`$, $`Q_m0.18`$. With this fraction, and equating the two SFR calibrations (eqs and ), we find that the luminosity densities at 1.4 GHz ($`L_{1.4GHz}`$) and $`2100`$Å ($`L_{UV}`$) are very nearly equal. We therefore assume that the rest frame 1.4 GHz flux density is directly given by the flux density at $`2100`$Å as specified by the synthesized model spectrum. In general terms, the observed radio flux (in the same units as the synthesised UV spectrum) can be written:
$$f_\nu (1.4\mathrm{GHz})_{obs}=(1+z)^{1\alpha }f_\nu (2100\mathrm{\AA })_{rest},$$
(7)
where $`\alpha 0.8`$ is the radio spectral index (Condon 1992) and $`f_\nu (2100\mathrm{\AA })_{rest}`$ is the rest frame (unreddened) UV spectral flux. We must emphasise that this relative radio flux is only that associated with the star-formation process. Possible additional contributions, such as contaminating AGN, are not considered in this model.
We calculated flux ratios involving the radio, near-infrared and optical bands using evolutionary synthesis models for ellipticals (E/SO), early (Sab/Sbc) and late (Scd/Sdm) type spirals, and “very blue” starbursts (SB). These are meant to represent the possible contributions to the sub-mJy radio sources, and each class is defined by a characteristic star formation rate as a function of time. As supported by local observations (eg. Gavazzi & Scodeggio 1996), we assumed that E/SO and Sab/Sbc galaxies have an exponentially decaying SFR of the form $`\psi (t)\tau ^1\mathrm{exp}(t/\tau )`$, where $`\tau `$ is the e-folding time. Values of $`\tau =1`$ and $`\tau =8`$ Gyr were adopted for the E/SO and Sab/Sbc galaxies respectively. For late-type spirals (Scd/Sdm) and young starbursts (SB), we assumed constant SFRs with different ages. All models used to generate the spectral energy distributions (SEDs) are summarised in Table 4. The models assume $`\mathrm{H}_0=50\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1`$, $`q_0=0.5`$ and a galaxy formation redshift $`z_f=10`$, which corresponds to an age of 12.7 Gyr at $`z=0`$.
To explore the effects of dust on our flux ratios and colors, each model SED was reddened in the source rest frame with an extinction curve $`\xi (\lambda )A_\lambda /A_V`$ characteristic of the Small Magellanic Cloud (SMC), which appears to be a good approximation for the ISM of nearby galaxies (Calzetti et al. 1994). This approximation is most accurate for the reddest wavelengths of starburst galaxy spectra ($`7000`$Å-$`3\mu `$m), although breaks down severely at $`\lambda 2500`$Å(Calzetti et al. 1999). We have used the analytical fit for $`\xi (\lambda )`$ as derived by Pei (1992) for the SMC. For simplicity, we assumed that the dust is distributed in a homogeneous foreground screen at the source redshift.
### 5.2 Data and Model Comparions
In Figure 9 we plot the radio–to–optical flux ratio, $`R(1.4/r)`$, defined as
$$R(1.4/r)=\mathrm{log}(S_{1.4}/\mathrm{mJy})+0.4r,$$
(8)
where $`S_{1.4}`$ and $`r`$ are the radio flux and optical $`r`$-band magnitude respectively. The distribution seen in observed values of $`R(1.4/r)`$ is larger at the faintest optical magnitudes $`r>21.5`$. There are few galaxies however with $`r<21.5`$, and nonetheless, the scatter at $`r>21.5`$ is consistent with that found at $`r19`$ in a larger follow-up study of sub-mJy radio sources by Georgakakis et al. (1999).
Figure 10 shows $`R(1.4/r)`$ as a function of $`rH`$ color for all sources with available optical and near-infrared data. Our reason for using $`rH`$ color is that the $`H`$ band yielded more “definite” detections than the other near-infrared bandpasses. Although the numbers are still relatively small, this facilitates the best comparison with the synthesis models. The predictions for four galaxy types (see Table 4 and §5.1) are shown for no dust reddening (thin curves) and a rest-frame extinction $`A_V=2.5`$ mag. The morphological mix of data shows a relatively large dispersion in $`rH`$ color that is more consistent with the range predicted by the models that include dust. This suggests that on average, the optical–to–near-infrared continua of most sources in Fig. 10 are reddened by a uniform (possibly “optically thin”) dust component with $`A_V22.5`$ mag absorption. This measure is consistent with spectroscopic studies of nearby starbursts by Calzetti et al. (1996), Meurer et al. (1997) and photometric modelling by Nakata et al. (1999).
We must emphasise that our models only account for radio emission produced from star-formation processes. The sources labelled as elliptical (or early type) in Fig. 10 are not expected to lie on any of the star-formation derived locii. An AGN most likely dominates their radio emission. We include them here merely as a comparision, and their relationship to normal starbursts is discussed further below.
The sources in Fig. 10 appear more-or-less consistent with the dusty “0.1 Gyr starburst” model. This could in principle apply to the two sources with spiral/disk-like morphology (labelled “S”), but is unconventional for the five elliptical morphologies (see below). A comparison with radio–to–near-infrared flux ratios further constrains the underlying properties of these sources. Figure 11 shows the radio–to–near-infrared flux ratio $`R(1.4/H)`$, defined analagous to eq., as a function of $`rH`$ color. The near-infrared emission is dominated by old stars and is less affected by dust than the optical. The radio–to–near-infrared flux ratio should therefore be relatively insensitive to dust. Given the simplicity of our models, the two disk-like sources may not necessarily represent “0.1 Gyr starbursts” as indicated in Fig. 11. They could also belong to the Sab/Sbc or Scd/Sdm classes. For this to be true, the following additional components may play an important role in more ‘realistic’ models: ‘optically-thick’ dust that completely obscures both the $`H`$ and $`r`$ emission without causing appreciable reddening in $`rH`$ color, or, contamination by at least an order of magnitude times more radio emission from a central AGN than that produced purely by supernovae. The second possibility is favored by radio observations of a number of luminous infrared galaxies by Norris et al. (1988), where some showed evidence for significant radio emission from compact Seyfert-like nuclei.
The large discrepancy between the four sources with elliptical morphologies (labelled “E”) and predictions from the early-type E/SO models suggests the importance of a significant AGN contribution to the radio emission. Appreciable amounts of optically-thick dust suppressing the optical and near-infrared light (except for extinction by diffuse, optically-thin dust) is not favored by observations (eg. Goudfrooij & de Jong 1995). Most, if not all of these ellipticals are likely to be radio-powered by AGN. At the limiting sensitivity of our radio survey ($`0.3`$ mJy), a nominal FR-I galaxy (with $`L_{1.4}=10^{24}\mathrm{WHz}^1`$) could be detected to $`z1.3`$, and indeed the spread in $`rH`$ color for the ellipticals in Fig. 11 is consistent with the E/SO (optically-thin dust) model to this redshift. A comparison between the model $`R(1.4/H)`$ values with actual observed values implies that such AGN will contribute a factor $`10^2`$ times more radio emission than that produced by any underlying star formation activity in these systems. It is important to note that the ratio of AGN-to-stellar powered radio activity has a huge spread for the elliptical population in general, and that the factor $`10^2`$ only illustrates a property specific to the ellipticals in our radio sample.
To summarise, our use of a simple synthesis model that includes radio emission and dust reddening to analyse the properties of sub-mJy radio sources has shown the following: first, the presence of dust with extinctions $`A_V2`$ mag and possibly greater, consistent with previous more direct determinations, and second, that the level of radio emission from non-stellar processes such as AGN could be easily inferred and constrained. This will be particularly important for starbursts hosting Seyfert nuclei where a comparison with more sophisticated dust models may be required to infer the relative contributions.
## 6 A Method to Select “ULIGS” via Radio/Optical Color
Since the emission (and dust absorption) properties from normal galaxy populations are reasonably well known, a color-color diagram such as Fig. 10 could provide a potential diagnostic for selecting ultraluminous infrared galaxies (ULIGS) to high redshift. The relatively low sensitivity of the Infrared-Astronomical Satellite (IRAS) has primarily confined ULIG selection to the local Universe (Sanders & Mirabel 1996), although there is some speculation that recently discovered faint “SCUBA” sources at sub-millimetre wavelengths could represent their high-redshift counterparts (eg. Blain et al. 1999). Approximately 80% of local ULIGS are believed to be powered by starbursts and the remainder show evidence for an AGN contribution (Genzel et al. 1998; Lutz et al. 1998). Far-infrared observations have shown that dust and molecular gas in local ULIGS is concentrated in compact regions $`1`$kpc (Okumura et al. 1991, Bryant 1996) and that a large fraction of the optical/UV emission is hidden by optically-thick dust (Sanders et al. 1988). A study of their properties and importance to galaxy evolution therefore requires observations at wavelengths virtually immune to dust absorption. Radio frequencies provide an excellent window of opportunity.
Figure 12 illustrates the predicted locus in color-color space using the synthetic SEDs of three local far-IR selected systems: Arp 220 ($`L_{IR}1.6\times 10^{12}L_{\mathrm{}}`$) - a ULIG undergoing a powerful starburst as seen via high resolution radio observations by Smith et al. (1998); M82 ($`L_{IR}6\times 10^{10}L_{\mathrm{}}`$) - a system undergoing a weak-to-moderate starburst, and Mrk 273 ($`L_{IR}2.6\times 10^{12}L_{\mathrm{}}`$) - a ULIG whose bolometric emission is believed to be dominated by a hidden central AGN from the presence of strong Seyfert-2 lines and moderately strong hard X-ray (2-10 keV) emission (Turner et al. 1997). We have used the synthetic SEDs generated by Devriendt et al. (1999) to model the starburst emission. These authors used a self-consistent modelling approach to predict the stellar optical/UV/near-IR emission, its reprocessing into the mid-IR–to–sub-mm by dust, and the nonthermal stellar-powered radio emission based on the empirical radio–to–far-IR luminosity correlation. Due to its strong AGN-dominated nature, the starburst synthetic SED predicted by Devriendt et al. for Mrk 273 differs appreciably from that observed in the radio. For this source, we therefore used the Devriendt et al. SED at wavelengths $`\lambda <1`$mm and extrapolated into the radio using its actual observed radio–to–1mm spectral slope and fluxes (obtained from the NASA/IPAC Extragalactic Database <sup>3</sup><sup>3</sup>3The NASA/IPAC Extragalactic Database (NED) is operated by the Jet Propulsion Laboratory, California Institute of Technology, under contract with the National Aeronautics and Space Administration.).
Figure 12 shows that a galaxy characteristic of the (low-IR luminous) M82 system will occupy a region similar to that occupied by normal galaxies in this study (and also their predicted synthetic colors in Fig. 10). Luminous systems classified as ULIGS however (Arp 220 and Mrk 273), will tend have higher radio–to–optical flux ratios which could be easily selected. This can be explained by the well-observed correlation between far-infrared luminosity and far-IR ($`60100\mu `$m)–to–optical spectral slope (Soifer et al. 1987). Consequently, the most IR-luminous systems with the largest far-IR–to–optical ratios are also likely to have a high level of radio-emission due to its strong correlation with IR luminosity. This will lead to a larger than average radio–to–optical flux ratio for ULIGS in general as shown in Fig. 12.
The existence of systems with either larger rest-frame optical/UV extinction or excess AGN contribution to the radio than the ULIGS considered here will be shifted further upwards on this plot. Diagnostics to distinguish between AGN and starburst dominated ULIGS using radio-to-optical color alone will not be trivial and is left to a future study. The three ULIGS in Fig. 12 represent a range of known ULIGS and their location on this plot simply serves as a diagnostic to pre-select ULIG candidates for further study.
A system like Arp 220 (with $`\nu L_\nu (1.4\mathrm{GHz})2.5\times 10^6L_{\mathrm{}}`$) could be observed to redshift $`z1.6`$<sup>4</sup><sup>4</sup>4Assumes $`q_0=0.5`$, $`H_0=50\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1`$ if initially selected from a radio survey limited to $`S_{1.4\mathrm{GHz}}50\mu `$Jy. Thus, to limiting sensitivities reached by existing 1.4 GHz surveys, such a method may not probe the highest redshifts. Nonetheless, as shown in Fig. 12, such systems could still be well separated from normal galaxies to this redshift. Assuming an Arp 220-like SED and moderate luminosity evolution ($`L_{60\mu \mathrm{m}}(1+z)^{2.5}`$), the surface density of ULIGS to $`z1.6`$ is expected to be of order $`150\mathrm{deg}^2`$ at $`S_{1.4\mathrm{GHz}}50\mu `$Jy, or about $`6\%`$ of the integral count to this sensitivity (Richards 2000). They should therefore exist in significant numbers in deep large-area radio surveys.
## 7 SUMMARY AND CONCLUSIONS
We have used the VLA radio telescope to image a contiguous $`33\times 33\mathrm{arcmin}^2`$ area to a (mean) limiting ($`5\sigma `$) sensitivity of $`0.35`$ mJy. From a total of 62 detections, the results of optical and near-infrared photometry are reported for 43 sources. Our optical photometry is more sensitive than previous optical follow-up studies of radio surveys of similar depth. Our main findings are:
(1) We have used a robust, likelihood-ratio method for determining optical identifications and their reliability. This method is seldom used in identification studies and is insensitive to assumptions concerning fluctuations in background source density and Gaussian error distributions. We assigned optical candidates to 26 radio sources with reliability $`80\%`$. Nine radio sources are uncertain and/or ambiguous, and eight are empty fields. Near-infrared photometry from the 2MASS database was reported for 7 sources.
(2) The eight optical empty field sources all display compact and symmetric radio morphologies and most probably represent compact starbursts at $`z0.3`$ strongly obscured by dust. They may require at least 4 magnitudes optical extinction to account for their large radio–to–optical flux ratio compared to the identified population. Our conclusion for them being ‘compact starbursts’ is very tentative as it is purely based on starburst versus AGN number statistics expected from sub-mJy radio surveys. Further deep infrared/optical imaging and spectroscopy will be necessary.
(3) Consistent with previous studies, our deep ($`r25`$) optical imaging shows that the optical appearence can be divided into two classes according to radio flux-density: elliptical-like morphologies for $`2`$ mJy, and peculiar or disturbed for $`2`$ mJy.
(4) Using a stellar synthesis model which includes radio emission and dust reddening, we find that the near-infrared–to–optical emission in a small, bright sub-sample is reddened by ‘optically thin’ dust with $`A_V22.5`$ mag, regardless of morphological type. This appears consistent with other more direct determinations. Consistent with previous studies, the radio emission from early-type systems seems to be powered by AGN rather than star-formation to account for their anomalously large radio–to–optical(–near-infrared) ratios.
(5) Our analysis shows that a radio/optical or radio/near-IR color selection technique could provide a potential means for detecting ULIG-type objects to $`z1.6`$.
Despite the lack of spectroscopic information, our study of a homogeneous population of faint radio sources has stressed the importance of dust on studies of intrinsic galaxy properties and their evolution at optical wavelengths. A future goal would be to obtain spectra, or multi-color optical/near-infrared photometry to better explore these sources and the validity of the simple stellar synthesis models presented in this paper. The ever improving resolution (and sensitivity) capabilities of optical/near-IR detectors over those feasible at (the longest) radio wavelengths requires robust identification techniques to better ascertain their properties. Likelihood-ratios provide one such technique. The present study complements other deep optical studies of faint radio sources to constitute a statistically significant sample for inferring their nature and importance to galaxy evolution.
FJM thanks Glenn Morrison and JoAnn O’Linger for valuable assistance with the data reduction and Rosalie Ewald for assistance with radio/optical image overlays. We thank the staff at Palomar Observatory for technical assistance during the observing run. This publication makes use of data products from the Two Micron All Sky Survey, which is a joint project of the University of Massachusetts and the Infrared Processing and Analysis Center/California Institute of Technology, funded by the National Aeronautics and Space Administration and the National Science Foundation. The National Radio Astronomy Observatory is operated by Associated Universities, Inc., under cooperative agreement with the National Science Foundation. This research has made use of the NASA/IPAC extragalactic database (NED) which is operated by the jet propulsion laboratory, caltech, under contract with the national aeronautics and space administration. FJM acknowledges support from a JPL/NASA postdoctoral fellowship grant. |
warning/0002/hep-ex0002025.html | ar5iv | text | # 𝑅-values in Low Energy 𝑒⁺𝑒⁻ Annihilation
## 1 Introduction
According to the quark-parton model, hadrons produced via $`e^+e^{}`$ collisions are characterized by the annihilation of $`e^+e^{}`$ pairs into a virtual $`\gamma `$ or $`Z^0`$ boson. In the lowest order, one defines the ratio of the rate of hadron production to that for muon pairs as
$$R\frac{\sigma (e^+e^{}\text{hadrons})}{\sigma (e^+e^{}\mu ^+\mu ^{})}=3\underset{f}{}Q_f^2,$$
(1)
where $`Q_f`$ is the fractional charge of the quark, and the factor of 3in front counts the three colors for each flavor. $`\sigma (e^+e^{}\mu ^+\mu ^{})=4\pi \alpha ^2/3s`$ is the cross section of the pure QED process. The value of $`R`$, which counts directly the number of quarks, their flavor and colors, is expected to be constant so long as the center-of-mass (cm) energy of the annihilated $`e^+e^{}`$ does not overlap with resonances or thresholds for the production of new quark flavors. One has
$`R`$ $`=`$ $`3[(2/3)^2+(1/3)^2+(2/3)^2]=2\text{for}u,d,s`$
$`=`$ $`2+3(2/3)^2=10/3\text{for}u,d,s,c,`$
$`=`$ $`10/3+3(1/3)^2=11/3\text{for}u,d,s,c,b.`$
These values of $`R`$ are only based on the leading order process $`e^+e^{}q\overline{q}`$. However, one should also include the contributions from diagrams where the quark and anti-quark radiate gluons. The higher order QCD corrections to $`R`$ have been calculated in complete 3rd order perturbation theory , and the results can be expressed as
$$R=3\underset{f}{}Q_f^2[1+(\frac{\alpha _s(s)}{\pi })+1.411(\frac{\alpha _s(s)}{\pi })^212.8(\frac{\alpha _s(s)}{\pi })^3+\mathrm{}],$$
(2)
where $`\alpha _s(s)`$ is the strong coupling constant. Precise measurement of $`R`$ at higher energy can be employed to determine $`\alpha _s(s)`$ according to eq.2, which exhibits a QCD correction known to $`O(\alpha _s^3)`$ . In addition, non-perturbative corrections could be important at low cm energy, particularly in the resonance region.
$`R`$ has been measured by many laboratories in the energy region from hadron production threshold to the $`Z^0`$ pole and recently to the energy of W pair production . The experimental $`R`$-values are in general consistent with theoretical predictions, which is an impressive confirmation of the hypothesis of the three color degrees of freedom for quarks. The measurements of $`R`$ in the low energy region were performed 15 to 20 years ago in Novosibirsk, Orsay, Frascati, SLAC and Hamburg . Fig. 1 shows the $`R`$-values for cm energies up to 10 GeV, including resonances. For cm energies below 5 GeV, the uncertainties in $`R`$-values are about 15% on average; and the structure in the charm threshold region is not well determined. The DASP group inferred the existence of narrow resonances at 4.04 GeV and 4.16 GeV. In addition to the resonance at 3.77 GeV, Mark I data shows broad enhancements at 4.04, 4.2 and 4.4 GeV. The resonance at 4.4 GeV was also observed by PLUTO , but the height and width of the resonance were reported differently. A new cross section measurement in the charm threshold region is needed to clarify the structure, which is important not only for the precision determination of $`\alpha (M_Z^2)`$ and the interpretation of the $`(g2)_\mu `$ measurement of E821 at BNL but also for the understanding of charmonium itself.
Between the charm and bottom thresholds, i.e., about 5-10.4 GeV, $`R`$ was measured by Mark I, DASP, PLUTO, Crystal Ball, LENA, CLEO, CUSB, and DESY-Heidelberg collaborations. Their systematic normalization uncertainties were about 5-10%. Above bottom threshold, the measurements were from PEP, PETRA and LEP with uncertainties of 2-7%.
Remarkable progress has been made in the precision test of the Standard Model (SM) during the last decade. The electroweak data from LEP is so copious and precise that one can make use of the radiative correction effect to test the SM. In particular, the indirect determination of $`m_H`$ depends critically on the precision of $`\alpha (M_Z^2)`$. Recently, for example, there has been an increasing interest in electroweak phenomenology to reduce the uncertainty in $`\alpha (M_Z^2)`$ which seriously limits further progress in the determination of the Higgs mass from radiative corrections to the SM . The uncertainty in $`\alpha (M_Z^2)`$ arises from the contribution of light quarks to the photon vacuum polarization $`\mathrm{\Delta }\alpha (s)=_\gamma ^{^{}}(s)`$ at the Z mass scale. They are independent of any particular initial or final states and can be absorbed in $`\alpha (s)\alpha /[1\mathrm{\Delta }\alpha (s)]`$, where the fine-structure constant $`\alpha `$=1/137.035 989 5(61). $`\mathrm{\Delta }\alpha =\mathrm{\Delta }\alpha _{lepton}+\mathrm{\Delta }\alpha _{had}`$, of which the leptonic part is precisely calculated analytically according to perturbation theory because free lepton loops are affected by small electromagnetic corrections . Whereas, the hadronic part $`\mathrm{\Delta }\alpha _{had}`$ cannot be entirely calculated from QCD because of ambiguities in defining the light quark masses $`m_u`$ and $`m_d`$ as well as the inherent non-perturbative nature of the problem at small energy scale since the free quark loops are strongly modified by strong interactions at low energy. An ingenious way to handle this is to relate $`\mathrm{\Delta }\alpha _{had}`$ from the quark loop diagram to $`R`$, making use of unitarity and analyticity,
$$\mathrm{\Delta }\alpha _{had}^{(5)}(s)=\frac{s}{4\pi ^2\alpha }P(_{4m_\pi ^2}^{E_{cm}^2}𝑑s^{}\frac{R^{data}(s^{})}{s^{}(s^{}s)}+_{E_{cm}^2}^{\mathrm{}}𝑑s^{}\frac{R^{PQCD}(s^{})}{s^{}(s^{}s)}),$$
(3)
where $`R(s)=\sigma (e^+e^{}hadrons)/\sigma (e^+e^{}\mu ^+\mu ^{})`$ and P is the principal value of the integral.
Much independent work has recently been done to evaluate $`\alpha (s)`$ at the energy of the Z pole. So far, the uncertainty of $`\mathrm{\Delta }\alpha (s)`$ is dominanted by the $`R`$-values at low energy ($`E_{cm}<5GeV`$) measured with an average uncertainty of $``$15%, as indicated in Fig. 1.
The anomalous magnetic moment of the muon $`a_\mu (g2)/2`$ receives radiative contributions that can in principle be sensitive to new degrees of freedom and interactions. Theoretically, $`a_\mu `$ is sensitive to large energy scales and very high order radiative corrections. It therefore provides an extremely clean test of electroweak theory and may give us hints on possible deviations from the SM. The experimental and the theoretical predictions on $`a_\mu `$ are well reviewed by Dr. Robert Lee in his talk given at LP99.
One can decompose $`a_\mu `$ as
$$a_\mu ^{SM}=a_\mu ^{QED}+a_\mu ^{had}+a_\mu ^{weak},$$
(4)
where the largest term, the QED contribution $`a_\mu ^{QED}`$, has been calculated to $`O(\alpha ^5)`$, including the contribution from $`\tau `$ vacuum polarization. $`a_\mu ^{weak}`$ includes the SM effects due to virtual W, Z and Higgs particle exchanges. $`a_\mu ^{had}`$ denotes the virtual hadronic (quark) contribution determined by QCD, part of which corresponds to the effects representing the contribution of the running of $`\alpha (s)`$ from low energy to high energy scale. It cannot be calculated from first principles but can be related to the experimentally determined $`R(s)`$ through the expression
$$a_\mu ^{had}=(\frac{\alpha m_\mu }{3\pi })^2_{4m_\pi ^2}^{\mathrm{}}𝑑s\frac{R(s)K(s)}{s^2},$$
(5)
where $`K(s)`$ is a kernel varying from 0.63 at $`s=4m_\pi ^2`$ to 1.0 at $`s=\mathrm{}`$.
The hadronic vacuum polarization is the most uncertain of all the SM contributions to $`a_\mu `$, the uncertainty is presently $`156\times 10^{11}`$. For several scenarios, it has been claimed that “the physics achievement of the effort to measure the cross section of $`e^+e^{}\text{hadrons}`$ that brings down the uncertainty of $`a_\mu `$ to $`60\times 10^{11}`$ is equivalent to that of LEP2 or even the LHC” .
From equations 3 and 5 one finds that $`a_\mu ^{had}`$ is more sensitive to lower energies than higher ones. Further measurement in the energy region of 0.5-1.5 GeV from VEPP-2M in Novosibirsk and DA$`\mathrm{\Phi }`$NE in Frascati will contribute to the interpretation of the $`a_\mu `$ measurement at Brookhaven and the luminosity measurement at CERN . However, their contribution to the precision determination of $`\alpha (M_Z^2)`$ is limited. The improved $`R`$ value from BESII at BEPC in the energy region of 2-5 GeV will make the major contribution to evaluate $`\alpha (M_Z^2)`$, and also partly contribute to the interpretation of $`a_\mu `$.
## 2 Recent measurements of $`R`$ in low energy $`e^+e^{}`$
There are two different approaches to the measurement of $`R`$. One is to study the exclusive hadronic final states, i.e. to measure the production cross section of each individual channel $`\sigma ^{exp}(e^+e^{}\text{hadrons})_j`$. The value of $`R`$ can then be obtained by summing over the measured hadron production cross section of all individual channels. This method demands that the detector has good particle identification and requires the understanding of each channel. It is usually used for cm energies below 2 GeV.
Another method treats the hadronic final states inclusively. It measures $`R`$ by dealing with all the hadronic events simultaneously and is suitable in an energy region where a reliable event generator for hadron production is available. With an improved Lund Model , we may be able to extend this region down to 2 GeV.
The typical features of hadron production below 5 GeV are:
$``$ many resonances in this energy region, such as, $`\rho `$, $`\omega `$, $`\phi `$, $`\rho ^{}`$, $`\omega ^{}`$, $`\phi ^{}`$, $`c\overline{c}`$ and charmed mesons $`J/\psi `$, $`\psi (2S)`$, $`D^+D^{}`$, $`D_s^+D_s^{}`$ and $`\tau ^+\tau ^{}`$, baryon-antibaryon pair production
$``$ a small number of final states and low charged multiplicity, usually $`N_{ch}6`$.
The experimental challenge here is how to subtract the beam-associated background and select $`N_{had}`$.
In the following section, I will first discuss some new measurements done by CMD-2 and SND at VEPP-2M in Novosibirsk, which are based on an exclusive analysis of the hadron production in the energy region around 0.4-1.4 GeV. Then I will concentrate on discussing the $`R`$ scan done by BESII at BEPC in Beijing in the energy region from 2-5 GeV, which measures $`R`$ values by dealing with hadronic final states inclusively.
### 2.1 Recent results from VEPP-2M
VEPP-2M, the $`e^+e^{}`$ collider with maximum luminosity of $`5\times 10^{30}`$ cm<sup>-2</sup>s<sup>-1</sup> at $`E_{beam}`$=510 MeV, has been operating since 1974 in the energy region $`E_{cm}=0.41.4`$ GeV ($`\rho `$, $`\omega `$, $`\varphi `$-meson region). SND and CMD-2 are the two detectors carrying out experiments at VEPP-2M. Since 1994, VEPP-2M performed a series of scans from 0.38 GeV to 1.38 GeV . With this data both SND and CMD-2 have measured the cross sections for the channels $`\pi ^+\pi ^{}`$, $`\pi ^+\pi ^{}\pi ^0`$, $`\pi ^+\pi ^{}\pi ^+\pi ^{}`$, $`\pi ^+\pi ^{}\pi ^0\pi ^0`$, $`\pi ^+\pi ^{}\pi ^+\pi ^{}\pi ^0`$ as well as $`K_LK_S`$ and $`K^+K^{}`$.
$``$ Study of $`e^+e^{}\pi ^+\pi ^{}\pi ^+\pi ^{},\pi ^+\pi ^{}\pi ^0\pi ^0`$
The four pion final states produced via $`e^+e^{}`$ annihilation in the $`E_{cm}=12`$ GeV energy region dominate and determine the main part of the hadronic contribution to $`a_\mu `$ and the QCD sum rules. Besides, these processes are important sources of information for the understanding of hadron spectroscopy, in particular for the study of the $`\rho `$-meson radial excitation . These processes were studied at the VEPP-2M, DCI and ADONE colliders . The statistical errors of these measurements were $`5\%`$, and the systematic errors were $`15\%`$. There was about a 20% discrepancy among the different experiments.
For the process of $`e^+e^{}\pi ^+\pi ^{}\pi ^+\pi ^{}`$, the systematic error from the recent SND results is $`7\%`$, mainly coming from the event selection and the luminosity determination. The measured total cross sections from CMD-2 for $`e^+e^{}2\pi ^+2\pi ^{}`$ are also illustrated in Fig. 3. Only the statistical errors are shown. The systematic uncertainties are $`7\%`$, attributed to the luminosity measurement, the event reconstruction and selection and the radiative correction as well.
For SND, backgrounds to the $`e^+e^{}\pi ^+\pi ^{}\pi ^0\pi ^0`$ channel are mainly from $`e^+e^{}K^+K^{}`$, QED processes $`e^+e^{}e^+e^{}e^+e^{},e^+e^{}\gamma \gamma `$, as well as cosmic rays and beam associated background. The systematic error is $`7\%`$, of which $`5\%`$ arise out of the variation of the detection efficiency shown from the simulation of the intermediate states like $`\omega \pi ^0`$, $`\rho ^0\pi ^0\pi ^0`$ and Lorentz-invariant phase space simulation(LIPS).
The total cross section for $`e^+e^{}\pi ^+\pi ^{}2\pi ^0`$ process measured by CMD-2 also is plotted in Fig. 3. The error bars indicates only the statistical error. The systematic uncertainties mainly come from event reconstruction, radiative corrections and the luminosity determination. The overall systematic uncertainty is estimated to be 7%. The cross section measured by this experiment is consistent with what was measured by OLYA and a recent result from SND . However, the cross section from all three measurements is apparently lower than that given by ND . For comparison, the results from Orsay and Frascati above 1.4 GeV are also shown in the figure.
CMD-2 finds that the dominant contribution to the cross section of the process $`e^+e^{}\pi ^+\pi ^{}2\pi ^0`$ comes from $`\omega \pi ^0`$ and $`\rho ^\pm \pi ^{}\pi ^0`$ intermediate states, whereas the $`\rho ^\pm \pi ^0\pi ^0`$ state is not observed. The $`\rho ^\pm \pi ^{}\pi ^0`$ states are saturated completely by the $`a_1(1260)\pi `$ intermediate state. This is also the dominant contribution to the cross section for the process $`e^+e^{}2\pi ^+2\pi ^{}`$. The theoretical predictions for the differential distributions and the total cross sections can be dramatically changed if one takes into account the interference of different amplitudes with various intermediate states but identical final states.
The cross section of $`e^+e^{}4\pi `$ can be related to the the four $`\pi `$ decays of the $`\tau `$-lepton through the hypothesis of CVC . This has been experimentally tested to be valid within an accuracy of 3-5% . The observed $`a_1(1260)\pi `$ dominance, if it is true, should be taken into account in $`\tau `$ decays.
$``$ The investigation of $`e^+e^{}\pi ^+\pi ^{}\pi ^0`$
This process was measured by ND at VEPP-2M in the energy region up to 1.1 GeV . The measured cross section is significantly higher than that predicted by the Vector Dominance Model (VDM). However, it is well known that the VDM is able to well describe the cross section near the $`\omega `$ and $`\varphi `$ resonances for the processes $`e^+e^{}\omega ,\varphi \rho \pi \pi ^+\pi ^{}\pi ^0`$. It is therefore necessary to perform new precise measurements in the non-resonance region to investigate the limitation of the VDM and determine possible contributions from heavier intermediate states like $`\omega (1120)`$ or $`\omega (1600)`$.
SND also measured this channel. The systematic errors from the detection efficiency, luminosity measurement and the background subtraction are 10%, 5% and 5% respectively, giving a total systematic error of $`12\%`$. The results from the new measurement agrees with the old ND data.
Invariant masses of $`\pi `$-meson pairs in the final $`3\pi `$ state were measured to investigate the intermediate state in the $`e^+e^{}\pi ^+\pi ^{}\pi ^0`$ process. The intermediate states might be $`\rho \pi `$ and the much less probable $`\omega \pi `$ with decay $`\omega 2\pi `$. Comparing the mass spectrum of $`\pi ^+\pi ^{}`$ with that of $`\pi ^0\pi ^\pm `$, one can observe the interference between the two intermediate states. The clear peak shown in $`\pi ^+\pi ^{}`$ mass spectrum of the experimental data proves the $`\rho \omega `$ interference in $`3\pi `$ final state, and the phase measured to be zero agrees with the VDM prediction. Fig. 5 plots the cross section for the production of $`\pi ^+\pi ^{}\pi ^0`$.
$``$ Cross section measurement for $`e^+e^{}K_SK_L`$
The cross section of $`e^+e^{}K_SK_L`$ reaction was measured in 1982 by OLYA in Novosibirsk and DM1 in Orsay in the energy region $`E_{cm}=1.061.40`$ GeV and $`E_{cm}=1.402.20`$ GeV respectively. It is desirable to re-measure this channel since the accuracy reached by both experiments is poor. So far, 1.8 pb<sup>-1</sup> data has been analyzed by SND, utilizing $`K_S\pi ^0\pi ^0`$ from $`e^+e^{}K_SK_L`$. $`e^+e^{}\omega \pi ^0\pi ^0\pi ^0\gamma `$ is the main background source. In addition, cosmic rays and beam associated background also contribute. The cross section measured by SND is illustrated in Fig. 5 with solid dots.
$``$ $`e^+e^{}\omega \pi ^+\pi ^{}`$ ($`\omega \pi ^+\pi ^{}\pi ^0`$), $`\eta \pi ^+\pi ^{}`$ ($`\eta \pi ^+\pi ^{}\pi ^0`$, or $`\gamma \gamma `$)
Figs. 7 and 7 plot the cross section measured by CMD-2 for $`\omega \pi ^+\pi ^{}`$ and $`\eta \pi ^+\pi ^{}`$, together with the measurement by DM2. The new measurement significantly reduced the uncertainties though the systematic errors are still as high as 15%.
$``$ $`e^+e^{}\pi ^+\pi ^{}`$
The cross section of the process $`e^+e^{}\pi ^+\pi ^{}`$ is given by
$$\sigma =\frac{\pi \alpha ^2}{3s}\beta _\pi ^3|F_\pi (s)|^2,$$
(6)
where $`F_\pi (s)`$ and $`\beta _\pi `$ are the pion form factor at the cm energy $`\sqrt{s}`$ and the velocity of the pion respectively. A precision measurement of the pion form factor is necessary to determine the $`R`$-values via an exclusive method. The relative uncertainty contribution to $`a_\mu `$ is dominated by the $`e^+e^{}\pi ^+\pi ^{}`$ channel with $`\sqrt{s}<2`$ GeV . The famous E821 experiment at BNL has measured $`a_\mu `$ to a precision of $`5`$ ppm and will further improve the accuracy to about 1 ppm. In order to compare a measurement with such a high accuracy with theory, the uncertainty in $`R`$ should be below 0.5% in this energy region. Therefore a new measurement of the pion form factor with smaller uncertainty is important for the interpretation of the E821 measurement.
The pion form factor was measured by the OLYA and CMD groups at VEPP-2M about twenty years ago . Twenty-four points from 360 to 820 MeV were studied by CMD with a systematic uncertainty of about 2%. The OLYA measurement scanned from 640 to 1400 MeV with small steps, giving a systematic uncertainty from 4% at the $`\rho `$-meson peak to 15% at 1400 MeV.
The pion form factor is one of the major experiments planned at CMD-2. A total of 128 energy points were scanned in the whole VEPP-2M energy region (0.36-1.38 GeV) in six runs performed from 1994 to 1998 . The discussion here is based on data taken from the first 3 runs with 43 energy points ranging from 0.61-0.96 GeV. The small energy scan step, 0.01 GeV, in this energy region allows the calculation of the hadronic contribution in a model-independent way. In order to investigate the $`\omega `$-meson parameters and the $`\rho \omega `$ interference, the energy steps were 2-6 MeV in the energy region near the $`\omega `$-meson. The beam energy was measured with the resonance depolarization technique for almost all the energy points, which significantly reduced the systematic error arising from the energy uncertainty. The charged trigger makes use of the information from the drift chamber and the Z-chamber and requires at least one track. There was an additional trigger criteria for the energy points between 0.81 and 0.96 GeV, which asks for the total energy deposited in the calorimeter to be greater than 20-30 MeV. The neutral trigger, served for monitoring the trigger efficiency, is based on the information only from the calorimeter.
The background is mainly from cosmic muons. Bhabha and dimuon production are also background sources. The shape of the energy deposition was carefully studied for the event separation and selection. An event vertex cut was applied to reject cosmic muons effectively.
To account for the fact that the radiative correction for $`e^+e^{}\pi ^+\pi ^{}`$ depends on the energy behavior of the cross section of $`e^+e^{}\pi ^+\pi ^{}`$ itself, the radiative correction factor was calculated iteratively. The existing $`|F_\pi (s)|^2`$ data were used as the first iteration for the calculation. The values of $`|F_\pi (s)|^2`$ were found to be stable after three iterations.
The corrections for the pion losses due to decays in flight and nuclear interaction, as well as the background from $`\omega 3\pi `$ were done using Monte Carlo simulations.
The total systematic uncertainty was estimated to be currently 1.5% and 1.7% for the energy region 0.78-0.784 GeV and 0.782-0.94 GeV respectively, and 1.4% for all the other points.
Other than the leading contribution from $`\rho (770)`$ and $`\omega (782)`$, the resonances $`\rho (1450)`$ and $`\rho (1700)`$ should be taken into account to describe the data for the determination of the pion form factor. In addition, the model based on the Hidden Local Symmetry (HLS), which predicts a point-like coupling $`\gamma \pi ^+\pi ^{}`$, can well describe the experimental data below 1 GeV. Both the Gounaris-Sakurai (GS) parameterization and the HLS parameterization approaches were used to fit the form factor. Only the higher resonance $`\rho (1450)`$ was taken into account in fitting the pion form factor in the relatively narrow energy region 0.61-0.96 GeV. Fig. 8 shows the fit of the pion form factor with CMD-2 94, 95 data according to GS and HLS models. Both theoretical curves are indistinguishable.
The pion form factor fit of CMD-2 94, 95 data is summarized by Table 2. PDG data is also shown for comparison.
With the remaining data collected, CMD-2 hopes to reduce the systematic error presented here by a factor of two. To achieve this goal, a new approach for the calculation of the radiative correction must be developed.
### 2.2 $`R`$ scan with BESII at BEPC in Beijing
With the upgraded machine and detector , the BES collaboration performed two scans to measure $`R`$ in the energy region of 2-5 GeV in 1998 and 1999. The first run scanned 6 energy points covering the energy from 2.6 to 5 GeV in the continuum. Separated beam running at each energy point was carried out in order to subtract the beam associated background from the data .
The second run scanned about 85 energy points in the energy region of 2-4.8 GeV. To subtract beam associated background, separated beam running was done at 26 energy points and single beam running for both $`e^{}`$ and $`e^+`$ was done at 7 energy points distributed over the whole scanned energy region. Special runs were taken at the $`J/\psi `$ to determine the trigger efficiency. The $`J/\psi `$ and $`\psi (2S)`$ resonances were scanned at the beginning and at the end of the $`R`$ scan for the energy calibration.
The $`R`$ values from the BESII scan data are measured by observing the final hadronic events inclusively, i.e. the value of $`R`$ is determined from the number of observed hadronic events ($`N_{had}^{obs}`$) by the relation
$$R=\frac{N_{had}^{obs}N_{bg}_lN_{ll}N_{\gamma \gamma }}{\sigma _{\mu \mu }^0Lϵ_{had}ϵ_{trg}(1+\delta )},$$
(7)
where $`N_{bg}`$ is the number of beam associated background events; $`_lN_{ll},(l=e,\mu ,\tau )`$ and $`N_{\gamma \gamma }`$ are the numbers of misidentified lepton-pairs from one-photon and two-photon processes events; $`L`$ is the integrated luminosity; $`\delta `$ is the radiative correction; $`ϵ_{had}`$ is the detection efficiency for hadronic events and $`ϵ_{trg}`$ represents the trigger efficiency.
The trigger efficiencies are measured by comparing the responses to different trigger requirements in special runs taken at the $`J/\psi `$ resonance. From the trigger measurements, the efficiencies for Bhabha, dimuon and hadronic events are determined to be 99.96%, 99.33% and 99.76%, respectively. As a cross check, the trigger information from the 2.6 and 3.55 GeV data samples is used to provide an independent measurement of the trigger efficiencies. This measurement is consistent with the efficiencies determined from the $`J/\psi `$ data. The errors in the trigger efficiencies for Bhabha and hadronic events are less than 0.5%.
The task of the hadronic event selection is to identify one photon multi-hadron production from all other possible contamination mechanisms. The event selection makes full use of all the information from each sub-detector of BESII, namely, the vertex position, the measured charged-particle momentum and the energy loss due to ionization, the related time of flight, the associated pulse height and pulse shape of the electromagnetic calorimeter, and the hits in $`\mu `$ counter.
The backgrounds involved in the measurement are mainly from cosmic rays, lepton pair production ($`e^+e^{}`$, $`e^+e^{}`$, $`\mu ^+\mu ^{}`$, $`\tau ^+\tau ^{}`$), two-photon processes, and beam associated processes. The cosmic rays and part of the lepton pair production events are directly removed by the event selection. The remaining background from lepton pair production and two-photon processes is then subtracted out statistically according to a Monte Carlo simulation.
The beam associated background sources are complicated. They may mainly come from beam-gas, beam-wall interaction, synchrotron radiation, and lost beam particles. The salient features of the beam associated background are that their tracks are very much along the beam pipe direction, the energy deposited in BSC is small, and most of the tracks are protons.
Separated-beam runs were performed for the subtraction of beam associated background. Most of the beam associated background events are rejected by vertex and energy cuts. Applying the same hadronic events selection criteria to the separated-beam data, one can obtain the number of separated-beam events $`N_{sep}`$ surviving these criteria. The number of beam associated events $`N_{bg}`$ in the corresponding hadronic events sample is given by $`N_{bg}=fN_{sep}`$, where $`f`$ is the ratio of the product of the pressure at the collision region times the integrated beam currents for colliding beam runs and that for the separated beam runs. To subtract beam associated background in this way, the variation of the pressure in the collision region and the beam current must be recorded for both colliding and separated-beam runs at each energy to be measured.
JETSET7.4 is used as the hadronic event generator to determine the detection efficiency for hadronic events. Parameters in the generator are tuned using a $`4\times 10^4`$ hadronic event sample collected near 3.55 GeV for the tau mass measurement done by the BES collaboration . The parameters of the generator are adjusted to reproduce distributions of kinematic variables such as multiplicity, sphericity, transverse momentum, etc.
The parameters have also been obtained using the 2.6 GeV data ($`5\times 10^3`$ events). The difference between the two parameter sets and between the data and the Monte Carlo data based on these parameter sets is used to determine a systematic error of 1.9-3.2% in the hadronic efficiency.
The Monte Carlo simulation packet JETSET was not designed to fully describe few body states produced by $`e^+e^{}`$ annihilation in the few GeV energy region, though the event shapes are consistent with that from the Monte Carlo simulation with tuned parameters at 3.5 GeV. A great effort has been made by the Lund group and BES collaboration to develop the formalism using the basic Lund Model area law directly for the Monte Carlo simulation, which is expected to describe the data better .
Radiative corrections determined using four different schemes agreed with each other to within 1% below charm threshold. Above charm threshold, where resonances are important, the agreement is within $`13\%`$. The major uncertainties common to all models are due to errors in previously measured $`R`$-values and in the choice of values for the resonance parameters. For the measurements reported here, we use the formalism of ref. and include the differences with the other schemes in the systematic error of 2.2-4.1%.
The $`R`$ values obtained at the six energy points scanned in 1998 are shown in Table 3 and graphically displayed in Fig. 9 with solid dots. The largest systematic error is due to the hadronic event selection and is determined to be 3.8-6.0% by varying the selection criteria. The systematic errors on the measurements below 4.0 GeV are similar and are a measure of the amount of error common to all points. The BES collaboration has also performed the analysis including only events with greater than two charged tracks; although the statistics are smaller, the results obtained agree well with the results shown here.
The $`R`$ values for $`E_{cm}`$ below 4 GeV are in good agreement with results from $`\gamma \gamma 2`$ and Pluto but are below those from Mark I . Above 4 GeV, our values are consistent with previous measurements.
Preliminary $`R`$-values at 2.4, 2.5, 2.6, 2.7, 2.8, 2.9, 3.0, 4.6 and 4.8 GeV are plotted with solid squares in Fig. 10. The preliminary errors, which add the statistical and systematical errors in quadrature, are all conservatively assigned to be 10%. However, it is believed that these errors can be decreased to be comparable to the error bars of the solid dots, i.e. $`7\%`$ for the energy points below 3.6 GeV and $`10\%`$ for energies above 4.5 GeV. The first scan repeated 3.4 GeV and the second scan repeated 2.6 and 4.6 GeV data points measured in the first scan. In all case, $`R`$-values obtained are consistent with each other at the same energy.
## 3 Prospects and Concluding Remarks
SND and CMD-2 at VEPP-2M have significantly improved the measurements of the hadron production cross section via $`e^+e^{}`$ collisions for some of the important exclusive channels in the energy region of 0.36-1.38 GeV. Further improvement with the analysis of the existing data is forthcoming. A major advance would be possible if the energy region could be extended to 2 GeV, which would link up to the lowest energy of BEPC.
CMD-2 and SND at VEPP-2M are planning to scan from threshold to 1.4 GeV in 1999-2000. A $`R`$ scan between 2 to 10 GeV with KEDR at VEPP-4 is proposed. A scan covering such a wide energy region with the same machine and detector would be very important if the measurement could be performed with a $`1\%`$ precision.
The $`R`$ scan performed with BESII at BEPC in Beijing can significantly reduce the uncertainties in $`R`$ in the energy region 2-5 GeV. The $`R`$-values from the first run data have already reduced the uncertainties in $`R`$ from 15-20% to 7%. BESII at BEPC in Beijing is analyzing the second run $`R`$ scan data. The preliminary $`R`$ values in the whole energy region of 2-5 GeV are expected to be shown in the Spring of 2000, and the final results will be presented in the summer of 2000.
The new $`R`$ ratio results in $`e^+e^{}`$ annihilations presented from Novosibirsk and Beijing have had a great impact on the value of $`\alpha (M_Z^2)`$. Using these new (albeit, preliminary) results, A.D. Martin et. al. re-evaluate $`\alpha (M_Z^2)`$ and find $`\alpha (M_Z^2)^1=128.973\pm 0.035`$ or $`128.934\pm 0.040`$, according to whether inclusive or exclusive cross sections are used. The uncertainties here are already decreased more than half if we compare with the previous value of $`\alpha (M_Z^2)^1=128.89\pm 0.09`$ evalued by using the old experimental $`R`$-values . With the final results from the Beijing inclusive measurement in the whole 2-5 GeV region and the more precise results from Novosibirsk, $`\alpha (M_Z^2)`$ and $`a_\mu `$ will be even more precisely determined from the experimental data.
A dedicated energy scan, aimed at a 1% precision direct measurement below 1.4 GeV, planned by KLOE at DAPNE is not possible in the short term since the DAFNE machine is tuned for the $`\varphi `$ resonance. However a machine upgrade is foreseen which will hopefully permit such an energy scan around 2004. Another method to measure the hadronic cross section is to measure events with Initial State Radiation (ISR). In this case one of the electrons or positrons of the beam radiates a hard photon and the cm energy of the hadronic system in the final state, mostly pions coming from the $`\rho `$-resonance, is lowered. KLOE has already started the analysis of those events .
Being one of the most fundamental parameters in particle physics, the $`R`$-value plays an important role in the development of the theory of particle physics and in testing the Standard Model. Experimental efforts to precisely measure $`R`$-values at low energies are crucial for the future electroweak precision physics. The measurements are not only important for the evaluation of $`\alpha (M_Z^2)`$ and for the interpretation of $`a_\mu `$, but also necessary for the understanding of the hadron production mechanism via $`e^+e^{}`$ annihilation.
A real breakthrough in electroweak theory physics with regard to the $`R`$-values at low energy would be possible only by measuring $`\sigma (e^+e^{}\text{hadrons})`$ at $`1\%`$ accuracy. Such a level of precision requires significant improvement to both machine and detector, and needs better theoretical calculation of the radiative correction and the event generator for hadron production.
Once the $`R`$-values has been measured with a precision of 1% in the energies covered by VEPP-2M and DAPNE, the central question will be how to further decrease the uncertainties of $`R`$-values measured by BESII in the energy region of 2-5 GeV, particularly from 2-3.7 GeV. Such a measurement would be extremely important for the interpretation of the $`a_\mu `$ experiment carrying out by E821 at BNL and for the precision determination of $`\alpha (M_Z^2)`$. This measurement would then be an important and attractive physics program for a $`\tau `$-c factory.
I would like to thank the many people who have at some point or another helped realize this presentation. In particular, my very special thanks are due Dr. Boris I. Khazin, Dr. Sergey Serednjakov for supplying me useful informations concerning the results from CMD-2 and SND. I’m also gratful to Dr. W. Dunwoodie, Dr. A. F. Haris and Dr. D. Kong who gave me helpful comments and suggestions. |
warning/0002/cond-mat0002345.html | ar5iv | text | # Control of Squeezed States
## 1 Introduction
In this paper we consider the problem of squeezing of harmonic oscillators from the point of view of control theory. Squeezing has been suggested as a method for reducing noise in quantum systems below the standard quantum limit. This can be achieved by using laser pulses and in that sense may be viewed as a quantum control problem, although the classical squeezing problem is also of interest. In the latter case one is interested in reducing noise induced by random perturbations.
The quantum control problem has been of great interest recently, see for example Brockett and Khaneja , Lloyd and Warren et. al. and references therein.
Here we consider squeezing as a control problem in both the classical and quantum setting. In the classical case we consider a system subject to thermal noise while in the quantum case we consider a system at zero temperature and in the presence of noise. In both cases the control is given by an external electromagnetic field and enters the control equations multiplicatively. In this sense the setting is similar to the NMR control problems analyzed by Brockett and Khaneja.
A key feature of squeezing is that it results in a redistribution of uncertainty between observables.
In this paper we consider a model for phonon squeezing in solids following the work of Garret at. al. , but one can equally well consider the case of photons in quantum optics. The control is via a single pulse on a large ensemble of oscillators and this sense we are considering under-actuated control systems in both the classical and quantum case.
We also model the effect of dissipation on the classical system and the effect of coupling to a heat bath in the quantum setting. This causes the squeezing effect to gradually moderate.
## 2 The Control Setting
In the classical setting we consider bilinear control systems of the form
$$\dot{x^i}=\underset{j=1}{\overset{n}{}}a_{ij}x^j+\underset{j=1}{\overset{m}{}}b_{ij}(x)u^j+\underset{j=1}{\overset{r}{}}g_{ij}\dot{w^k}$$
(2.1)
where $`a_{ij}`$ is a constant matrix (i.e. the free dynamics is linear), $`b_{ij}`$ is linear in $`x`$, i.e. the control $`u^j`$ enters bilinearly, $`\dot{w^k}`$ is white noise and the state space is $`^n`$.
In the quantum context we want to consider a similar equation but defined on an appropriate Hilbert space:
$$ih\frac{\psi }{t}=H_s\psi +\underset{j}{}u_jH^j\psi +\underset{j}{}\dot{w^j}H_j\psi $$
(2.2)
where $`H_s`$ is the Schroedinger operator, the $`H_j`$ are (linear) input operators, $`u^j`$ are functions of time, and the $`\psi `$ is a vector in the Hilbert space.
## 3 Classical Squeezing of the Harmonic Oscillator
In this section we consider classical squeezing of a set of identical coupled harmonic oscillators. Denote the position of each oscillator by $`u^i`$.
The Hamiltonian for the system is of the form:
$$H=\underset{i}{}\frac{p_i^2}{2}+\underset{i}{}\frac{\omega _i^2}{2}(q^iq^{i+1})^2$$
(3.1)
where the oscillators are assumed to have unit mass and $`p_i=\dot{u^i}`$.
In order to analyze the system we decompose it into its normal modes. Denoting the normal mode coordinates by $`Q^i`$ we thus obtain a system of uncoupled harmonic oscillator equations of the form $`\ddot{Q^i}+\mathrm{\Omega }_i^2Q^i=0`$.
The main control mechanism we consider here is squeezing by pulses. In this case each oscillator is forced by a pulse at time $`t=0`$ which is proportional to its displacement, i.e. we have equations of the form:
$$\ddot{Q^i}+\mathrm{\Omega }_i^2Q^i=2\lambda Q^i\delta (t)$$
(3.2)
where $`\delta (t)`$ is the Dirac delta function and $`\lambda `$ is a constant which is proportional to the frequency $`\mathrm{\Omega }`$.
Thus we obtain
$$\dot{Q}^i(0^+)=\dot{Q}^i(0^{})+2\lambda Q^i(0).$$
(3.3)
Thus, if one considers the system subject to white noise,
$$\ddot{Q^i}+\mathrm{\Omega }_i^2Q^i=2\lambda Q^i\delta (t)+\alpha \dot{w}^i,$$
(3.4)
one sees that while one starts with a spherical equilibrium distribution which is invariant in time, after the pulse one has an elliptical distribution which rotates in time at twice the harmonic frequency (by the $`_2`$ symmetry of the ellipse). (A precise analysis is given below in the course of our treatment of the quantum mechanical case.) Noise reduction is then achieved by viewing the system “stroboscopically” when the noise is low.
Actually the above is an idealization: in actuality the oscillator should be viewed as in equilibrium with a heat bath which dissipates energy. In the classical setting one can model this by simple linear dissipation (in the quantum setting one has to introduce a heat bath – see below).
Thus we have a system of the form
$$\ddot{Q^i}+\mathrm{\Omega }_i^2Q^i=\eta _i\dot{Q}^i+U_i(t)+\alpha \dot{w}^i$$
(3.5)
where $`\eta _i`$ is a dissipation constant and $`U_i(t)`$ is the control which we can choose to be a single pulse or a sequence of pulses. Depending on the dissipation strength an initial squeezing effect will decay away and we need a continual sequences of pulse to keep the system in a squeezed state.
It is worthwhile remarking on the how the control enters in our setting: the control is a single pulse applied overall (and in this sense the system is under-actuated) while the effect on each (normal mode) of oscillation is to apply a pulse proportional to displacement (minus the mean displacement which is of course zero for each oscillator). This is effected by the type of interaction of the oscillators with the field that the pulse induces. We note also that in the full nonlinear setting the mean displacement may not be zero and must be taken into account.
We shall return to the classical squeezing of oscillator by pulses, and in particular a computation of mean square displacement, after a discussion of the quantum case below.
We note also that parametric resonance control can achieve similar squeezing effects in the classical case. In this case we consider oscillator motion in the presence of a modulating drive:
$$\ddot{Q}\omega ^2(1+ϵ\mathrm{cos}2\omega t)Q=0$$
(3.6)
where $`ϵ`$ parameterizes the strength of the drive. We omit details of this approach here.
## 4 Squeezing of the Quantum Harmonic Oscillator
We now turn to the quantum setting.
Consider the following Hamiltonian
$$H=\frac{P^2}{2m}+\frac{m\omega ^2}{2}Q^2+\lambda \delta (t)Q^2,$$
(4.1)
which reflects an impulsive change in the spring constant and where $`\omega =\sqrt{K/m}`$, $`K`$ being the original spring constant.
The variables $`P`$ and $`Q`$, which are operators in the quantum case, obey canonical commutation rules $`[P,Q]=i\mathrm{}`$. We can rewrite the above Hamiltonian in terms of creation operators $`a`$ and $`a^{}`$ defined through
$$Q=\sqrt{\frac{\mathrm{}}{2m\omega }}(a+a^{}),P=i\sqrt{\frac{\mathrm{}m\omega }{2}}(a^{}a),$$
(4.2)
with $`[a,a^{}]=1`$. Written in terms of the new variables, the Hamiltonian is
$$H=\mathrm{}\omega (a^{}a+1/2)+\lambda \delta (t)(a+a^{})^2.$$
(4.3)
The ground state of the system, for $`t0`$, $`|0`$, corresponds to the vacuum of $`a`$, ($`a|0=0`$), and the excited states are of the form $`(a^{})^2|0`$.
We now want to study the behavior of the system at $`t>0`$, given that the system is in its ground state at $`t<0`$. The wave function at $`t=0^+`$ is of the form $`|\psi (t=0^+)=\mathrm{exp}(i\lambda Q^2)|0,`$ and for longer times the system evolves with the “unperturbed” Hamiltonian: $`|\psi (t>0)=\mathrm{exp}(iH_0t)e^{i\lambda Q^2}|0.`$ Our first quantity of interest is $`\psi (t)|Q^2|\psi (t)Q^2(t)`$. Let us compute it using the general method of coherent states. We find
$$Q^2(t)=0|e^{i\lambda Q^2}(ae^{i\omega t}+a^{}e^{i\omega t})^2e^{i\lambda Q^2}|0,$$
(4.4)
where we have used the fact that $`e^{iH_0t}ae^{iH_0t}=ae^{i\omega t}`$, which states that $`a^{}`$ and $`a`$ respectively destroy and create eigenstates of $`H_0`$, and where $`Q`$ is defined in units of $`\sqrt{\mathrm{}/(2m\omega )}`$.
Now we introduce a basis of coherent states $`|z`$, which satisfy $`a|z=z|z,z|a^{}=z|z^{},`$ and form an overcomplete set of states:
$$1=\frac{1}{2\pi i}𝑑z𝑑z^{}e^{zz^{}}|zz|.$$
(4.5)
Inserting (4.5) in (4.4) we find
$`Q^2(t)={\displaystyle \frac{1}{2\pi i}}{\displaystyle 𝑑z𝑑z^{}e^{zz^{}}}`$
$`(z^2e^{2i\omega t}+z^2e^{2i\omega t}+2zz^{}1)|0|e^{i\lambda x^2}|z|^2.`$
In order to evaluate the last term we need the position representation of the ground state (note that at this point $`Q`$ is a real number)
$$0|Q=\frac{1}{\pi ^{\frac{1}{4}}}e^{Q^2/2}$$
(4.6)
and that of the coherent state
$$Q|z=\frac{1}{\pi ^{\frac{1}{4}}}e^{Q^2/2+\sqrt{2}zQz^2/2}.$$
(4.7)
A simple integration gives
$`0|e^{i\lambda Q^2}|z`$ $`=`$ $`{\displaystyle 𝑑x0|QQ|ze^{i\lambda Q^2}}`$ (4.8)
$`=`$ $`{\displaystyle \frac{1}{\sqrt{1i\lambda }}}e^{i\lambda z^2/2(1i\lambda )}.`$ (4.9)
Changing to the variables $`z=u+iv`$ we have
$$e^{zz^{}}|0|e^{i\lambda Q^2}|z|^2=\frac{1}{\sqrt{1+\lambda ^2}}e^{[v^2+(2\lambda ^2+1)u^2+2\lambda uv]/(1+\lambda ^2)},$$
(4.10)
and
$`Q^2(t)={\displaystyle \frac{4}{\pi \sqrt{1+\lambda ^2}}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑u{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑v`$ (4.11)
$`\left(u^2\mathrm{cos}^2\omega t+v^2\mathrm{sin}^2\omega t+uv\mathrm{sin}2\omega t{\displaystyle \frac{1}{4}}\right)`$
$`\times e^{[v^2+(2\lambda ^2+1)u^2+2\lambda uv]/(1+\lambda ^2)}`$
$`=`$ $`1+4\lambda ^2\mathrm{sin}^2\omega t+2\lambda \mathrm{sin}2\omega t`$ (4.12)
It is interesting to compare this with an ensemble of classical oscillators with initial conditions taken from a heat bath. For simplicity let us take $`\omega =m=k_B=T=1`$ ($`k_B`$ is Boltzman’s constant). An arbitrary oscillator will evolve as $`Q(t)=u\mathrm{cos}t+v\mathrm{sin}t,`$ with $`u`$ and $`v`$ its initial position and velocity. If a pulse is applied at $`t=0`$ of the form treated above: $`Q(t)=u\mathrm{cos}t+(v+2\lambda u)\mathrm{sin}t.`$ Now let us average over initial conditions taken from a measure given by (a thermal bath)
$`Q^2(t)`$ $``$ $`{\displaystyle 𝑑u𝑑v[u\mathrm{cos}t+(v+2\lambda u)\mathrm{sin}t]^2e^{(u^2+v^2)}}`$ (4.13)
$`=`$ $`1+4\lambda ^2\mathrm{sin}^2t+2\lambda \mathrm{sin}2t.`$
It is interesting to note that the two expressions for, respectively, the quantum oscillator at zero temperature and the classical oscillator at finite temperature, are exactly the same. The general time dependence of the variance for a squeezed harmonic oscillator with frequency $`\omega `$ can thus be written in the following form
$$[Q(t)]^2=\frac{ϵ_0}{K}\left[1+\left(\frac{2\lambda }{\omega }\right)\mathrm{sin}2\omega t+\left(\frac{2\lambda }{\omega }\right)^2\mathrm{sin}^2\omega t\right]$$
(4.14)
with $`ϵ_0=\mathrm{}\omega /2`$ for the quantum case and $`ϵ_0=k_BT`$ for the classical oscillator at a temperature $`T`$.
The method of coherent states presented above has the advantage of being suitable for calculating other quantities. For example, if the oscillators are atoms within a solid, the scattering amplitude for an X-ray is decreased by a factor (called the Debye-Waller factor – see Ziman ) $`\mathrm{exp}ikQ(t)`$, with $`k`$ the wave-vector of the X-ray. We now ask ourselves what is the time evolution of the Debye-Waller factor for a squeezed phonon. This means that we need to compute the following expression
$`I(\lambda ,t)=0|e^{i\lambda Q^2}e^{(ae^{i\omega t}+a^{}e^{i\omega t})}e^{i\lambda Q^2}|0`$ (4.15)
$`={\displaystyle \frac{1}{\sqrt{e}}}{\displaystyle \frac{1}{\sqrt{1+\lambda ^2}}}{\displaystyle \frac{1}{\pi }}{\displaystyle 𝑑u𝑑v}`$
$`e^{2u\mathrm{cos}\omega t+2v\mathrm{sin}\omega t\frac{[v^2+(2\lambda ^2+1)u^2+2\lambda uv]}{(1+\lambda ^2)}}`$
$`=`$ $`e^{1+4\lambda ^2\mathrm{sin}^2\omega t+2\lambda \mathrm{sin}2\omega t}.`$
For the Debye-Waller factor, we obtain the following time dependence
$$e^{ikQ(t)}=e^{k^2Q^2(t)}$$
(4.16)
Measurement of the Debye-Waller factor may provide a practical method of detecting the squeezing phenomenon experimentally.
## 5 Squeezing and dissipation
In this section we consider the squeezing of a quantum oscillator coupled to a an infinite number of oscillators representing a “heat” bath. We show that this causes a decay in the squeezing oscillation for small time and true damping in the limit of a continuum of oscillators. This damping effect of the heat bath is similar to that analyzed classically in Lamb , Komech , Sofer and Weinstein and Hagerty, Bloch and Weinstein . We stress that we are considering a zero temperature case, and the damping effects appear due to a) the coupling of a single variable with a continuum of variables and b) an “asymmetry” in the initial conditions. The applied pulse on the oscillator generates outgoing waves on the continuum system which in turn gives rise to a positive damping (for a detailed discussion of negative versus positive damping see Keller and Bonilla ).
We start with a general formulation, and at the end of this section discuss a specific continuum example.
The Hamiltonian of the system consists of three parts: $`H_0`$ describing the original oscillator:
$$H_0=\frac{p_0^2}{2m}+\frac{m\omega _0^2}{2}q_0^2,$$
(5.1)
the Hamiltonian $`H_e`$ of the environment:
$$H_e=\underset{\alpha }{}\left[\frac{p_\alpha ^2}{2m}+\frac{m\omega _\alpha ^2}{2}q_\alpha ^2\right],$$
(5.2)
and a linear coupling between the two
$$H_{\mathrm{int}}=\underset{\alpha }{}\xi _\alpha q_\alpha q_0.$$
(5.3)
Formally, the total Hamiltonian $`H=H_0+H_e+H_{\mathrm{int}}`$ can be written in terms of its normal mode coordinates $`X_\nu `$ and $`P_\nu `$:
$$H=\underset{\nu }{}\left[\frac{P_\nu ^2}{2m}+\frac{m\omega _\nu ^2}{2}X_\nu ^2\right],$$
(5.4)
and we will consider a situation in which the initial (before the pulse) wave function corresponds to all the modes in the ground state:
$$\mathrm{\Psi }_0=\underset{\nu }{}\left(\frac{\omega _\nu }{\pi \mathrm{}}\right)^{1/4}e^{\omega _\nu X_\nu ^2/2\mathrm{}}.$$
(5.5)
At $`t=0`$ a pulse is applied to the (original) oscillator, the wave function immediately after the pulse given by:
$`\mathrm{\Psi }_0(t=0^+)`$ $`=`$ $`e^{i\lambda q_0^2}\mathrm{\Psi }_0`$ (5.6)
$`=`$ $`e^{i\lambda _{\mu \nu }U_{0\mu }U_{0\nu }X_\mu X_\nu }\mathrm{\Psi }_0,`$ (5.7)
where $`U_{\mu \nu }`$ is the matrix transforming from the original (uncoupled) modes to the coupled system ($`q_0=_\nu U_{0\nu }X_\nu `$).
As in previous sections, we are interested in the fluctuations of the variance of $`q_0`$, given in this case by
$$q_0^2(t)=\underset{\mu \nu }{}U_{0\mu }U_{0\nu }X_\mu X_\nu (t),$$
(5.8)
and that we will compute by solving the equation of motion obeyed by the correlations $`X_\mu X_\nu (t)`$. Since $`X_\mu `$ and $`X_\nu `$ correspond to harmonic coordinates, using the quantum mechanical commutation relations we compute the equations of motion to be:
$`{\displaystyle \frac{d}{dt}}X_\mu X_\nu ={\displaystyle \frac{1}{m}}(P_\mu X_\nu +P_\nu X_\mu )`$
$`{\displaystyle \frac{d^2}{dt^2}}X_\mu X_\nu =(\omega _\mu ^2+\omega _\nu ^2)X_\mu X_\nu +{\displaystyle \frac{2}{m^2}}P_\mu P_\nu `$
$`{\displaystyle \frac{d}{dt}}P_\mu P_\nu =m\left(\omega _\mu ^2X_\mu P_\nu +\omega _\nu ^2X_\nu P_\mu \right)`$
$`{\displaystyle \frac{d^2}{dt^2}}P_\mu P_\nu =(\omega _\mu ^2+\omega _\nu ^2)P_\mu P_\nu +2m^2\omega _\mu ^2\omega _\nu ^2X_\mu X_\nu .`$
Note that the above equations are identical to those of classical harmonic oscillators for the quantities $`X_\mu (t)X_\nu (t)`$ etc., with initial conditions given by the values of the correlations evaluated for the quantum wave function:
$`X_\mu X_\nu (0^+)=\delta _{\mu \nu }{\displaystyle \frac{\mathrm{}}{2m\omega _\mu }},`$
$`P_\mu P_\nu (0^+)=\delta _{\mu \nu }{\displaystyle \frac{\mathrm{}m\omega _\mu }{2}}`$
$`+2\mathrm{}^2\lambda ^2(1+\delta _{\mu \nu }){\displaystyle \frac{U_{0\mu }}{m\omega _\mu }}{\displaystyle \frac{U_{0\nu }}{m\omega _\nu }}q_0^2`$
$`(X_\mu P_\nu +P_\nu X_\mu )(0^+)=4\lambda \mathrm{}U_{0\mu }U_{0\nu }{\displaystyle \frac{\mathrm{}}{2m}}({\displaystyle \frac{1}{\omega _\mu }}+{\displaystyle \frac{1}{\omega _\nu }})`$
with $`q_0^2q_0^2(0^{})=_\alpha \mathrm{}U_{0\alpha }^2/2m\omega _\alpha `$.
Collecting the above equations we obtain
$$q_0^2(t)=q_0^2\left\{1+4\lambda ^2S^2(t)+\frac{\lambda }{q_0^2}C(t)S(t)\right\},$$
(5.9)
with
$$S(t)=\underset{\mu }{}\frac{\mathrm{}U_{0\mu }^2}{m\omega _\mu }\mathrm{sin}\omega _\mu tC(t)=\underset{\mu }{}\frac{\mathrm{}U_{0\mu }^2}{m\omega _\mu }\mathrm{cos}\omega _\mu t.$$
All the information of the evolution of the variance is contained in the function $`J(\omega )`$, the physical interpretation of which is that of a local density of states of the oscillator, defined as
$$J(\omega )=\underset{\mu }{}\frac{\mathrm{}U_{0\mu }^2}{m\omega _\mu }\delta (\omega \omega _\mu ),$$
(5.10)
from which
$$S(t)=𝑑\omega J(\omega )\mathrm{sin}\omega t,C(t)=𝑑\omega J(\omega )\mathrm{cos}\omega t.$$
(5.11)
Note that $`J(\omega )`$ is a sum over delta functions, giving rise to a superposition of oscillations with the frequencies $`\omega _\nu `$ for both $`S(t)`$ and $`C(t)`$. In the limit of an infinite system, and when the modes are spatially extended over all space $`J(\omega )`$ becomes a continuous function. In that case the oscillatory behavior acquires a damped component, the detailed time dependence being given by the frequency spectrum of $`J(\omega )`$. A lorenzian shape for $`J(\omega )`$ will give an exponentially damped oscillation for both $`S(t)`$ and $`C(t)`$. As an illustration of this point we consider a model for which $`J(\omega )`$ can be computed explicitly – see the classical analysis in Lamb Komech . Consider a one-dimensional string coupled to our oscillator. The string is described by a “transverse” displacement $`u(x,t)`$. The classical equations of motion of the system are
$`u_{tt}(x,t)`$ $`=`$ $`c^2u_{xx}(x,t)`$
$`Md^2q_0(t)/dt^2`$ $`=`$ $`Vq_0(t)+T[u_x(0+,t)u_x(0,t)]`$
$`q_0(t)`$ $`=`$ $`u(0,t).`$ (5.12)
The normal modes consist of even and odd (in $`x`$) solutions. The odd solutions do not involve $`q_0`$ and are of the form $`u_{q,o}(x,t)=e^{icqt}\mathrm{sin}qx`$, whereas the even solutions are of the form $`u_{q,e}(x,t)=e^{icqt}\mathrm{cos}(q|x|+\delta _q),`$ with $`\delta _q`$ a phase shift (to be found). The wave vectors $`q`$ label the normal modes, and play the role of the index $`\mu `$ in the above discussion: $`\omega _\mu =cq`$, and $`U_{\mu 0}^2=\mathrm{cos}^2(\delta _q)`$ (up to a normalization constant) in the present case. Substituting this expression in (5.12) we obtain ($`\omega _0^2=V/M`$)
$$\mathrm{tan}\delta _q=\frac{Mc}{2T}\frac{(\omega _0^2\omega _q^2)}{\omega _q},$$
(5.13)
from which $`U_{\mu 0}^2=\mathrm{cos}^2\delta _q`$ is given by
$$U_{\mu 0}^2=\frac{\alpha ^2\omega _q^2}{\alpha ^2\omega _q^2+(\omega _q^2\omega _0^2)^2}U_q^2,$$
(5.14)
where we have defined $`\alpha =2T/Mc`$. Note that $`U_q`$ represents the transformation matrix that has to be normalized and since the frequencies form a continuum we normalize $`U_q(\omega _q)`$ to its integral over $`\omega _q`$. Omitting the index $`q`$ in $`\omega _q`$, we obtain
$$U(\omega )=\frac{2\alpha }{\pi }\frac{\omega ^2}{\alpha ^2\omega ^2+(\omega ^2\omega _0^2)^2}=\frac{m\omega }{\mathrm{}}J(\omega ).$$
(5.15)
Substituting (5.15) in (5.11) we obtain
$$S(t)=\frac{\mathrm{}}{m\omega _0}e^{\mathrm{\Gamma }t}\mathrm{sin}\mathrm{\Omega }_0t,C(t)=\frac{\mathrm{}}{m\omega _0}e^{\mathrm{\Gamma }t}\mathrm{cos}\mathrm{\Omega }_0t,$$
(5.16)
with
$`\mathrm{\Omega }_0`$ $`=`$ $`\omega _0\left(1+\left[\alpha /\omega _0\right]^2\right)^{1/4}\mathrm{cos}\delta /2,`$ (5.17)
$`\mathrm{\Gamma }`$ $`=`$ $`\omega _0\left(1+\left[\alpha /\omega _0\right]^2\right)^{1/4}\mathrm{sin}\delta /2,`$ (5.18)
where $`\delta =\mathrm{tan}^1\alpha /\omega _0`$.
In the realistic limit $`\alpha \omega _0`$ which corresponds to a “weak” coupling to the environment) this expressions take the form: $`S(t)(\mathrm{}/(m\omega _0)\mathrm{exp}(Tt/Mc)\mathrm{sin}\omega _0t,C(t)(\mathrm{}/(m\omega _0)\mathrm{exp}(Tt/Mc)\mathrm{cos}\omega _0t.`$
Note that in this model, and in the limit of weak coupling, the initial variance $`q_0^2`$ of the reference oscillator is unchanged due to the coupling to the environment, and is given by $`q_0^2=\mathrm{}/2m\omega _0`$. Our final result for this section is then
$`q_0^2(t)q_0^2\{1+e^{2(T/Mc)t}`$
$`[\left({\displaystyle \frac{2\lambda \mathrm{}}{m\omega _0}}\right)\mathrm{sin}2\omega _0t+\left({\displaystyle \frac{2\lambda \mathrm{}}{m\omega _0}}\right)^2\mathrm{sin}\omega _0^2t]\},`$ (5.19)
which reduces simply to (4.13) in the uncoupled case of $`T=0`$.
In summary we have shown in this section that the coupling to the environment can be included in general, giving rise to dissipation, and that the squeezing effect in the presence of dissipation can be computed explicitly for the Lamb model.
Additional details of the analysis here, extensions to the squeezing of a nonlinear oscillator, and a treatment of the quantum measurement issue will appear in forthcoming publications.
Acknowledgement: We would like to thank Roger Brockett for useful discussions.
References
Brockett, R. and N. Khaneja On the stochastic control of quantum ensembles, preprint.
Garret, G.A., A.G. Rojo, A.K. Sood, J.F. Whitaker and R. Merlin Vacuum Squeezing of solids: Macroscopic quantum states driven by light pulses, Science 275, 1638-1640.
Hagerty, P. A.M. Bloch and M.I Weinstein Radiation induced instability in interconnected systems, Proc. 37th Conference on Decisions and Control, IEEE, 651-656.
Keller, J.B. and L.B. Bonilla Irreversibility and nonrecurrance, Journal of Statistical Physics 42 1115.
Komech, A. I. On the stabilization of string-oscillator interaction. Russian Journal of Mathematical Physics 3, 2, 227-247.
Lamb, H. On the peculiarity of the wave-System due to the free vibrations of a nucleus in an extended medium, Proceeding of the London Math. Society, 32, 208-211.
Lloyd, S. Universal quantum simulators, Science 273, 1073-1078.
Negele, J. W. and H. Orland Quantum Many Particle Physics, Addison-Wesley, New York.
Soffer, A. and Weinstein, M.I. Resonances, Radiation Damping and Instability in Hamiltonian Nonlinear Wave Equations, Invent. Math. 136 9–74.
Warren, W., H. Rabitz, and M. Dahleh Coherent control of quantum dynamics: The dream is alive, Science 259, 1581-1589.
Ziman, J. Principles of The Theory of Solids, Second Edition, Cambridge University Press. |
warning/0002/cond-mat0002128.html | ar5iv | text | # Statistical Physics of Structural Glasses
## 1 Introduction
While the experimental and phenomenological knowledge on glasses has improved a lot in the last decades, the progress on a first principle, statistical mechanical study of the glass phase has turned out to be much more difficult.
Take any elementary textbook on solid state physics. It deals with a special class of solid state, the crystalline state, and usually avoids to elaborate on the possibility of amorphous solid states. The reason is very simple: there is no theory of amorphous solid states. Schematically, the first elementary steps of the theory of crystals are the following. One computes the ground state energy of all the crystalline structures. The small vibrations around these structures are easily handled, either using the simple Einstein approximation of independent atoms in harmonic traps, or computing the phonon dispersion relations and going to the Debye theory. Then one can study the one electron problem and compute the band structure. The basic thermodynamic properties are already well reproduced by these elementary computations. Anharmonic vibrations, electron-phonon and electron-electron interactions can then be added to these basic building blocks.
Until very recently, none of the above computations, even in the simplest-minded approximation, could be done in the case of the glass state. The reason is obvious: all of them are made possible in crystals by the existence of the symmetry group. The absence of such a symmetry, which is a defining property of the glass state, forbids the use of all the solid state techniques. If one takes a snapshot of a glass state, an instantaneous configuration of atoms, it looks more like a liquid configuration. In fact the techniques which we shall use are often borrowed from the theory of the liquid state. But while the liquid phase is ergodic (which means that the probability distribution of positions is translationally invariant), the glass phase is not. The problem is to describe a non-ergodic phase without a symmetry: an amorphous solid state.
The work which we report on here has been elaborated during the last year and aims at building the first steps of a first principle theory of glasses. The fact that this is being made possible now is not fortuitous, but rather results from a conjunction of several sets of ideas, and the general progress of the last two decades on the theory of amorphous systems.
The oldest ingredients are the phenomenological ideas, originating in the work of Kauzmann , and developed among others by by Adam, Gibbs and Di-Marzio , which identify the glass transition as a ‘bona fide’ thermodynamic transition blurred by some dynamical effects. As we shall discuss below, in this scenario the transition is associated with an ‘entropy crisis’, namely the vanishing of the configurational entropy of the thermodynamically relevant glass states.
A very different, and more indirect, route, was the study of spin glasses. These are also systems which freeze into amorphous solid states, but one of their constitutive properties is very different from the glasses we are interested in here: there exists in spin glasses some ‘quenched disorder’: the exchange-interaction coupling constants between the spin degrees of freedom are quenched (i.e. time independent on all experimental time scales) random variables. Anyhow, a few years after the replica symmetry breaking (RSB) solution of the mean field theory of spin glasses , it was realized that there exists another category of mean-field spin glasses where the transition is due to an entropy crisis . These are now called discontinuous spin glasses because their phase transition, although it is of second order in the Ehrenfest sense, has a discontinuous order parameter, as first shown in . Another name often found in the literature is ‘one step RSB’ spin glasses, because of the special pattern of symmetry breaking involved in their solution. The simplest example of these is the Random Energy Model , but many other such discontinuous spin glasses were found subsequently, involving multispin interactions .
The analogy between the phase transition of discontinuous spin glasses and the thermodynamic glass transition was first noticed by Kirkpatrick, Thirumalai and Wolynes in a series of inspired papers of the mid-eighties . While some of the basic ideas of the present development were around at that time, there still missed a few crucial ingredients. On one hand one needed to get more confidence that this analogy was not just fortuitous. The big obstacle was the existence (in spin glasses) versus the absence (in structural glasses) of quenched disorder. The discovery of discontinuous spin glasses without any quenched disorder provided an important new piece of information: contrarily to what had been believed for long, quenched disorder is not necessary for the existence of a spin glass phase (but frustration is). A second confirmation came very recently from the developments on out of equilibrium dynamics of the glass phase. Initiated by the exact solution of the dynamics in a discontinuous spin glass by Cugliandolo and Kurchan , this line of research has made a lot of progress in the last few years. It has become clear that, in realistic systems with short range interactions, the pattern of replica symmetry breaking can be deduced from the measurements of the violation of the fluctuation dissipation theorem . Although these difficult measurements are not yet available, numerical simulations performed on different types of glass forming systems have provided an independent and spectacular confirmation of their ‘one step rsb’ structure on the (short) time scales which are accessible. The theory was then facing the big challenge: understanding what this replica symmetry breaking could mean, in systems void of quenched disorder, in which there is thus no a priori reason to introduce replicas. The recent progress has brought the answer to this question and turned it into a computational method, allowing for a first principle computation of the equilibrium thermodynamics of glasses .
In the context of glasses, the words ‘equilibrium thermodynamics’ call for some comments. First, it is not obvious whether the glass phase is an equilibrium phase of matter. It might be a metastable phase, reachable only by some fast enough quench, while the ‘true equilibrium’ phase would always be crystalline. The answer depends on the interaction potential. Numerically it is known that the frustration induced by considering for instance binary mixtures of soft spheres of different radii strongly inhibits crystallisation. But what is the true equilibrium state is unknown, and not very relevant. One can study crystals without having proven that they are stable phases of matter (by the way, simply proving that the fcc-hcp is the densest packing of hard spheres in 3 dimensions, a simple zero temperature statement, has resisted the efforts of scientists for centuries ), and one can study the properties of diamond, even though it is notoriously unstable. The point is to have reproducible properties, which is certainly the case. Letting aside the crystal, a more interesting question is how to reach equilibrium glass states. Experimentally nobody knows how to achieve this. In a ferromagnet, one can reach an equilibrium state and eliminate domain walls by using an external magnetic field. In a glass there is no such field conjugate to the order parameter, and the fate is an out of equilibrium situation. The same is true in spin glasses, and in fact in all kind of glass phases. Why study the equilibrium thermodynamics then? The answer is twofold. First principle computations are certainly much easier as far as the equilibrium is concerned, therefore it is natural to start with these in order to first get some detailed understanding of the free energy landscape, which will be useful in the more realistic dynamical studies. Secondly, we have strong indications, and some general arguments, in favour of a close relationship between the equilibrium properties and the observable out of equilibrium dynamical observations . Let us also mention here the recent developments of some phenomenological theory of the out of equilibrium theory of glasses .
In this paper we shall introduce the main ideas of the recent elaboration of the equilibrium theory of glasses. We shall not present the details which can be found in the literature. The general replica strategy can be found in . The explicit computations have been done first for soft spheres in , and then generalized to binary mixtures of soft spheres or Lennard Jones particles .
## 2 Hypotheses on the glass phase
The general framework of our approach is a familiar one in physics: we shall start from a number of basic hypotheses on the glass phase, derive some quantitative properties starting from these hypotheses, and then compare them with numerical, and hopefully, in the future, experimental results. We work with a simple glass former, $`N`$ undistinguishable particles move in a volume $`V`$ of a d-dimensional space, and we take the thermodynamic limit $`N,V\mathrm{}`$ at fixed density $`\rho =N/V`$. The interaction potential is a two body one, defined by a short range function $`v(x)`$ (for instance one may consider a soft spheres system where $`v(x)=1/x^{12}`$).
Let us introduce a free energy functional $`F(\rho )`$ which depends on the density $`\rho (x)`$ and on the temperature. We suppose that at sufficiently low temperature this functional has many minima (i.e. the number of minima goes to infinity with the number $`N`$ of particles). Exactly at zero temperature these minima, labelled by an index $`\alpha `$, coincide with the mimima of the potential energy as function of the coordinates of the particles. A more detailed discussion of the valleys and their relationship to the inherent structures will be given in sect. 6. To each valley we can associate a free energy $`F_\alpha `$ and a free energy density $`f_\alpha =F_\alpha /N`$. The number of free energy minima with free energy density $`f`$ is supposed to be exponentially large:
$$𝒩(f,T,N)\mathrm{exp}(N\mathrm{\Sigma }(f,T)),$$
(1)
where the function $`\mathrm{\Sigma }`$ is called the complexity or the configurational entropy (it is the contribution to the entropy coming from the existence of an exponentially large number of locally stable configurations). This function is not defined in the regions $`f>f_{max}(T)`$ or $`f<f_{min}(T)`$, where $`𝒩(f,T,N)=0`$, it is convex and it is supposed to go to zero continuously at $`f_{min}(T)`$, as found in all existing models so far (see fig.1). In the low temperature region the total free energy of the system, $`\mathrm{\Phi }`$, can be well approximated by:
$$e^{\beta N\mathrm{\Phi }}\underset{\alpha }{}e^{\beta Nf_\alpha (T)}=_{f_{min}}^{f_{max}}𝑑f\mathrm{exp}\left(N[\mathrm{\Sigma }(f,T)\beta f]\right),$$
(2)
where $`\beta =1/T`$. The minima which dominate the sum are those with a free energy density $`f^{}`$ which minimizes the quantity $`\mathrm{\Phi }(f)=fT\mathrm{\Sigma }(f,T)`$. At large enough temperatures the saddle point is at $`f>f_{min}(T)`$. When one decreases $`T`$ the saddle point free energy decreases. The Kauzman temperature $`T_K`$ is that below which the saddle point sticks to the minimum: $`f^{}=f_{min}(T)`$. It is a genuine phase transition, the ‘ideal glass transition’.
This scenario for the glass transition is precisely the one which is at work in discontinuous spin glasses, and can be studied there in full details. The transition is of a rather special type. It is of second order because the entropy and internal energy are continuous. When decreasing the temperature through $`T_K`$ there is a discontinuous decrease of specific heat, as seen experimentally. On the other hand the order parameter is discontinuous at the transition, as in first order transitions. To show this we have to provide a definition of the order parameter in our framework of equilibrium statistical mechanics. This is not totally trivial because of the lack of knowledge on the valleys themselves. The best way is to introduce two identical copies of the system. We have one system of undistinguishable ‘red’ particles, interacting between themselves through $`v(x)`$, another system of undistinguishable ‘blue’ particles, interacting between themselves through $`v(x)`$, and we turn on a small interaction between the blue and red particles, which is short range. We take the thermodynamic limit first, and then send this red-blue coupling to zero. If the position correlations between the red and blue particles disappear in this double limit, the system is in a liquid phase, otherwise it is in a solid phase. Clearly, the order parameter, which is the red-blue pair correlation function, is discontinuous at the transition: there is no correlation in the liquid phase, while in the solid phase one gets an oscillating pair correlation, similar to that of a dense liquid, but with an extra peak at the origin. In some sense, in this framework, the role of the unknown conjugate field, needed in order to polarize the system into one state, is played by the coupling to the second copy of the system. The small red-blue coupling is here to insure that the two systems will fall into the same glass state.
The above scenario, relating the glass transition to the vanishing of the configurational entropy, is the main hypothesis of our work. Clearly it is in agreement with the phenomenology of the glass transition, and with the old ideas of Kauzman, Gibbs and Di-Marzio. It is also very interesting from the point of view of the dynamical behaviour.
In discontinuous mean field spin glasses, the slowing down of the dynamics takes a very special form. There exist a dynamical transition temperature $`T_c>T_K`$. When T decreases and gets near to $`T_c`$, the correlation function relaxes with a characteristic two step forms: a fast $`\beta `$ relaxation leading to a plateau takes place on a characteristic time which does not grow, while the $`\alpha `$ relaxation from the plateau takes place on a time scale which diverges when $`TT_c`$. This dynamic transition is exactly described by the schematic mode coupling equations. The existence of a dynamic relaxation at a temperature above the true thermodynamic one is possible only in mean field, and the conjecture is that in a realistic system like a glass, the region between $`T_K`$ and $`T_c`$ will have instead a finite, but very rapidly increasing, relaxation time, as shown in fig. 2.
On this figure we see the existence of several temperature regimes:
-a relatively high temperature regime where mode coupling theory applies
\- an intermediate region, extending from $`T_k`$ up to the temperature above $`T_c`$ where mode coupling predictions start to be correct. This is the region of activated processes, where one can identify some traps in phase space in which the system stays for a long time, and then jumps.
-the low temperature, glass phase $`T<T_K`$.
The dynamics of the glass is expected to show aging effects in the glass region, but also in the intermediate region provided the laboratory time is smaller than the relaxation time.
Here we shall focus onto the equilibrium study of the low temperature phase. One main reason is that the direct study of out of equilibrium dynamics is more difficult, and that one might be able to make progress by a careful analysis of the landscape . Another motivation is to go into a more quantitative test of the basic scenario: while it agrees qualitatively with several observations, as we just discussed, it should also be able to help make more quantitative predictions.
Our strategy will be to start from this set of hypotheses and derive the quantitative predictions which can be checked independently. We shall be able to compute for instance the configurational entropy versus free energy within some well controlled approximations, and compare it to the results of some numerical simulations.
## 3 Replicas
In order to cope with the degeneracy of glass states and the existence of a configurational entropy, a choice method is the replica method. Initially replicas were introduced in order to study systems with quenched disorder, in which one needs to compute the disorder average of the logarithm of the partition function . It took a few years to realize that a large amount of information is encoded in the distribution of distances between replicas. This is true again in structural glasses. The simplest example was given above when we explained the use of two replicas in order to define the order parameter. A much more detailed information can be gained if one studies in general a set of $`m`$ replicas, sometimes named ‘clones’ in this context, coupled through a small extensive attraction which will eventually go to zero . In the glass phase, the attraction will force all $`m`$ systems to fall into the same glass state, so that the partition function is:
$$Z_m=\underset{\alpha }{}e^{\beta Nmf_\alpha (T)}=_{f_{min}}^{f_{max}}𝑑f\mathrm{exp}\left(N[\mathrm{\Sigma }(f,T)m\beta f]\right)$$
(3)
In the limit where $`m1`$ the corresponding partition function $`Z_m`$ is dominated by the correct saddle point $`f^{}`$ for $`T>T_K`$. The interesting regime is when the temperature is $`T<T_K`$, and the number $`m`$ is allowed to become smaller than one. The saddle point $`f^{}(m,T)`$ in the expression (3) is the solution of $`\mathrm{\Sigma }(f,T)/f=m/T`$. Because of the convexity of $`\mathrm{\Sigma }`$ as function of $`f`$, the saddle point is at $`f>f_{min}(T)`$ when $`m`$ is small enough, and it sticks at $`f^{}=f_{min}(T)`$ when $`m`$ becomes larger than a certain value $`m=m^{}(T)`$, a value which is smaller than one when $`T<T_K`$. The free energy in the glass phase, $`F(m=1,T)`$, is equal to $`F(m^{}(T),T)`$. As the free energy is continuous along the transition line $`m=m^{}(T)`$, one can compute $`F(m^{}(T),T)`$ from the region $`mm^{}(T)`$, which is a region where the replicated system is in the liquid phase. This is the clue to the explicit computation of the free energy in the glass phase. It may sound a bit strange because one is tempted to think of $`m`$ as an integer number. However the computation is much clearer if one sees $`m`$ as a real parameter in (3). As one considers low temperatures $`T<T_K`$ the $`m`$ coupled replicas fall into the same glass state and thus they build some molecules of $`m`$ atoms, each molecule being built from one atom of each ’colour’. Now the interaction strength of one such molecule with another one is basically rescaled by a factor $`m`$ (this statement becomes exact in the limit of zero temperature where the molecules become point like). If $`m`$ is small enough this interaction is small and the system of molecules is liquid. When $`m`$ increases, the molecular fluid freezes into a glass state at the value $`m=m^{}(T)`$. So our method requires to estimate the replicated free energy, $`F(m,T)=\mathrm{log}(Z_m)/(\beta mN)`$, in a molecular liquid phase, where the molecules consist of $`m`$ atoms and $`m`$ is smaller than one. For $`T<T_K`$, $`F(m,T)`$ is maximum at the value of $`m=m^{}`$ smaller than one, while for $`T>T_K`$ the maximum is reached at a value $`m^{}`$ is larger than one. The knowledge of $`F_m`$ as a function of $`m`$ allows to reconstruct the configurational entropy function $`Sc(f)`$ at a given temperature $`T`$ through a Legendre transform, using the parametric representation (easily deduced from a saddle point evaluation of (3)):
$$f=\frac{\left[mF(m,T)\right]}{m};\mathrm{\Sigma }(f)=\frac{m^2}{T}\frac{F(m,T)}{m}.$$
(4)
The Kauzmann temperature (’ideal glass temperature’) is the one such that $`m^{}(T_K)=1`$. For $`T<T_K`$ the equilibrium configurational entropy vanishes. Above $`T_K`$ one obtains the equilibrium configurational entropy $`\mathrm{\Sigma }(T)`$ by solving (4) at $`m=1`$.
More explicitly, one must thus introduce $`m`$ clones of each particle, with positions $`x_i^a,a1,\mathrm{},m`$. The replicated partition function is:
$`Z_m={\displaystyle \frac{1}{N!^m}}{\displaystyle }{\displaystyle \underset{i=1}{\overset{N}{}}}{\displaystyle \underset{a=1}{\overset{m}{}}}dx_i^a\mathrm{exp}(\beta {\displaystyle \underset{1i<jN}{}}{\displaystyle \underset{a=1}{\overset{m}{}}}v(x_i^ax_j^a)`$
$`\beta ϵ{\displaystyle \underset{i,j=1}{\overset{N}{}}}{\displaystyle \underset{1a<bm}{}}w(x_i^ax_j^b)),`$ (5)
where $`v`$ is the original interparticle potential and $`w`$ is an attractive potential. This attractive potential must be of short range (the range should be less than the typical interparticle distance in the solid phase), but its precise form is irrelevant. Assuming that $`w`$ is equal to $`1`$ at very small distances, and zero at large distances (notice that the scale of the inter-replica interaction is fixed by the parameter $`ϵ`$), the coupling $`w`$ can be used to define an overlap between two configurations, in a way similar to the crucial concept of overlaps in spin glasses. Taking two configurations $`x_i`$ and $`y_i`$ of the $`N`$ particles, one defines the overlap between the configurations as $`q(x,y)=1/N_{i,k=1,N}w(x_iy_k)`$, or the distance as $`1q`$. The replicated partition function with $`m`$ clones is thus (in more compact notations where $`dx=_{i=1}^Ndx_i/N!`$ and $`H(x)_{i<j}v(x_ix_j)`$ is the total energy of the system):
$$Z_m=\underset{a}{}dx^a\mathrm{exp}\left(\beta \underset{a}{}H(x_a)+\beta ϵN\underset{a,b}{}q(x_a,x_b)\right).$$
(6)
This can be defined also for non integer $`m`$ using an analytic continuation (if our hypothesis of the glass transition being of the same nature as the one step rsb in spin glasses is correct, there is no replica symmetry breaking between the clones, and the continuation is straightforward). Alternatively, one can define it through the formula
$$Z_m𝑑\mu (\varphi )Z(\varphi )^m$$
(7)
where $`\varphi `$ is a quenched random potential defined in the full space, which has a Gaussian distribution with moments:
$$d\mu (\varphi )=1,d\mu (\varphi )\varphi (x)=0,d\mu (\varphi )\varphi (x)\varphi (y)=c^tw(xy),$$
(8)
and $`Z(\varphi )`$ is the partition function of one system in the external potential $`\varphi `$:
$$Z(\varphi )=𝑑x\mathrm{exp}\left(\beta H(x)\sqrt{\beta ϵ}\underset{i=1}{\overset{N}{}}\varphi (x_i)\right).$$
(9)
## 4 The molecular liquid
The explicit computation of $`Z_m`$ in the regime $`m<m^{}(T)`$ is a complicated problem of dense molecular liquids, which requires some approximate treatments. Several types of approximations have been developed recently, leading to fully consistent results. Focusing onto the low temperature regime, where the molecules have a small radius, it is natural to write the partition function in terms of the center of mass and relative coordinates $`\{r_i,u_i^a\}`$, with $`x_i^a=r_i+u_i^a`$ and $`_au_i^a=0`$, and to expand the interaction in powers of the relative displacements $`u`$. After a proper renumbering of the particles, so that particles in the same molecule have the same $`i`$ index, one gets:
$`Z_m`$ $`=`$ $`{\displaystyle \frac{1}{N!}}{\displaystyle }dr{\displaystyle \underset{a=1}{\overset{m}{}}}du^a{\displaystyle \underset{i=1}{\overset{N}{}}}\left(m^3\delta ({\displaystyle \underset{a=1}{\overset{m}{}}}u_i^a)\right)\mathrm{exp}(\beta {\displaystyle \underset{i<j,a}{}}[v(r_ir_j)`$ (10)
$`+{\displaystyle \underset{p=2}{\overset{\mathrm{}}{}}}(u_i^au_j^a)^p{\displaystyle \frac{v^{(p)}(r_ir_j)}{p!}}]{\displaystyle \frac{ϵ}{4}}{\displaystyle \underset{i,a,b}{}}(u_i^au_i^b)^2).`$
The last term is the small inter-replica coupling ($`ϵ`$ will be sent to zero in the end), which we have approximated for convenience by its quadratic approximation. The expression (10) can be expanded, at low temperatures, in the following ways:
* ‘Harmonic resummation’: One keeps only the $`p=2`$ term. The action is quadratic in $`u`$, and after performing the exact $`u`$ integral one obtains an effective interaction for the center of mass degrees of freedom, which we shall detail below. The parameter $`m`$ appears as a coupling constant, the analytic continuation in $`m`$ is thus trivial, and the whole problem reduces to treating the liquid of center of masses, interacting through the effective interaction.
* ‘Small cage expansion’: One expands the exponential in powers of the relative variables $`u`$, keeping only the $`ϵ`$ term in the exponent. Again, the $`u`$ integrals can be done exactly to each order of the approximation. In this way one generates an expansion of the free energy in powers of $`1/ϵ`$. This function can be Legendre transformed with respect to $`ϵ`$, leading to a generalized free energy expressed as a series in terms of the ‘cage radius’, $`A=2/(3m(m1))_{a,b}<(u_i^au_i^b)^2>`$. Notice that the $`1/ϵ`$ expansion is just an intermediate step in order to generate the small $`A`$ expansion of the potential (the same can be done for instance when computing the Gibbs potential of an Ising model in terms of the magnetization $`M`$ at low temperatures: even if one is interested in the zero magnetic field case, one can introduce the field as an intermediate device and first expand in powers of $`\mathrm{exp}(\beta h)`$, before turning the result into an expansion in $`1M`$).
The two methods are complementary. They both lead to the study of a liquid of center of mass positions. The small cage expansion is simpler because the result is expressed in terms of various correlation functions of the pure liquid of center of masses at the effective temperature $`T/m`$, which can be handled using traditional liquid state techniques. On the other hand the leading ($`p=2`$) term at low temperatures is not treated exactly. In the harmonic resummation scheme the interaction potential of the center of masses is modified: one gets
$$Z_m=Z_m^0𝑑r\mathrm{exp}\left(\beta mH(r)\frac{m1}{2}Tr\mathrm{log}M\right)$$
(11)
where $`Z_m^0=m^{Nd/2}\sqrt{2\pi T}^{Nd(m1)}/N!`$, and the matrix $`M`$, of dimension $`dN\times dN`$, is given by:
$$M_{(i\mu )(j\nu )}=\frac{^2H(r)}{r_i^\mu r_j^\nu }=\delta _{ij}\underset{k}{}v_{\mu \nu }(r_ir_k)v_{\mu \nu }(r_ir_j)$$
(12)
and $`v_{\mu \nu }(r)=^2v/r_\mu r_\nu `$ (the indices $`\mu `$ and $`\nu `$ denote space directions). The effective interaction contains the complicated ‘$`Tr\mathrm{log}M`$’ piece which is not a pair potential. Because of this term, in the whole glass phase where one is interested in the $`m<1`$ regime, the partition function receives some contributions only from those configurations $`r_i`$ such that all eigenvalues of $`M`$ are positive: these are locally stable glass configurations. In order to handle this additional constraint, we used so far the following (rather crude) approximate treatment, which consists of two steps. First, a ’quenched approximation’, which amounts to neglecting the feedback of vibration modes onto the centers of masses, substitutes $`\mathrm{exp}\left(\frac{m1}{2}Tr\mathrm{log}M\right)`$ by $`\mathrm{exp}\left(\frac{m1}{2}Tr\mathrm{log}M\right)`$, where $`.`$ is the Boltzmann expectation value at the effective temperature $`T/m`$. One is then left with the computation of the spectrum of $`M`$ in a liquid. This is an interesting problem in itself. The treatment done in corresponds to keeping the leading term in a high density limit. Further recent progress should allow for a better controlled approximation of the spectrum.
We shall not review here the details of these computations, which can be found in as far as the simple glass former with the ‘soft sphere’ $`1/x^{12}`$ potential is concerned, in for the mixtures of soft spheres and in for mixtures of Lennard-Jones particles. Once one has derived an expression for the replicated free energy, one can deduce from it the whole thermodynamics, as described above. In all three cases, one finds an estimate of the Kauzman temperature which is in reasonable agreement with simulations, with a jump in specific heat, from a liquid value at $`T>T_K`$ to the Dulong-Petit value $`C=3/2`$ (we have included only positional degrees of freedom) below $`T_K`$. This is similar to the experimental result, where the glass specific heat jumps down to the crystal value when one decreases the temperature (Our approximations so far are similar to the Einstein approximation of independent vibrations of atoms, in which case the contribution of positional degrees of freedom to the crystal specific heat is $`C=3/2`$). The parameter $`m^{}(T)`$ and the cages sizes are nearly linear with temperature in the whole glass phase. This means, in particular, that the effective temperature $`T/m`$ is always close to $`T_K`$, so in our theoretical computation we need only to evaluate the expectation values of observables in the liquid phase, at temperatures where the HNC approximation for the liquid still works quite well.
A more detailed numerical checks of these analytical predictions involves the measurement of the configurational entropy. We shall review these checks in sect. 6, but we first wish to present some alternative derivation of the low temperature results.
## 5 Without replicas
For those who do not appreciate the beauty and efficacy of the replica approach, it may be useful to derive some of the above results without resorting to the replica method . Specifically, we shall study the simplest case of the zero temperature limit in the harmonic approximation through a direct approach, and reinterpret the above results. At low temperatures, the critical value $`m^{}`$ of the parameter $`m`$ goes to zero linearly with $`T`$. We thus write $`\gamma =\beta m`$ and take the $`T,m0`$ limit of (11) at fixed $`\gamma `$. This gives:
$$Z_m\left(\frac{\gamma }{2\pi }\right)^{Nd/2}_C𝑑r\sqrt{detM(r)}\mathrm{exp}\left(\gamma H(r)\right),$$
(13)
where $`_C`$ is restricted to configurations in which all eigenvalues of $`M`$ are positive. A direct derivation of this formula, making all hypotheses explicit, is the following. At zero temperature one is interested in configurations where every particle is in equilibrium: $`i,\mu ,H/x_i^\mu =0`$. The number of such configurations at energy $`NE`$,
$$\mu (E)=𝑑x|detM(x)|\delta (NEH(x))\underset{i,\mu }{}\delta \left(\frac{H}{x_i^\mu }\right),$$
(14)
can be approximated at low enough energy, where most extrema are minima , by the expression
$$\nu (E)=_C𝑑x𝑑etM(x)\delta (NEH(x))\underset{i,\mu }{}\delta \left(\frac{H}{x_i^\mu }\right).$$
(15)
Within this approximation $`\nu (E)`$ is related to the configurational entropy through $`\nu (E)=\mathrm{exp}(N\mathrm{\Sigma }(E))`$, and one can compute its Laplace transform:
$$\zeta (\gamma )𝑑E\nu (E)\mathrm{exp}\left(\gamma NE\right)=𝑑E\mathrm{exp}\left(N\left[\mathrm{\Sigma }(E)\gamma E\right]\right).$$
(16)
Using an exponential representation of the ground state constraints, this effective partition function is:
$$\zeta (\gamma )=\left(\frac{\gamma }{2\pi }\right)^{Nd}𝑑x\underset{k}{}d\lambda _k^\mu detM(x)\mathrm{exp}\left(\gamma H(x)+i\gamma \underset{k,\mu }{}\lambda _k^\mu \frac{H}{x_k^\mu }\right)$$
(17)
One can change variables from $`x_k`$ to $`y_k=x_ki\lambda _k`$. At low temperatures it is reasonable to assume that the only configurations which contribute are those in the neighborhood of the minima. Expanding in powers of $`\lambda `$, and neglecting anharmonic terms, one writes:
$`H(x)`$ $``$ $`H(y)+i{\displaystyle \underset{k,\mu }{}}\lambda _k^\mu {\displaystyle \frac{H(y)}{y_k^\mu }}{\displaystyle \frac{1}{2}}{\displaystyle \underset{k,\mu ,l,\nu }{}}\lambda _k^\mu \lambda _l^\nu {\displaystyle \frac{^2H(y)}{y_k^\mu y_l^\nu }}`$
$`{\displaystyle \frac{H(x)}{x_k^\mu }}`$ $``$ $`{\displaystyle \frac{H(y)}{y_k^\mu }}+{\displaystyle \underset{l,\nu }{}}\lambda _l^\nu {\displaystyle \frac{^2H(y)}{y_k^\mu y_l^\nu }}.`$ (18)
The $`\lambda `$ integral in (17) is then quadratic, and one gets:
$$\zeta (\gamma )=\left(\frac{\gamma }{2\pi }\right)^{Nd/2}_C𝑑y\sqrt{detM(y)}\mathrm{exp}(\gamma H(y)),$$
(19)
a result identical to the low $`T`$ limit (13) of the replica approach within the harmonic approximation.
## 6 Configurational entropy: theory and simulations
The configurational entropy (sometimes called also complexity) is a key concept in the theory of glasses. There is no difficulty of principle in defining a valley and its entropy in the low temperature phase $`T<T_K`$. As we have seen, we can take a thermalized configuration as a reference system, add a small attraction to this configuration, and take the thermodynamic limit before the limit of a vanishing attraction. This procedure defines the restricted partition function in the valley containing the reference configuration $`y`$, and therefore the free energy of the valley. Computing $`S_c(f,T)`$ is thus in principle doable, but it is still a formidable challenge to get equilibrated configurations $`y`$ in this temperature range.
On the other hand in the intermediate temperature regime $`T_K<T<T_c`$, the valleys and the configurational entropy remain well defined in the mean field theory. The existence of a decoupling of time scales points to the possibility of defining metastable valleys in the whole region where activated (’hopping’) processes are found. This region is particularly interesting, both because of the rapid change of relaxation times, and because part of this region can be studied experimentally or numerically. It often happens that different authors use different definitions of the configurational entropy, which should be hopefully be equivalent at low temperature but behave rather differently at high temperatures. Therefore it seems to us appropriate to start this section with a comparison of the various definitions of configurational entropies which have been introduced and studied so far.
If we consider the configurational entropy versus temperature, which is non-zero for $`T>T_K`$, in a first approximation we can distinguish three different types of definitions:
* A first definition is based on the presence of many minima of the Hamiltonian, i.e. inherent structures.
* A second definition is based on the fact that the phase space at sufficient low energy may be decomposed in many disconnected region (let us call it the microcanonical one).
* A third definition is based on the thermodynamics. One starts from the definition
$$S(T)=\mathrm{\Sigma }(T)+S_{valley}(T)$$
(20)
where $`S(T)`$ is the total entropy and $`S_{valley}`$ is the entropy of the generic valley at temperature $`T`$. In this case the problem consists in finding a precise definition of $`S_{valley}`$.
In this paper we have used the third definition, however we think useful to recall the other two definitions in order to avoid possible misunderstanding.
### 6.1 The inherent structure entropy
Given the Hamiltonian $`H(x)`$ of a system with $`N`$ particles, we can consider the solution $`x(t)`$ of the equation
$$\frac{dx}{dt}=\frac{H}{x}$$
(21)
as function of the initial conditions $`x(0)`$. At large time $`x(t)`$ will go to one of the minima of the Hamiltonian, called an inherent structure. We label by $`a`$ each coherent structure and we call $`𝒟_a`$ the set of those configurations which for large times go to the coherent structure labeled by $`a`$. The union of all the sets $`𝒟_a`$ is the whole phase space. The probability of finding the system at a temperature $`T`$ inside a given inherent structure is proportional to
$$P(a)=Z(a)/\underset{b}{}Z(b);Z(a)_{x𝒟_a}dx\mathrm{exp}(\beta H(x)).$$
(22)
The configurational entropy density, $`\mathrm{\Sigma }_{is}`$, is defined by
$$N\mathrm{\Sigma }_{is}(T)=\underset{a}{}P(a)\mathrm{ln}(P(a)).$$
(23)
This definition makes sense at all temperatures. In the limit of large $`T`$ one finds
$$\underset{T\mathrm{}}{lim}\mathrm{\Sigma }_{is}(T)=\underset{a}{}V(a)\mathrm{ln}(V(a)),$$
(24)
where $`V(a)`$ is proportional to the volume in phase space of the region $`𝒟_a`$, normalized in such a way that $`_aV(a)=1`$. It is reasonable to expect that this inherent-structures configurational entropy starts to decrease when the temperature is decreased around $`T=T_c`$ and vanishes at $`T=T_K`$.
### 6.2 Microcanonical entropy
We consider the hypersurface of constant energy density, $`H(x)=EN`$, and decompose this energy surface in connected components which we label by $`a`$. The number of connected components clearly depends on $`E`$.
Calling $`V_a`$ the normalized phase space volume of each connected component, we define the microcanonical configurational entropy density as
$$N\widehat{\mathrm{\Sigma }}_m(E)=\underset{a}{}V(a)\mathrm{ln}(V(a))$$
(25)
The microcanonical configurational entropy density as function of the temperature is naturally defined as
$$\mathrm{\Sigma }_m(T)=\widehat{\mathrm{\Sigma }}_m(E(T))$$
(26)
where $`E(T)`$ is the internal energy density as function of the temperature. It is clear that at high energies the configuration space contains only one connected component and therefore
$$\underset{T\mathrm{}}{lim}\mathrm{\Sigma }_m(T)=0$$
(27)
The two configurational entropies introduced so far, $`\mathrm{\Sigma }_{is}(T)`$ and $`\mathrm{\Sigma }_m(T)`$ certainly differ at high temperature and many hands must be waved in order to argue that both entropies behave in a similar way at low temperature and vanish together at $`T_K`$.
### 6.3 The thermodynamic configurational entropy
As we have already stated the thermodynamic configurational entropy can be defined by the relation
$$\mathrm{\Sigma }_t=S(T)S_{valley}(T)$$
(28)
The main difficulty is the precise definition of the valleys, and of $`S_{valley}(T)`$, in the regime $`T>T_K`$ where the system is still ergodic. The basic idea is to take a generic equilibrium configuration ($`y`$) at temperature $`T`$ and to define $`S_{valley}(T)`$ as the thermodynamic entropy of the system constrained to stay at a distance not too large from the equilibrium configuration $`y`$. If we impose a strong constraint (i.e. $`x`$ too near to $`y`$) the entropy will depend on the constraint, but the constraint cannot be taken vanishingly small because the system is ergodic.
One may be worried that this method contains an unavoidable ambiguity. It turns out that there exists a way to modify this method slightly in order to get rid of this ambiguity. The modified method was introduced in and called the potential method. Let us summarize it here briefly. Given two configurations $`x`$ and $`y`$ we define their overlap as before as $`q(x,y)=1/N_{i,k=1,N}w(x_iy_k),`$ where $`w(x)=1\text{for}x`$ small, $`w(x)=0\text{for}x`$ larger than the typical interatomic distance. Instead of adding a strict constraint we add an extra term to the Hamiltonian: we define
$`\mathrm{exp}(N\beta F(y,ϵ))={\displaystyle 𝑑x\mathrm{exp}(H(x)+\beta ϵNq(x,y))},`$
$`F(ϵ)=F(y,ϵ),`$ (29)
where $`f(y)`$ denotes the average value of $`f`$ over equilibrium configurations $`y`$ thermalized at temperature $`\beta ^1`$.
We introduce the Legendre transform $`W(q)`$ of the free energy $`F(ϵ)`$:
$$W(q)=F(ϵ)+ϵq;q=\frac{F}{ϵ}.$$
(30)
Analytic computation in mean field models , as well as in glass forming liquids using the replicated HNC approximation , show that the behaviour of $`W(q)`$ is qualitatively given by the graphs of fig. 3.
Fig. 4 shows the expectation value of $`q`$ as function of $`ϵ`$ in the corresponding four temperature ranges. The results for the potential $`W(q)`$ in the unstable region where its second derivative is negative and $`q`$ is a decreasing function of $`ϵ`$ are a clear artefact of the mean field approximation, while the results in the metastable region correspond to phenomena that can be observed on time scales shorter than the lifetime of the metastable state.
The thermodynamic configurational entropy is the value of the potential $`W(q)`$ at the secondary minimum with $`q0`$ , and it can be defined only if the minimum do exist (i.e. for $`T<T_c`$). It is evident that the secondary minimum for $`T>T_k`$ is always in the metastable region. However if one would start from a large value of $`ϵ`$ and would decrease $`ϵ`$ to zero not too slowly, the system would not escape from the metastable region and one obtains a proper definition of the thermodynamic configurational entropy in this region $`T>T_K`$. In a similar way one could compute $`q(ϵ)`$ in the region ($`ϵ>ϵ_c`$) where the high $`q`$ phase is thermodynamically stable and extrapolate it to $`ϵ0`$. The ambiguity in the definition of the thermodynamic configurational entropy at temperatures above $`T_k`$ becomes larger and larger when the temperature increases. It cannot be defined for $`T>T_c`$.
### 6.4 Numerical estimates of the configurational entropy
Most attempts at estimating numerically the thermodynamic configurational entropy start from the decomposition (28). The liquid entropy is estimated by a thermodynamic integration of the specific heat from the very dilute (ideal gas) limit. It turns out that in the deeply supercooled region the temperature dependence of the liquid entropy is well fitted by the law predicted in : $`S_{liq}(T)=aT^{2/5}+b`$, which presumably allows for a good extrapolation at temperatures $`T`$ which cannot be simulated. As for the ’valley’ entropy, it can be estimated as that of an harmonic solid. One needs however the vibration frequencies of the solid. These have been approximated by several methods, which are all based on some evaluation of the Instantaneous Normal Modes (INM) in the liquid phase, and the assumption that the spectrum of frequencies does not depend much on temperature below $`T_K`$. Starting from a typical configuration of the liquid, one can look at the INM around it. In general there exist some negative eigenvalues (the liquid is not a local minimum of the energy) which one must take care of. Several methods have been tried: either keep only the positive eigenvalues, or one considers the absolute values of the eigenvalues . Alternatively one can also consider the INM around the nearest inherent structure which has by definition a positive spectrum . The computation of the thermodynamic entropy, using its definition as a system coupled to a reference thermalized configuration, has also been studied in .
The results for the configurational entropy as a function of temperature are shown in fig. 5, for binary mixtures of soft spheres and of Lennard-Jones particles. The agreement with the analytical result obtained from the replicated fluid system is rather satisfactory, considering the various approximations involved both in the analytical estimate and in the numerical ones.
In a recent work, Sciortino Kob and Tartaglia have computed the configurational entropy of inherent structures, $`\mathrm{\Sigma }_{is}(T)`$, defined in (23), in binary Lennard-Jones system. Assuming that the free energy $`T\mathrm{log}Z(a)`$ of an inherent structure $`a`$ ($`Z(a)`$ is defined in (22)) can be approximated by $`E_a+\delta F(T)`$, with a correction $`\delta F`$ which is nearly independent of $`E_a`$, then the logarithm of the probability of finding an inherent structure with a given energy $`E_{IS}`$ is given by $`\beta E_{IS}+\mathrm{\Sigma }_{is}(E_{IS})+c^t`$. One can thus deduce the $`E_{IS}`$ dependence of $`\mathrm{\Sigma }_{IS}`$. Shifting the curves vertically in order to try to superimpose them with the thermodynamic configurational entropy, they have checked that all these curves coincide in the region of small enough energy, confirming thus that these two definitions of the configurational entropy agree at low enough energy or temperature. In fig. 6 we compare their result for the configurational entropy of inherent structures to the one obtained analytically, using the description of the molecular fluid of binary Lennard-Jones particles of . Apart from a small shift in the ground state energy which may have several origins (finite size effects, small uncertainties in the description of the correlation in the molecular fluid), the figures are in rather good agreement.
## 7 Remarks
We believe that we have now a consistent scheme for computing the thermodynamic properties of glasses at equilibrium. What is needed is on the one hand some better approximations of the molecular liquid state, on the other hand some precise numerical results in the glass phase at equilibrium, as well as measurements of the fluctuation dissipation ratio in the out of equilibrium dynamics (which should give the value of $`m`$ ). Another obvious direction is to study, with the present methods, various types of interaction potentials, including some which are characteristic of strong glasses. Eventually, one would like to proceed to a first principle study of the out of equilibrium dynamics.
## 8 Acknowledgments
We wish to thank W. Kob for providing the data discussed in the last section, and for giving us the energy shift of the truncated Lennard-Jones problem, used in the comparison of fig. 6. We wish to thank P. Verrocchio for providing the analytic prediction shown in fig. 6.
## 9 References |
warning/0002/hep-th0002094.html | ar5iv | text | # 1 Introduction
## 1 Introduction
These four lectures aim at providing a summary of –and some guidance through– the existing literature dealing with the so-called pre-big bang (PBB) scenario, a new cosmological model largely based on the new symmetries underlying superstring cosmology. The lectures will be pedagogical in nature and will not presuppose an advanced knowledge either of modern inflationary cosmology or of superstring/M-theory. Elements of both will be included in the lectures in order to make them reasonably self-contained. More exhaustive treatments of pre-big bang cosmology are (or will soon be ) available elsewhere, while a homepage on the PBB scenario is being kept updated on the Web .
The four lectures roughly correspond to the four forthcoming sections and deal, respectively, with:
* BASIC MOTIVATIONS AND IDEAS
* HOW COULD IT HAVE STARTED?
* PHENOMENOLOGICAL CONSEQUENCES
* HOW COULD IT HAVE STOPPED?
In particular, lecture II (Section 3) contains a discussion of the initial conditions, lecture III (Section 4) discusses the phenomenological virtues and shortcomings of the model, while lecture IV (Section 5) deals with the most important open theoretical issues.
## 2 Basic Motivations and Ideas
### 2.1 Why string cosmology?
The first question that comes to one’s mind when thinking about cosmology and string theory is: Why bother? Indeed, even if string/M-theory is the correct theory of nature, only its effective (low-energy) quantum field theory description appears to be relevant to most of the history of our Universe, i.e. since a very short time after the big bang. This is certainly the case for the standard (hot-big-bang) cosmological model, but it is also true for the standard models of inflation, provided we confine our attention to what happened during the last $`70`$ e-fold of inflation and later (i.e. to what happened after our present horizon reached the size of the inflationary Hubble radius). In both instances, one is only confronting situations in which curvatures are very small with respect to the fundamental scale of string theory.
On the other hand, both the hot-big bang model and its inflationary variant suffer from initial condition problems. In the former case, these are just the well-known homogeneity and flatness problems that motivated inflation. In the latter case, although the problems look less severe, it is still a matter of heated discussion whether or not one should naturally expect a quasi-homogeneous inflaton field highly displaced from the minimum of its potential to emerge from the Planck era. In either case, the question of how to get physically appealing initial conditions lies in the realm of Planck-scale quantum gravity.
At present, the only candidate for a consistent synthesis of general relativity (GR) and quantum mechanics (QM) is superstring theory (see for a recent review, as well as for a non-specialized introduction), or, if we prefer, the mysterious M-theory that reduces to various superstring theories in appropriate limits. It thus seems mandatory to ask whether the above questions on initial conditions do –or do not– find an answer within string theory. Although most string theorists would certainly agree with the above statements –this being after all one of the most selling ads for string theory– many of them would still object to tackling these problems now. The “excuse” is that our understanding of string theory, especially at large curvatures, is still largely incomplete. Furthermore, most of the recent progress in non-perturbative string theory has been achieved in the context of “vacua” (i.e. classical solutions to the field equations) that respect a large number of supersymmetries. By definition, a cosmological background (a fortiori one that evolves rapidly in time) breaks (albeit spontaneously) supersymmetry. This is why the Planckian regime of cosmology appears to be intractable for the time being.
There is however a pleasant surprise. About ten years of work on string cosmology have led naturally to considering a scenario –the so-called pre-big bang (PBB) scenario– in which the Universe enjoyed a long perturbative “life” before the big bang. Starting from an almost trivial state (asymptotic past triviality, see Section 3), the Universe would have evolved towards stronger and stronger curvature and coupling, thereby inflating, until it entered the non-perturbative phase that replaces the big bang singularity of more standard cosmological models.
The situation is very much reminiscent of QCD and strong interactions. Perturbative QCD has been very successful in predicting a huge number of observables for short-distance-dominated hard processes. Successes in the non-perturbative, large-distance regime have been meagre, by comparison: we still lack a definitive proof of confinement, of spontaneous chiral symmetry breaking, of explicit $`U(1)_A`$ breaking, etc. Yet, we do believe that QCD is the correct description of hadronic physics down to scales of $`10^{15}\mathrm{cm}`$ or so. This is largely based on the belief that large- and short-distance physics “decouple”, e.g. on the assumption that the soft hadronization process does not affect certain infrared-safe quantities computed at the quark–gluon level. Fortunately, we did not wait until the confinement problem was solved, to take QCD seriously!
A very similar attitude will be defended here in the case of string cosmology, with one amusing twist: large- and short-distance physics get somehow swapped as we go from QCD to gravity/cosmology. Figure 1 (from Ref. ) illustrates this point. The easy regime for gravity is at large distance/small curvatures; the tough one turns out to be the high-curvature regime that replaces here the big bang singularity. Yet, we shall argue that some consequences of string cosmology, those related to scales that were very large with respect to the string scale in the high-curvature regime, should not be affected, other than by a trivial kinematical red-shift, by the details of the pre- to post-big bang transition $`\mathrm{}`$ provided, of course, that such a transition does indeed take place (the counterpart to assuming that confinement does occur in QCD).
The above reasoning does not imply, of course, that one should not address the hard questions now. On the contrary, the easy part of the game will give precious information on what the relevant hard questions are (for cosmology) and on how to formulate them. I have already mentioned an example of what I mean: insisting too much on (extended) SUSY vacua appears to be an unacceptable limitation for the problems at hand. Another example is that of demanding stability of an acceptable string vacuum: we shall see (in Section 4) that inflationary string vacua lead to tachyonic, i.e. to growing rather than to oscillating, modes. Such modes appear to horrify most string theorists; however, they are just what inflationary cosmologists happily use all the time in order to generate large-scale structure (LSS), and what PBB cosmology uses to generate heat and entropy from an initially cold Universe (see Section 5).
A completely different criticism of string cosmology comes from the cosmology end: for someone accustomed to a data-driven “bottom-up” approach, string cosmology is too much “top-down”. There is certainly a point here. I do not believe that a good model of cosmology is likely to emerge from theoretical considerations alone. Input from the data will be essential in the selection among various theoretical alternatives. We shall see explicit examples of what I mean in Section 5. Yet, it appears that a combination of top-down and bottom-up would be highly desirable. If past history can teach us something in this respect, the construction of the standard model of particle physics (and of QCD in particular) is a perfect example of a fruitful interplay of theoretically sound ideas and beautiful experimental results. Cosmology today resembles the particle physics of the sixties: interesting new data keep coming in at a high pace, while compelling theoretical pillars on which to base our understanding of those data are still missing.
As a final remark, let me turn things around and claim that cosmology could be the only hope that we have for testing string theory in the foreseeable future by using the cosmos itself as the largest conceivable accelerator. The cosmological red-shift since the big bang has kindly brought down Planck-scale physics to a macroscopic scale, thus opening for us a window on the very early Universe. As we shall see in Subsection 2.3, even in this respect, standard and PBB inflation are markedly different.
### 2.2 Why/Which inflation?
The reasons why the standard hot-big-bang model is unsatisfactory have been repeatedly discussed in the literature. For details, we refer to two excellent reviews . Let me briefly summarize here the basic origin of those difficulties with the simplest Friedmann–Robertson–Walker (FRW) cosmology.
In the FRW framework the size of the (now observable) Universe was about $`10^2\mathrm{cm}`$ at the start of the classical era, say at $`t\mathrm{a}\mathrm{few}\mathrm{times}t_P`$, where $`t_P10^{43}\mathrm{s}`$ is the so-called Planck time. This is of course a very tiny Universe w.r.t. its present size ($`10^{28}\mathrm{cm}`$), yet it is huge w.r.t. the horizon (the distance travelled by light) at that time, i.e. to $`l_P=ct_P10^{33}\mathrm{cm}`$. In other words, a few Planck times after the big bang, our observable Universe was much too large! It consisted of $`(10^{30})^3=10^{90}`$ Planckian-size, causally disconnected regions. There had not been, since the beginning, enough time for the Universe to become homogeneous (e.g. to thermalize) over its entire size. Also, soon after $`t=t_P`$, the Universe was characterized by a huge hierarchy between its Hubble radius on one side and its spatial-curvature radius on the other. The relative factor of (at least) $`10^{30}`$ appears as an incredible amount of fine-tuning on the initial state of the Universe, corresponding to a huge asymmetry between time and space derivatives. Was this asymmetry really there? And, if so, can it be explained in any, more natural way?
It should be stressed that, while the above unexplained ratio becomes larger and larger as we approach the Planck time (and would go to infinity at $`t=0`$ if we could trust the equations throughout), it represents the ratio of two classical length scales. It so happens that one of the two lengths becomes the (quantum) Planck scale at $`t=t_P`$, but the ratio is still huge at much later times when both scales have nothing to do with (and are much larger than) $`t_P`$. This comment will be very relevant to the discussion of fine-tuning issues given in Subsection 3.6.
It is well known that a generic way to wash out inhomogeneities and spatial curvature consists in introducing, in the history of the Universe, a long period of accelerated expansion, called inflation . This still leaves two alternatives: either the Universe was generic at the big bang and became flat and smooth because of a long post-bangian inflationary phase; or it was already flat and smooth at the big bang as a result of a long pre-bangian inflationary phase.
Assuming, dogmatically, that the Universe (and time itself) started at the big bang, leaves only the first alternative. However, that solution has its own problems, in particular those of fine-tuned initial conditions and inflaton potentials. Besides, it is quite difficult to base standard inflation in the only known candidate theory of quantum gravity, superstring theory. Rather, as we shall argue in a moment, superstring theory gives strong hints in favour of the second (pre-big bang) possibility through two of its very basic properties, the first in relation to its short-distance behaviour, the second from its modifications of GR even at large distance.
### 2.3 Superstring-inspired cosmology
As just mentioned, two classes of properties of string theory are relevant for cosmology. Let us discuss them in turn.
A) Short-distance properties
Since the classical (Nambu–Goto) action of a string is proportional to the area $`A`$ of the surface it sweeps, its quantization must introduce a quantum of length $`\lambda _s`$ through:
$$S/\mathrm{}=A/\lambda _s^2.$$
(1)
This fundamental length, replacing Planck’s constant in quantum string theory , plays the role of a minimal observable length, of an ultraviolet cut-off. Thus, in string theory, physical quantities are expected to be bound by appropriate powers of $`\lambda _s`$, e.g.
$`H^2RG\rho <\lambda _s^2`$
$`k_BT/\mathrm{}<c\lambda _s^1`$
$`R_{comp}>\lambda _s.`$ (2)
In other words, in quantum string theory, relativistic quantum mechanics should solve the singularity problems in much the same way as non-relativistic quantum mechanics solved the singularity problem of the hydrogen atom by keeping the electron and the proton a finite distance apart. By the same token, string theory gives us a rationale for asking daring questions such as: What was there before the big bang? Certainly, in no other present theory can such a question be meaningfully asked.
B) Large-distance properties
Even at large distance (low-energy, small curvatures), superstring theory does not automatically give Einstein’s GR. Rather, it leads to a scalar–tensor theory of the JBD variety. The new scalar particle/field $`\varphi `$, the so-called dilaton, is unavoidable in string theory, and gets reinterpreted as the radius of a new dimension of space in so-called M-theory . By supersymmetry, the dilaton is massless to all orders in perturbation theory, i.e. as long as supersymmetry remains unbroken. This raises the question: Is the dilaton a problem or an opportunity? My answer is that it could be both; and while we can try to avoid its potential dangers, we may try to use some of its properties to our advantage … Let me discuss how.
In string theory, $`\varphi `$ controls the strength of all forces , gravitational and gauge alike. One finds, typically:
$$l_P^2/\lambda _s^2\alpha _{gauge}e^\varphi ,$$
(3)
showing the basic unification of all forces in string theory and the fact that, in our conventions, the weak-coupling region coincides with $`\varphi 1`$. In order not to contradict precision tests of the equivalence principle, and of the constancy of the gauge and gravitational couplings in the “recent” past, we require the dilaton to have a mass (see, however, for an amusing alternative) and to be frozen at the bottom of its own potential today. This does not exclude, however, the possibility of the dilaton having evolved cosmologically (after all, the metric did!) within the weak coupling region where it was practically massless. The amazing (yet simple) observation is that, by so doing, the dilaton may have inflated the Universe!
A simplified argument, which, although not completely accurate, captures the essential physical point, consists in writing the Friedmann equation (for a spatially flat Universe):
$$3H^2=8\pi G\rho ,$$
(4)
and in noticing that a growing dilaton (meaning through (3) a growing $`G`$) can drive the growth of $`H`$ even if the energy density of standard matter decreases in an expanding Universe. This new kind of inflation (characterized by growing $`H`$ and $`\varphi `$) has been termed dilaton-driven inflation (DDI). The basic idea of pre-big bang cosmology is thus illustrated in Fig. 2: the dilaton started at very large negative values (where it was practically massless), ran over a potential hill, and finally reached, sometime in our recent past, its final destination at the bottom of its potential ($`\varphi =\varphi _0`$). Incidentally, as shown in Fig. 2, the dilaton of string theory can easily roll-up —rather than down— potential hills, as a consequence of its non-standard coupling to gravity.
DDI is not just possible. It exists as a class of (lowest-order) cosmological solutions thanks to the duality symmetries of string cosmology , , . Under a prototype example of these symmetries, the so-called scale-factor duality (SFD) , , a FRW cosmology evolving (at lowest order in derivatives) from a singularity in the past is mapped into a DDI cosmology going towards a singularity in the future. Of course, the lowest order approximation breaks down before either singularity is reached. A (stringy) moment away from their respective singularities, these two branches can easily be joined smoothly to give a single non-singular cosmology, at least mathematically. Leaving aside this issue for the moment (see Section 5 for more discussion), let us go back to DDI. Since such a phase is characterized by growing coupling and curvature, it must itself have originated from a regime in which both quantities were very small. We take this as the main lesson/hint to be learned from low-energy string theory by raising it to the level of a new cosmological principle, that of “Asymptotic Past Triviality”, to be discussed in the next Lecture.
### 2.4 Explicit solutions
Many explicit exact PBB-type solutions to the low-energy effective action equations have been constructed and discussed in the literature. For an excellent review, see . Exact solutions can only be obtained in the presence of symmetries (isometries) and, although they are heuristically very important, they are too special from the point of view of an inflationary cosmology, which, as such, should not accept fine-tuned initial conditions. This is why we shall not go into an exhaustive discussion of explicit solutions here. Instead, in Section 3, we will adress the general problem of the evolution of asymptotically trivial initial data.
Here we shall limit our attention to the simplest Bianchi I-type solutions and to their quasi-homogeneous generalizations, after recalling that many more solutions can be obtained from the former by using the non-compact $`O(d,d)`$ symmetry of the low-energy string-cosmology equations when the Kalb–Ramond (KR) field $`B_{\mu \nu }`$ is turned on, or by S-duality transformations (see e.g. ) generating a homogeneous axion field (related to $`B_{\mu \nu }`$ by yet another duality transformation).
The generic homogeneous Bianchi I solution with $`B_{\mu \nu }=0`$ reads, for $`t<0`$,
$`ds^2`$ $`=`$ $`dt^2+{\displaystyle \underset{i}{}}(t)^{2\alpha _i}dx^idx^i,`$
$`\varphi `$ $`=`$ $`(1{\displaystyle \underset{i}{}}\alpha _i)\mathrm{log}(t)`$
$`1`$ $`=`$ $`{\displaystyle \underset{i}{}}\alpha _i^2.`$ (5)
i.e. represents a generalization of the well-known Kasner solutions (see e.g. ) in which one of the two Kasner constraints (the one linear in the $`\alpha _i`$) is replaced by the equation giving the time dependence of $`\varphi `$ ($`\varphi `$ is absent, or constant, for Kasner, hence the second constraint).
Note that, unlike Kasner’s, (5) allows for isotropic solutions ($`\alpha _i=\pm 1/\sqrt{d}`$ for all $`i`$). Also, the quadratic Kasner constraint automatically has $`2^d`$ SFD-related branches, obtained by changing the sign of any subset of the $`\alpha ^{}s`$. Also note that the so-called shifted dilaton defined by:
$$\overline{\varphi }=\varphi \frac{1}{2}\mathrm{log}(\mathrm{det}g_{ij}),$$
(6)
which is invariant under the full $`O(d,d)`$ group, is always given by:
$$\overline{\varphi }=\mathrm{log}(t).$$
(7)
A quasi-homogeneous generalization of (5) was first discussed in (see also ) and reads:
$`ds^2`$ $`=`$ $`dt^2+{\displaystyle \underset{a}{}}e_i^a(x)e_j^a(x)(t)^{2\alpha _a(x)}dx^idx^j,`$
$`\varphi `$ $`=`$ $`(1{\displaystyle \underset{i}{}}\alpha _i(x))\mathrm{log}(t)`$
$`1`$ $`=`$ $`{\displaystyle \underset{i}{}}\alpha _i^2(x),t<0,`$ (8)
where $`x`$ stands for the space coordinates. Equation (8) can be shown to be a generic asymptotic solution of the full PDEs near the $`t=0`$ singularity where spatial gradients become less and less important w.r.t. time derivatives, justifying the validity of the so-called gradient expansion . Note that Eq. (7) is not modified in the quasi-homogeneous solutions. Besides allowing isotropic cosmologies in the homogeneous case, the presence of the dilaton also removes the necessity of a chaotic (BKL-type ) behaviour near the singularity .
### 2.5 Phase diagrams and Penrose-style overview
It is useful to visualize the PBB scenario with the help of some diagrams. Since the actual phase space of the model is multidimensional, each of these diagrams necessarily represents just a cross section of the complete picture.
A very commonly used diagram (Fig. 3) is the flow-diagram in the $`\dot{\overline{\varphi }},H`$ plane (time being just a parameter along the flow lines). Since, at lowest order, $`\ddot{\overline{\varphi }}0`$, the flow is always from left to right near the origin. The four straight lines represent the four (isotropic for simplicity) solutions connected by SFD and time-reversal. The product of the two transformations represents the physically interesting case, since it maps ordinary decelerating FRW cosmology (top left) to dilaton-driven inflation (top right). Clearly, our scenario needs a high-curvature phase during which the left-to-right flow is inverted (as shown by the dotted line joining the two perturbative branches). This can only happen as the result of higher-order corrections (see Section 5).
A second useful diagram (Fig. 4) is the $`e^\varphi ,H`$ plot, i.e. the curvature (energy) coupling plane. The fully perturbative domain (where evolution starts according to the APT postulate) lies, in a log-log plot, to the far left-bottom corner. Sticking again, for simplicity, to the isotropic case, DDI evolution is represented by parallel lines distinguished by different initial values of the dilaton (i.e. of the coupling). It is clear that all these solutions run, eventually, into strong curvature or strong coupling (shown as thick solid lines), which one is hit first being determined by the above-mentioned initial coupling. A discussion of what might happen afterwards is given in Section 5.
As a third possibility, let us use a Carter–Penrose style plot (Fig. 5) to represent, on a finite piece of paper, the entire evolution of the Universe. Unlike in ordinary cosmology, where the CP diagram is truncated by the (space-like) hypersurface of the big-bang singularity, here the whole CP diagram, going from past to future time-like and null infinities, is physically meaningful because of our assumption that finite-string-size effects remove the big bang singularity. This diagram will be discussed and used in the following sections.
Finally, let us represent the basic difference between the standard inflation scenario and that of PBB cosmology by plotting, for each cosmological model, the Hubble horizon ($`H^1`$) and the physical scale that coincides with it today, as functions of cosmic time. This gives rise to two “wine glasses” (Fig. 6), which are very similar in their upper parts (corresponding to recent epochs) but differ markedly at very early times. The most salient difference appears in the early behaviour of the Hubble horizon, an increasing function of time in the standard inflation, a decreasing one in the PBB case. The figure allows me to stress one phenomenological advantage of PBB inflation: Planck- (or string)-scale physics, being no longer washed out by a long, subsequent inflationary phase, becomes accessible to present (or near-future) experiments at the millimetre ($`100`$ GHz) scale. At the same time, larger-scale experiments (such as those on small-angle CMB anisotropies) will test (sub-Planckian-energy) physics during the pre-bangian phase. By contrast, as we have already mentioned in the Introduction, in standard inflation large-scale data probe the Universe as it was seventy e-folds or so before the end of inflation, while shorter scales tells us about more recent epochs. Since we know that, seventy e-folds before the end of inflation, $`H_{infl}`$ was less than $`10^5M_P`$ (or else excessive large scale anisotropies are created, see Section 4), and that such a scale slowly decreases during (slow-roll) inflation, it is clear that, according to standard inflation, physics at energies larger than $`10^5M_P`$ remains unaccessible.
## 3 How could it have started?
### 3.1 Generic asymptotically-trivial past
We have already mentioned that, in standard non-inflationary cosmology, initial conditions have to be fine-tuned to incredible accuracy in the far past (i.e. at $`tt_P10^{43}`$ s). What does this fine-tuning problem look like if we accept hints from scale-factor duality and assume asymptotically trivial, yet generic, initial conditions?
The concept of asymptotic past triviality (APT) is quite similar to that of “asymptotic flatness”, familiar from general relativity . The main differences consist in making only assumptions concerning the asymptotic past (rather than future or space-like infinity) and in the additional presence of the dilaton. It seems physically (and philosophically) satisfactory to identify the beginning with simplicity (see e.g. the entropy-related arguments given in Subsection 5.7). What could be simpler than a trivial, empty and flat Universe? Nothing, of course! The problem is that such a Universe, besides being uninteresting, is also non-generic. By contrast, asymptotically flat/trivial Universes are initially simple, yet generic, in a precise mathematical sense that we shall now discuss.
From the point of view of space-time (taken here, for simplicity, to be $`(3+1)`$-dimensional) the generic solution depends upon four arbitray functions of three coordinates related to the metric, plus two more each for the dilaton and the KR field $`B_{\mu \nu }`$. Amusingly, there is an exact correspondence between this “target-space” counting and a “world-sheet” counting. In the latter, those eight arbitrary functions correspond to eight arbitrary functions of three-momentum entering the most general physical (i.e. on shell) vertex operator describing gravitons, dilatons, and the KR field (which, in four dimensions, is equivalent to a pseudoscalar, the KR axion). We will see in Subsection 3.4 how these arbitrary functions appear in the asymptotic expansion of our fields.
Can a very rich and complicated Universe, like our own, emerge from such extremely simple initial conditions? This would look much like a miracle. However, as I shall argue below, this is precisely what should be expected, owing to well-known classical and quantum gravitational instabilities.
### 3.2 The asymptotic past’s effective action and different (conformal) frames
The APT postulate implies that the early-time evolution of the Universe can be described in terms of the low-energy tree-level action of string theory. Taking a generic closed superstring theory, this reads:
$$\mathrm{\Gamma }_{eff}=\lambda _s^{1d}d^{d+1}x\sqrt{|g|}e^\varphi \left(R+g^{\mu \nu }_\mu \varphi _\nu \varphi \frac{1}{12}(dB)^22\mathrm{\Lambda }\right),$$
(9)
where $`dB`$ is the (three-form) field strength associated with $`B_{\mu \nu }`$.
A further simplification comes from assuming to be dealing with so-called critical superstring theory, the case in which the tree-level (and actually the all-order perturbative) cosmological constant $`\mathrm{\Lambda }`$ vanishes. This requires a total of $`D=10`$ space-time dimensions. If $`D10`$ there will be an effective cosmological constant $`O(\lambda _s^2)`$ preventing any low-curvature solution of the field equations to exist. A similar conclusion is reached if we consider critical, but non-supersymmetric, string theories (see Subsection 3.3).
Equation (9) receives corrections when curvatures become $`O(\lambda _s^2)`$ or when the coupling $`e^\varphi `$ becomes $`O(1)`$. If such corrections are both negligible, it sometimes becomes useful to perform a change of variable by going to the so-called Einstein frame (not to be confused with different frames in GR). This is done by defining:
$$g_{\mu \nu }=g_{\mu \nu }^{(E)}e^{\frac{2}{d1}(\varphi \varphi _0)}.$$
(10)
It is relatively easy to rewrite the action (9) using the Einstein metric. The result is simply:
$$\mathrm{\Gamma }_{eff}^E=l_P^{1d}d^{d+1}x\sqrt{|g^{(E)}|}\left(R\frac{1}{d1}_\mu \varphi ^\mu \varphi \frac{1}{12}e^{\frac{4}{d1}\varphi }(dB)^2\right),$$
(11)
where $`l_P^{d1}=e^{\varphi _0}\lambda _s^{d1}`$ is the present value of the Planck length.
Although the use of the Einstein frame could simplify some calculations, and we shall see examples of this below, it should be kept in mind that the form of the corrections is no longer so simple. For instance, higher-derivative corrections become important when the Einstein-frame curvature is $`O(l_P^2e^{\frac{2}{d1}\varphi }=\lambda _s^2)`$, i.e. reaches a dilaton-dependent critical value. Similarly, having a constant Newton “constant” in this frame is a mere illusion because (even tree-level) string masses do now depend upon $`\varphi `$. For these reasons, although physical results are frame-independent, we shall always describe them with reference to the original string-frame metric in which the stringh length $`\lambda _s`$ is constant.
Let us finally remark that the two frames have been made to coincide today, with the dilaton fixed at its present value $`\varphi _0`$. Similarly, the assumption of APT would also allow the identification of the two frames in the far past, since the dilaton approaches a constant as $`t\mathrm{}`$. However, the two Einstein frames that coincide with the string frame at $`t=\pm \mathrm{}`$ differ from each other by an enormous conformal factor, i.e. by a huge blowing-up of all physical scales.
### 3.3 Classical asymptotic symmetries: the importance of SUSY
The classical equations that follow from varying (9) or (11), besides being generally covariant, are also invariant under a two-parameter group of (global) transformations acting as follows:
$`\varphi `$ $``$ $`\varphi +c,`$
$`g_{\mu \nu }`$ $``$ $`\lambda ^2g_{\mu \nu }.`$ (12)
Indeed (9), (11) are simply rescaled by a constant factor under this group. These two symmetries depend crucially on the validity of the tree-level low-energy approximation and on the absence of a cosmological constant. Loop corrections clearly spoil invariance under dilaton shifts, while lower derivatives (a cosmological constant) or higher derivatives ($`\alpha ^{}`$) corrections spoil invariance under a rescaling of the metric. Note that, using general covariance, the latter symmetry is equivalent to an overall rescaling of all the coordinates. The relevance of the two classical symmetries on the issue of fine-tuning will become obvious in the next two subsections.
The importance of dealing with critical superstring theory now becomes evident: if one would consider non-supersymmetric string theories, a cosmological constant would almost certainly be generated at some finite order of the loop expansion: this would change completely the large-distance properties and spoil the symmetries of the field equations.
### 3.4 Dilaton-driven inflation as gravitational collapse
For simplicity, we will only illustrate here the simplest case of gravi-dilaton system already compactified to four space-time dimensions. Through the field redefinition (10), our problem is reduced to the study of a massless scalar field minimally coupled to gravity. It is well known that such a form of matter cannot give inflation (since it has positive pressure). Instead, it can easily lead to gravitational collapse (GC). Thus, in the Einstein frame, the problem becomes that of finding out under which conditions gravitational collapse occurs if asymptotically-trivial initial data are assigned. Gravitational collapse usually means that the (Einstein) metric (hence the volume of 3-space) shrinks to zero at a space-like singularity. However, typically, the dilaton blows up at that same singularity. Given the relation (10) between the Einstein and the (physical) string metric, we can easily imagine that the latter blows up near the singularity, as implied by DDI.
How generically does GC happen? Let us recall the singularity theorems of Hawking and Penrose , which state that, under some general assumptions, singularities are inescapable in GR. Looking at the validity of those assumptions in the case at hand, one finds that all but one are automatically satisfied. The only condition to be imposed is the existence of a closed trapped surface (CTS) (a closed surface from which future light cones lie entirely in the region inside the surface). Rigorous results show that this condition cannot be waived: sufficiently weak initial data do not lead to closed trapped surfaces, to collapse, or to singularities. Sufficiently strong initial data do. But where is the border-line? This is not known in general, but precise criteria do exist for particularly symmetric space-times, e.g. for those endowed with spherical symmetry (see Subsection 3.6).
However, no matter what the general collapse/singularity criterion will eventually turn out to be, we do know, from the classical symmetries described in the previous subsection, that such a criterion cannot depend
* on an over-all additive constant in $`\varphi `$, or
* on an over-all multiplicative factor in $`g_{\mu \nu }`$.
A characterization of APT initial data can be made following the pioneering work of Bondi, Sachs, Penrose, and others. Since our initial quanta are assumed to consist of massless gravitons and dilatons, their past infinity is null: it is the famous $`^{}`$ of the Penrose diagram (Fig. 5). APT means that dilaton and metric can be expanded near $`^{}`$ in inverse powers of $`r\mathrm{}`$, while advanced time $`v`$ and two angular variables, $`\theta `$ and $`\phi `$, are kept fixed. We shall thus write:
$$\varphi (x^\lambda )=\varphi _0+\frac{f(v,\theta ,\phi )}{r}+o\left(\frac{1}{r}\right),$$
(13)
$$g_{\mu \nu }(x^\lambda )=\eta _{\mu \nu }+\frac{f_{\mu \nu }(v,\theta ,\phi )}{r}+o\left(\frac{1}{r}\right).$$
(14)
The null wave data on $`^{}`$ are: the asymptotic dilatonic wave form $`f(v,\theta ,\phi )`$, and two polarization components, $`f_+(v,\theta ,\phi )`$ and $`f_\times (v,\theta ,\phi )`$, of the asymptotic gravitational wave form $`f_{\mu \nu }(v,\theta ,\phi )`$, whose other components can be gauged away. The three functions $`f`$, $`f_+`$, $`f_\times `$ of $`v,\theta ,\phi `$ are equivalent to six functions of $`r,\theta ,\phi `$ with $`r0`$, because the advanced time $`v`$ ranges over the full line $`(\mathrm{},+\mathrm{})`$. This is how the six arbitrary functions of the generic solution to the gravi-dilaton system are recovered.
Of particular interest here are the so-called News functions, simply given by
$$N(v,\theta ,\phi )_vf(v,\theta ,\phi ),N_+_vf_+,N_\times _vf_\times ,$$
(15)
and the “Bondi mass” given by:
$`M_{}(v)`$ $`=`$ $`{\displaystyle \frac{1}{4\pi }}{\displaystyle d^2\mathrm{\Omega }M_{}(v,\theta ,\phi )},`$
$`g_{vv}`$ $`=`$ $`\left(1{\displaystyle \frac{2M_{}(v,\theta ,\phi )}{r}}\right)+o\left({\displaystyle \frac{1}{r}}\right).`$ (16)
The Bondi mass and the News are connected by the energy–momentum conservation equation, which tells us that the advanced-time derivative of $`M_{}(v)`$ is positive-semidefinite and related to incoming energy fluxes controlled by the News:
$$dM_{}(v)/dv=\frac{1}{4}d^2\mathrm{\Omega }\left(N^2+N_+^2+N_\times ^2\right).$$
(17)
The physical meaning of $`M_{}(v)`$ is that it represents the energy brought into the system (by massless sources) by advanced time $`v`$. In the same spirit one can define the Bondi mass $`M_+(u)`$ at future null infinity $`^+`$. It represents the energy still present in the system at retarded time $`u`$. If only massless sources are present, the so-called ADM mass is given by
$$M_{}(+\mathrm{})=M_+(\mathrm{})=M_{ADM},$$
(18)
while $`M_{}(\mathrm{})=0`$, and $`M_+(+\mathrm{})=M_C`$ represents the mass that has not been radiated away even after waiting an infinite time, i.e. the mass that underwent gravitational collapse . Collapse (resp. no-collapse) criteria thus aim at establishing under which initial conditions one expects to find $`M_C>0`$ (resp. $`M_C=0`$).
Since, as we shall see in the particular case of spherical symmetry, collapse criteria i) do not involve any particularly large number, and ii) do not contain any intrinsic scale but just dimesionless ratios of various classical scales, we expect i) gravitational collapse to be quite a generic phenomenon and ii) that nothing, at the level of our approximations, will be able to fix either the size of the horizon or the value of $`\varphi `$ at the onset of collapse. Generically, and quite randomly and chaotically, some regions of space will undergo gravitational collapse, will form horizons and singularities therein. When this is translated into the string frame, the region of space-time within the horizon undergoes a period of DDI in which both the initial value of the Hubble parameter and that of $`\varphi `$ are left arbitrary. In the next subsection we shall see that such arbitrariness provides an answer to the fine-tuning allegations that have been recently moved to the PBB scenario. This section will be concluded with a discussion of how more precisely the case of spherical symmetry can be dealt with.
### 3.5 Fine-tuning issues
The two arbitrary parameters discussed in the previous subsection are very important, since they determine the range of validity of our description. In fact, since both curvature and coupling increase during DDI, the low-energy and/or tree-level description is bound to break down at some point. The smaller the initial Hubble parameter (i.e. the larger the initial horizon size) and the smaller the initial coupling, the longer we can follow DDI through the effective action equations and the larger the number of reliable e-folds we shall gain.
This does answer, in my opinion, the objections raised recently to the PBB scenario according to which it is fine-tuned. The situation here actually resembles that of chaotic inflation . Given some generic (though APT) initial data, we should ask which is the distribution of sizes of the collapsing regions and of couplings therein. Then, only the “tails” of these distributions, i.e. those corresponding to sufficiently large, and sufficiently weakly coupled, regions will produce Universes like ours, the rest will not. The question of how likely a “good” big bang is to take place is not very well posed and can be greatly affected by anthropic considerations .
In conclusion, we may summarize recent progress on the problem of initial conditions by saying that :
Dilaton-Driven Inflation in String Cosmology
is as generic as
Gravitational Collapse in General Relativity.
Furthermore, asking for a sufficiently long period of DDI amounts to setting upper limits on two arbitrary moduli of the classical solutions.
Figure 7 (from Ref. ) gives a $`(2+1)`$-dimensional sketch of a possible PBB Universe: an original “sea” of dilatonic and gravity waves leads to collapsing regions of different initial size, possibly to a scale-invariant distribution of them. Each one of these collapses is reinterpreted, in the string frame, as the process by which a baby Universe is born after a period of PBB inflationary “pregnancy”, the size of each baby Universe being determined by the duration of the corresponding pregnancy, i.e. by the initial size of (and coupling in) the corresponding collapsing region. Regions initially larger than $`10^{13}\mathrm{cm}`$ can generate Universes like ours, smaller ones cannot.
A basic difference between the large numbers needed in (non- inflationary) FRW cosmology and the large numbers needed in PBB cosmology should be stressed. In the former, the ratio of two classical scales, e.g. of total curvature to its spatial component, which is expected to be $`O(1)`$, has to be taken as large as $`10^{60}`$. In the latter, the above ratio is initially $`O(1)`$ in the collapsing/inflating region, and ends up being very large in that same region, thanks to DDI. However, the (common) order of magnitude of these two classical quantities is a free parameter, and it is taken to be much larger than the classically irrelevant quantum scale. Indeed, the smallness of quantum corrections (which would introduce a scale in the problem) was explicitly checked in .
We can visualize analogies and differences between standard and pre-big bang inflation by looking again at Figs. 6a and 6b. The common feature in the two pictures is that the fixed comoving scale corresponding to the present horizon was “inside the horizon” for some time during inflation, possibly very deeply inside at its onset. The difference between the two scenarios is just in the behaviour of the Hubble radius during inflation: increasing in standard inflation (a), decreasing in string cosmology (b). Thus, while standard inflation is still facing the initial-singularity question and needs a non-adiabatic phenomenon to reheat the Universe (a kind of small bang), PBB cosmology faces the singularity problem later, combining it with the exit and heating problems (see Section 5).
### 3.6 The spherically symmetric case
In the spherically symmetric case many authors have studied the problem of gravitational collapse of a minimally coupled scalar field both numerically and analytically. In the former case I will only mention the well-known results of Choptuick , pointing at mysterious universalities near critical collapse (i.e. at the border-line situation in which the collapse criteria are just barely met). In this case, a very small black hole forms. This is not the case we are really interested in for the reasons we just explained. We shall thus turn, instead, to what happens when the collapse criteria are largely fulfilled. For this we make use of the rather powerful results due to Christodoulou over a decade of beautiful work , , , .
There are no gravitational waves in the spherically symmetric case so that null wave data consist of just an angle-independent asymptotic dilatonic wave form $`f(v)`$, with the associated scalar News $`N(v)=f^{}(v)`$.
A convenient system of coordinates is the double null system, $`(u,v)`$, such that
$$\varphi =\varphi (u,v),$$
(19)
$$ds^2=\mathrm{\Omega }^2(u,v)dudv+r^2(u,v)d\omega ^2,$$
(20)
where $`d\omega ^2=d\theta ^2+\mathrm{sin}^2\theta d\phi ^2`$. The field equations are conveniently re-expressed in terms of the three functions $`\varphi (u,v)`$, $`r(u,v)`$ and $`m(u,v)`$, where the local mass function $`m(u,v)`$ is defined by:
$$1\frac{2m}{r}g^{\mu \nu }(_\mu r)(_\nu r)=\frac{4}{\mathrm{\Omega }^2}\frac{r}{u}\frac{r}{v}.$$
(21)
One gets the following set of evolution equations for $`m`$, $`r`$ and $`\varphi `$
$`2{\displaystyle \frac{r}{u}}{\displaystyle \frac{m}{u}}=\left(1{\displaystyle \frac{2m}{r}}\right){\displaystyle \frac{r^2}{4}}\left({\displaystyle \frac{\varphi }{u}}\right)^2,`$ (22)
$`2{\displaystyle \frac{r}{v}}{\displaystyle \frac{m}{v}}=\left(1{\displaystyle \frac{2m}{r}}\right){\displaystyle \frac{r^2}{4}}\left({\displaystyle \frac{\varphi }{v}}\right)^2,`$ (23)
$`r{\displaystyle \frac{^2r}{uv}}={\displaystyle \frac{2m}{r2m}}{\displaystyle \frac{r}{u}}{\displaystyle \frac{r}{v}},`$ (24)
$`r{\displaystyle \frac{^2\varphi }{uv}}+{\displaystyle \frac{r}{u}}{\displaystyle \frac{\varphi }{v}}+{\displaystyle \frac{r}{v}}{\displaystyle \frac{\varphi }{u}}=0.`$ (25)
The quantity
$$\mu (u,v)\frac{2m(u,v)}{r}$$
(26)
plays a crucial rôle in the problem. If $`\mu `$ stays everywhere below 1, the field configuration will not collapse but will finally disperse at infinity as outgoing waves. By contrast, if the mass ratio $`\mu `$ can reach anywhere the value 1, this signals the formation of an apparent horizon $`𝒜`$. The location of this apparent horizon is indeed defined by the equation
$$𝒜:\mu (u,v)=1.$$
(27)
The above statements are substantiated by some rigorous inequalities stating that:
$`{\displaystyle \frac{r}{u}}<0,{\displaystyle \frac{m}{v}}>0,`$ (28)
$`{\displaystyle \frac{r}{v}}\left(1\mu \right)>0,{\displaystyle \frac{m}{u}}\left(1\mu \right)<0.`$ (29)
Thus, in weak-field regions ($`\mu <1`$), $`_vr>0`$, while, as $`\mu >1`$, $`_vr<0`$, meaning that the outgoing radial null rays (“photons”) emitted by the sphere $`r=\mathrm{const}`$ become convergent, instead of having their usual behaviour. This is nothing else but the signature of a CTS!
In the case of spherical symmetry, it has been possible to prove that the presence of trapped surfaces implies the existence of a future singular boundary $``$ of space-time where a curvature singularity occurs. Furthermore, the behaviour of various fields near the singularity is just that of a quasi-homogeneous DDI as described by Eqs. (8)! This highly non-trivial result strongly supports the idea that PBB inflation in the string frame is the counterpart of gravitational collapse in the Einstein frame.
Reference gives the following sufficient criterion on the strength of characteristic data, considered at some finite retarded time $`u`$
$$\frac{2\mathrm{\Delta }m}{\mathrm{\Delta }r}\left[\frac{r_1}{r_2}\mathrm{log}\left(\frac{r_1}{2\mathrm{\Delta }r}\right)+\frac{6r_1}{r_2}1\right],$$
(30)
where $`r_1r_2`$, with $`r_23r_1/2`$, are two spheres, $`\mathrm{\Delta }r=r_2r_1`$ is the width of the “annular” region between the two spheres, and $`\mathrm{\Delta }m=m_2m_1m(u,r_2)m(u,r_1)`$ is the mass “contained” between the two spheres, i.e. more precisely the energy flux through the outgoing null cone $`u=`$ const, between $`r_1`$ and $`r_2`$. Note the absence of any intrinsic scale (in particular of any short-distance cut-off) in the above criterion. The theorem proved in is not exhausted in the above statement. It contains various bounds as well, e.g.
* an upper bound on the retarded time at which the CTS (i.e. a horizon) is formed,
* a lower bound on the mass, i.e. on the radius of the collapsing region.
The latter quantity is very important for the discussion of the previous subsection since it gives, in the equivalent string-frame problem, an upper limit on the Hubble parameter at the beginning of DDI. Such an upper limit depends only on the size of the advanced-time interval satisfying the CC; since the latter is determined by the scale-invariant condition (30), the initial scale of inflation will be classically undetermined.
The above criterion is rigorous but probably too conservative. It also has the shortcoming that it cannot be used directly on $`^{}`$, since $`u\mathrm{}`$ on $`^{}`$. In Ref. a less rigorous (or less general) but simpler criterion directly expressible in terms of the News (i.e. on $`^{}`$) was proposed on the basis of a perturbative study. It has the following attractive form:
$$\underset{\genfrac{}{}{0pt}{}{v_1,v_2}{v_1v_2}}{sup}\mathrm{Var}(N(x))_{x[v_1,v_2]}>C=O(1/4),$$
(31)
where:
$$\mathrm{Var}(N(x))_{x[v_1,v_2]}(N(x)N_{[v_1,v_2]})^2_{x[v_1,v_2]}.$$
(32)
Thus $`\mathrm{Var}(g)_{[v_1,v_2]}`$ denotes the “variance” of the function $`g(x)`$ over the interval $`[v_1,v_2]`$, i.e. the average squared deviation from its mean value.
According to this criterion the largest interval satisfying (31) determines the size of the collapsing region and thus, through the collapse inflation connection, the initial value of the Hubble parameter. It would be interesting to confirm the validity of the above criterion and to determine more precisely the value of the constant appearing on its r.h.s. through more analytic or numerical work. Actually, numerical studies of spherically symmetric PBB cosmologies have already appeared , while more powerful numerical codes should soon be available .
## 4 Phenomenological Consequences
### 4.1 Cosmological amplification of vacuum fluctuations: general properties
I will start by recalling the basic physical mechanism underlying particle production in cosmology (for a nice review, see ) and by introducing the corresponding (and by now standard) jargon. By the very definition of inflation ($`\ddot{a}>0`$) physical wavelengths are stretched past the Hubble scale ($`H^1`$) during inflation. After the end of inflation each wavelength grows slower than $`H^1`$ and thus “re-enters” the horizon. Obviously, the larger the scale the earlier it crosses the horizon outward and the later it crosses it back inward. Hence larger scales “spend” more time “outside the horizon” than smaller ones.
The attentive reader may worry at this point about the way this description applies when distances are measured using the Einsten-frame metric. As we have seen in the previous section, PBB inflation corresponds to accelerated contraction in the Einstein frame. Nonetheless, one can show that physical quantities (that is, typically, dimensionless ratios of physical quantities) do not depend on the choice of the frame: after all, changing frame is nothing more than a local field-redefinition, which is known not to affect the physics. It is amusing to notice, for instance, that physical wavelengths go outside the horizon during the Einstein-frame equivalent of DDI. Indeed, although physical EF scales shrink during the collapse, the horizon $`H^1`$ shrinks even faster! I refer to the first paper in for further discussion on this point.
Consider now a generic perturbation $`\mathrm{\Psi }`$ on top of a homogeneous background, which includes a cosmological-type metric, a dilaton, and, possibly, other fields, such as another inflaton field, an axion, etc. Since $`\mathrm{\Psi }=0`$ is, by definition of a perturbation, a classical solution, $`\mathrm{\Psi }`$ intself enters the effective low-energy action quadratically. Soon after the beginning of inflation the background itself becomes homogeneous, isotropic, and spatially flat, so that the perturbed action takes the generic form:
$$I=\frac{1}{2}𝑑\eta d^3xS(\eta )\left[\mathrm{\Psi }^2(\mathrm{\Psi })^2\right].$$
(33)
Here $`\eta `$ is conformal-time ($`ad\eta =dt`$), and a prime denotes $`/\eta `$. The function $`S(\eta )`$ (sometimes called the “pump” field) is, for any given $`\mathrm{\Psi }`$, a given function of the scale factor $`a(\eta )`$, and of other scalar fields (four-dimensional dilaton $`\varphi (\eta )`$, moduli $`b_i(\eta )`$, etc.), which may appear non-trivially in the background.
While it is clear that a constant $`S`$ may be reabsorbed by rescaling $`\mathrm{\Psi }`$, and is thus ineffective, a time-dependent $`S`$ couples non-trivially to $`\mathrm{\Psi }`$ and leads to the production of pairs of quanta (with equal and opposite momenta). In order to see this, it is useful to go over to a Hamiltonian description of the perturbation and of its canonically conjugate momentum $`\mathrm{\Pi }`$:
$$\mathrm{\Pi }=\frac{\delta I}{\delta \mathrm{\Psi }^{^{}}}=S\mathrm{\Psi }^{}.$$
(34)
The Hamiltonian corresponding to (33) is thus given by
$$H=\frac{1}{2}d^3x\left[S^1\mathrm{\Pi }^2+S(\mathrm{\Psi })^2\right],$$
(35)
and the first-order Hamilton equations read
$$\mathrm{\Psi }^{}=\frac{\delta H}{\delta \mathrm{\Pi }}=S^1\mathrm{\Pi },\mathrm{\Pi }^{}=\frac{\delta H}{\delta \mathrm{\Psi }}=S^2\mathrm{\Psi },$$
(36)
leading to the decoupled second order equations
$$\mathrm{\Psi }^{\prime \prime }+\frac{S^{}}{S}\mathrm{\Psi }^{}^2\mathrm{\Psi }=0,\mathrm{\Pi }^{\prime \prime }\frac{S^{}}{S}\mathrm{\Pi }^{}^2\mathrm{\Pi }=0.$$
(37)
In Fourier space the Hamiltonian (35) is given by
$$H=\frac{1}{2}\underset{\stackrel{}{k}}{}\left(S^1\mathrm{\Pi }_\stackrel{}{k}\mathrm{\Pi }_\stackrel{}{k}+Sk^2\mathrm{\Psi }_\stackrel{}{k}\mathrm{\Psi }_\stackrel{}{k}\right),$$
(38)
where $`\mathrm{\Psi }_\stackrel{}{k}=\mathrm{\Psi }_\stackrel{}{k}^{}`$ and $`\mathrm{\Pi }_\stackrel{}{k}=\mathrm{\Pi }_\stackrel{}{k}^{}`$. The equations of motion become
$$\mathrm{\Psi }_\stackrel{}{k}^{}=S^1\mathrm{\Pi }_\stackrel{}{k},\mathrm{\Pi }_\stackrel{}{k}^{}=Sk^2\mathrm{\Psi }_\stackrel{}{k},$$
(39)
where $`k=|\stackrel{}{k}|`$. The transformation
$$\mathrm{\Pi }_\stackrel{}{k}\stackrel{~}{\mathrm{\Pi }}_\stackrel{}{k}=k\mathrm{\Psi }_\stackrel{}{k},\mathrm{\Psi }_\stackrel{}{k}\stackrel{~}{\mathrm{\Psi }}_\stackrel{}{k}=k^1\mathrm{\Pi }_\stackrel{}{k},S\stackrel{~}{S}=S^1$$
(40)
leaves the Hamiltonian, Poisson brackets, and equations of motion unchanged. This symmetry of linear perturbation theory, and its physical consequences, was discussed in under the name of S-duality, since it contains the usual strong–weak coupling (electric–magnetic) duality in the special case of gauge perturbations.
In order to solve the perturbation equations, and to normalize the spectrum, it is convenient to introduce the normalized (but no longer canonically conjugate) variables $`\widehat{\mathrm{\Psi }}`$, $`\widehat{\mathrm{\Pi }}`$, whose Fourier modes are defined by
$$\widehat{\mathrm{\Psi }}_k=S^{1/2}\mathrm{\Psi }_k,\widehat{\mathrm{\Pi }}_k=S^{1/2}\mathrm{\Pi }_k,$$
(41)
so that the Hamiltonian density takes the canonical form:
$$H=\frac{1}{2}\underset{\stackrel{}{k}}{}\left(|\widehat{\mathrm{\Pi }}_k|^2+k^2|\widehat{\mathrm{\Psi }}_k|^2\right).$$
(42)
Under S-duality, these new variables transform as the original ones. They satisfy the Schrödinger-like equations
$$\widehat{\mathrm{\Psi }}_k{}_{}{}^{\prime \prime }+[k^2(S^{1/2}){}_{}{}^{\prime \prime }S_{}^{1/2}]\widehat{\mathrm{\Psi }}_k=0,\widehat{\mathrm{\Pi }}_k{}_{}{}^{\prime \prime }+[k^2(S^{1/2}){}_{}{}^{\prime \prime }S_{}^{1/2}]\widehat{\mathrm{\Pi }}_k=0.$$
(43)
The amplification of perturbations is typically associated with a transition from an inflationary phase in which the pump field is accelerated to a post-inflationary phase in which the pump field is decelerated or constant. In such a class of backgrounds, the “effective potentials”, $`V_\mathrm{\Psi }=(S^{1/2}){}_{}{}^{\prime \prime }S_{}^{1/2}`$ and $`V_\mathrm{\Pi }=(S^{1/2}){}_{}{}^{\prime \prime }S_{}^{1/2}`$, grow during the phase of accelerated evolution, and decrease in the post-inflationary, decelerated epoch, vanishing asymptotically both for very early times, $`\eta \mathrm{}`$, and for very late times, $`\eta +\mathrm{}`$.
The initial evolution of perturbations, for all modes with $`k^2>|V_\mathrm{\Psi }|`$, $`|V_\mathrm{\Pi }|`$, may be described by the WKB-like approximate solutions of Eqs. (43)
$`\widehat{\mathrm{\Psi }}_k(\eta )`$ $`=`$ $`\left(k^2V_\mathrm{\Psi }\right)^{1/4}e^{i\underset{\eta _0}{\overset{\eta }{}}𝑑\eta ^{}\left(k^2V_\mathrm{\Psi }\right)^{1/2}},`$
$`\widehat{\mathrm{\Pi }}_k(\eta )`$ $`=`$ $`k\left(k^2V_\mathrm{\Pi }\right)^{1/4}e^{i\underset{\eta _0}{\overset{\eta }{}}𝑑\eta ^{}\left(k^2V_\mathrm{\Pi }\right)^{1/2}},`$ (44)
which we have normalized to a vacuum fluctuation, and where the extra factor of $`k`$ in the solution for $`\widehat{\mathrm{\Pi }}_k`$ comes from consistency with the first order equations (39). We have ignored a possible relative phase in the solutions. Solutions (44) manifestly preserve the S-duality symmetry of the equations, since the potentials $`V_\mathrm{\Psi }`$, $`V_\mathrm{\Pi }`$ get interchanged under $`SS^1`$.
Let us now discuss two opposite regimes:
* When the perturbation is deeply inside the horizon ($`k/aH`$) we find “adiabatic” behaviour, i.e.
$$k\mathrm{\Phi }_kS^{1/2},\mathrm{\Pi }_kS^{1/2},$$
(45)
implying, through (35), that the contribution to the Hamiltonian of modes inside the horizon stays constant.
* When the perturbation is far outside the horizon ($`k/aH`$), it enters the so-called freeze-out regime in which $`\mathrm{\Psi }`$ and $`\mathrm{\Pi }`$ stay constant (better have a constant solution, see ). Such a behaviour implies, again through (35), that the contribution of super-horizon modes to the Hamiltonian grows in time. If $`\dot{S}>0`$, the growth of $``$ is due to $`\mathrm{\Psi }`$, while, for $`\dot{S}<0`$, it is due to $`\mathrm{\Pi }`$. In either case the growth is due to particle production in squeezed states , i.e. states in which one canonical variable is very sharply defined and the conjugate one is largely undetermined. Although, strictly speaking, quantum coherence is not lost, in practice the sub-fluctuating variable cannot be measured with unlimited precision (coarse graining) and therefore entropy is produced (see Subsection 4.6).
It is not too hard to join the two extreme regimes mentioned above and to find the qualitative and quantitative features of the solutions. For lack of space we refer the reader to the original literature (see, e.g. ).
The above considerations were very general. What is instead typical of the PBB scenario? There are at least two features that are quite unique to string cosmology:
* Pump fields, and in particular their contributions to the evolution equations (43), grow during PBB inflation, while they tend to decay in standard inflation.
* The richer set of backgrounds and fluctuation present in string theory allows for the amplification of new kinds of perturbations.
One can easily determine the pump fields for each one of the most interesting perturbations appearing in the PBB scenario. The result is:
$`\mathrm{Gravity}\mathrm{waves},\mathrm{dilaton}`$ $`:`$ $`S=a^2e^\varphi `$
$`\mathrm{Heterotic}\mathrm{gauge}\mathrm{bosons}`$ $`:`$ $`S=e^\varphi `$
$`\mathrm{Kalb}\mathrm{Ramond},\mathrm{axions}`$ $`:`$ $`S=a^2e^\varphi .`$ (46)
In the following subsections we shall briefly describe the characteristics of these four perturbations after their original vacuum fluctuations are amplified by PBB inflation. For further details, see also .
### 4.2 Tensor perturbations: an observable cosmic gravitational radiation background (CGRB)?
It is not surpising to find that, for tensor and dilaton perturbations, the pump field is nothing but the scale factor in the Einstein frame ($`a_E=ae^{\varphi /2}`$) since, in this frame, the action for gravity and for the dilaton take the canonical form. The Einstein-frame scale factor corresponds to a collapsing Universe (see Section 3), hence to the decreasing pump field $`a_E(\eta )\eta ^{1/2}`$ during DDI. For scales that go outside the horizon during DDI, this implies a Raileigh–Jeans-like spectrum, $`d\mathrm{\Omega }/d\mathrm{log}kk^3`$, up to logarithmic corrections .
When the curvature scale reaches the string scale we expect DDI to end, and a high (string scale) curvature phase to follow, before the eventual exit to the FRW phase takes place (see Section 5). Not much is known about the string phase, but, using some physical arguments as well as some quantitative estimates, it can be argued that such a phase will lead to copious GW production at frequencies corresponding to the string scale at the time of exit. After the transition to the FRW phase, all particle production switches off. This is why our GW spectrum has an end point that corresponds to the string/Planck scale at the beginning of the FRW phase. If no inflation takes place after, the end-point frequency corresponds, today, to $`\omega =\omega _1100`$ GHz.
As illustrated in Fig. 8, the GW spectrum can be rather flat below the end point, up to the frequency $`\omega _s`$, the last scale that went out of the horizon during DDI. Further below $`\omega _s`$ we get the above-mentioned steep $`\omega ^3`$ spectrum. It thus looks as if the best chances for the detection of our stochastic background lie precisely near $`\omega _s`$, where a kink (or knee) is expected.
Unfortunately, the position of the knee and the value of $`\mathrm{\Omega }_{GW}`$ at that point depend on two background parameters that are, so far, difficult to predict. One corresponds to the duration of (better, the total red-shift during) the string phase, the other to the value of $`l_P/\lambda _s`$ at the end of DDI (hence to the value of the dilaton at that time). As shown in Fig. 8, values of $`\mathrm{\Omega }_{GW}`$ in the range of $`10^6`$$`10^7`$ are possible in some regions of parameter space, which, according to some estimates of sensitivities reported in the same figure for $`\omega _s10^2`$Hz, could be inside detection capabilities in the near future.
The signal is predicted to consist of randomly distributed standing waves, a feature that has been argued to further help detection. In any case, cross-correlation experiments are mandatory here in order to disentangle this stochastic signal from real noise. Sensitivities to a CGRB of this type have been estimated for a variety of two-detector combinations . A comprehensive review of GW experiments and of their relevance to the early Universe can be found in .
### 4.3 Dilaton perturbations
Since the dilaton is, after all, the inflaton of PBB cosmology, its fluctuations are the most natural source of adiabatic scalar perturbations. We recall that, in standard cosmology, inflaton fluctuations naturally lead to a quasi scale-invariant, Harrison-Zeldovich (HZ) spectrum of adiabatic perturbations, something highly desirable both to explain CMB anisotropy and for models of LSS formation. Can we get something similar from the dilaton? The answer, unfortunately, is no! Let me spend a moment explaining why.
Unlike tensor perturbations, which do not couple to the scalar field to linear order and are gauge-invariant by themselves, scalar perturbations are contained, a priori, in five functions defined by:
$`ds^2`$ $`=`$ $`a^2(\eta )\left[(1+2\mathrm{\Phi })d\eta ^2+\left((12\mathrm{\Psi })\delta _{ij}+_i_jE\right)dx^idx^j2_iBdx^id\eta \right]`$
$`\varphi `$ $`=`$ $`\varphi _0(\eta )+\chi (\eta ,\stackrel{}{x}).`$ (47)
The five functions $`\mathrm{\Phi },\mathrm{\Psi },B,E,\chi `$ are not separately gauge-invariant. However, the following “Bardeen” combinations are gauge-invariant:
$`\mathrm{\Phi }_B`$ $`=`$ $`\mathrm{\Phi }+{\displaystyle \frac{1}{a}}\left[a(BE^{})\right]^{^{}},`$
$`\mathrm{\Psi }_B`$ $`=`$ $`\mathrm{\Psi }{\displaystyle \frac{a^{}}{a}}(BE^{}),`$
$`\chi _{GI}`$ $`=`$ $`\chi +{\displaystyle \frac{\varphi _0^{}a}{a^{}}}\mathrm{\Psi }.`$ (48)
Introducing the variable $`va\chi _{GI}`$, the scalar field enters the quadratic action “canonically” i.e.:
$$S_{eff}(v)=\frac{1}{2}𝑑\eta d^3x\left[v^2(\stackrel{}{}v)^2+(z^{^{\prime \prime }}/z)v^2\right],z\frac{\varphi _0^{}a^2}{a^{}},$$
(49)
giving the evolution equation
$$v_k^{^{\prime \prime }}+\left(k^2(z^{^{\prime \prime }}/z)\right)v_k=0.$$
(50)
In the DDI background, $`za`$ and thus the canonical scalar field obeys the same equation as the canonical graviton field, therefore giving identical spectra (as far as the dilaton remains massless, of course). This strongly suggests that adiabatic perturbations in PBB cosmology have a Raleigh–Jeans, rather than HZ, spectrum and that they are unsuitable for generating CMBA or LSS. Before being sure of that, however, we have to analyse the scalar fluctuations of the metric itself in terms of the above-mentioned Bardeen potentials $`\mathrm{\Phi }_B,\mathrm{\Psi }_B`$.
A popular gauge (particularly advertised in ) is the so-called longitudinal gauge, defined by $`B=E=0`$, where $`\mathrm{\Phi }_B=\mathrm{\Phi }`$ and $`\mathrm{\Psi }_B=\mathrm{\Psi }`$. In this gauge one of the constraints simply reads $`\mathrm{\Phi }=\mathrm{\Psi }`$, while a second constraint relates either one of them to $`v`$:
$$\mathrm{\Psi }_k=\frac{\varphi _0^{}}{4k^2}(v_k/a)^{}k^{3/2}\frac{1}{|k\eta |^2}(k/a)_{HC},$$
(51)
where we have inserted the small $`k`$ behaviour of $`v_k`$, which is identical to that of tensor perturbations.
Unfortunately, Eq. (51) leads to very large fluctuations of $`\mathrm{\Psi }=\mathrm{\Phi }`$ at small $`k\eta `$, so large that one leaves the linear-perturbation regime for the expansion (47) of the metric much before the high-curvature scale is reached. Does this mean that the metric becomes very inhomogeneous? It would look to be the case … unless the growth of $`\mathrm{\Psi }`$ and $`\mathrm{\Phi }`$ is in some way a gauge artefact. But how can it be a gauge artefact if $`\mathrm{\Psi }`$ and $`\mathrm{\Phi }`$ correspond, in this gauge, to the gauge-invariant Bardeen potentials? The answer to this question was provided in . By going from the longitudinal gauge to an “off-diagonal” gauge with $`\mathrm{\Psi }=E=0`$, or, even better, to one in which only $`\mathrm{\Phi }`$ and $`E`$ appear, one finds that perturbations of the metric remain small at all $`\eta `$ till Planckian/string-scale curvatures are reached.
This is easy to see, for instance, in a gauge with $`\mathrm{\Psi }=B=0`$, where $`\mathrm{\Psi }_B(a^{}/a)E^{}`$. Clearly this gives $`E\eta ^2\mathrm{\Psi }_B`$ and, since $`E`$ enters the metric with two spatial derivatives, this implies that $`h_{ij}(k\eta )^2\mathrm{\Psi }_B`$, which is sufficiently small at small $`k\eta `$ for linear perturbation theory to be valid. One can then look for physical effects of these scalar perturbations (e.g. for contributions to CMBA) and find that they actually remain as small as the tensor contributions. In conclusion, once gauge artefacts are removed, it seems that adiabatic scalar perturbations, as well as their tensor counterparts, remain exceedingly small at large scales.
On the other hand, the rather large yields at short scales also apply to dilatons. This allows for a possible source of scalar waves if the dilaton is very light. However, as recently discussed by Gasperini , it is very unlikely that such a signal will be observable, given the constraints on the dilaton mass due to tests of the equivalence principle (see Section 2). Other restrictions on the dilaton mass come from the possibility that their density may become overcritical and close the Universe. This and other possible interesting windows in parameter space are discussed in , and will not be reported in any detail here.
### 4.4 Gauge-field perturbations: seeds for $`\stackrel{}{B}_{gal}`$?
In standard inflationary cosmology there is no amplification of the vacuum fluctuations of gauge fields. This is a straightforward consequence of the fact that inflation makes the metric conformally flat, and of the decoupling of gauge fields from a conformally flat metric precisely in $`D=3+1`$ dimensions.
As a very general remark, apart from pathological solutions, the only background field that can amplify, through its cosmological variation, e.m. (more generally gauge-field) quantum fluctuations is the effective gauge coupling itself . By its very nature, in the pre-big bang scenario the effective gauge coupling inflates together with space during the PBB phase. It is thus automatic that any efficient PBB inflation brings together a huge variation of the effective gauge coupling, and thus a very large amplification of the primordial e.m. fluctuations . This can possibly provide the long-sought for origin for the primordial seeds of the observed galactic magnetic fields.
To be more quantitative, since the pump field for electromagnetic perturbations is the effective (four-dimensional) gauge coupling itself (see Eq. (46)), the total amplification of e.m. perturbations on any given scale $`\lambda `$ is given by $`\alpha _0/\alpha _{ex}`$, i.e. by the ratio of the fine structure constant now and the fine structure constant at the time of exit of the scale $`\lambda `$ during DDI. It turns out that, in order to produce sufficiently large seeds for the galactic magnetic fields, such a ratio has to be enormous for the galactic scale, i.e. about $`10^{66}`$. Taken at face value, this would be a very strong indication in favour of the PBB scenario, more particularly of DDI. Indeed, only in such a framework is it natural to expect that the effective gauge coupling grew during inflation by a factor whose logarithm is of the same order as the number of inflationary e-folds.
Notice, however, that, unlike GW, e.m. perturbations interact quite considerably with the hot plasma of the early (post-big bang) Universe. Thus, converting the primordial seeds into those that may have existed at the protogalaxy formation epoch is by no means a trivial exercise (see, e.g. ). The question of whether or not the primordial seeds generated in PBB cosmology can evolve into the observed galactic magnetic fields thus remains, to this date, an unsolved, yet very interesting, problem.
### 4.5 Axion perturbations: seeds for CMBA and LSS?
In four dimensions the curl of $`B_{\mu \nu }`$, $`H_{\mu \nu \rho }`$, is equivalent to a pseudoscalar field, the (KR) axion $`\sigma `$, through
$$H_{\mu \nu \rho }=e^\varphi ϵ_{\mu \nu \rho \tau }^\tau \sigma .$$
(52)
It is easy to see that, while the pump field for $`B_{\mu \nu }`$ is $`a^2e^\varphi `$, that for $`\sigma `$ is $`a^2e^\varphi `$. Indeed their respective perturbations are related by the duality of perturbations discussed in Subsection 4.1. We can use either description with identical physical results. Note that, while $`a`$ and $`\varphi `$ worked in opposite directions for tensor and dilaton perturbations, generating strongly tilted (blue) spectra, the two work in the same direction for axions, so that spectra can be flat or even tilted towards large scales (red spectra) . An interesting fact is that, unlike the GW spectrum, that of axions is very sensitive to the cosmological behaviour of internal dimensions during the DDI epoch. On one side, this makes the model less predictive. On the other, it tells us that axions represent a window over the multidimensional cosmology expected generically from string theories, which must live in more that four dimensions. Parametrizing the spectrum by:
$$\mathrm{\Omega }_{ax}(k)=\left(\frac{H_{max}}{M_P}\right)^2(k/k_{max})^\alpha ,$$
(53)
and considering the case of three non-compact and six compact dimensions with separate isotropic evolution, one finds:
$$\alpha =\frac{3+3r^22\sqrt{3+6r^2}}{1+3r^2},$$
(54)
where
$$r\frac{1}{2}\frac{\dot{V_6}V_3}{V_6\dot{V_3}}$$
(55)
is a measure of the relative evolution of the internal and external volumes. Equations (54), (55) show that the axion spectrum becomes exactly HZ (i.e. scale-invariant) when $`r=1`$, i.e. when all nine spatial dimensions of superstring theory evolve in a rather symmetric way . In situations near this particularly symmetric one, axions are able to provide a new mechanism for generating large-scale CMBA and LSS.
Calculation of the effect gives , for massless axions:
$$l(l+1)C_lO(1)\left(\frac{H_{max}}{M_P}\right)^4(\eta _0k_{max})^{2\alpha }\frac{\mathrm{\Gamma }(l+\alpha )}{\mathrm{\Gamma }(l\alpha )},$$
(56)
where $`C_l`$ are the usual coefficients of the multipole expansion of $`\mathrm{\Delta }T/T`$
$$\mathrm{\Delta }T/T(\stackrel{}{n})\mathrm{\Delta }T/T(\stackrel{}{n}^{})=\underset{l}{}(2l+1)C_lP_l(\mathrm{cos}\theta ),\stackrel{}{n}\stackrel{}{n}^{}=\mathrm{cos}\theta ,$$
(57)
and $`\eta _0k_{max}10^{30}`$. In string theory, as repeatedly mentioned, we expect $`H_{max}/M_PM_s/M_P1/10`$, while the exponent $`\alpha `$ depends on the explicit PBB background with the above-mentioned HZ case corresponding to $`\alpha =0`$. The standard tilt parameter $`n=n_s`$ ($`s`$ for scalar) is given by $`n=1+2\alpha `$ and is found, by COBE, to lie between $`0.9`$ and $`1.5`$, corresponding to $`0<\alpha <0.25`$ (a negative $`\alpha `$ leads to some theoretical problems). With these inputs we can see that the correct normalization ($`C_210^{10}`$) is reached for $`\alpha 0.2`$, which is just in the middle of the allowed range. In other words, unlike in standard inflation, we cannot predict the tilt, but when this is given, we can predict (again unlike in standard inflation) the normalization.
With some extra work one can compute the $`C_l`$ in the acoustic-peak region adding vector and tensor contributions from the seeds. It turns out that the acoustic-peak structure is very sensitive to $`\alpha `$, hence to the behaviour of the internal dimensions during the DDI phase. The above-mentioned value, $`\alpha =1`$, does not give peaks at all and, as such, looks ruled out by the data. Values of $`\alpha `$ in the range $`0.3`$$`0.4`$ appear to be preferred (especially in the presence of a cosmological constant with $`\mathrm{\Omega }_\mathrm{\Lambda }0.7`$). We saw, however, that the overall normalization was very sensitive to the value of $`\alpha `$. For $`\alpha `$ in the $`0.3`$$`0.4`$ range, the normalization is off (way too small) by many orders of magnitude. Therefore, if present indications are confirmed, as they seem to be from the recent release of the Boomerang 1997 data analysis , one will be forced to a $`k`$-dependent $`\alpha `$, meaning different phases in the evolution of internal dimensions during DDI.
### 4.6 Heating up the Universe
Before closing this section, I wish to recall how one sees the very origin of the hot big bang in this scenario. One can easily estimate the total energy stored in the quantum fluctuations, which were amplified by the pre-big bang backgrounds (for a discussion of generic perturbation spectra, see . The result is, roughly,
$$\rho _{quantum}N_{eff}H_{max}^4,$$
(58)
where $`N_{eff}`$ is the effective number of species that are amplified and $`H_{max}`$ is the maximal curvature scale reached around $`t=0`$. We have already argued that $`H_{max}M_s=\lambda _s^1`$, and we know that, in heterotic string theory, $`N_{eff}`$ is in the hundreds. Yet, this rather huge energy density is very far from critical, as long as the dilaton is still in the weak-coupling region, justifying our neglect of back-reaction effects. It is very tempting to assume that, precisely when the dilaton reaches a value such that $`\rho _{quantum}`$ is critical, the Universe will enter the radiation-dominated phase. This PBBB (PBB bootstrap) constraint gives, typically:
$$e^{\varphi _{exit}}1/N_{eff},$$
(59)
i.e. a value for the dilaton close to its present value.
The entropy in these quantum fluctuations can also be estimated following some general results . The result for the density of entropy $`S`$ is, as expected,
$$SN_{eff}H_{max}^3.$$
(60)
It is easy to check that, at the assumed time of exit given by (59), this entropy saturates recently proposed holography bounds. The discussion of such bounds is postponed to Subsection 5.7 since is has also interesting implications for the exit problem.
## 5 How could it have stopped?
We have argued that, generically, DDI, when considered at lowest order in derivatives and coupling, evolves towards a singularity of the big bang type. Similarly, at the same level of approximation, non-inflationary solutions of the FRW type emerge from a singularity. Matching these two branches in a smooth, non-singular way has become known as the (graceful) exit problem in string cosmology . It is, undoubtedly, the most important theoretical problem the PBB scenario is facing today.
Of course, one would not only like to know that a graceful exit does take place: one would also like to describe the transition between the two phases in a quantitative way. Achieving this goal would amount to nothing less than a full description of what replaces the big bang of standard cosmology in the PBB scenario. As mentioned in Section 1, this difficult problem is the analogue, in string cosmology, of the (still not fully solved) confinement problem of QCD. The exit problem is particularly hard because, by its very nature, and by the existing no-go theorems , it must occur, if at all, at large curvature and/or coupling and, because of fast time-dependence, must break (spontaneously) supersymmetry. The phenomenological predictions made in the previous section were based on the assumption that i) a graceful exit does take place; ii) sufficiently large scales are only affected by it kinematically, i.e. through an overall red-shift of all scales.
In this section, after recalling some no-go theorems for the exit, we will review various proposals that circumvent those theorems starting from the mathematically simplest, but physically least realistic, proposals and ending with the physically favoured, but harder to analyse, suggestions. The latter proposals suggest possible lines along which a quantitative description of the exit might eventually emerge. Needless to say, in spite of the many encouraging results, much work remains to be done: perhaps new techniques, and/or a deeper understanding of string theory in its non-perturbative regimes through the construction of the still largely unknown M-theory , need to be developed before a full quantitative description can be hoped for.
I should also mention that there have been suggestions that the BB singularity can be avoided if the DDI phase is highly anisotropic. While this is an interesting suggestion, with isotropization taking place later in the non-inflationary regime, we will stick here to the simplest case, in which DDI has already prepared a very homogeneous Universe before exit takes place. This is why our discussion on the exit problem is limited to the case of homogeneous cosmologies. Also, for lack of space, I shall refer the reader to the literature for most of the details.
### 5.1 No-go theorems
Under some restrictive conditions , it was shown that one cannot have a change of branch, i.e. that the Universe cannot make a permanent transition from the inflationary pre-big bang to a FRW post-big bang solution. Perhaps the best way to convey the physical meaning of those theorems is in terms of the necessary conditions for exit recently given by Brustein and Madden . These authors give necessary conditions for two subsequent events to occur: firstly, a branch change in the string frame should take place: this imposes the violation of some energy conditions; secondly, a bounce should occur in the E-frame metric since, as we have seen in Section 3, DDI represents a collapse in the E-frame. This latter transition requires further violation of energy conditions.
Before the reader gets too worried about these violations, I should point out that these refer to the equations of state satisfied by some “effective” sources, which include both higher-derivative and higher-loop corrections. It is well known that such sources generically do lead to violations of the standard energy conditions satisfied by normal matter or radiation-like classical sources.
### 5.2 Exit via a non-local $`V`$
This is perhaps the simplest example of an exit. It was first discussed in . The reason why this is not considered an appealing mechanism for the exit is that the potential it employs depends on $`\overline{\varphi }`$ (instead of $`\varphi `$), in order to preserve SFD. By general covariance such a potential, if non-trivial, must be non-local. Unfortunately, there has been no convincing proposal to explain how such non-local potentials might arise within a superstring theory framework.
### 5.3 Exit via $`B_{ij}`$
The antisymmetric KR field may lead to violation of the energy conditions and thus induce an exit. Some amusing examples were given in , where a non-trivial $`B_{ij}`$ is introduced through $`O(d,d)`$ transformations acting on a pure metric-dilaton cosmology of the type described in Subsection 2.5. It was found that these so-called “boosted” cosmologies were less singular than the original ones. In some cases they were even completely free of singularity and provided examples of exit, albeit in not-so-realistic situations.
It is tempting to speculate that this softening of singularities, due to a non-trivial $`B_{ij}`$ field, could be related to recent developments in the field of non-commutative geometry induced by a $`B_{ij}`$ field. Work along these lines is in progress.
### 5.4 Exit via quantum tunnelling
Several groups have attempted to describe the transition from the pre- to the post-big bang without modifying the low-energy tree-level effective action, by exploiting the quantum cosmology approach based on the Wheeler–De Witt (WDW) equation. In Refs. an $`O(d,d)`$-invariant WDW equation was derived in the $`(d^2+1)`$-dimensional mini-superspace consisting of a homogeneous Bianchi I metric, the antisymmetric tensor, and the dilaton. The $`O(d,d)`$ symmetry helps avoiding the ordering ambiguities which usually plague the WDW equation. For the time being only the mathematically simpler case of an $`O(d,d)`$-invariant potential $`V(\overline{\varphi })`$ has been analysed since, in that case, $`d^2`$-conserved charges can be defined and the “radial” part of the WDW equation reduces to a one-dimensional Schroedinger equation for a scattering problem.
It is amusing that, from such a point of view, the initial state of the Universe is described by a right-moving plane wave, which later encounters a potential, giving rise to both a transmitted and a reflected (i.e. left-moving) wave. The transmission coefficient gives the probability that the Universe ends up in the pre-big-bang singularity, while the reflection coefficient gives the probability of a successful exit into the post-big-bang decelerating expansion.
For certain forms of $`V(\overline{\varphi })`$ the wave is classically reflected and the WDW approach just confirms this expectation by giving a $`100\%`$ probability for the exit. However, even when there is no classical exit, the probability of wave-reflection is non-zero because of quantum tunnelling. The quantum probability of a classically forbidden exit turns out to be exponentially suppressed in the coupling constant $`e^\varphi `$, which is just fine. Unfortunately, it is also exponentially suppressed in the total volume of 3-space (in string units) after the pre-big-bang. Thus, only tiny regions of space have a reasonable chance to tunnel.
### 5.5 Higher-derivative corrections
While the examples of exit given in the previous subsections are theoretically interesting, they do look somewhat artificial and non-generic. In this and in the following subsection we shall describe two mechanisms for exit that involve very general properties of the lowest order solutions and of string theory. The present feeling is that, if graceful exit occurs, it should be maily induced by some combination of higher-derivative and higher-loop effects. Let us start with the former.
Toy examples have shown that DDI can flow, thanks to higher-curvature corrections, towards a de-Sitter-like phase, i.e. into a phase of constant $`H`$ (curvature) and constant $`\dot{\varphi }`$. This phase is expected to last until loop corrections become important and give rise to a transition to a radiation-dominated phase (see the next subsection). The idea is to justify the strong curvature transition from the dilatonic to the string phase by proving the existence of an exact de Sitter-like solution to the field equation, which acts as a late time attractor for the perturbative DDI branch. As shown in , the existence of such attractors depends on the existence of (non-trivial) solutions for a system of $`n`$ algebraic equations in $`n`$ unknowns. In general, we may expect a discrete number of solutions to exist. If at least one of them has some qualitative characteristics, it will act as a late-time attractor for solutions approaching DDI in the far past. An explicit example of this phenomenon was constructed in . In this connection, it is worth mentioning that solutions connecting duality-related low-energy branches through a high-curvature CFT were already proposed in .
It was recently pointed out that the reverse order of events is also possible. The coupling may become large before the curvature does. In this case, at least for some time, the low-energy limit of M-theory should be adequate: this limit is known to give $`D=11`$ supergravity and is therefore amenable to reliable study. It is likely, though not yet clear, that, also in this case, strong curvatures will have to be reached before the exit can be completed.
### 5.6 Loop corrections and back reaction
The idea here is to invoke the back reaction from particle production as the relevant mechanism. Since the back reaction is an $`O(e^\varphi \alpha ^{}H^2)`$ correction, its effect is contained in one-loop $`O(R^2)`$ contributions to the effective action. A recent calculation shows that, indeed, loop corrections to DDI work in the right direction and become relevant precisely when expected according to the exit criterion (59).
A class of such contributions was analysed some time ago by Antoniadis et al. in the case of a spatially flat ($`k=0`$) cosmology and by Easther and Maeda in the case of a closed Universe ($`k=1`$). Both groups find non-singular solutions to the loop-corrected field equations. However, neither group is actually able to obtain solutions that start in the dilaton-driven superinflationary regime and later evolve through a branch change.
More recently, several examples of full exit have been constructed . Although they are based on $`\alpha ^{}`$ and loop-corrected actions, which have not been derived from reliable calculations, they seem to indicate, at least, that the BM conditions for exit may turn out to be not just necessary but also sufficient. It also appears that exit occurs when the entropy bound becomes threatened by the entropy in the amplified/squeezed quantum fluctuations, as we shall now discuss.
### 5.7 Entropy considerations
Entropy-related considerations have recently led to model-independent arguments in favour of the occurrence of a graceful exit in string cosmology. As we shall see, those are physically quite close to the arguments based on back-reaction and loop corrections, which we have just discussed in the previous subsection.
Almost twenty years ago Bekenstein suggested that, for a limited gravity system of energy $`E`$, and whose size $`R`$ is larger than its gravitational radius, $`R>R_g2G_NE`$, entropy is bounded by $`S_{BEB}`$:
$$S_{BEB}=ER/\mathrm{}=R_gRl_P^2.$$
(61)
Holography suggests that maximal entropy is bounded by $`S_{HOL}`$,
$$S_{HOL}=Al_P^2,$$
(62)
where $`A`$ is the area of the space-like surface enclosing the region of space whose entropy we wish to bound. For systems of limited gravity, since $`R>R_g,A=R^2`$, (61) implies the holography bound (62).
Can these entropy bounds be applied to the whole Universe, i.e. to cosmology? A cosmological Universe is not a system of limited gravity, since its large-distance behaviour is determined by the gravitational effect of its matter content through Friedmann’s equation (4). Furthermore, the holography bound obviously fails for sufficiently large regions of space since, for a given temperature, entropy grows like $`R^3`$ while area grows like $`R^2`$. The generalization of entropy bounds to cosmology turned out to be subtle.
In 1989, Bekenstein himself gave a prescription for a cosmological extension by choosing $`R`$ in Eq. (61) to be the particle horizon. Amusingly, he arrived at the conclusion that the bound is violated sufficiently near the big-bang singularity, implying that the latter is fake (if the bound is always valid). About a year ago, Fischler and Susskind (FS) proposed a similar extension of the holographic bound to cosmology, arguing that the area of the particle horizon should bound entropy on the backward-looking light cone according to (62). It was soon realized, however, that the FS proposal requires modifications, since violations of it were found to occur in physically reasonable situations. An improvement of the FS bound applicable to light-like hypersurfaces was later made by Bousso .
Of more interest here are the attempts made at deriving cosmological entropy bounds on space-like hypersurfaces . These identify the maximal size of a spatial region for which holography works: the Hubble radius , the apparent horizon , or, finally, a so-called causal connection (Jeans) scale .
For our purposes here, we do not need to enter into the relative merits of these various proposals. Rather, we will only outline the physical idea behind them. Consider, inside a quasi-homogeneous Universe, a sphere of radius $`H^1`$. We may consider “isolated” bodies, in the sense of ref. , fully contained in the sphere, i.e. with radius $`R<H^1`$. For such systems, the usual BB holds and is saturated by a black hole of size $`R`$. We may consider next several black holes inside our Hubble volume, each carrying an entropy proportional to the square of its mass. If two, or more, of these black holes merge, their masses will add up, while the total entropy after the merging, being quadratic in the total mass, will exceed the sum of the initial entropies. In other words, in order to maximize entropy, it pays to form black holes as large as possible.
Is there a limit to this process of entropy increase? The suggestion made in , which finds support in old results by several groups , is that a critical length of order $`H^1`$ is the upper limit on how large a classically stable black hole can be. If we accept this hypothesis, the upper bound on the entropy contained in a given region $``$ of space will be given by the number of Hubble volumes in $``$, $`n_H=VH^3`$ times the Bekenstein–Hawking entropy of a BH or radius $`H^1`$, $`H^2l_P^2`$. The two factors can be combined in the suggestive formula:
$$S()<l_P^2_{}d^3x\sqrt{h}\stackrel{~}{H}S_{HB},$$
(63)
where $`_{}d^3x\sqrt{h}`$ is the volume of the space-like hypersurface whose entropy we wish to bound, and $`\stackrel{~}{H}`$ differs from one proposal to another , but is, roughly, of the order of the Hubble parameter. Actually, since $`H`$ is proportional to the trace of the second fundamental form on the hypersurface, Eq. (63) reminds us of the boundary term that has to be added to the gravitational action in order to correctly derive Einstein’s equations from the usual variational principle. This shows that the bound (63) is generally covariant for $`\stackrel{~}{H}=H`$. It can also be written covariantly for the identification of $`\stackrel{~}{H}`$ made in .
For the qualitative discussion that follows, let us simply take $`\stackrel{~}{H}=H`$ and let us convert the bound to string-frame quantities, taking into account the relation between $`l_P`$ and $`\lambda _s`$, given in Eq. (3). We obtain :
$$S()<(VH^3)(H^2\lambda _s^2e^\varphi )=e^{\overline{\varphi }}H\lambda _s^2,$$
(64)
where we have fixed an arbitrary additive constant in the definition (6) of $`\overline{\varphi }`$. Equation (64) thus connects very simply the entropy bound of a region of fixed comoving volume to the most important variables occurring in string cosmology (see, e.g., the phase diagram of Fig. 3).
An immediate application of the bound (64) was pointed out in . Noting that the bound is initially saturated in the BDV picture of collapse/inflation, the bound itself cannot decrease without a violation of the second law. This gives immediately:
$$\dot{\overline{\varphi }}\frac{\dot{H}}{H}.$$
(65)
It is easy to check that this inequality is obeyed, but just so, during DDI, in the sense that it holds with the equality sign. In other words, the HEB is saturated initially and throughout DDI in the BDV picture. The bound also turns out to give a physically acceptable value for the entropy of the Universe just after the big bang: a large entropy on the one hand (about $`10^{90}`$); a small entropy for the total mass and size of the observable Universe on the other, as often pointed out by Penrose . Thus, PBB cosmology neatly explains why the Universe, at the big bang, looks so fine-tuned (without being so) and provides a natural arrow of time in the direction of higher entropy .
What happens in the mysterious string phase, where we are desperately short of reliable techniques? It is quite clear that Eq. (65) does not allow $`H`$ to reach saturation ($`\dot{H}=0`$) in the first quadrant of Fig. 3 since $`\dot{\overline{\varphi }}>0`$ there. Instead, saturation of $`H`$ in the second quadrant (where $`\dot{\overline{\varphi }}0`$) is perfectly all right. But this implies having attained the sought for branch change!
Let us finally look at the loop corrections. Since, physically, these correspond to taking into account the back-reaction from particle production, we may check when the entropy in the cosmologically produced particles starts to threaten the bound. As discussed in Subsection 4.6, the entropy density in quantum fluctuations is given by $`\sigma N_{eff}H^3`$, which equals the bound $`\sigma _{HEB}Hl_P^2`$ precisely when $`l_P^2H^2N_{eff}=O(1)`$. But, as already pointed out, this is just the line on which the energy density in quantum fluctuations becomes critical (see Eq. (59)) and where, according to , the back-reaction becomes $`O(1)`$. Similar conclusions are reached by applying generalized second law arguments .
The picture that finally emerges from all these considerations is best illustrated with reference to the diagram of Fig. 4. Two lines are shown, representing boundaries for the possible evolution. The horizontal boundary is forced upon by the large-curvature corrections, while the tilted line in the first quadrant corresponds to the equation $`l_P^2H^2N_{eff}=O(1)`$ that we have just discussed. Amusingly, this line was also suggested by Maggiore and Riotto as a boundary beyond which copious production of $`0`$-branes would set in. Thus, depending on initial conditions, the PBB bubble corresponding to our Universe would hit first either the high-curvature or the large-entropy boundary and initiate an exit phase. Hopefully, a universal late-time attractor will emerge guiding the evolution into the FRW phase of standard cosmology (shown as a vertical line in Fig. 10).
Needless to say, all this has to be considered, at best, as having heuristic value. If taken seriously, it would suggest that the Universe will never enter the strong-coupling, strong-curvature regime, where the largely unknown M-theory should be used. The low-energy limit of the latter (the much better understood 11-D supergravity) could suffice to deal with the fundamental exit problem of string cosmology. We refer to the literature for several other attempts at M-cosmology .
## 6 Outlook
The outlook for the pre-big bang scenario, as formulated at present, is not necessarily an optimistic one. I am not sure I would bet a lot of money on it being right! But this is not really the issue. We have to remember that the PBB scenario is a top–down approach to cosmology. As stressed in the introduction, it would be quite a miracle if the correct model could be guessed without extensive feed-back from the data. The good news here is that new data are coming in all the time, and will continue to do so with more and more precision in the coming years!
Rather, we should draw some lessons from this new attempt at very early cosmology, whether it succeeds or it fails. As I can see, the main lessons to be drawn are the following:
* Our Universe did not have to emerge, together with space and time, from a singularity; in string theory, the singularity should be fake, either because it is tamed by finite-string-size effects, or because it simply signals the need for new degrees of freedom in the description of physics at very short distances;
* Because string theory is an extension of GR, inflation is possible in that context even in the absence of potential energy (i.e. of an effective cosmological constant); actually, inflation is very natural and easy to achieve, being a consequence of the duality symmetries of the string-cosmology equations;
* Inflation in string cosmology can be related, mathematically, to the problem of gravitational collapse in GR; as such, it is a generic phenomenon, once the assumption of asymptotic past triviality is made; furthermore, the curvature scale and the coupling at the onset of PBB inflation are arbitrary classical moduli;
* The Universe did not have to start hot! A hot Universe can emerge from a cold one thanks to quantum particle production in inflationary backgrounds;
* PBB cosmology predicts a rich spectrum of perturbations with different spectra depending on each perturbation’s “pump” field and on its evolution in the PBB era; observable relics of these perturbations may serve as a window on physics in the pre-bangian Universe all the way down to the string/Planck scale;
* The simplest PBB models either predict too small perturbations at large scales, or a spectrum of isocurvature perturbations which may be already “experimentally challenged” (as Rocky Kolb would kindly say);
* The exit problem still remains the hardest theoretical challenge to the whole idea of PBB cosmology;
* Hopefully, the combination of the above-mentioned experimental and theoretical challenges will be able to tell us whether the PBB idea is just doomed, or whether parts of it should be kept while searching for a better scenario; it should also suggest new avenues for physics-driven research in string/M-theory;
* Last but not least, the PBB idea has taught us that we need not lock ourselves into preconceived ideas in cosmology (cf. “the big bang is the beginning of time”, “inflation needs a scalar potential”); rather, we should contemplate as wide a range of theoretically sound possibilities as we can in order for Nature to choose, at best, one of them.
ACKNOWLEDGEMENTS
I am very grateful to Pierre Binétruy and Richard Schaeffer for having invited me to such a pleasant and stimulating school, to François David for the perfect organization, and to all the students for their patience in listening (after 5 weeks of courses!) and for their interesting questions. |
warning/0002/cond-mat0002140.html | ar5iv | text | # Thermodynamics of the Spin-1/2 Antiferromagnetic Uniform Heisenberg Chain
## Abstract
We present a new application of the traditional thermodynamic Bethe ansatz to the spin-1/2 antiferromagnetic uniform Heisenberg chain and derive exact nonlinear integral equations for just two functions describing the elementary excitations. Using this approach the magnetic susceptibility $`\chi `$ and specific heat $`C`$ versus temperature $`T`$ are calculated to high accuracy for $`5\times 10^{25}T/J5`$. The $`\chi (T)`$ data agree very well at low $`T`$ with the asymptotically exact theoretical low-$`T`$ prediction of S. Lukyanov, Nucl. Phys. B 522, 533 (1998). The unknown coefficients of the second and third lowest-order logarithmic correction terms in Lukyanov’s theory for $`C(T)`$ are estimated from the $`C(T)`$ data.
The spin $`S=1/2`$ antiferromagnetic (AF) uniform Heisenberg chain has a long and distinguished history in condensed matter physics and exhibits unusual static and dynamical properties unique to one-dimensional spin systems. It has been used as a testing ground for many theoretical approaches. The Hamiltonian is $`=J_i𝑺_i𝑺_{i+1}`$, where $`J>0`$ is the AF Heisenberg exchange interaction between nearest-neighbor spins. In this paper we usually set $`k_\mathrm{B}=1`$ and $`g\mu _\mathrm{B}=1`$ where $`k_\mathrm{B}`$ is Boltzmann’s constant, $`g`$ is the spectroscopic splitting factor of the spins and $`\mu _\mathrm{B}`$ is the Bohr magneton; also, the reduced temperature $`tT/J`$ where $`T`$ is the absolute temperature.
The $`S=1/2`$ Heisenberg chain is known to be exactly solvable, i.e. all eigenvalues can be obtained from the so-called Bethe ansatz equations. Despite the amazing property of being integrable, the Heisenberg chain has defied many attempts to calculate physical observables including thermodynamic quantities. A rather direct evaluation of the partition function was constructed in and is known as the “thermodynamic Bethe ansatz” (TBA), but this did not allow for high accuracy calculations especially in the low temperature region. The fundamental problem in is the necessity to deal with infinitely many coupled nonlinear integral equations for which the truncation procedures are difficult to control.
The possibility to accurately calculate the physical properties of the $`S=1/2`$ Heisenberg chain improved following the development of the path integral formulation of the transfer matrix treatment of quantum systems. On the basis of a Bethe ansatz solution to the quantum transfer matrix, Eggert, Affleck and Takahashi in 1994 obtained numerically exact results for the magnetic susceptibility $`\chi (t)`$ down to much lower temperatures than before and compared these with their low-$`t`$ results from conformal field theory. They found, remarkably, that $`\chi (t0)`$ has infinite slope: their conformal field theory calculations showed that the leading order $`t`$ dependence is $`\chi (t0)=\chi (0)\{1+1/[2\mathrm{ln}(t_0/t)]\}`$, where the value of $`t_0`$ is not predicted by the field theory. Such log terms are called “logarithmic corrections”. From their comparison of their field theory and Bethe ansatz calculations which extended down to $`t=0.003`$, Eggert, Affleck and Takahashi estimated $`t_07.7`$. Their numerical $`\chi (t)`$ values are up to $`10`$% larger than the former Bonner-Fisher extrapolation for $`t0.25`$.
Lukyanov has recently presented an exact asymptotic field theory for $`\chi (t)`$ and the specific heat $`C(t)`$ at low $`t`$, including the exact value of $`t_0`$ . These results are claimed to be exact in the sense of a renormalization group treatment close to a fixed point where only few operators are responsible for perturbations. Questions arising about such calculations are whether these operators have been correctly identified and whether the effective theory has been properly evaluated. A meaningful test of Lukyanov’s theory is only possible using numerical data of very high accuracy and at extremely low temperatures, such as we have attained in our numerical calculations to be presented below.
In this Letter we present a new application of the traditional TBA to the spin-1/2 Heisenberg chain and derive exact nonlinear integral equations \[Eqs. (37) below\] involving just two functions describing the elementary excitations. Our derivation evolved from earlier work by one of us using the powerful lattice approach. By means of a lattice path integral representation of the finite temperature Heisenberg chain and the formulation of a suitable quantum transfer matrix, a set of numerically well-posed expressions for the free energy was derived. A serious disadvantage of this approach lies in the complicated and physically non-intuitive mathematical constructions, which strongly inhibits generalizations to other integrable, notably itinerant fermion models. The present work is a new analytic derivation of the finitely-many integral equations of by means of the intuitive TBA approach. Our Eqs. (37) are identical to those obtained in by a rigorous, however much more involved method. In our new construction, we assume that magnons (on paths $`C_\pm `$) are elementary excitations and contain all information about the thermodynamics. Bound states are implicitly taken into account by use of the exact scattering phase probed in the analyticity strip. The a posteriori success of our reasoning is important for two reasons. First, our construction is as simple as the standard TBA, however avoiding the problems of dealing with density functions for (up to) infinitely many bound states. This may be of great advantage in the study of more complicated systems. Second, we have a simple particle approach to the Heisenberg chain which will allow for a study of transport properties like the Drude weight which has not been possible within the path integral approach .
We also demonstrate here that using our integral equations one can improve the accuracy and extend the temperature range of numerical calculations of $`\chi (t)`$ and $`C(t)`$ for the $`S=1/2`$ Heisenberg chain on the lattice far beyond those of previous calculations. We find agreement of our data with the above theory of Lukyanov for $`\chi (t)`$ to high accuracy ($`1\times 10^6`$) over a temperature range spanning 18 orders of magnitude, $`5\times 10^{25}t5\times 10^7`$; the agreement in the lower part of this temperatures range is much better, $`𝒪(10^7)`$. For $`C(t)`$, the logarithmic correction in Lukyanov’s theory is insufficient to describe our numerical data accurately even at very low $`t`$, so we estimate the coefficients of the next two logarithmic correction terms in his theory from our $`C(t)`$ data.
Derivation of integral equations
We start with the partially anisotropic Hamiltonian $`=J_i(S_i^xS_{i+1}^x+S_i^yS_{i+1}^y+\mathrm{cos}(\gamma )S_i^zS_{i+1}^z)h_iS_i^z`$ with $`0<\gamma <\pi /2`$ and magnetic field $`h`$. The dynamics of the magnons, i.e. the elementary excitations above the ferromagnetic state, constitute the Bethe ansatz. Momentum $`p`$ and energy $`ϵ`$ are suitably parametrized in terms of the spectral parameter $`x`$
$$p(x)=i\mathrm{log}\frac{\mathrm{sinh}(xi\gamma /2)}{\mathrm{sinh}(x+i\gamma /2)},ϵ(x)=J\frac{\mathrm{sin}\gamma }{2}p^{}(x)h,$$
(1)
where real values are obtained for Im$`x=0`$ and $`\pi /2`$, defining magnon bands of type “$`+`$” and “$``$”.
Any two magnons with spectral parameters $`x`$ and $`y`$ scatter with phase shift $`\mathrm{\Theta }(xy)`$ where
$$\mathrm{\Theta }(z)=i\mathrm{log}\frac{\mathrm{sinh}(zi\gamma )}{\mathrm{sinh}(z+i\gamma )}.$$
(2)
Next we apply the standard TBA just to the magnons and ignore bound states! However, the magnons on band $``$ are considered for spectral parameter $`x`$ with Im $`x`$ $`=\gamma `$ hence avoiding the branch cut in the scattering phase.
The density functions for particles $`\rho _j`$ and holes $`\rho _j^h`$ for the bands $`j=+`$, $``$ give rise to the definition of the ratio function $`\eta _j=\rho _j^h/\rho _j`$. Our analysis shows that $`\eta _+`$ and $`\eta _{}`$ are analytic continuations of each other. Quantitatively we find $`\eta _{}(x+i\gamma )=\eta _+(x)=:\eta (x)`$ subject to the non-linear integral equation
$$\mathrm{log}\eta (x)=\frac{ϵ(x)}{T}+_C\kappa (xy)\mathrm{log}(1+\eta ^1(y))𝑑y$$
(3)
where $`\kappa (x)=\frac{1}{2\pi }\mathrm{\Theta }^{}(x)`$ and $`C`$ is a contour consisting of the paths $`C_+`$ and $`C_{}`$ with Im $`y=0`$ and $`\gamma `$ encircled in clockwise manner. Substituting $`\mathrm{log}(1+\eta ^1)=\mathrm{log}(1+\eta )\mathrm{log}\eta `$ on the contour $`C_+`$ and resolving for $`\mathrm{log}\eta `$ we find
$`\mathrm{log}\eta (x)={\displaystyle \frac{\overline{ϵ}(x)}{T}}`$ $`+`$ $`{\displaystyle _{C_+}}\overline{\kappa }(xy)\mathrm{log}(1+\eta (y))𝑑y`$ (4)
$``$ $`{\displaystyle _C_{}}\overline{\kappa }(xy)\mathrm{log}(1+\eta ^1(y))𝑑y`$ (5)
with
$`\overline{ϵ}(x)`$ $`=`$ $`J{\displaystyle \frac{\mathrm{sin}\gamma }{2}}e_0(x)+{\displaystyle \frac{\pi }{2(\pi \gamma )}}h,e_0(x)={\displaystyle \frac{\frac{\pi }{\gamma }}{\mathrm{cosh}\frac{\pi }{\gamma }x}},`$ (6)
$`\overline{\kappa }(x)`$ $`=`$ $`{\displaystyle \frac{1}{2\pi }}{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \frac{\mathrm{sinh}(\frac{\pi }{2}\gamma )k}{2\mathrm{cosh}\frac{\gamma }{2}k\mathrm{sinh}\frac{\pi \gamma }{2}k}}\mathrm{e}^{ikx}𝑑k.`$ (7)
Finally, ignoring $`T`$ and $`h`$ independent contributions we obtain the free energy as
$$f=\frac{T}{2\pi }_{\mathrm{}}^{\mathrm{}}e_0(x)\mathrm{log}\left[(1+\eta (x))(1+\eta ^1(xi\gamma ))\right]𝑑x.$$
(8)
Numerical study of low-$`T`$ behavior
Lukyanov’s low-$`t`$ asymptotic expansion of $`\chi (t)`$ is
$`\chi _{\mathrm{lt},\mathrm{g}}(t)J={\displaystyle \frac{1}{\pi ^2}}\{1`$ $`+`$ $`{\displaystyle \frac{g}{2}}+{\displaystyle \frac{3g^3}{32}}+𝒪(g^4)`$ (9)
$`+`$ $`{\displaystyle \frac{\sqrt{3}}{\pi }}t^2[1+𝒪(g)]\},`$ (10)
where $`g(t/t_0)`$ obeys the transcendental equation $`\sqrt{g}\mathrm{exp}(1/g)=t_0/t,`$ with a unique value of $`t_0`$ given by $`t_0=\sqrt{\pi /2}\mathrm{exp}(\gamma +1/4)2.866`$ where $`\gamma `$ is Euler’s constant. His expansion for the free energy per spin at $`h=0`$ yields the specific heat per spin as
$`C_{\mathrm{lt},\mathrm{g}}(t)={\displaystyle \frac{2t}{3}}[1`$ $`+`$ $`{\displaystyle \frac{3}{8}}g^3+𝒪(g^4)]`$ (11)
$`+`$ $`{\displaystyle \frac{(2)3^{5/2}t^3}{5\pi }}[1+𝒪(g)],`$ (13)
where the exact prefactor $`2t/3`$ was found by Affleck in 1986, and the prefactor 3/8 in the logarithmic correction term agrees with .
Numerical data for $`\chi (t)`$ and $`C(t)`$ were obtained using our free energy expression (8). These data are considerably more accurate than those presented previously in . Our $`\chi (t)J`$ data, and the exact value $`1/\pi ^2`$ at $`t=0`$ , are plotted in Fig. 1. The calculations have an absolute accuracy of $`1\times 10^9`$. The data show a maximum at a temperature $`t^{\mathrm{max}}=\mathrm{0.6\hspace{0.17em}408\hspace{0.17em}510}(4)`$ with a value $`\chi ^{\mathrm{max}}J=\mathrm{0.146\hspace{0.17em}926\hspace{0.17em}279}(1)`$, yielding the $`J`$-independent product $`\chi ^{\mathrm{max}}T^{\mathrm{max}}=\mathrm{0.0\hspace{0.17em}941\hspace{0.17em}579}(1)`$. These values are consistent within the errors with those found by Eggert, Affleck and Takahashi, but are much more accurate.
The differences between our low-$`t`$ Bethe ansatz $`\chi (t)J`$ calculations and Lukyanov’s theoretical $`\chi _{\mathrm{lt},\mathrm{g}}(t)J`$ prediction in Eq. (10) are shown in Fig. 2. The error bar on each data point is the estimated uncertainty in $`\chi _{\mathrm{lt},\mathrm{g}}J`$ arising from the presence of the unknown $`𝒪(g^4)`$ and higher-order terms in Eq. (10), which was arbitrarily set to $`g^4(t)/\pi ^2`$; the uncertainty in the $`t^2`$ contribution, $`\sqrt{3}t^2g(t)/\pi ^3`$, is negligible at low $`t`$ compared to this. At the lower temperatures, the data agree extremely well with the prediction of Lukyanov’s theory. At the highest temperatures, higher order $`t^n`$ terms also become important. Irrespective of these uncertainties in the theoretical prediction at high temperatures, we can safely conclude directly from Fig. 2 that our numerical $`\chi (t)`$ data are in agreement with the theory of Lukyanov to within an absolute accuracy of $`1\times 10^6`$ (relative accuracy $`10`$ ppm) from $`t=5\times 10^{25}`$ to $`t=5\times 10^7`$. The agreement at the lower temperatures, $`𝒪(10^7)`$, is much better than this.
Our $`C(t)`$ data for $`t2`$ are shown in the inset of Fig. 3 and have an estimated accuracy of $`3\times 10^{10}C(t)`$. The data show a maximum with a value $`C^{\mathrm{max}}=\mathrm{0.3\hspace{0.17em}497\hspace{0.17em}121\hspace{0.17em}235}(2)`$ at a temperature $`t_C^{\mathrm{max}}=\mathrm{0.48\hspace{0.17em}028\hspace{0.17em}487}(1)`$. The electronic specific heat coefficient $`C(t)/t`$ is plotted in Fig. 3. These data exhibit a maximum with a value $`(C/t)^{\mathrm{max}}=\mathrm{0.8\hspace{0.17em}973\hspace{0.17em}651\hspace{0.17em}576}(5)`$ at $`t_{\mathrm{C}/\mathrm{t}}^{\mathrm{max}}=\mathrm{0.30\hspace{0.17em}716\hspace{0.17em}996}(2)`$. The existence of low-$`t`$ log corrections to $`C(t)`$ is revealed in the top plot of $`\mathrm{\Delta }C(t)/t`$ in Fig. 4, where $`\mathrm{\Delta }C(t)=C(t)2t/3`$ and $`2t/3`$ is the low-$`t`$ limit of $`C(t)`$. The influence of the $`g^3`$ log correction term in Eq. (13) is evaluated by subtracting it in the plot of $`\mathrm{\Delta }C(t)/t`$ as shown by the middle curve in Fig. 4. The $`t=0`$ singularity is still present but with reduced amplitude; this demonstrates that additional logarithmic correction terms are important within the accuracy of the data.
We estimate the unknown coefficients of the next two logarithmic correction ($`g^4,g^5`$) terms in Eq. (13) from our $`C(t)`$ data as follows. From Eq. (13), if we plot the data as $`[C(t)/t(2/3)(1+3g^3/8)]/g^4`$ vs $`g`$ and fit the lowest-$`t`$ data by a straight line, the $`y`$-intercept gives the coefficient of the $`g^4`$ term and the slope gives the coefficient of the $`g^5`$ term. We fitted a straight line to the data in such a plot for $`5\times 10^{25}t5\times 10^9`$ as shown by the weighted linear fit in Fig. 5 where the parameters of the fit are given in the figure. By subtracting the influences of these two logarithmic correction terms from the middle data set as shown in the bottom data set in Fig. 4, the singular behavior as $`t0`$ is largely removed, leaving a behavior which is close to a $`t^2`$ dependence as predicted by the last term in Eq. (13). Further discussion of the predictions of , and high-accuracy fits ($`0t5`$) to our $`C(t)`$ and $`\chi (t)`$ data and the respective exact $`t=0`$ values, will be presented elsewhere.
In conclusion, we have presented an analytic approach to the thermodynamics of the $`S=1/2`$ AF Heisenberg chain on the basis of a finite number of elementary excitations. We envisage that this approach can be generalized to study a variety of other systems such as Hubbard and $`t`$-$`J`$ models, quantum spin chains with higher symmetries and systems with orbital degrees of freedom. Our free energy expression has allowed numerical calculations of $`\chi (t)`$ and $`C(t)`$ for the Heisenberg chain to be carried out to much higher accuracy and to much lower temperatures than heretofore attained. Our $`\chi (t)`$ data are in excellent agreement with the theory of Lukyanov at low $`t`$. The logarithmic correction in Lukyanov’s theory for $`C(t)`$ is found insufficient to describe our $`C(t)`$ data accurately even at very low $`t`$. However, the $`t`$ dependence of the deviation agrees with the form of his theory, which enabled us to estimate the unknown coefficients of the next two logarithmic correction terms in his theory for $`C(t)`$ from our $`C(t)`$ data. Thus we have verified Lukyanov’s theory of a critical system perturbed by marginal operators and have given evidence that his asymptotic expansion can be systematically extended to higher order.
The authors acknowledge valuable discussions with U. Löw and K. Fabricius. Comparison of our results with their numerical data for the thermodynamics of finite systems proved essential to achieve high accuracy in the treatment of the nonlinear integral equations. D.C.J. thanks the University of Cologne and the Stuttgart Max-Planck-Institut für Festkörperforschung for their hospitality. A.K. acknowledges financial support by the Deutsche Forschungsgemeinschaft under grant No. Kl 645/3 and by the research program of the Sonderforschungsbereich 341, Köln-Aachen-Jülich. Ames Laboratory is operated for the U.S. Department of Energy by Iowa State University under Contract No. W-7405-Eng-82. The work at Ames was supported by the Director for Energy Research, Office of Basic Energy Sciences. |
warning/0002/astro-ph0002262.html | ar5iv | text | # Structural and Photometric Classification of Galaxies – I. Calibration Based on a Nearby Galaxy Sample
## 1 Introduction
It is now well-established that a large fraction of galaxies discovered at intermediate and high redshift have unusual morphologies, and thus cannot be classified in terms of the nominal Hubble-Sandage system (Driver et al. 1995, 1998; Abraham et al. 1996a, 1996b). The Hubble classification scheme is also difficult to apply to many local galaxies, dubbed ‘peculiar,’ or any galaxies imaged at low signal-to-noise ($`S/N`$), or apparently small size (relative to the point-spread function). The Hubble-Sandage classification system was predicated on the study of nearby, “normal” galaxies – luminous and relatively quiescent objects(\[Sandage 1961\], \[Sandage & Tamman 1987\], \[Sandage & Bedke 1993\]). While the classification system developed by de Vaucouleurs et al. (1976) makes an attempt to push the framework to “later” types, it still suffers from the above shortcomings. Fundamentally, these traditional classification schemes are based on the concept of ‘pigeon-holing’ galaxies based on a reference-set, or archetypes. These archetypes are selected from samples in the local universe, and are preferentially axisymmetric systems. Since our local census is undoubtedly incomplete, and, since galaxies evolve, such reference sets by their very definition are incomplete. Thus it is not surprising that these systems are of marginal utility in the study of dwarf galaxies, interacting galaxies, or galaxies at high redshift.
An alternative classification scheme could be based on quantitative indices, the inter-relation of which is not predetermined by a finite reference set. This would permit galaxies to be classified, for example, in different stages of their evolution; albeit the classification would be different but the basis set of indices would be the same. The goal of this paper is to define such a set of indices that can be used as quantitative, objective classifiers of galaxies (i) over a wide range in redshift, and (ii) for wide range of galaxy types. In particular, we desire classifiers that are well suited to typing both “normal” galaxies and the compact galaxies that are the focus of a companion study (Jangren et al. 2000; hereafter, paper II). We anticipate that such a classification scheme is both necessary and enabling for the exploration of the physical mechanisms driving galaxy evolution (Bershady 1999).
What are the desired characteristics of classification parameters? They should be physically interesting (closely related to underlying physical properties of galaxies), model-independent, and measurable for all galaxy types. It also should be possible to accurately determine the parameters chosen for a wide range of image resolution and signal-to-noise ratios.
From Hubble’s classification a posteriori we have learned that a strong correlation exists between galaxy spectral type and apparent morphological features – at least for the galaxy types which fit well within his scheme. This correlation – noted by Hubble as early as 1936 (Hubble 1936) – can loosely be termed a ‘color-morphology’ relation, although the correlation is not necessarily limited to broad-band color. This is a triumph of Hubble’s classification explicitly because it is not part of the classification. Furthermore, the correlation yields clues about the physical connection of the present matter distribution and the star-formation histories in galaxies. But while morphology (or ‘form’) and spectral type are correlated, there is also significant dispersion in this correlation. Some of the more notable deviations from the nominal color-morphology relation are found in the plethora of forms for spectrally ‘late’ type galaxies, the presence of ‘E+A’ galaxies (Dressler & Gunn 1993), and the compact, luminous, blue, emission-line galaxies studied in paper II (as we shall show). This points to the importance of form and spectral type as key, yet independent axes of a revised classification system.
However, the only example of such a revised classification system is that of Morgan (1958, 1959), where central light concentration is used as the primary classification parameter. Morgan was motivated by the fact that (i) a salient criterion used in classifying galaxies in the Hubble-Sandage system is the degree of central concentration of light; (ii) there was a significant dispersion in spectral type and Hubble type (\[Humason et al. 1956\]); and (iii) spectral type appeared to correlate more strongly with light concentration. In this way, Morgan hoped to wed the classification of stellar populations to the classification of galaxies. Nonetheless, he was compelled to introduce a secondary parameter, i.e. the ‘Form Family,’ because there was still a dispersion of morphological forms within each of his spectral types. Today, one should be able to improve upon Morgan’s scheme by introducing quantitative measures of image concentration and other indices of form, and by independently assessing the spectral type via colors or spectra.
A number of subsequent attempts have been made to construct quantitative classification system that could replace or modify the current Hubble scheme. Yet these schemes are generally based purely either on photometric form (e.g. \[Elmegreen & Elmegreen 1982\]; \[Okamura et al. 1984\]; \[Watanabe et al. 1985\]; \[Doi et al. 1993\]; \[Abraham et al. 1994\]; \[Odewahn 1995\]; \[Han 1995\]) or spectral type (e.g. \[Bershady 1995\]; \[Connolly et al. 1995\]; \[Zaritsky et al. 1995\]; \[Folkes et al. 1996\]; \[Bromley et al. 1998\]; \[Ronen et al. 1999\]). In essence, they have relied implicitly on an assumed correlation between galaxy spectral type and apparent morphology. Related attempts have been made to use artificial neural networks to reproduce the Hubble scheme in an objective way (e.g. \[Burda & Feitzinger 1992\]; \[Storrie-Lombardi et al. 1992\]; \[Spiekermann 1992\]; \[Serra-Ricart et al. 1993\]; \[Naim et al. 1995\]; \[Odewahn 1995\]; \[Odewahn et al. 1996\]). Yet these go no further in differentiating between spectral type and form. Only in Whitmore’s (1984) scheme are spectral and structural parameters combined, i.e., $`BH`$ color, size, and bulge-to-total ratio are used to define two principal classification axes of scale and form. But again, the correlation(s) between galaxy spectral type, scale and form are not explicit.
Here we attempt to expand on Morgan’s program by fully quantifying the classification of form via image concentration and several other structural parameters, and explicitly using color as an indicator of spectral type. In this study we choose to use only a single color ($`BV`$), but we anticipate that a more desirable, future development would be to include broad-wavelength coverage, multi-color data and spectroscopic line-indices. Spectroscopic line-indices would be required, for example, to identify E+A galaxies. While such galaxies are not the focus of the present work, a comprehensive classification scheme should be able to isolate these systems and determine the range of their morphology (cf. Dressler & Gunn 1992, Couch et al. 1994, and Wirth et al. 1994). Nonetheless, broad-band colors are a cost-effective way to characterize the spectral continuum (cf. Bershady, 1995, and Connolly et al. 1995). Of more direct relevance to the study at hand, a future elaboration of including $`UV`$ and $`VK`$ would enhance the ability to distinguish between spectral types particularly for galaxies with extremely blue, optical colors (e.g. Aaronson 1978, Bershady 1995).
We have also chosen to quantify form and scale via non-parametric measures, such as luminosity, half-light size and surface-brightness, asymmetry, and image concentration. An alternative, model-dependent approach is to decompose a galaxy’s light profile into a disk and bulge. The traditional one-dimensional decompositions are fraught with technical problems such that decompositions can only be achieved reliably for about half of all disk galaxies (\[Kent 1985\]). The newer two-dimensional decomposition techniques are superior (e.g. \[de Jong 1996b\]), and have been shown to successfully reproduce observed light profiles for faint galaxies (e.g. \[Simard et al. 1999\]). Indeed, one can argue that two-dimensional model fitting to imaging data is optimum in terms of using the available information, and for minimizing random error. At high $`S/N`$ and high angular resolution, however, even the most “normal” galaxies exhibit peculiarities (as discussed in more detail in §3.3.2) such that simple bulge-plus-disk models cannot reproduce these frequently observed peculiarities in light distributions with high fidelity. The situation worsens for “peculiar” galaxies. For this reason we have some concerns about the uniqueness of the observationally derived model parameters, and hence their interpretation. We anticipate future developments which use the models and non-parametric measurements in a hybrid scheme optimal for characterizing galaxy light distributions both in terms of random and systematic errors.
It is worth noting again that bright galaxy samples are notorious for missing or under-representing certain galaxy types – particularly dwarfs and low-surface-brightness galaxies. The samples used here are no exception. While this was one of our complaints about the classical Hubble scheme, there are two key differences with our approach: (i) the classification parameters we develop are objective; and (ii) these parameters do not assume the presence of basic axi-symmetry, disk-plus-bulge structure, or spiral patterns which underly the Hubble scheme. As we will show, the galaxies examined here are sufficiently diverse to establish the parameter space for a comprehensive classification scheme, although not the comprehensive classification itself. By developing an initial classification of these galaxies, however, we intend to use it as a foil against which we can begin to compare the classification of more distant samples: How are the classifications different? Do the nearby and distant samples occupy the same regions of parameter space? If not, do the differences represent continuous extensions of these parameters, or are they physically disjoint? These are the types of questions one can address given the limitations of current local samples. Note that we must stop short of identifying differences as “new,” epoch-specific classes of galaxies. Without a complete census of both the nearby and distant universe, it is not possible to establish whether there are different “classes” of galaxies at different redshifts; apparent differences could simply be artifacts of the presently limited samples. With such a complete census, in the future we may hope to address deeper issue of how the comoving space-densities of different classes evolve.
Towards the goal of establishing a comprehensive classification scheme of utility to distant galaxy studies, in this paper we assemble a robust set of non-parametric, photometric and structural properties for a range of nearby, lumimous galaxies. We define a multivariate, photometric parameter space that forms an initial classification scheme for these galaxies. This classification can be used reliably to identify comparable samples in other surveys and at higher redshift. In the accompanying paper (II) we measure these properties for compact, luminous emission-line galaxies at intermediate redshift, compare them to the “normal,” nearby galaxies studied here, and demonstrate that our classification parameter space distinguishes between these two samples. We discuss the implications for the evolution of this intermediate-redshift sample therein. In future papers in this series we intend to extend our analysis (a) to more representative samples of the local volume that include dwarf and emission-line galaxies (e.g., the University of Michigan Objective Prism Survey (\[Salzer et al. 1989\]); (b) to more comprehensive samples of distant galaxies, e.g magnitude-limited samples from the Hubble Deep Field; and (c) to studies of the morphological evolution of these distant samples. The classification scheme which we propose here is intended as a framework for these future studies.
The data sets are presented in §2; the analysis is described in §3. The results are presented in §4, and summarized in §5. Throughout this paper we adopt $`H_0`$ = 50 km s<sup>-1</sup> Mpc<sup>-1</sup>, $`q_0=0.1`$, $`\mathrm{\Lambda }=0`$.
## 2 Nearby galaxy samples
As a primary reference sample, 101 of the 113 local Hubble-type galaxies from the catalog of Frei et al. (1996) were analyzed. This sample will define what we mean by “normal” galaxies in this paper. This catalog is the only digital, multi-band, sample publicly available that is reasonably comprehensive; it consists of ground-based CCD images of bright galaxies, all apparently large (most have diameters of 4 \- 6) and well resolved. As a result, the sample contains mostly luminous and physically large galaxies: out of the 101 objects we used in our analysis, only seven have $`L<0.1L^{}`$. We excluded 12 objects whose apparent sizes were larger than the CCD field of view (thus their image structure parameters could not be well estimated). Two of the excluded objects are early-type galaxies (E–S0), seven are intermediate (Sa–Sb), and three are late-type (Sc–Irr).<sup>1</sup><sup>1</sup>1The excluded objects are: NGC 2403, 2683, 3031, 3079, 3351, 3623, 4406, 4472, 4594, 4826, 5746, and 6503. The majority of the remaining sample are spirals and S0 galaxies. Frei et al. have removed foreground stars from the images of the nearby galaxies, in a few cases leaving visible “scars;” except in the case of NGC 5792, these residuals did not cause noticeable problems when determining the structural parameters (§3.3).
In several instances in the present analysis we reference the sample of Kent (1984, 1985), which is composed of 53 nearby, luminous and physically large galaxies similar to the Frei et al. sample. We find Kent’s sample useful for comparison of both photometric and structural parameters. We also reference the sample of 196 normal (non-active) Markarian galaxies studied by Huchra (1977a). Relevant characteristics of the above three samples are summarized in Table 1, including an enumeration of the effective filter systems used in each study. Further details on these photometric systems are found in the studies listed in the Table and references therein.
### 2.1 Comparison of reference samples to emission-line galaxy samples
Both the Frei et al. and Kent samples are under-representative of dwarf galaxies, and contain neither HII galaxies nor low surface-brightness galaxies. The latter objects have been shown to make up a significant fraction of the local galaxy population (de Jong 1995, 1996a). Clearly our reference samples do not constitute a representative template of the local population. Here we estimate where these samples may be particularly un-representative with an eye towards the study of faint galaxy samples in future papers. In Figures 1 and 2 we compare the Frei et al. samples photometric properties of color and luminosity to (i) the normal Markarian galaxies (Huchra 1977a), (ii) dwarf spheroidals (as described in the following section), and (iii) the intermediate redshift samples presented in paper II.
Since the Markarian galaxies were selected from objective prism plates based on their strong UV continua, the sample is biased toward bluer colors than the Frei et al. galaxies and is thus likely more representative of star-forming galaxies. Huchra’s sample contains fainter galaxies that extend the magnitude range down to $`M_B14`$ and the color-color locus blue-ward of $`BV=0.4`$.
The intermediate-redshift galaxies, also selected in part due to their blue color (see paper II), have blue luminosities comparable to the brighter half of the Frei et al. sample, but with bluer colors. This places most of them in a distinct region of the color-luminosity plot from the Frei et al. sample. In contrast, the distribution of the Markarian galaxies extends into the region occupied by the intermediate-redshift objects. In the color-color diagram, again the intermediate-redshift galaxies largely overlap with the Markarian sample in the region corresponding to extreme blue colors not occupied by the Frei et al. or Kent samples.
In short, the Frei et al. sample is spectro-photometrically disjoint from extreme samples of blue, star-forming galaxies at intermediate-redshift (e.g., paper II), even though both contain intrinsically luminous and moderate-to-high surface-brightness systems. Yet clearly there ar local examples (e.g., from Markarian) which are as blue and luminous as these intermediate-redshift, star-forming galaxies. These sources are simply missing from the Frei et al. sample. The comparison of the global properties of the intermediate-redshift, compact, star-forming galaxies in paper II to those of local galaxies from Frei et al. (here) is then an initial step in mapping the range of galaxy types at any redshift. Further investigation of the nature and evolution of these types of extreme, star-forming systems will be greatly facilitated by future work quantifying the image structure of local counterparts $`BV`$$`<0.4`$ and $`M_B<19`$.
### 2.2 Comparison of reference samples to dwarf spheroidals
We have made some attempt, where possible, to access the photometric and structural properties of other key dwarf populations. We schematically indicate the locus of dwarf ellipticals/spheroidals in Figures 1 and 2 using data from \[Caldwell 1983\], \[Bingelli & Cameron 1991\], and \[Bingelli & Jerjen 1998\]. The dwarf spheroidals occupy a virtually unpopulated region of the color-luminosity diagram at relatively red colors and low luminosity. The absence of such objects from most surveys is attributed typically to a selection bias since these sources are at low surface-brightness. It is interesting to note that in the color-color diagram the dwarf spheroidals occupy a region over-lapping with the early-to-intermediate type spirals. Hence the integrated broad-band light of these systems are unusual compared to our reference samples only with respect to their luminosity. We refer to the dwarf spheroidal properties extensively in future papers where we also explore their image structural properties.
## 3 Analysis
As noted in the Introduction, many galaxies are sufficiently unusual that they cannot be classified in terms of the normal Hubble scheme. This becomes increasingly true at intermediate redshifts. The compact, luminous emission-line galaxies in paper II are such an example. This is not due to poor spatial resolution, but to truly unusual morphological properties, e.g., off-centered nuclei, tails, asymmetric envelopes, etc. To compare such objects morphologically to “normal” galaxies, we define here six fundamental parameters of galaxy type that are quantitative, can be reliably determined over a range in redshift, and are physically meaningful.
Two of these parameters are photometric, derived from existing ground-based imaging and estimated $`k`$-corrections: rest-frame color $`BV`$, and absolute blue luminosity $`M_B`$. Two are image structure parameters, derived from multi-aperture photometric analysis presented below: physical half-light radius $`R_e`$, and image concentration $`C`$. One is a combined photometric-structural parameter: average rest-frame surface brightness $`SB_e`$ within $`R_e`$. Of the three parameters luminosity, half-light radius, and surface brightness, any one can obviously be derived from the other two. (We consider all three since in any given range of, e.g., luminosity, there is significant dispersion in both $`SB_e`$ and $`R_e`$.) The sixth parameter, a 180-rotational asymmetry index ($`A`$), utilizes the multi-aperture photometry indirectly through definition of the extraction radius for rotation; we refer to $`A`$ as a structural parameter. Table 2 contains all individual measurements for the Frei et al. sources. Luminosities and all image-structure parameters are measured in the rest-frame $`B`$ band.
### 3.1 Photometric parameters: restframe color and luminosity
While the Frei et al. (1995) data set contains $`B_J`$ images for 75% of the sample and $`g`$ band images for the remaining objects, there is no blue bandpass in which observations are available for all galaxies (see Table 1). We have made a comparison of the apparent (uncorrected) $`B`$ magnitudes listed in the Third Reference Catalogue of Bright Galaxies (\[de Vaucouleurs et al. 1991\], RC3) to those derived from our photometry of the Frei et al. images, appropriately transformed to the $`B`$-band using the tabulated corrections of Frei and Gunn (1994). This comparison shows that while the two magnitude estimates do not differ in the mean, there is a 0.25 mag (rms) scatter. To avoid the uncertainty associated with SED-dependent color transformations (see also §A.3) we use the RC3 uncorrected $`B`$ magnitudes and $`BV`$ colors instead of the values from our own photometry.
We apply $`k`$-corrections and corrections for galactic extinction to the $`BV`$ colors and apparent $`B`$ band magnitudes of the nearby galaxy sample in the manner described in RC3. The heliocentric velocities $`v_{hc}`$ of the galaxies are small (no greater than 3000 km s<sup>-1</sup> for any object); the average velocity is $`v_{hc}1000`$ km s<sup>-1</sup>. Hence the associated $`k`$-corrections are $`<0.05`$ mag. We use the distances given in the Nearby Galaxies Catalogue (\[Tully 1988\]), recalculated to $`H_0`$=50 km s<sup>-1</sup> Mpc<sup>-1</sup>, to derive absolute magnitudes $`M_B`$ from the corrected apparent magnitudes. Note we do not correct for internal extinction since the suitability and procedure for applying such corrections may be ill-defined for higher-redshift galaxies.
### 3.2 Multi-aperture photometry
To characterize the light distributions of the galaxies, we performed multi-aperture photometry on all images. The apertures are centered at the intensity-weighted centroid of each object. Since much of the profile shape information is contained in the central parts of the image, logarithmically spaced apertures are used. For the photometry of the Hubble Space Telescope ($`HST`$) images in paper II, the smallest aperture corresponds to 0.<sup>′′</sup>05, the largest to 15<sup>′′</sup>; for the nearby galaxy photometry here, the aperture radii are scaled to correspond to similar linear sizes. The apertures are circular to accommodate the irregular morphology of the intermediate redshift galaxies in paper II that would be difficult to fit with another geometrical figure. The efficacy of this approach is addressed in more detail below and in the Appendix (§A.3).
### 3.3 Structural parameters
Image structure is most commonly quantified via bulge–disk decomposition, yielding a bulge-to-total ratio, $`B/T`$. We refrain from this approach here, for reasons which we alluded to in the Introduction. For example, $`B/T`$ parameter may be poorly defined for asymmetric and compact galaxies. Irregularities in the surface brightness profiles, which can be caused by asymmetric structure, rings, or lenses, also cause problems for bulge–disk decompositions. While Kent showed that the concentration parameter correlates well with the bulge-to-total ratio, this holds only for objects with $`B/T<0.63`$. At larger values of $`B/T`$, bulge-disk decomposition fails for several objects in Kent’s sample, resulting in galaxies of type S0 – Sa being given extremely high values of $`B/T`$. Bulge–disk decomposition also becomes unreliable when galaxy disks are fainter than the bulges.
It is worth noting again that these problems mainly arise from older, one-dimensional methods of decomposition. The newer two-dimensional decomposition techniques are clearly successful at reproducing the observed light profiles, with remarkably small residuals (Schade et al. 1995, 1996; \[Simard 1998\]; \[Marleau & Simard 1998\]). Still, there are physical situations where bulge–disk decomposition techniques in general become problematic, namely where the astrophysical reality is more complex than simple bulge–disk models. Some galaxies have central condensations better described by an exponential profile rather than an $`r^{1/4}`$-law (\[Wyse et al. 1997\]); many galaxies have strong bi-symmetries, such as bars; virtually all galaxies have varying degrees of asymmetry due to star formation, dust, or large-scale gravitational perturbations and lopsidedness. All of these features represent details that decomposition into bulge and disk components do not address correctly. Simple disk and bulge decomposition is also inadequate for disk galaxies where the luminosity profile deviates from a pure exponential (\[Freeman 1970\]), e.g. type I and type II disks. (Type I disk profiles have an added component which contributes to the light just outside the bulge region; the surface brightness of a type II profile shows the opposite behavior (an inner truncation), and drops below the level of an exponential profile in the region near the center.)
Given the astrophysical complexity of real galaxies, the physical interpretation of the derived model parameters of disk-bulge fits remains uncertain. Nonetheless, such profile-fitting methods should be useful for estimating non-parametric structural and photometric parameters (e.g. characteristic sizes, surface brightness, image concentration, and ellipticity) in a way that uses the data in an optimal manner. In the current effort, however, we have taken a completely non-parametric approach of measuring sizes, surface-brightness, image concentration and asymmetry using multi-aperture photometry rather than deriving a model-dependent $`B/T`$ parameter.
#### 3.3.1 Half-light radii and surface-brightness
We define first our working definition of a total magnitude since it represents the critical zeropoint for measurement of the half-light radius and surface-brightness. We use the dimensionless parameter $`\eta `$ to define the total aperture of the galaxies – a limiting radius which is not based on isophotes.<sup>2</sup><sup>2</sup>2Isophotal radii introduce redshift-dependent biases unless careful consideration and corrections are made for dimming due to the expansion ($`(1+z)^3`$ in broad-band photon counts) and $`k`$-corrections. While such redshift-dependent biases are not an issue for the samples studied in this paper, in future papers in the series this would be an issue were we not to avoid isophotes. The concept of defining the size of a galaxy based on the rate of change in the enclosed light as a function of radius was first introduced by Petrosian (1976). In terms of intensity, $`\eta `$ can be defined as the ratio of the average surface brightness within radius $`r`$ to the local surface brightness at $`r`$ (\[Djorgovski & Spinrad 1981\]; \[Sandage & Perelmuter 1990\]). Like Wirth et al. (1994), we follow Kron’s (1995) suggestion to use the inverted form, $`\eta (r)I(r)/I(r)`$, which equals one at the center of the galaxy and approaches zero at large galactic radii. The radius $`r(\eta =0.5)`$ corresponds roughly to the half-light radius $`r_e`$.
Since $`\eta `$ is defined as an intensity ratio, it is not affected by the surface brightness dimming effect that makes the use of isophotes problematic. Moreover, $`\eta `$ is only dependent on the surface brightness within a given radius and not on any prior knowledge of total luminosity or the shape of the light profile. These properties make it advantageous for faint object photometry. We defined the “total” aperture of the intermediate-redshift objects as twice the radius $`r(\eta =0.2)`$. The apparent total magnitudes are then defined within this aperture. For ideal Gaussian or exponential profiles, the magnitude $`m_{0.2}`$ within the radius $`2r(\eta =0.2)`$ is approximately equal to the true total magnitude $`m_{tot}`$; more than 99$`\%`$ of the light is included with the radius $`r(\eta =0.2)`$. For an $`r^{1/4}`$-law profile, there is a difference $`m_{0.2}m_{tot}0.13`$ mag; this is due to the slow decline in luminosity at large radii that characterizes this profile. The radius $`r(\eta =0.2)`$ was chosen based on visual inspection of the curves of growth, derived from the aperture photometry, out to large radii.
For reference, the theoretical value for the ratio of $`r(\eta =0.2)`$ to half-light radius is 2.16, 1.95, and 1.82 for three standard profiles: Exponential, Gaussian, and $`r^{1/4}`$-law, respectively. The observed ratio is 2.3 $`\pm `$ 0.3 for $`BV`$$`<0.85`$ (with little trend with color), but rises slightly (2.6 $`\pm `$ 0.25) for the reddest galaxies with $`BV`$$`>0.85`$. A contributing cause to this rise is that for about half of the reddest objects, $`r_{1/2}`$ has been underestimated by $`20\%`$ because of their higher ellipticity. As we show in the Appendix (§ A.3.2), the half-light radii of early-type galaxies with axis ratio $`a/b>2`$ are systematically underestimated by up to 30$`\%`$. This effect will also cause small changes to the measured image concentration (§ 3.3.2) of these galaxies.
A weak downward trend can be seen from blue towards red colors; this is what we expect since bluer objects tend to have exponential luminosity profiles, and redder objects are better described by $`r^{1/4}`$-law profiles. However, this trend is broken by the reddest objects ($`BV>0.87`$), which have higher values of $`r(\eta =0.2)/r_{1/2}`$ than what is expected for an $`r^{1/4}`$-law profile.
Finally, the angular half-light radii $`r_e`$ were determined from the normalized curves of growth. Based on $`M_B`$ and (corrected) $`R_e`$ we calculated the photometric-structural parameter $`SB_e`$, the average blue surface brightness within the half-light radius, for all objects. For the nearby galaxy sample, the Tully catalog distances (as described in § 3.1) were used to determine $`R_e`$ (kpc).
#### 3.3.2 Image concentration
We use the image concentration parameter $`C`$ as defined by Kent (1985), which is based on the curve of growth. This parameter was shown to be closely correlated with Hubble type for “normal” galaxies:
$`C5log(r_o/r_i)`$
In the above equation, $`r_o`$ and $`r_i`$ are the outer and inner radii, enclosing some fraction of the total flux. In contrast, the concentration parameter defined by Abraham et al. (1994) is not based on curve of growth radii, but on a flux ratio within two isophotal radii.
However, in practice Kent also uses isophotes: He replaces the outer radius $`r_o`$, which encloses 80$`\%`$ of the total light, by the radius of the 24th mag/arcsec<sup>2</sup> isophote. He has demonstrated that this radius encloses $`79\%`$ of the total light for all galaxy types in the restrictive confines of his sample (\[Kent 1984\]). Because of the surface brightness dimming effect that becomes important for non-local galaxies, we instead use a method that is independent of isophotes. The total aperture of the galaxy, which determines the curve of growth, is defined based on the $`\eta `$-radius as described in § 3.3.1.
We have also explored the possibility of using $`\eta `$-radii to define a concentration parameter. However, a concentration parameter based on the curve of growth was ultimately found to be the more robust measure: the curve of growth increases monotonically with galactic radius for all objects, while the $`\eta (r)`$-function will be non-monotonic for a “bumpy” light profile (like that of a well-resolved spiral galaxy). As a consequence, image concentration defined by the curve of growth rather than $`\eta `$ exhibits less scatter when plotted against other correlated observables (e.g. color, surface-brightness) than an image concentration parameter based on the $`\eta `$-function.
Anticipating our need to measure image concentration for small galaxies in paper II and future papers in this series, we have studied the effects of spatial resolution and $`S/N`$ on $`C`$. Here we focus primarily on resolution, as this was the dominant effect. The importance of resolution is demonstrated by the comparison of Schade et al. (1996) of decompositions of compact objects in ground-based and $`HST`$ images: the cores of the blue nucleated galaxies are not resolved in ground-based imaging, and hence they are frequently misclassified as having much lower $`B/T`$-ratios than what is revealed by $`HST`$-imaging. In paper II we analyze this sample of galaxies, and hence this illustration is of particular relevance.
Resolution effects on image concentration were estimated by block-averaging the images of nearby galaxy sample over a range of values until the spatial sampling (as measured in pixels per half-light radius) was comparable to that of the compact galaxies at intermediate redshift observed with the WFPC2. The details of these simulations are presented in the Appendix (§A.1). In short, as the objects’ half-light radii get smaller, the scatter in the measured concentration indices increases. While larger inner radii or a smaller outer radii decrease this scatter (due to improved resolution and $`S/N`$, respectively), such choices decrease the dynamic range of the concentration index.
Based on these simulations, we chose a definition of $`C`$ that is, to first order, sufficiently robust to allow a direct comparison of the image concentration of the local and the higher-redshift samples studied here and in paper II, and furthermore gives a large dynamic range: $`C=5log(r(80\%)/r(20\%))`$. This concentration index is remarkably stable: The mean concentration does not deviate from that measured in the original image by more than 0.2, or $`8\%`$ of the dynamic range in $`C`$, down to resolution of five pixels per half-light radius.
Our definition is sufficiently close to that of Kent’s (1985) so that it is meaningful to compare our values directly to those he determined from photometric analysis of a sample of nearby galaxies. With this choice of radii, a theoretical $`r^{1/4}`$-law profile has $`C=5.2`$, an exponential profile has $`C=2.7`$, and a Gaussian has $`C=2.1`$. These values agree well with the results of Kent’s analysis: he finds that elliptical galaxies have $`C5.2`$, and late-type spirals have $`C3.3`$.
Lastly, since we use circular apertures, the measured image concentration may be affected by the ellipticity of the galaxy. Based on the comparison between our results for the Frei et al. sample and those of Kent’s elliptical aperture photometry, we believe this to be a negligible effect in all cases but the earliest, must elliptic galaxies. Wirth, Koo and Kron (1994) found that for an $`r^{1/4}`$ law profile with axis ratio $`b/a=0.2`$, the change in $`C`$ is less than 5$`\%`$. The effect appears to be larger in our study. A more detailed description of this possible systematic is given in the Appendix.
#### 3.3.3 Image asymmetry
The last image structure parameter is rotational asymmetry, $`A`$, as defined by Conselice, Bershady & Jangren (2000). This definition differs from earlier methods in that the asymmetry is determined within a constant $`\eta `$-radius of $`\eta =0.2`$, a noise correction is applied, and an iterative procedure which minimizes $`A`$ is used to define the center of rotation. This algorithm was tested to be robust to changes in spatial resolution and signal-to-noise by Conselice et al. (1999) using simulations similar to those described here for the concentration parameter $`C`$; the systematics with resolution are below 10$`\%`$ of the original value for galaxies in paper I and II here.
#### 3.3.4 Morphological $`k`$corrections
To obviate the issue of ‘morphological’ k-corrections, image structural parameters should ideally be measured at the same rest-frame wavelength for all objects. Anticipating our needs to derive the structural parameters for intermediate-redshift objects in paper II (and future papers in this series), we have adopted the following protocol: (i) For the nearby galaxy sample we use the images in the $`B_J`$ and $`g`$ bands to derive the primary local image structure parameters. (The rest-frame wavelengths sampled by the $`R,r`$ band images correspond to bands redshifted into the near-infrared for the intermediate-redshift galaxies.) (ii) We use the multi-band images of the Frei et al. sample to determine corrections to compensate for the wavelength dependence of asymmetry, concentration, and half-light radius – as described in the Appendix (§ A.2). For example, $`HST`$ Wide Field Planetary Camera-2 (WFPC-2) images in the $`I_{814}`$ band of objects between $`0.3<z<0.8`$ correspond to first order to the rest-wavelength range of the $`B_J`$ and $`g`$ bands. Nonetheless, the effective rest-wavelength for such intermediate-redshift galaxies is typically slightly redward of rest-frame $`B_J`$ and $`g`$ bands. The corrections in § A.2 are suitable for such samples, as well as higher redshift samples imaged in redder bands.
## 4 Results
### 4.1 Mean Properties, Distributions and Correlations
The mean properties for our six parameters (M<sub>B</sub>, $`BV`$, $`R_e`$, $`SB_e`$, $`C`$, and $`A`$) are listed in Table 3, as a function of Hubble Type. While we would like to move away from using ‘Hubble Types,’ they are so ingrained in the astronomical culture that they are a useful point of departure. For clarity in the following discussion, we group these types together into “Early” (E-S0), “Intermediate” (Sa-Sb), and “Late” (Sc-Irr). These names are potentially misleading, of course, and so we encourage the reader to treat them as labels which evoke, at best, a well-conceived galaxy type, but not necessarily an evolutionary state. Clearly further sub-division could be made, but our current purposes are illustrative, not definitive.
A typical approach to exploring the correlations in (and dimensionality of) a multivariate parameter space is principal component analysis. While this is valuable, it is not particularly instructive for a first understanding of the distribution of different types of objects in the parameter space. We are interested both in correlations between observables and in trends as a function of the qualitative Hubble-type. These correlations and trends need not be one and the same. For example, two observables can be uncorrelated but still exhibit a distribution segregated by Hubble type. To develop such an understanding, we therefore inspected the 15 possible 2-dimensional projections of our 6-dimensional parameter space.
To distill this information further, we considered that there are in fact three types of physically-distinct parameters:
1. spectral index (color): this parameter is purely photometric, by which we mean there is no information about the shape of the light profile. There is also no scale information, i.e. the amplitude and size of the light profile is also unimportant. In the balance of this paper we will use “color” and “spectral index” interchangeably.
2. form ($`A`$, $`C`$): these parameters are purely structural, by which we mean that they do not depend – to first order – on the amplitude or the shape of the spectral energy distribution, nor on the physical scale of the light distribution; they reflect only the shape of the light profile.<sup>3</sup><sup>3</sup>3We consider image concentration to be a form, in contrast to Morgan who used it as a surrogate for spectral index.
3. scale ($`R_e`$, $`L`$, and $`SB_e`$): these parameters are physically distinct. Luminosity is purely photometric (by our above definition). Size, which we also refer to as a structural parameter, is influenced by image shape, i.e., depending on the definition of size, two galaxies with different light profile shapes can have relatively different sizes (see §3.3.1, for example). Surface-brightness is a hybrid, photometric and structural parameter; it is a function of size and luminosity. While surface-brightness is a ratio of luminosity to surface area, it is still a measure of “scale” – in this case, the luminosity surface-density. <sup>4</sup><sup>4</sup>4A fourth scale parameter which we do not consider here is line-width, or some measure of the amplitude of the internal dynamics.
This reduces the types of combinations (by parameter-type) to 6, i.e. between color, form, and scale.
We find the strongest and physically most interesting correlations are between color, form, and the one scale parameter, $`SB_e`$ (Figures 3-5). We focus on these for the remainder of the paper. Before turning to them, for completeness we first summarize our observations of the other types of correlations:
Color-color correlations are strong and well known (e.g. Figure 2). Effectively they add higher-order information about spectral type. Here we consider only $`BV`$ as a simple spectral index which effectively represents the first-order information of spectral type. In general, one might adopt several spectral indices, e.g. $`UV`$ and $`VK`$, or a single index based on multi-colors.
Color-scale correlations also have been explored in detail elsewhere, e.g. color-luminosity relationships, known to exist for all galaxy types in both the optical and near-infrared (Huchra, 1977b; Mobasher et al. 1986; Bershady 1995). The limited dynamic range of the Frei et al. sample in size and luminosity (they are mostly large and luminous systems) preclude useful results being drawn here in this regard. For example, the correlation of color with size in this sample is subtle and depends in detail on how size is defined, as noted above. Form-scale correlations including size and luminosity are also difficult to assess for this sample for the same reasons of limited dynamic range in scale. However, scale versus scale is an interesting diagnostic because, for example, size and luminosity allow one to probe the range of surface-brightness in the sample. We explore this in paper II.
#### 4.1.1 Spectral index versus form and scale
Strong correlations exist in all three plots of color versus form parameters $`C`$ and $`A`$ and scale parameter $`SB_e`$ (Figure 3). Early-type galaxies are redder, more concentrated, high-surface-brightness, and more symmetric than Intermediate- and Late-type systems. The best correlation is between color and concentration in the sense that there is a smooth change in both quantities with Hubble Type. This is expected from a simple interpretation of the Hubble Sequence as a sequence parameterized by the relative dominance of a red, concentrated bulge (or spheroid) versus a bluer, more diffuse disk.<sup>5</sup><sup>5</sup>5A few of the local galaxies have values of $`C`$ that are lower than the theoretical concentration for an exponential disk (the errors in $`C`$ are $`<0.02`$ for all of them). The majority of these objects are late-type spiral galaxies with prominent, bright regions of star formation in the spiral arms. The star-forming regions cause the image profiles to become less centrally concentrated than a simple disk profile. In contrast, the distinction between Hubble Types in $`SB_e`$ and $`A`$ is most pronounced between Early-types and the remainder; Intermediate- and Late-types galaxies are not well distinguished by either of these parameters.
A more complete local sample will likely include a larger fraction of objects that do not follow these trends. For example, amorphous galaxies have surface brightnesses comparable to elliptical galaxies but are generally quite blue in color (\[Gallagher & Hunter 1987\]; \[Marlowe et al. 1997\]). Nonetheless, what is physically compelling about these color-form correlations is that each axis carries distinct information, respectively, on the integrated stellar population and its spatial distribution.
#### 4.1.2 Form versus form and scale
There are clear trends present in the two plots of form versus $`SB_e`$ (scale) in Figure 4 as well as the plot of form parameters along in Figure 5. More centrally concentrated galaxies have higher average surface-brightnesses and lower asymmetry; more symmetric objects have higher surface-brightness. In general, the concentrated, high surface-brightness galaxies are Early-type, while the Late-type galaxies are less-concentrated, have lower surface-brightness, and are more asymmetric. While there is substantial scatter in the form and scale parameters for Early and Late types, these two extreme groups still are well-separated in the above three plots. The Intermediate-type galaxies, however, are not well separated from these extremes, and tend to overlap substantially with the Late-type galaxies, consistent with what is found in plots of color versus form and scale: Intermediate- and late-type galaxies have comparable degrees of asymmetry, and similar surface-brightness.
One should be cautious in concluding the relative merits of form-scale and form-form and versus color-form and color-scale correlations based on the relative separation of Hubble Types. Using Hubble Types may be unfair if, for example, they were designed to correlate well with color but not necessarily with the quantitative form and scale parameters explored here. Since the form-form and form-scale correlations themselves are comparable, and nearly as strong as for color-form and color-scale, we are inclined to consider both as part of a general classification scheme. Certainly the form and scale parameters will each have different sensitivity to stellar evolution than color and so are advantageous to consider in isolation.
#### 4.1.3 Comparisons to previous work
The correlation between image concentration and mean surface-brightness within the effective radius (Figure 4) has been explored by several groups in the context of galaxy classification (\[Okamura et al. 1984\]; \[Watanabe et al. 1985\]; \[Doi et al. 1993\]; \[Abraham et al. 1994\]). We focus here, however, on Kent’s (1985) $`r`$-band study since his definition of image concentration and effective surface-brightness are the most similar to our own. While similar, nonetheless the slope of the correlation is steeper for our sample, albeit with much larger scatter, as illustrated in the top panel of Figure 6. As the middle and bottom panels reveal, the cause of the steeper slope in our sample is due to a smaller dynamic range in image concentration. This is likely due to the fact that we use circular apertures when performing surface photometry, whereas Kent used elliptical apertures. We attempt to quantify the systematics due to differences in aperture shape in §A.3. While the dynamic range in image concentration is reduced using circular apertures for the Frei et al. sample, there does appear to be a somewhat smaller scatter in $`C`$ as a function of $`BV`$.
The nature of the large scatter in the top two panels of Figure 6 for the Frei et al. sample is also discussed further in §A.3. In short, we believe much of this scatter is due to uncertainties in the $`R`$\- and $`r`$-band zeropoints of the Frei et al. sample. These uncertainties adversely affect only the surface-brightness values in Figure 6. Robust estimators of the scatter about a mean regression (i.e., iterative, sigma-clipping of outlying points) eliminate the outlying points, but still yield 50% larger scatter in $`R`$-band $`SB_e`$ for the Frei et al. sample as a function of either image concentration or $`BV`$. A plausible additional source contributing to this larger scatter is that Kent’s observed surface brightnesses were converted to face-on values, while ours were not “corrected” in this way. We conclude that if accurate and appropriate inclination corrections are possible to apply to all galaxies in a given study, this would be desirable. Since such corrections cannot be performed for the intermediate-redshift objects in paper II (and in general, if such corrections are not possible for a critical subset of the data), we believe it is best not make such corrections at any redshift.
The asymmetry–concentration plane has also been explored for galaxy classification purposes by, e.g., Abraham et al. (1994, 1996a) and Brinchmann et al. (1998). Our methods of measuring these parameters differ from theirs, and thus our quantitative results cannot be directly compared. However, a qualitative comparison to the $`AC`$ plot of Brinchmann et al. shows that both methods yield very similar results: the distribution of galaxies can be subdivided into sectors where early-type, intermediate-type, and late-type objects dominate. Brinchmann et al. also use the local sample from the Frei catalog to define these bins, but note however that the points they plot represent a sample of intermediate-redshift galaxies. The $`AC`$ correlation in the Brinchmann et al. diagram is not as clear as that seen here for the local sample in Figure 5; the scatter in their diagram is comparable to the dynamic range of the parameters. This is probably due to the different properties of the samples, rather than to the differences in how we determine the parameters. For a more direct comparison, we plot $`B`$ band asymmetry and concentration versus rest-frame $`BV`$ color for 70 galaxies from the Frei et al. sample (Figure 7), using both the $`A`$, $`C`$ values from this study and those found by Brinchmann et al. It can be seen that the distributions are overall quite similar; however, the separation in asymmetry of the different Hubble types is more apparent in this study, and the scatter in concentration is somewhat smaller. The conclusion here, then, is that our methodology offers typically modest, but sometimes significant improvements over previous work.
### 4.2 Classification
The above results point to how we can most effectively define a parameter demarcation to isolate, identify, and classify normal galaxies. In the four-dimensional parameter space of ($`BV`$, $`A`$, $`SB_e`$, $`C`$), we define boundaries (“cuts”) in the 6 two-dimensional projections between galaxies classified in the Hubble Sequence as Early/Intermediate and Intermediate/Late. These boundaries, selected by eye on the basis of the distribution of Hubble types, are listed in Table 4 and illustrated in Figures 3-5. Segregation by higher-dimensional hyper-surfaces are likely to be more effective (galaxies appear to be distributed on a ‘fundamental’ hyper-surface – the subject of a future paper), but the projected boundaries here are meant as illustrative, and practical for application when all of the parameters are not available. We stress that these boundaries are not definitive in some deeper physical sense. For example, in terms of formal Hubble Types cuts involving color are clearly “best;” however, as noted above, this may not be physically significant.
It would be uninteresting if all of the cuts provided the same classification. Moreover, one expects there will be discrepancies for objects near boundaries. We find that 49% of the sample matches in all cuts, while 64%, 87%, and 99% of the sample matches in at least 5, 4, or 3 cuts, respectively. (Hereafter, we refer to cases where 5 out of 6 cuts match as “5/6,” etc.) This degree of consistency seems reasonable so we have not tried to fine-tune the boundaries (such fine-tuning would not be sensible anyway since the details of the classification self-consistency are likely to be sample-dependent): The preponderance of objects are classifiable by a simple majority of the classifications based on the 6 cuts; 13% of the objects have a more ambiguous classification.
Of interest are the discrepancies within and between cuts in different combinations of color, form, and scale. We found that it is useful to group the six cuts into two groups of three. The first consists of the cuts in Figure 3 between color, form and scale, which we refer to as color-form/scale. The second consists of the cuts in Figures 4 and 5 between form and scale, which we refer to as scale/form-form. For example, 64% of the variance in the 5/6 cases comes from cuts in $`C`$$`SB_e`$, whereas cuts in $`C`$$`(BV)`$ and $`A`$$`C`$ are always consistent with the majority classification. More generally, scale/form-form cuts are internally mis-matched 40% of the time, while color-form/scale cuts are internally mis-matched only 21% of the time (and two-thirds of these color-form/scale mis-matches are also present in scale/form-form mis-matches). In other words, the color-form/scale cuts tend to be more consistent; much of the variance in the scale/form-form cuts again comes from $`C`$$`SB_e`$.
Only two galaxies pose a substantial problem for classification: NGC 4013 and NGC 4216. They are classified by various cuts to be in all categories (Early, Intermediate, and Late), and have no majority classification. However, both are highly inclined (4013 is edge on), which appears to give them unusual observed properties. Indeed, they are extreme outliers in several of the projections in Figures 3 and 4 (see also A.3.2 and figures therein). Hence such problem cases are likely to be easy to identify. Three other sources classified in all three categories (NGC 4414, 4651, and 5033) are not a problem: They have 4/6 consistent classifications. Two of these (NGC 4651 and 5033) have Seyfert nuclei, and are outliers only in plots with image concentration; they are highly concentrated for their color. NGC 4414 is not an outlier in any of the plots.
Finally, it is interesting to note that 23% of the sample has inconsistent majority classifications in color-form/scale versus scale/form-form cuts. This is true for 100% of the 3/6 cases, and 55% of the 4/6 cases. However, we believe this is for different reasons. In the latter cases (only) we find that the galaxies are predominantly at high inclination ($``$50% excess in the top half and top quartile of the sample distribution in inclination). Moreover, the color-form/scale classifications in these cases are all earlier than the majority scale/form-form classifications. We surmise this is due to the effects of reddening on $`BV`$.<sup>6</sup><sup>6</sup>6Inclination will also cause changes in other measured parameters. Changes in $`C`$, however, appear to be small (see §A.3.2). Surface-brightness will tend to increase at modest inclinations, and then decrease at high inclinations if a prominent dust lane obscures the bulge. Likewise, $`A`$ may increase due to a dust lane until the galaxy is directly edge-on. As a consequence of these changes and the distributions and cuts, $`C`$$`SB_e`$ tends to mimic the color-form cuts in the high-inclination cases, while $`A`$$`SB_e`$ and $`A`$$`C`$ do not. While the color-form/scale classifications tend to be earlier for the 3/6 cases, because there is no apparent inclination dependence, these differences are due likely to other physical effects. Two possibilities include low star-formation rates or high metallicity for galaxies of their form. Both of these conjectures are testable via spectroscopic observation.
We suggest then, as a practical, simple prescription, that the majority classification for all 6 cuts be taken as the classifier, except in the situation where the galaxy in question is highly inclined. In the latter case, the majority classification of the scale/form-form cuts should be adopted. When galaxies have only 3/3 consistent classifications, (13% of the Frei et al. sample), the adopted classifier should be intermediate between the two most common classifications. It also may be of interest to note if the color-form/scale and scale/form-form majority classifications differ. However, further elaboration based on these two-dimensional projections of a higher-dimensional distribution is not likely to be warranted.
#### 4.2.1 Discussion
We note that there are no distance-dependent scale parameters in our classification. By this we mean specifically that the classification parameters do not depend on knowledge of the distance modulus. Hence this classification is both quantitative and independent of the cosmological distance-scale and its change with cosmological epoch (i.e. no a priori knowledge is needed about H<sub>0</sub> or q<sub>0</sub>). The effects of the expansion do change the observed classification parameters. However, with knowledge of galaxy redshifts and judicious choice of “redshifted” photometric bands, surface-brightness dimming can be corrected and band-shifting either eliminated or corrected via the protocol described in the Appendix. Galaxy evolution, of course, will also modify the values of the parameters, but this is precisely the utility of the classification systems as applied to such a study: In what way do the parameters and their correlations evolve? How do the scale parameters change for a fixed range in classification parameters? These are issues which we intend to explore in subsequent papers in this series.
We also comment on the efficacy of using the four-dimensional parameter space of color, concentration, surface-brightness, and asymmetry for the classification of distant galaxies. As noted earlier, Abraham et al. (1996a) and Brinchmann et al. (1998) have explored the use of the asymmetry–concentration plane as a tool for distant galaxy classification. The use of the additional parameters of color and surface-brightness are clearly advantageous; they offer substantially more information, particularly as a diagnostic of the stellar population age and surface-density. The reasoning behind using $`A`$ and $`C`$ alone has been that to first order, they can be estimated without redshift information. Yet the wavelength dependence of both parameters (i.e., what is referred to as ‘morphological k-corrections’) can lead to measurement systematics. These systematics, if not corrected, in turn result in objects over a range in redshift to be systematically misclassified. For example, Brinchmann estimates that at $`z=0.9`$, $`25\%`$ of spiral galaxies are mis-classified as peculiar objects in the $`AC`$ plane. This fraction is expected to increase at larger redshifts. Hence, for high-$`z`$ studies of galaxy morphology, redshift information is crucial even when using asymmetry and concentration. Therefore, since redshift information is crucial no matter what, there is no reason not to use the four-dimensional classification we have outlined in future studies. The recent refinements and calibration of the technique of estimating redshifts photometrically make this all the more tractable.
Finally, we note that while the classification we have proposed here is practical and useful, there are five areas where we anticipate it can be improved or elaborated: (i) As we have mentioned before, the spectral-index parameter could have much greater leverage at distinguishing between different stellar populations by adding pass-bands that expand the wavelength baseline (e.g. the $`U`$ and $`K`$ bands in the near-UV and near-IR, respectively), or by increasing the spectral resolution (e.g. line-strengths and ratios). A further step of elaboration would be to explore spatially resolved spectral indices (gradients) and determine their correlation with form parameters. (ii) Internal kinematics should be considered. Ideally, the kinematic information would include estimates of both the random and ordered motion (rotation) so that the dynamical temperature could be assessed, in addition to the overall scale. Kinematics are relatively expensive to obtain (compared to images), but with modern spectrographs on large telescopes, the absolute cost is minimal at least for nearby galaxies. (iii) Higher-dimensional correlations are worthy of exploration to determine, for example, whether “fundamental” hyper-planes can adequately describe the entirety of the galaxy distribution. (iv) It is worth considering whether there are additional form parameters of value for classification that have not been included here. (v) The classification scheme needs to be tested against much larger, and more volume-representative samples of galaxies.
## 5 Summary and Conclusions
We have presented results from a study of the photometric and image-structural characteristics and correlations of a sample of local, bright galaxies (Frei et al. 1996). We find it illuminating to distinguish between parameters which characterize spectral-index (color), form (image concentration and asymmetry), and scale (size, luminosity, and surface-brightness). In this context, we arrive at the following main results and conclusions.
* We find that a combination of spectral-index, form and scale parameters has the greatest discriminatory power in separating normal Hubble-types. The strongest correlation is found between color and image concentration. However, there are equally strong correlations between form parameters (e.g. $`A`$ and $`C`$), but here the Hubble-types are not as well distinguished. As an indicator of classification utility, we suggest that the strength of the correlation between parameters is likely more important than the separation of Hubble Types within the correlation.
* It is possible to define a quantitative classification system for normal galaxies based on a four-parameter sub-set of spectral-index, form and scale: rest-frame $`BV`$ color, image concentration, asymmetry, and average surface brightness within the half-light radius. We propose a specific classification that distinguishes between “normal” galaxies as Early, Intermediate, and Late based on cuts in these four parameters. The classification is successful for 99% of the Frei et al. sample. Nonetheless, we designate this as “preliminary” until larger, more comprehensive samples of galaxies are needed than analyzed in the present study.
* Distance-dependent scale parameters are not part of this preliminary classification.
* These classification parameters can be measured reliably over a broad range in $`S/N`$ and image resolution, and hence should be applicable to reliably distinguishing between a wide variety of galaxies over a large range in redshift.
* Redshift information is needed to estimate reliably both the photometric properties (rest-frame color and surface brightness) as well as the structural parameters asymmetry and concentration at a fixed ($`B`$-band) rest-frame wavelength. In terms of redshift independence, asymmetry and concentration alone thus offer no advantages over the additional parameters classifiers proposed here. Indeed, incorporating the full suite of parameters defined here is advantageous for the purposes of classification.
The authors wish to thank Greg Wirth for his highly-refined algorithm for calculating $`\eta `$-radii user here and in paper II; our collaborators David Koo and Rafael Guzmán for their comments on the manuscript and input on assembling dwarf galaxy samples; and Jarle Brinchmann for providing us with his measurements of structural parameters for the local galaxy sample. We also gratefully acknowledge Jay Gallagher and Jane Charlton for a critical reading of the original version of this paper (I and II), and for useful discussions on this work. Most importantly, we thank Zolt Frei, Puragra Guhathakurta, James Gunn, and Anthony Tyson for making their fine set of digital images publicly available. Funding for this work was provided by NASA grants AR-07519.01 and GO-7875 from STScI, which is operated by AURA under contract NAS5-26555. Additional support came from NASA/LTSA grant NAG5-6032.
APPENDIX
## A A. Corrections for Measurement Systematics
Here we establish the measurement systematics due to changes in image resolution for half-light radius and image concentration, and for band-shifting effects on half-light radius and image concentration, and asymmetry. Systematic effects of images resolution and noise on asymmetry are quantified in Conselice et al. (1999).
### A.1 A.1. Resolution dependence of observed size and image concentration
To maximize the dynamic range of the measured concentration index, $`C`$, the inner radius should be small, and the outer radius large relative to the half-light radius. In this way, one samples the light profile gradients in both the central and outer regions of a galaxy where the bulge and disk contribute quite differently. This strategy maximizes the leverage for discriminating between different profiles, e.g. exponential and $`r^{1/4}`$-law. In the presence of noise and limited spatial resolution, however, the choice of radii determines the robustness of the concentration index: As noted by Kent (1985), the inner radius should be large enough to be relatively insensitive to seeing effects, and the outer radius should not be so large that it is affected by uncertainties in the sky background and $`S/N`$. In the current study, the sources are resolved, and the images are at moderately high signal-to-noise: within the half-light radius, the sample of local galaxies have $`600<S/N<3000`$. The intermediate redshift galaxies in paper II have $`S/N`$ in the range 40 to 90, with a mean of $`55`$. This is sufficiently high that we focus our attention here on the effects of spatial sampling and resolution.
Even in the absence of significant image aberrations, an additional limiting factor is the number of resolution elements sampling the inner radius. This is likely to become a limiting factor when the half-light radii is only sampled by a few pixels. To understand this potential systematic, we have calculated six concentration indices $`C=5log(r_o/r_i)`$ for several different choices of inner and outer radii. We use $`r_i`$ enclosing 20 and 30$`\%`$ of the light, and $`r_o`$ enclosing 50, 70, an 80$`\%`$ of the light. The radii were measured for nearby galaxies that were block-averaged by factors 2, 4, and 6 to simulate coarser spatial sampling, as shown in Figure 8. The six different concentration indices are plotted as a function of sampling in Figure 9. These simulations span sufficient dynamic range in size to cover most galaxies observed, for example, in the Hubble Deep Field. With factors of 4 and 6, we measure radii with pixel sampling similar to that observed in the HST/WFPC-2 images of the intermediate-redshift objects of paper II. Typically, these galaxies have half-light radii of 0.3 – 0.7<sup>′′</sup>. For the Planetary Camera, the scale is 0.046<sup>′′</sup>/pixel, and for the Wide Field, 0.10<sup>′′</sup>/pixel; hence the half-light radii are of order 3 to 15 pixels.
The half-light radius $`r(50\%)`$ is remarkably stable, even with poor sampling. Unfortunately, the dynamic range given by concentration indices with $`r_o=r(50\%)`$ is too small to be useful. As expected, the 30$`\%`$ radius was more stable than the 20$`\%`$ radius to decreased spatial resolution. However, the concentration indices using $`r(30\%)`$ were less sensitive to the differences between galaxy types, and gave a smaller dynamic range than indices using $`r(20\%)`$. The inner radius dominated the effect on the amplitude of the systematics; changing the outer radius from $`70\%`$ to 80$`\%`$ decreased the scatter only marginally. With a block-averaging factor of 6, where the half-light radii are typically only $`5`$ pixels, the scatter becomes large for all choices of concentration indices.
Based on these simulations, we decided to use the radii enclosing 80$`\%`$ and 20$`\%`$ of the total light (as did Kent, 1985), even though $`r_o=r(70\%)`$ gives concentration indices with slightly smaller scatter at poor resolution. For objects with half-light radii of only 7 pixels, the mean differences in concentration (relative to the original image) are $`\mathrm{\Delta }C_{80:20}=0.10_{0.60}^{+0.20}`$ and $`\mathrm{\Delta }C_{70:20}=0.10_{0.50}^{+0.15}`$. Even at a resolution of only five pixels per half-light radius, the concentration index only deviates by 0.2 relative to the original image; this is $`8\%`$ of the dynamic range in $`C_{80:20}`$. Thus we consider this parameter to be robust enough to useful in the comparison of local and intermediate-redshift samples.
### A.2 A.2. Systematics with wavelength
Observations at different wavelengths sample preferentially different stellar populations in a galaxy. Since these populations are not always spatially homogeneous, the image-structural characteristics (concentration, asymmetry, and half-light radius) will will have some wavelength dependence \[e.g., see de Jong’s (1995) study of disk scale-lengths\]. Hence, when comparing one of these parameters for different galaxies, the parameter ideally should be measured at the same rest-frame wavelength for all objects. This is not possible in general for studies over a wide range in redshifts employing a finite number of observed bands. To determine the amplitude of the wavelength-dependence for the measured structural parameters, we therefore compare the $`B_J`$ and $`R`$ structural parameters for 72 of the Frei et al. galaxies. The differences between the red and blue structural parameters versus the rest-frame color, $`BV`$, are shown in Figure 10. For comparison, all intermediate-redshift objects in paper II, except two, fall in the bluest bin ($`BV<0.62`$).
The plot of $`\mathrm{\Delta }C=C_BC_R`$ shows that in most cases, the values are slightly negative, i.e. the majority of objects are more highly concentrated in the red band than in the blue, as expected because of the redness of the central bulge. Only the bluest galaxies have comparable image concentration in both bands. There is a weak trend towards more negative values for the redder (early-type) objects, which also show a larger scatter than the bluer objects.
In the plot of $`\mathrm{\Delta }A=A_BA_R`$ it is clear that most galaxies have at positive values, i.e. their image structures are more asymmetric in the blue band than in the red, as shown by Conselice (1997). The difference in asymmetry is only seen for late- and intermediate-type objects; red objects are generally very symmetric in both bands, and have $`A_BA_R0`$. This trend was also noted by Brinchmann et al. (1998).
The plot of half-light radii ($`\mathrm{\Delta }R_e=R_{e,B}R_{e,R}`$) shows that most values are slightly positive, with a larger scatter for redder objects. No other trend with color is seen. The fact that the galaxies have slightly larger half-light radii in the blue band is consistent with their image concentration being higher in the red band, as a bulge profile generally has a much smaller scale length than an exponential profile.
In summary, the average differences ($`\pm 1\sigma `$) between parameters for galaxies with $`BV<0.62`$, determined from the $`B_J`$ and $`R`$ bands, are: $`\mathrm{\Delta }C=0.15\pm 0.30`$, $`\mathrm{\Delta }A=0.013\pm 0.044`$, and $`\mathrm{\Delta }R_e=0.21\pm 0.80`$ kpc.
#### A.2.1 A.2.1. Corrections for wavelength systematics
Based on the mean values above, we correct the measured structural parameters for galaxies at non-zero redshift to the rest-frame $`B`$ band values as follows, where for clarity we use the intermediate redshift galaxies in paper II as an example. The structural parameters of these intermediate-redshift objects generally were measured at rest-frame wavelengths between $`B_J`$ and $`R`$, i.e., in the observed WFPC2 $`I_{814}`$ band for $`z0.6`$. For “normal” galaxies, this would cause us to overestimate $`C`$, and underestimate $`A`$ and $`R_e`$. Hence we use the differences listed above and the redshift of the objects to linearly interpolate the correction to the measured values. Specifically, for a given parameter and color bin, we use the mean difference between values measured in the $`B_J`$ and $`R`$ bands, and the position of the rest-frame wavelength relative to the $`B_J`$ band, to make corrections to the measured parameters. (Note that the correction made to $`R_e`$ also affects the value of $`SB_e`$ in general.) For some objects, the combination of observed band-pass and redshift corresponds to rest-frame wavelengths slightly blueward of $`B_J`$. When computing the corrections for these objects, we assumed that the wavelength trends continue outside the $`B_J`$$`R`$ wavelength range. Overall, these corrections are small for objects in paper II, while for higher-redshift objects we expect band-shifting effects to become increasingly important.
We add a final, cautionary note that it is not certain the corrections for intermediate-redshift objects should be made based on the correlations we see for the nearby sample. When comparing the observations in the bluer bands ($`B_{450}`$ or $`V_{606}`$) to those in, e.g. the $`I_{814}`$ used in paper II, we find that most objects are more concentrated in the blue band, and slightly larger in the red band – this is the opposite of what we see for the Frei et al. sample.<sup>7</sup><sup>7</sup>7Indeed, Huchra noted that the Markarian galaxies get bluer toward their centers, reminiscent of the blue “bulges” seen in the blue nucleated galaxies of paper II, yet in contrast to the color gradients found for “normal” galaxies. This type of color-aperture relation was also noted by de Vaucouleurs (1960, 1961) for the latest Hubble-type galaxies (Sm, Im). For asymmetry, the trend is the same for both samples (higher $`A`$ in bluer bands). The trends are not directly comparable, however, to what we see in the local sample, as the observations in the bluer bands correspond to rest-frame wavelengths in the UV region for most intermediate-redshift galaxies. For this reason, and since the small sample of intermediate-$`z`$ objects poorly defines the variation in image structure with wavelength, we adopt the more well-determined trends seen for the Frei et al. sample to calculate the band-shifting corrections. These corrections based on local galaxy trends tend to make the intermediate-redshift objects somewhat less “extreme”; their half-light radii become larger, their surface brightnesses fainter, and their image concentrations lower. If instead we had based our corrections on the trends seen within the intermediate-$`z`$ sample of paper II, then this sample would be even more extreme relative to the local galaxy sample. The corrections would then tend to shift the positions of the intermediate-$`z`$ objects even farther from the nearby galaxies in diagrams that include any of the parameters $`R_e`$, $`SB_e`$, and $`C`$.
### A.3 A.3. Systematics with aperture shape
#### A.3.1 A.3.1. Comparison to elliptical aperture photometry
Circular-aperture surface-photometry will yield systematic differences in the measured structural parameters when compared to those derived from elliptical-aperture surface-photometry. To assess this, we compared our results for the Frei et al. catalog in $`R`$ and $`r`$ bands to the results of Kent (1985) for a sample of local, Hubble-type galaxies (Figure 6). Kent used elliptical apertures tailored to fit the axis ratio and position angle of each isophote in galaxy images to determine $`r`$ band image concentration and average surface brightness within the half-light radius.
As we detail in the figure caption, we have attempted to transform all of the surface-brightness values to the Cousins $`R`$ band ($`R_c`$). For each of the relations in Figure 6 we have characterized the slopes and scatter about a mean regression using a simple linear, least-squares algorithm with an iterative, sigma-clipping routine to remove outlying points. Given the nature of the data, such an algorithm is not statistically correct (see, e.g., Akritas and Bershady, 1996). However, given the potential photometric uncertainties (discussed below) and the need for robust estimation, it is not possible to formally implement more appropriate algorithms. Nonetheless, the relative characterization of the slopes and scatter between Frei et al. and Kent samples is useful.
As discussed in §4.1.3, the slope of the correlation between average surface-brightness and image concentration is steeper for our study than for Kent’s because of a decreased range in image concentration in our study. The effect (bottom panel of Figure 6) is such that the bluest galaxies have comparable image concentration values in both studies while the image concentration of the reddest galaxies differ by as much as 1 unit in the mean (Kent’s values are larger). We interpret this as likely to be the effect of different aperture shapes. The results of our study of systematics with axis ratio (below) support this conclusion. Surprisingly, there is no indication that elliptical apertures give significantly different results than circular apertures for intermediate- and late-type (disk dominated) galaxies.
The larger scatter in the Frei et al. (1996) sample in the top two panels of Figure 6 might lead one to conclude that the elliptical apertures provide a superior measurement of effective surface-brightness. However, much of the scatter is due to the subset of the Frei et al. sample observed at Palomar Observatory. We believe that zeropoint problems are the cause of much of this scatter, consistent with discussion in Frei et al. concerning the difficulty of photometric calibration. The bulk of the objects observed at Lowell Observatory are consistent with independent $`R_c`$-band photometry from Buta and Williams (1996), although there are some points that are very discrepant. In general, the overlap is excellent in $`SB_e`$ and $`BV`$ between the Kent sample, the Lowell subset of the Frei et al. sample, and the subset of the Frei et al. sample with Buta and Williams’ photometry.
#### A.3.2 A.3.2. Systematics with axis ratio
A second approach to determine the systematic effects of aperture shape on measured structural parameters was also used: we quantify the degree to which “normal” galaxies with the same intrinsic morphology but with different axial ratios $`a/b`$ will have different $`C`$ when measured with circular apertures. The galaxies in the Frei et al. catalog were divided into early-, intermediate- and late-type objects (using the same bins as elsewhere in the paper), and we plot concentration and half-light radius versus the logarithm of the axis ratio (taken from the RC3 catalog).
In the image concentration plots (Figure 11), a weak trend can be seen for the late-type galaxies (top panel), with slightly higher values of $`C`$ for the more inclined objects. This effect, if caused by the shape of the apertures, will lead us to overestimate the image concentration by at most $`0.1`$ ($`3\%`$) for the nearly edge-on galaxies. We do not expect this to be a problem for our analysis. The two labeled objects have unusually high values of $`C`$ for their morphological type. One of them, NGC 5033, is known to be a Seyfert 1 galaxy; the other, NGC 4651, is a suspected “dwarf-Seyfert” galaxy (Ho et al. , 1997). For intermediate-type objects (middle panel), no trend is observed. The lowest $`C`$-value, which belongs to NGC 4013, could be caused by the prominent dust lane in this object: the central light distribution is divided into two parts, making it difficult to determine the position of the center. Effects like these will likely be more problematic for objects with high values of $`a/b`$. The highest $`C`$-value in this plot is that of NGC 4216, which also is highly inclined and has spiral arm dust lanes superimposed on the bulge. In the bottom panel, a trend is observed for the early-type galaxies: the concentration is lower for objects with higher $`a/b`$ratio. This effect will cause us to underestimate the image concentration of these objects by $`0.5`$, or 10–15$`\%`$. This result agrees well with what was seen in the comparison of the Frei et al. sample to Kent’s image concentration measurements, as described above. This leads us to conclude that our circular aperture photometry will underestimate the image concentration somewhat for elliptical/S0 galaxies. Again, there is no indication that the aperture shapes lead to different results for intermediate- and late-type galaxies.
In the plots of half-light radius $`R_e`$ versus $`a/b`$ (Figure 12), no trends are seen for the intermediate- and late-type objects. For the early-type objects, however, the measured half-light radii become progressively smaller for increasing values of $`a/b`$. The trend is weak; it will cause us to underestimate the half-light radii by at most 30 $`\%`$ for objects with $`a/b4`$. If this effect is real, the derived surface brightness will be too bright by $`<0.7`$ magnitudes for the most highly elliptical early-type galaxies.
Figure captions
Fig. 1.— Rest-frame $`BV`$ versus $`M_B`$ for the nearby galaxy sample of Frei et al. (E-S0, Sa-Sb, and Sc-Irr) and Huchra’s (1977a) sample of normal Markarian galaxies (pluses). The dotted outline indicates the approximate locus of dE/dSph galaxies. The intermediate redshift samples from paper II are also plotted for comparison: blue nucleated galaxies (BNGs) compact, narrow emssion-line galaxies (CNELGs), and small, blue galaxies (SBGs). (The two SBGs and the BNG that we ultimately determine not to be “Luminous Blue Compact Galaxies” in paper II are shown as hatched symbols.) Only a few Markarian galaxies and late-type galaxies from the Frei et al. catalog share the extreme color–magnitude properties of the intermediate-redshift objects. In this plot, and in Figures 2-6, the vigorously star-forming galaxy NGC 4449 is labeled. Characteristic random errors are indicated separately for the Frei et al. sample and the intermediate-$`z`$ objects.
Fig. 2.— Rest-frame $`UB`$ versus $`BV`$ for the sample samples as in Figure 1. The intermediate-redshift samples of paper II largely overlap with the bluest Markarian galaxies, which extend blueward the color-color relation seen for the “normal” galaxies from Kent.
Fig. 3.— Rest-frame $`B`$-band form and scale parameters versus spectral index for the Frei et al. sample. Top panel: Average surface-brightness within the half-light radius ($`SB_e`$) versus rest-frame $`BV`$. Middle panel: Image concentration ($`C`$) versus $`BV`$. Bottom panel: 180-degree rotational image asymmetry ($`A`$) versus $`BV`$. Characteristic errors are given in the top-left corner of each panel. Outlying objects are labeled and discussed in the text. Dashed lines demark Early, Intermediate, and Late types in our classification scheme. Symbols are by Hubble type, as defined in the key. Different Hubble types are well distinguished, particularly in color. Morphological types are also well separated in $`C`$, but only the earliest types are well separated in $`SB_e`$ and $`A`$.
Fig. 4.— Rest-frame $`B`$-band parameters of form versus scale for the Frei et al. sample. Top panel: Image asymmetry ($`A`$) versus average surface-brightness ($`SB_e`$). Bottom panel: image concentration ($`C`$) versus $`SB_e`$. Outlying objects are labeled and discussed in the text. Dashed lines demark Early, Intermediate, and Late types in our classification scheme. The separation of morphological types is less clear than in Figure 3, but the different Hubble types are reasonably segregated.
Fig. 5.— Form versus form parameters for the Frei et al. sample: Rest-frame $`B`$-band image asymmetry ($`A`$) versus image concentration ($`C`$). Outlying objects are labeled and discussed in the text. Dashed lines demark Early, Intermediate, and Late types in our classification scheme. The separation of morphological types is less clear than in Figure 3, but is comparable to figure 5 where the different Hubble types are reasonably segregated.
Fig. 6.— Comparison of form, scale, and spectral index correlations between Frei et al. and Kent samples. Top panel: average $`R`$ band (Kron-Cousins) surface brightness within the half-light radius, $`SB_e(R_c)`$, versus $`R`$\- or $`r`$-band image concentration, $`C(R)`$. Middle panel: $`SB_e(R_c)`$ versus rest-frame $`BV`$. Bottom panel: $`C(R)`$ versus rest-frame $`BV`$. Structural parameters: We measured half-light radius and image concentration for the Frei et al. sample using their $`R`$ or $`r`$-band CCD images and circular photometry apertures. Kent measured these structural parameters using elliptical apertures on $`F`$-band CCD images. Photometric parameters: The Frei et al. sample is subdivided between objects observed at (a) Lowell Observatory (filled squares), (b) Palomar Observatory (dotted-circles), and (c) an overlapping subset of the Frei et al. sample with existing $`R_c`$-band photometry from Buta and Williams (1996; outlined-triangles). For (a) and (b) we used the zeropoints from the Frei et al. image headers (DNATO\_BV), and transformations from Thuan-Gunn $`r`$ and Gullixson et al. $`R`$ to Cousins $`R_c`$ from Frei & Gunn (1994). We have transformed Kent’s photometry reported in the Thuan-Gunn $`r`$-band to $`R_c`$ again based on transformations in Frei & Gunn (1994); Kent corrected surface brightnesses to “face-on” values. Regressions: Lines indicate $`\pm 1\sigma `$ about linear least-squares fits to the correlations (dotted, Kent; dashed, Frei et al. ) using an iterative clipping method ($`\pm 2.5\sigma `$ clip; 10 iterations). In the top and middle panels only the Lowell subset of the Frei et al. sample was used in the regressions. The substantial scatter in the Frei et al. $`SB_e(R_c)`$ values we infer is due primarily to zeropoint difficulties; we detect no noticeable systematics effects with inclination in $`SB_e(R_c)`$. The difference in the correlation between $`SB_e(R_c)`$ and $`C(R)`$ is largely due to the shallower trend in $`C(R)`$ with $`BV`$ for the Frei et al. sample. This may be due to differences between circular versus elliptical apertures. While elliptical aperture photometry provides greater dynamic range in $`C(R)`$, the correlation of $`C(R)`$ with $`BV`$ has larger scatter.
Fig. 7.— Form parameters and spectral index for 70 galaxies from the Frei et al. sample as determined by Brinchmann et al. Top panel: $`B`$ band image concentration versus rest-frame $`BV`$. Middle panel: $`B`$ band asymmetry versus rest-frame $`BV`$. Bottom panel: $`B`$ band asymmetry versus concentration. The asymmetry parameter was determined in a very similar manner as our own and thus should have comparable dynamic range. Since our $`C`$ parameter is logarithmic, we plot the logarithm of the Brinchmann et al. $`C`$ values. These plots are displayed so that they may be directly comparable to Figures 3 and 5. The trend in asymmetry for the different Hubble types is more apparent in Figure 3 and 5. In the concentration–color plane, the distributions are similar for both studies, although we find a smaller scatter among the late-type galaxies, and a larger scatter among the early-type objects.
Fig. 8.— A representative subset of galaxy images from the Frei et al. catalog, block-averaged by factors 1, 2, 4, and 6 (top to bottom). While the apparent change in qualitative (visually-assessed) morphology is small, the effects on the quantitative parameters $`C`$ and $`A`$ can be substantial. Half-light radius and surface-brightness are only weakly affected.
Fig. 9.— Resolution dependence of image concentration, $`C`$, for the galaxies in the Frei et al. catalog: $`\mathrm{\Delta }C`$ versus the half-light radius $`R_e`$ (in pixel units of the block-averaged images). $`\mathrm{\Delta }C`$ is the difference between the concentration index for a given simulated value of $`R_e`$ relative to the original concentration value (i.e. that value measured on the observed image). Measurements for six definitions of the concentration index are plotted (two types per panel, labeled by line-type). The central line (bold) is the median value of this difference, and the bounding lines are the 25$`\%`$ and 75$`\%`$ values, i.e. 50$`\%`$ of the simulations are contained between the upper and lower lines for each index.
Fig. 10.— Wavelength dependence of structural parameters for galaxies in the Frei et al. sample, plotted versus galaxy rest-frame color. Dashed lines show the mean differences between blue and red bands and the error bars show the 1$`\sigma `$ dispersions for three bins in color: $`BV<0.62`$ (late-type), $`0.62<BV<0.87`$ (intermediate-type), and $`BV>0.87`$ (early-type). Top panel: Image concentration $`C_BC_R`$. Nearly all galaxies are more highly concentrated in the red band than in the blue, and thus fall below the dotted line at $`C_BC_R=0`$. This difference is slightly larger for galaxies with intermediate-type morphology. Middle panel: Image asymmetry $`A_BA_R`$. Late- and intermediate-type galaxies are more asymmetric in the blue band than in the red band. Red objects are generally very symmetric in both bands, and have $`A_BA_R0`$. This panel can be compared to Figure 2 in Conselice et al. (1999) where $`A_BA_R`$ is plotted versus $`A_R`$. Since asymmetry and color are strongly correlated for the Frei et al. sample (as seen in Figure 3), the trend in Conselice’s plot is similar to what is shown here. Bottom panel: Half-light radius $`R_{e,B}R_{e,R}`$. Although the scatter in this diagram is relatively large, it is clear that the half-light radius shows little wavelength dependence over this wavelength range (cf. de Jong 1995). Objects of intermediate $`BV`$ color tend to be slightly larger in the blue band than in the red, but this trend is not seen for either the bluest or the reddest objects.
Fig. 11.— Axis ratio dependence of image concentration $`C`$ for galaxies of different morphological types in the Frei et al. sample. The dotted line separates the sample into two bins at log$`{}_{10}{}^{}[a/b]=0.3`$, corresponding to an inclination of 60. The dashed lines and error bars show the mean and 1$`\sigma `$ dispersion for each morphological type and bin. Labeled objects are discussed in the text. Top panel: Late-type objects with high inclination have slightly higher measured $`C`$ than more face-on objects. Middle panel: For intermediate-type galaxies the measured concentration indices show no correlation with the axial ratio $`[a/b]`$. Bottom panel: For early-type galaxies, a tendency can be seen where objects with larger axial ratio $`[a/b]`$ are measured to have lower image concentration.
Fig. 12.— Axis ratio dependence of half-light radius $`R_e`$ for galaxies of different morphological types in the Frei et al. sample. The dotted line separates the sample into two bins at log$`{}_{10}{}^{}[a/b]=0.3`$, corresponding to an inclination of 60. The dashed lines and error bars show the mean and 1$`\sigma `$ dispersion for each morphological type and bin. Top panel: The measured $`R_e`$ are slightly larger for late-type objects with high axial ratios. Middle panel: Intermediate-type objects have somewhat smaller $`R_e`$ for high values of $`[a/b]`$. In both of these panels the scatter is large and the differences between the bins are small. Bottom panel: Early-type galaxies with larger axial ratio $`[a/b]`$ are measured to have $`<30\%`$ smaller half-light radii. |
warning/0002/hep-ph0002133.html | ar5iv | text | # Spin polarisabilities of the nucleon at NLO in the chiral expansion
## A Full amplitude
The full amplitude in the Breit frame for diagrams 2a-2h are as follows. The notation $`t_i`$ is used for the tensor structures which multiply the amplitudes $`A_i`$.
$`T_a`$ $`=`$ $`{\displaystyle \frac{g^2e^2}{4m_Nf_\pi ^2}}[(m^2\omega ^2(1+\mathrm{cos}\theta )){\displaystyle \frac{J_0(\omega ,m^2)}{\omega }}2\omega J_0(\omega ,m^2)]t_3(\omega \omega )`$ (A1)
$`T_b`$ $`=`$ $`{\displaystyle \frac{g^2e^2}{2m_Nf_\pi ^2}}[{\displaystyle \frac{2m^2}{\omega }}(J_2^{}(\omega ,m^2)J_2^{}(0,m^2))t_3+(1+\mathrm{cos}\theta ){\displaystyle \frac{J_2(\omega ,m^2)}{\omega }}t_3`$ (A3)
$`\omega J_2^{}(\omega ,m^2)t_5+\omega (t_52(1\mathrm{cos}\theta )t_3){\displaystyle _0^1}dxJ_2^{}(x\omega ,m^2)](\omega \omega )`$
$`T_c`$ $`=`$ $`{\displaystyle \frac{g^2e^2}{2m_Nf_\pi ^2}}\tau _3\omega J_0(\omega ,m^2)t_3(\omega \omega )`$ (A4)
$`T_d`$ $`=`$ $`{\displaystyle \frac{g^2e^2}{m_Nf_\pi ^2}}\tau _3\omega {\displaystyle _0^1}𝑑xJ_2^{}(x\omega ,m^2)t_3(\omega \omega )`$ (A5)
$`T_e`$ $`=`$ $`{\displaystyle \frac{g^2e^2}{2m_Nf_\pi ^2}}(1\tau _3){\displaystyle \frac{1}{\omega }}\left(J_2(\omega ,m^2)J_2(0,m^2)\right)t_3(\omega \omega )`$ (A6)
$`T_f`$ $`=`$ $`{\displaystyle \frac{g^2e^2}{4m_Nf_\pi ^2}}\omega \left(2(\mu _v\mu _s\tau _3)(t_3\mathrm{cos}\theta t_4)+(1\tau _3)t_6\right){\displaystyle _0^1}𝑑x(12x)J_2^{}(x\omega ,m^2)(\omega \omega )`$ (A7)
$`T_g`$ $`=`$ $`{\displaystyle \frac{g^2e^2}{4m_Nf_\pi ^2}}\omega \left(2(\mu _v+\mu _s\tau _3)(t_3\mathrm{cos}\theta +t_4t_5)+(1+\tau _3)t_6\right){\displaystyle _0^1}𝑑xJ_2^{}(x\omega ,m^2)(\omega \omega )`$ (A8)
$`T_h`$ $`=`$ $`{\displaystyle \frac{g^2e^2}{m_Nf_\pi ^2}}\omega ^2{\displaystyle _0^1}dy{\displaystyle _0^{1x}}dx[((7x1)(t_6t_5)+7(1xy)t_4){\displaystyle \frac{J_6^{\prime \prime }(\stackrel{~}{\omega },m^2xyt)}{\stackrel{~}{\omega }}}`$ (A12)
$`+(2V(x,y,\theta )(xt_6xt_5+(1xy)t_4)`$
$`(1xy)(9xyxy)t_7)\omega ^2{\displaystyle \frac{J_2^{\prime \prime }(\stackrel{~}{\omega },m^2xyt)}{\stackrel{~}{\omega }}}`$
$`xy(1xy)\omega ^4V(x,y,\theta )t_7{\displaystyle \frac{J_0^{\prime \prime }(\stackrel{~}{\omega },m^2xyt)}{\stackrel{~}{\omega }}}](\omega \omega )`$
where $`\stackrel{~}{\omega }=(1xy)\omega `$,
$$J_6(\omega ,m^2)=\frac{1}{d+1}\left((m^2\omega ^2)J_2(\omega ,m^2)\frac{\omega m^2}{d}\mathrm{\Delta }_\pi \right),$$
(A13)
$`J_0(\omega ,m^2)`$, $`J_2(\omega ,m^2)`$ and $`\mathrm{\Delta }_\pi `$ have their usual meanings, prime denotes differentiation with respect to $`m^2`$, and
$$V(x,y,\theta )=(2xyxy+1)\mathrm{cos}\theta x(1x)y(1y).$$
(A14) |
warning/0002/hep-th0002218.html | ar5iv | text | # 1 Introduction
## 1 Introduction
There is a well-known construction of correlation functions in two decoupled critical Ising models as correlation functions in a $`c=1`$ free boson theory . A natural question to ask is how to extend this to the boundary model and how to extend the analysis of Cardy and Lewellen to include non-local bulk fields such as the Ising disorder field.
We first recall the Ising model field content. The maximal set of primary fields one normally considers are the identity field 11 of weight $`(0,0)`$; the spin field $`\sigma `$ and disorder field $`\mu `$ of weights $`(1/16,1/16)`$; the energy operator $`ϵ`$ of weight $`(1/2,1/2)`$; and the fermion fields $`\psi `$ and $`\overline{\psi }`$ of weights $`(1/2,0)`$ and $`(0,1/2)`$ respectively. These fields are not mutually local, and the three maximal sets of mutually local fields are
$$\{\text{1}\text{1},\sigma ,ϵ\},\{\text{1}\text{1},\mu ,ϵ\},\{\text{1}\text{1},\psi ,\overline{\psi },ϵ\}.$$
We shall take the point of view that our primary fields are the local set 11, $`\sigma `$ and $`ϵ`$ and we can then add the disorder field $`\mu `$ at the expense of introducing disorder lines joining pairs of disorder fields. The effect of the disorder lines is that the correlation function changes sign when a spin field passes through a disorder line; without these disorder lines the correlation functions containing both spin and disorder fields would be double-valued. We also recall that bulk correlation functions are related by duality, under which $`\sigma \mu `$, $`ϵϵ`$.
A conformally-invariant boundary condition on a conformal field theory the upper half plane defines a set of correlation functions of the bulk fields satisfying $`T(x)=\overline{T}(x)`$ on the $`x`$–axis. The boundary theory of the local field theory containing $`\sigma `$ and $`ϵ`$ was investigated by Cardy who found that (under certain assumptions such as the uniqueness of the vacuum) there are three such boundary conditions, denoted by ‘$`+`$’, ‘$``$’ and ‘$`f`$’, which have the interpretation that the spin variable in the lattice realisation is fixed up, fixed down and free, respectively.
The fields on the boundary can be classified into primary and descendent fields. Under Cardy’s assumptions, one finds that the identity is the only primary field on the $`\pm `$ boundaries, whereas on the free b.c. there is a non-trivial primary field $`\sigma ^B(x)`$ of weight $`1/2`$ which has the interpretation of the boundary spin field. One can further consider fields $`\psi ^{(\alpha \beta )}`$ which interpolate different b.c.’s $`\alpha `$ and $`\beta `$, and which again fall into primary and descendent fields; $`\psi ^{(\pm )}`$ have weight 1/2 and $`\psi ^{(\pm f)},\psi ^{(f\pm )}`$ have weight 1/16.
To complete the picture one needs the structure constants in the bulk-boundary opes (expressing the expansion of a bulk field in boundary fields) and the opes of boundary fields. For the Ising model these were found by Cardy and Lewellen , and the bulk-boundary opes take the form
$$\begin{array}{cccccc}\hfill \sigma (z,\overline{z})|_\pm & & {}_{}{}^{(\pm )}B_{\sigma }^{1}|2y|^{1/8}\hfill & \hfill ϵ(z,\overline{z})|_\alpha & & {}_{}{}^{(\alpha )}B_{ϵ}^{1}|2y|^1\hfill \\ \hfill \sigma (z,\overline{z})|_f& & {}_{}{}^{(f)}B_{\sigma }^{\sigma ^B}|2y|^{3/8}\sigma ^B(x)\hfill & & & \end{array}$$
(1.1)
where $`{}_{}{}^{(\pm )}B_{\sigma }^{1}=\pm 2^{1/4},{}_{}{}^{(f)}B_{\sigma }^{\sigma ^B}=2^{1/4},{}_{}{}^{(\pm )}B_{ϵ}^{1}=1,{}_{}{}^{(f)}B_{ϵ}^{1}=1.`$
A natural question is how to include the bulk disorder field in this analysis. Roughly speaking, fixed and free b.c.’s are interchanged under duality. However, as noted in , this leads to an apparent contradiction between duality and the (conjectured) $`g`$–theorem for boundary renormalisation group flows. This conjecture states that the boundary entropy $`g`$ in a unitary theory decreases along a renormalisation group flow . The values for the free and fixed b.c.’s are $`g_{\mathrm{free}}=1`$, $`g_{\mathrm{fixed}}=2^{1/2}`$. When the free b.c.’s are perturbed by a boundary magnetic field, there is a flow to the fixed b.c. However, under duality there should be a similar flow from fixed to free induced by the boundary disorder field, which would apparently be forbidden by the $`g`$–theorem. A resolution of this contradiction was obtained in in which it was shown that a careful consideration of the two local field theories leads one to the result that for the local theory containing $`\{\text{1}\text{1},\mu ,ϵ\}`$ one has three different boundary conditions, ‘fixed’, ‘free+’ and ‘free$``$’ which are dual to ‘free’, ‘fixed+’ and ‘fixed$``$’. In terms of these new b.c.’s, the flow is from ‘fixed’ with $`g=1`$ to ‘fixed$`\pm `$’ with $`g=2^{1/2}`$, in agreement with the $`g`$–theorem. Here we suggest another resolution adapted to our situation in which we wish to remain within a theory containg the spin field. As we see below, the boundary disorder field is locally identical to the operator interpolating fixed $`+`$ and fixed $``$ b.c.’s, and so one can formulate a perturbation by the boundary disorder field only in a Hilbert space describing the direct sum of $`+`$ and $``$ b.c.’s. Since the actual flow is now from a direct sum of $`+`$ and $``$ b.c.’s with $`g=\mathrm{\hspace{0.17em}2}2^{1/2}=\mathrm{\hspace{0.17em}2}^{1/2}`$ to free with $`g=1`$, $`g`$ decreases along this flow, and again there is no contradiction with the $`g`$–theorem.
Using duality, the bulk-boundary ope of the disorder field in the free b.c. should be dual to that of the spin field in fixed b.c.’s; i.e. in the free b.c. the bulk disorder field only couples to the identity,
$$\mu (z,\overline{z})|_{\mathrm{free}}{}_{}{}^{(\mathrm{free})}B_{\mu }^{1}|2y|^{1/8},{}_{}{}^{(\mathrm{free})}B_{\mu }^{1}=\pm 2^{1/4},$$
(1.2)
where the sign ambiguity reflects the overall ambiguity in correlation functions involving the disorder field.
However, when a disorder field approaches a free boundary, it carries its associated disorder line with it. Since the leading term in the ope of the disorder field with the boundary is the identity operator, this means that a disorder line can end on a free boundary at an operator of weight 0. This explains the observation in that the boundary state describing a circular boundary in the plane with free b.c. with a disorder line is rotationally invariant – one might at first sight think that the presence of the disorder line ending on the boundary would lead to a dependence on the position of the disorder line on the boundary, but since it couples to field of weight 0 this is not the case. Alternatively, the disorder line simply represents a cut in the value of correlation functions and has no physical content, it does not break the rotational symmetry of the correpsonding boundary state.
We shall see this directly in section 3 where we construct the one-point function of the bulk disorder field on the cylinder with free b.c.’s which interpolates the frustrated and unfrustrated partition functions as the disorder field crosses the cylinder.
In fixed b.c.’s conversely, the bulk disorder field will only couple to a boundary field of weight 1/2, the boundary disorder field (also called the boundary freedom field ) $`\mu ^B(x)`$,
$$\mu (z,\overline{z})|_{\mathrm{fixed}}{}_{}{}^{(\mathrm{fixed})}B_{\mu }^{\mu ^B}|2y|^{3/8}\mu ^B(x),{}_{}{}^{(\mathrm{fixed})}B_{\mu }^{\mu ^B}=2^{1/4},$$
(1.3)
We have denoted the b.c. by ‘fixed’ rather than by ‘$`\pm `$’ since the presence of the disorder fields means that it is not possible to unambiguously identify a fixed b.c. as either ‘$`+`$’ or ‘$``$’. The ability to move a disorder line across a boundary at no cost (other than a sign change when it passes through a spin field) is also true for the fixed b.c., thereby changing the b.c. from $`+`$ to $``$ or vice versa. Using this property we see that the boundary disorder field on a fixed b.c. can be identified with a $`(\pm )`$ boundary condition changing operator connected to a disorder line, as shown in figure 1
$$\begin{array}{ccc}\text{}& & \\ \text{Figure 1: the equivalence of the boundary disorder field and the boundary changing operator}& & \end{array}$$
One can ‘pull’ the disorder line attached to the bulk disorder field down and to the left until it crosses the edge of the upper-half-plane leaving the bulk disorder field attached to a disorder line that ends directly on the boundary at a field of weight 0. When the bulk field then approaches the boundary, this disorder line shrinks and locally the leading term in the bulk-boundary ope is the boundary-condition-changing operator interpolating $`\pm `$ b.c.’s. This shows that the boundary disorder field is locally identical to a boundary-condition changing operator, the only difference being the position of disorder lines.
## 2 Bosonisation formulae
As explained in , one can realise two independent copies of the Ising model with local fields $`\{1,\sigma ,ϵ\}`$ as the $`_2`$ orbifold of a single free boson compactified on a circle of radius $`r=1`$. If we label the fields in the two Ising models by 1 and 2, this is generated by
$$\{\text{1}\text{1},\sigma _1,ϵ_1,\sigma _2,ϵ_2\},$$
(2.1)
and is a local field theory; we shall call this model A. One can consider a different $`c=1`$ theory containing (amongst others) the following alternative mutually local set of fields,
$$\{\text{1}\text{1},\sigma _1\sigma _2,\mu _1\mu _2,ϵ_1+ϵ_2,ϵ_1ϵ_2\},$$
(2.2)
which we shall call model B. This is the $`_2`$ orbifold of model A, i.e. the un-orbifolded free boson and contains only symmetric combinations of the fields in the two Ising models, although not all – for example it does not contain $`\sigma _1+\sigma _2`$ which would be non-local with respect to $`\mu _1\mu _2`$.
The advantage of considering model B is that it enables one to construct correlation functions $`𝒪`$ in a single Ising model as the square root of the correlation function $`𝒪_1𝒪_2`$ in a free boson theory, as explained in . The correlation functions in model B are single-valued, but taking the square–root leads to the (possibly) multi-valued correlation functions in the Ising model.
The bosonisation formulae on the plane are
$$\begin{array}{cccccc}\hfill \sigma _1\sigma _2& =& \sqrt{2}:\mathrm{cos}\frac{1}{2}\phi :,\hfill & \hfill \mu _1\mu _2& =& \sqrt{2}:\mathrm{sin}\frac{1}{2}\phi :,\hfill \\ \hfill ϵ_1ϵ_2& =& :\phi \overline{}\phi :,\hfill & \hfill ϵ_1(z,\overline{z})+ϵ_2(z,\overline{z})& =& 2:\mathrm{cos}\phi :,\hfill \end{array}$$
(2.3)
Note that the field representing $`\mu _1\mu _2`$ is absent from the orbifold model of .
We now consider this model on a cylinder of length $`L`$, circumference $`R`$, which we take to be the rectangle in the complex $`z`$ plane with vertices $`0,R,iL,R+iL`$. This can be mapped to an annulus in the $`\zeta `$ plane by $`\zeta =\mathrm{exp}(2\pi iz/R)`$. In this form, the boundary conditions on the two ends of the cylinder are represented by boundary states in the full bulk theory.
All the possible boundary states for the model A were found by Oshikawa and Affleck in . To use their results, we first give our conventions. We take the mode expansion of $`\phi (\zeta ,\overline{\zeta })`$ on the plane to be
$$\phi (\zeta ,\overline{\zeta })=qip\mathrm{ln}(\zeta \overline{\zeta })iw\mathrm{ln}(\zeta /\overline{\zeta })+i\underset{n0}{}(\frac{a_n}{n}\zeta ^n+\frac{\overline{a}_n}{n}\overline{\zeta }^n),$$
where $`w`$ is the winding number, $`[q,p]=i`$ and $`[a_m,a_n]=[\overline{a}_m,\overline{a}_n]=m\delta _{m,n}`$. With this choice of $`\phi `$, compactification on a circle of radius 1 in the sense of means we identify $`\phi \phi +4\pi `$, and hence $`w`$ takes integer and $`p`$ half-integer values (n.b. this is different to ). For the boundary states we consider, we only need the Dirichlet-type Ishibashi states of zero winding number, viz.
$$|k=\mathrm{exp}(\underset{n>0}{}\frac{a_n\overline{a}_n}{n})|k,$$
(2.4)
which satisfy
$$(a_n\overline{a}_n)|k=0.$$
(2.5)
The boundary states in the orbifold model also contain states in the twisted sector, but these will play no role for us.
Although the space of states in models A and B differ, one of the boundary conditions of is common to the two models, and that is the condition that the spins on both models are ‘free’. We denote this by $`|\mathrm{free}`$ and it is given by
$$|\mathrm{free}=\underset{k}{}(1)^k|k.$$
(2.6)
To check that this does indeed correspond to free boundary conditions, one can calculate the partition function for the system with this boundary state on the two ends. Given the cylinder Hamiltonian
$$H(R)=\frac{2\pi }{R}(L_0+\overline{L}_01/12)$$
the partition function with this ‘free’ boundary condition on the both ends of the cylinder is
$`Z_{\mathrm{free},\mathrm{free}}`$ $`=`$ $`\mathrm{free}|e^{LH(R)}|\mathrm{free}={\displaystyle \underset{k}{}}\stackrel{~}{q}^{k^2/2}/\eta (\stackrel{~}{\tau })=\theta _3(0|\stackrel{~}{\tau })/\eta (\stackrel{~}{\tau })=\theta _3(0|\tau )/\eta (\tau )`$ (2.7)
$`=`$ $`(\chi _0(q)+\chi _{1/2}(q))^2=(Z_{\mathrm{free},\mathrm{free}}^{\mathrm{Ising}})^2,`$
where $`\stackrel{~}{q}=\mathrm{exp}(4\pi L/R)`$, $`q=\mathrm{exp}(\pi R/L)`$, $`\stackrel{~}{\tau }=2iL/R`$, $`\tau =iR/(2L)`$, $`\chi _h(q)`$ is the character of the $`c=1/2`$ Virasoro algebra representation of weight $`h`$, and the definitions and properties of the $`\theta `$ functions can be found in e.g. chap. 10 of .
Hence the state $`|\mathrm{free}`$ naturally describes free boundary conditions with no disorder lines (or an even number of disorder lines) ending on the boundary. To find the boundary state $`|\mathrm{free}^{}`$ for the system with an odd number of disorder lines, we can use the fact that a bulk disorder operator carries with it a disorder line and has the bulk–boundary ope
$$\mu (x+iy)|_{\mathrm{free}}2^{1/4}(2y)^{1/8}\text{1}\text{1}+\mathrm{}.$$
Hence we can calculate $`|\mathrm{free}^{}`$ by acting with the bosonic expression for $`\mu _1\mu _2`$ on $`|\mathrm{free}`$.
Using the standard expression for the normal ordered vertex operator on the cylinder expressed in terms of the field on the plane,
$$:e^{i\alpha \phi _{\mathrm{cyl}.}(z,\overline{z})}:=(\frac{2\pi }{R})^{\alpha ^2}|\zeta |^{\alpha ^2}e^{i\alpha \phi _<(\zeta ,\overline{\zeta })}e^{i\alpha q}|\zeta |^{2\alpha p}e^{i\alpha \phi _>(\zeta ,\overline{\zeta })},$$
where $`\zeta =\mathrm{exp}(2\pi iz/R)`$ and
$$\phi _{\genfrac{}{}{0pt}{}{>}{<}}(\zeta ,\overline{\zeta })=i_{n\genfrac{}{}{0pt}{}{>}{<}0}(\frac{a_n}{n}\zeta ^n+\frac{\overline{a}_n}{n}\overline{\zeta }^n),$$
one easily finds that
$$:e^{i\alpha \phi _{\mathrm{cyl}.}(z,\overline{z})}:|k=(\frac{R}{\pi }\mathrm{sinh}\frac{2\pi y}{R})^{\alpha ^2}|\zeta |^{2\alpha p}e^{i\alpha (\phi _<(\zeta ,\overline{\zeta })\phi _<(1/\zeta ,1/\overline{\zeta }))}|k+\alpha .$$
(2.8)
Therefore, using the bosonisation formulae (2.3), we find that
$$\mu _1\mu _2(x+iy)|\mathrm{free}=i\sqrt{2}(2y)^{1/4}\underset{k}{}(1)^k|k+1/2+O(y^{3/4}),$$
(2.9)
and so the boundary state $`|\mathrm{free}^{}`$ of a free boundary condition with an odd number of disorder lines is given by
$$|\mathrm{free}^{}=i\underset{k}{}(1)^k|k+1/2.$$
(2.10)
As a check we can calculate the partition function for the Ising model with a disorder line ending on each end,
$`Z_{\mathrm{free}^{},\mathrm{free}^{}}`$ $`=`$ $`\mathrm{free}^{}|e^{LH(R)}|\mathrm{free}^{}={\displaystyle \underset{r+1/2}{}}\stackrel{~}{q}^{r^2/2}/\eta (\stackrel{~}{\tau })=\theta _2(0|\stackrel{~}{\tau })/\eta (\stackrel{~}{\tau })=\theta _4(0|\tau )/\eta (\tau )`$ (2.11)
$`=`$ $`(\chi _0(q)\chi _{1/2}(q))^2=(Z_{\mathrm{free}^{},\mathrm{free}^{}}^{\mathrm{Ising}})^2,`$
i.e. the square of the frustrated partition function as should be the case.
As a final check,
$$Z_{\mathrm{free},\mathrm{free}^{}}=\mathrm{free}|e^{LH(R)}|\mathrm{free}^{}=0,$$
(2.12)
reflecting the fact that there are no configurations with a single disorder line ending on one end of the cylinder and no disorder fields.
Having found the boundary states for the two Ising models in (free) and (free’), we must complete the discussion by considering the case of fixed boundary conditions. In model A, all possible combinations of up and down conditions on the two models are possible, and all such boundary states have been found in . In model B however, the spin fields only appear in the symmetric combination $`(\sigma _1\sigma _2)`$, and so one can only possibly distinguish the relative signs of the two spins. This means that the possible boundary conditions in model B are ‘fixed same’ and ‘fixed opposite’, giving the relative signs of the two spins. If we denote the corresponding boundary states by $`|\mathrm{same}`$ and $`|\mathrm{opp}.`$, they are given in terms of Oshikawa’s and Affleck’s boundary states as
$$\begin{array}{ccccc}\hfill |\mathrm{same}& =& \frac{|+++|}{\sqrt{2}}\hfill & =& _k|k/2\hfill \\ \hfill |\mathrm{opp}.& =& \frac{|++|+}{\sqrt{2}}\hfill & =& _k(1)^k|k/2\hfill \end{array}$$
(2.13)
Using (2.3) and (2.8) one finds that
$$\begin{array}{ccc}\hfill \sigma _1\sigma _2(x+iy)|\mathrm{same}& =& +\sqrt{2}(2y)^{1/4}|\mathrm{same}+O(y^{3/4}),\hfill \\ \hfill \sigma _1\sigma _2(x+iy)|\mathrm{opp}.& =& \sqrt{2}(2y)^{1/4}|\mathrm{opp}.+O(y^{3/4}).\hfill \end{array}$$
As a final comment, we note how the fixed and free boundary states are related by duality. If $`V\phi V=\phi `$, then duality $`\phi \pi \phi `$ is implemented by the operator $`D=Ve^{i\pi p}`$. Hence we have
$$D|k=e^{i\pi k}|k,D|\mathrm{free}=\frac{|\mathrm{same}+|\mathrm{opp}.}{\sqrt{2}},D|\mathrm{free}^{}=\frac{|\mathrm{same}|\mathrm{opp}.}{\sqrt{2}}.$$
## 3 Cylinder correlation functions
We can also use these bosonisation formulae and boundary states to calculate correlation functions of fields on the cylinder. One interesting case to consider is the expectation value of a single disorder field with free b.c.’s on the two sides of the cylinder.
To recall, the cylinder is given by the rectangle in the upper half $`z`$ plane with vertices $`0,R,iL,R+iL`$, with the vertical edges $`\mathrm{Re}(z)=0,R`$ identified. We must also choose on which end of the cylinder the disorder line attached to the bulk field ends. We choose it to end on the top edge $`\mathrm{Im}(z)=L`$.
Using standard free field techniques, we can explicitly evaluate
$`\mu _1\mu _2(iy)`$ $`=`$ $`\mathrm{free}^{}|e^{LH(R)}(\sqrt{2}:\mathrm{sin}\frac{1}{2}\phi (iy):)|\mathrm{free}`$ (3.1)
$`=`$ $`\sqrt{2}{\displaystyle \frac{\theta _3(\stackrel{~}{\nu }/2|\stackrel{~}{\tau })}{\eta (\stackrel{~}{\tau })}}\left[{\displaystyle \frac{1}{iR}}{\displaystyle \frac{\theta _1^{}(0|\stackrel{~}{\tau })}{\theta _1(\stackrel{~}{\nu }|\stackrel{~}{\tau })}}\right]^{1/4}`$
$`=`$ $`\sqrt{2}{\displaystyle \frac{\theta _3(\nu /2|\tau )}{\eta (\tau )}}\left[{\displaystyle \frac{1}{2L}}{\displaystyle \frac{\theta _1^{}(0|\tau )}{\theta _1(\nu |\tau )}}\right]^{1/4},`$
where $`\nu =y/L`$ and $`\stackrel{~}{\nu }=2iy/R`$. This gives the cylinder expectation value in the Ising model as
$$\mu (iy)_{\mathrm{free}^{},\mathrm{free}}=2^{1/4}\sqrt{\frac{\theta _3(\frac{\nu }{2}|\tau )}{\eta (\tau )}}\left|\frac{\theta _1^{}(0|\tau )}{2L\theta _1(\nu |\tau )}\right|^{1/8}.$$
(3.2)
This can also be found as the appropriate combination of the two chiral blocks given in (12.109) of .
As the disorder field approaches the $`x`$–axis, $`\nu 0`$, we find
$`\mu (iy)_{\mathrm{free},\mathrm{free}}`$ $`=`$ $`2^{1/4}(2\pi \nu )^{1/8}\left(\sqrt{{\displaystyle \frac{\theta _3(0|\tau )}{\eta (\tau )}}}+O(\nu )\right)=\mathrm{\hspace{0.17em}2}^{1/4}(2y)^{1/8}\left(\chi _0+\chi _{1/2}\right)+\mathrm{}`$ (3.3)
$`=`$ $`2^{1/4}(2y)^{1/8}Z_{\mathrm{free},\mathrm{free}}^{\mathrm{Ising}}+\mathrm{}.`$
This is the exactly what we would have expected from the bulk-boundary ope of the disorder field (1.2) at a free boundary. Hence we see explicitly that as the disorder field approaches the $`x`$–axis the disorder line shrinks to zero.
However, if the disorder field approaches the other side, the disorder line gets stretched across the cylinder. As it approaches the boundary $`\mathrm{Im}(z)=L`$ the leading term in the bulk-boundary ope will again be the identity operator, so that the leading behaviour of the correlation function should be the frustrated partition function, i.e. the partition function including a disorder line,
$$Z_{\mathrm{free}^{},\mathrm{free}^{}}^{\mathrm{Ising}}=\chi _0\chi _{1/2};.$$
Using (3.2) it is easy to check that this is in fact the case. Putting $`y=L\stackrel{~}{y}`$, $`\stackrel{~}{y}=Lv`$, $`\nu =1v`$, we find
$`\mu (i(L\stackrel{~}{y}))_{\mathrm{free},\mathrm{free}}`$ $`=`$ $`2^{1/4}\sqrt{{\displaystyle \frac{\theta _4(\frac{v}{2}|\tau )}{\eta (\tau )}}}\left|{\displaystyle \frac{\theta _1^{}(0|\tau )}{2L\theta _1(v|\tau )}}\right|^{1/8}`$ (3.4)
$`=`$ $`2^{1/4}(2Lv)^{1/8}\left(\sqrt{{\displaystyle \frac{\theta _4(0|\tau )}{\eta (\tau )}}}+O(v)\right)`$
$`=`$ $`2^{1/4}(2\stackrel{~}{y})^{1/8}\left(\chi _0\chi _{1/2}\right)+\mathrm{},`$
as expected. The two limits and the way the frustration line is stretched across the cylinder are shown below in Figure 2.
$$\begin{array}{c}\text{}\\ \text{Figure 2.}\end{array}$$
## 4 Correlation functions on the upper half plane
We finish with some comments on the structure of Ising correlation functions on the upper half plane (UHP). From , we can express the various boundary states in the B model in terms of various combinations of Dirichlet boundary states $`|D(\phi _0)`$, i.e. states for which $`\phi `$ takes the value $`\phi _0`$ on the boundary:
$$\begin{array}{cccccc}\hfill |\mathrm{free}& =& \frac{|D(\pi )+|D(\pi )}{\sqrt{2}}\hfill & \hfill |\mathrm{free}^{}& =& \frac{|D(\pi )|D(\pi )}{\sqrt{2}}\hfill \\ \hfill |\mathrm{same}& =& |D(0)\hfill & \hfill |\mathrm{opp}.& =& |D(2\pi )\hfill \end{array}$$
Note that our field $`\phi `$ is twice that of and hence the values of $`\phi _0`$ are also twice theirs.
It is easy to express a free boson $`\phi (\zeta ,\overline{\zeta })`$ satisfying $`\phi =\phi _0`$ on the $`x`$–axis in terms of a chiral boson $`\varphi (\zeta )`$ as
$$\phi (\zeta ,\overline{\zeta })=\phi _0+\varphi (\zeta )\varphi (\overline{\zeta }).$$
(4.1)
The only subtlety is that one cannot substitute this directly into vertex operators as the result will be ill-defined. We choose to define our normal ordering by
$$:\mathrm{exp}(i\alpha \phi (\zeta ,\overline{\zeta })):|_{D(\phi _0)}=(2\eta )^{\alpha ^2}:\mathrm{exp}(i\alpha (\phi _0+\varphi (\zeta )\varphi (\overline{\zeta })):,$$
(4.2)
where $`\zeta =\xi +i\eta `$ with $`\xi ,\eta `$.
If we also use the results of evaluating the bulk–boundary opes,
$$\begin{array}{ccc}\hfill \sqrt{2}(2\eta )^{1/4}:\mathrm{sin}\frac{1}{2}(\varphi (\zeta )\varphi (\overline{\zeta })):& =& \frac{1}{\sqrt{2}}(2\eta )^{3/4}\left[i\varphi (\xi )+O(\eta )\right],\hfill \\ \hfill \sqrt{2}(2\eta )^{1/4}:\mathrm{cos}\frac{1}{2}(\varphi (\zeta )\varphi (\overline{\zeta })):& =& \sqrt{2}(2\eta )^{1/4}[\mathrm{\hspace{0.17em}1}+\frac{1}{2}\eta ^2:(i\varphi (\xi ))^2:+O(\eta ^3)],\hfill \end{array}$$
we can further find the expressions for the boundary spin and disorder fields. We summarise the results in the following table of bosonisation formulae:
$$\begin{array}{cccc}& & & \\ & D\left(0\right)& D\left(\pi \right)& D\left(\pi \right)\\ & & & \\ \multicolumn{4}{c}{}\\ \sigma _1\sigma _2& \sqrt{2}\left(2\eta \right)^{{\scriptscriptstyle \frac{1}{4}}}:\mathrm{cos}\frac{\varphi \left(\zeta \right)\varphi \left(\overline{\zeta }\right)}{2}:& \sqrt{2}\left(2\eta \right)^{{\scriptscriptstyle \frac{1}{4}}}:\mathrm{sin}\frac{\varphi \left(\zeta \right)\varphi \left(\overline{\zeta }\right)}{2}:& \sqrt{2}\left(2\eta \right)^{{\scriptscriptstyle \frac{1}{4}}}:\mathrm{sin}\frac{\varphi \left(\zeta \right)\varphi \left(\overline{\zeta }\right)}{2}:\\ \mu _1\mu _2& \sqrt{2}\left(2\eta \right)^{{\scriptscriptstyle \frac{1}{4}}}:\mathrm{sin}\frac{\varphi \left(\zeta \right)\varphi \left(\overline{\zeta }\right)}{2}:& \sqrt{2}\left(2\eta \right)^{{\scriptscriptstyle \frac{1}{4}}}:\mathrm{cos}\frac{\varphi \left(\zeta \right)\varphi \left(\overline{\zeta }\right)}{2}:& \sqrt{2}\left(2\eta \right)^{{\scriptscriptstyle \frac{1}{4}}}:\mathrm{cos}\frac{\varphi \left(\zeta \right)\varphi \left(\overline{\zeta }\right)}{2}:\\ ϵ_1ϵ_2& :\varphi \left(\zeta \right)(\varphi )\left(\overline{\zeta }\right):& :\varphi \left(\zeta \right)(\varphi )\left(\overline{\zeta }\right):& :\varphi \left(\zeta \right)(\varphi )\left(\overline{\zeta }\right):\\ \sigma _1^B\sigma _2^B& \text{}& i\varphi \left(\xi \right)& i\varphi \left(\xi \right)\\ \mu _1^B\mu _2^B& i\varphi \left(\xi \right)& \text{}& \text{}\end{array}$$
Considering these formulae, we see that the only effect of taking the particular linear combinations in the free and free b.c.’s is to ensure that correlators including odd and even numbers of disorder operators vanish respectively. As a result, we can safely evaluate correlators in the ‘generic’ free boundary conditions (i.e. not paying attention to the number of disorder lines ending on the boundary) by using simply the bosonisation formulae for $`\phi _0=\pi `$. To summarise, we can calculate correlation functions in the UHP by using expressions for $`D(0)`$ to obtain fixed b.c.’s in the Ising models, and the expressions for $`D(\pi )`$ to obtain free b.c.’s.
Note that as we have defined them, the bosonisation formulae for $`D(0)`$ and $`D(\pi )`$ are not exactly related by the duality $`\phi \pi \phi `$ which should simply interchange the expressions for the spin and disorder fields but instead differ by an irrelevant sign in the spin operator.
As an example we use these formulae to calculate a correlation function which would be hard to find using standard conformal field theory techniques, being equivalent to a five point chiral block. Consider the expectation value of a bulk spin field, a bulk disorder field and a boundary disorder field in fixed boundary conditions. Using the formulae for $`D(0)`$ we have
$`f(\zeta )`$ $`=`$ $`\left(\sigma (\zeta )\mu (i)\mu ^B(1)_{\mathrm{fixed}}\right)^2=\sigma _1\sigma _2(\zeta )\mu _1\mu _2(i)\mu _1^B\mu _2^B(1)`$
$`=`$ $`2(2\eta )^{1/4}(2)^{1/4}:\mathrm{cos}(\frac{1}{2}(\varphi (\zeta )\varphi (\overline{\zeta })))::\mathrm{sin}(\frac{1}{2}(\varphi (i)\varphi (i))):i\varphi (0)`$
$`=`$ $`i\mathrm{\hspace{0.17em}2}^{1/2}\eta ^{1/4}`$
$`\left(({\displaystyle \frac{1}{\zeta }}{\displaystyle \frac{1}{\overline{\zeta }}}+{\displaystyle \frac{1}{i}}{\displaystyle \frac{1}{i}})\left[{\displaystyle \frac{(i\zeta )(i+\overline{\zeta })}{(i\overline{\zeta })(i+\zeta )}}\right]^{1/4}({\displaystyle \frac{1}{\zeta }}{\displaystyle \frac{1}{\overline{\zeta }}}{\displaystyle \frac{1}{i}}+{\displaystyle \frac{1}{i}})\left[{\displaystyle \frac{(i\zeta )(i+\overline{\zeta })}{(i\overline{\zeta })(i+\zeta )}}\right]^{1/4}\right)`$
This expression is single valued in $`\zeta `$, but when we take the square root to find the result for a single Ising model there are square-root branch points at $`\zeta =0`$ and at $`\zeta =i`$. Where we put this branch cut has implications for the identification of the operators in the correlation function, as shown in figure 3 where the functions plotted are identical up to a choice of the position of the branch cut.
$$\begin{array}{ccc}|f(\zeta )|^{1/2}=\sigma (\zeta )\mu ^B(0)\mu (i)_+& & \mathrm{sign}(\xi )|f(\zeta )|^{1/2}=\sigma (\zeta )\psi ^{(+|)}(0)\mu (i)_{+|}\\ \text{}& & \text{}\\ \text{}& & \text{}\\ \multicolumn{3}{c}{\text{Figure 3: two choices of the position of the disorder line showing the equivalent interpretation of the boundary disorder field }\mu ^B\text{ as the boundary-condition changing operator }\psi ^{(\pm )}\text{.}}\end{array}$$
In the first case we plot $`|f|^{1/2}`$, and with this choice the disorder line joins the bulk and boundary disorder fields, and the boundary condition appears to be uniform, fixed ‘up’.
In the second case we plot $`\mathrm{sign}(\xi )|f|^{1/2}`$, and the disorder line from the bulk spin field extends to infinity in the bulk, and the operator inserted at the origin appears to be a boundary changing operator interpolating fixed ‘up’ and fixed ‘down’ b.c.’s.
## 5 Conclusions
We have shown how the bulk bosonisation formulae of Di Francesco et al. can be extended to the Ising model on the upper half plane when the bosonic field satisfies Dirichlet boundary conditions. This has helped understand the nature of boundary disorder fields, and has also helped clarify the results of in which boundary states describing frustrated systems were found to be rotationally invariant – the result being that the disorder line ends on an operator of weight zero and hence is invariant under rotations.
The boundary conformal field theory we have been looking at is unusual on two counts. Firstly the bulk theory is non-local, and secondly it seems we are forced to consider boundary conditions with apparently two vacua; for the free case, these are vacua with even and odd numbers of disorder lines, and in the fixed these are the fixed up and fixed down vacua. Such boundary conditions fall outside Cardy’s classification on both counts and it would be interesting to understand more generally the relations between the non-locality of the bulk theory and these peculiarities of the boundary theory.
The perturbation by the boundary disorder operator also provides a very simple example of a unitary perturbation by a boundary-condition changing operator. Given the renewed interest in such perturbations , it will be of interest to generalise this to other models.
Acknowledgments
I would like to thank K. Graham, M.E. Ortiz, A. Recknagel, P. Ruelle and I. Runkel for many helpful discussions, and the organisers of the 1999 Oberwolfach meeting on “Mathematical Aspects of String Theory” for a very enjoyable meeting where this work was started. The work was supported in part by a TMR grant of the European Commission, contract reference ERBFMRXCT960012, and by an EPSRC advanced fellowship. |
warning/0002/hep-th0002213.html | ar5iv | text | # 1 Introduction
## 1 Introduction
There is a large class of massive 2D integrable quantum field theories (IQFTs), which can be considered as perturbed conformal field theories (CFTs) . The ultraviolet (UV) behavior of these IQFTs is encoded in the CFT data while their long distance properties are defined by the S-matrix data. If the basic CFT admits the representation of the primary fields of full symmetry algebra in terms of the exponential fields the CFT data include “reflection amplitudes”. These functions define the linear transformations between different exponential fields, corresponding to the same primary field. Reflection amplitudes play the crucial role for the calculation of the one point functions as well as for the description of the zero mode dynamics in integrable perturbed CFTs. In particular, the zero mode dynamics determines the UV asymptotics of the ground state energy $`E(R)`$ (or effective central charge $`c_{\mathrm{eff}}(R)`$) for the system on the circle of size $`R`$. The function $`c_{\mathrm{eff}}(R)`$ admits in this case the UV series expansion in the inverse powers of $`\mathrm{log}(1/R)`$. The solution of the quantization condition for the vacuum wave function (which can be written in terms of the reflection amplitudes), supplemented with the exact relations between the parameters of the action and the masses of the particles determines all logarithmic terms in this UV expansion.
The effective central charge $`c_{\mathrm{eff}}(R)`$ in IQFT can be calculated independently from the S-matrix data using the TBA method . At small $`R`$ its asymptotics can be compared with that following from the CFT data. In the case when the basic CFT is known the agreement of both approaches can be considered as nontrivial test for the S-matrix amplitudes in IQFT. The corresponding analysis based on the both approaches was previously done for the sinh-Gordon , super-symmetric sinh-Gorgon, Bullough-Dodd models and simply-laced affine Toda field theories (ATFTs) .
In this paper we study the UV behavior of the effective central charge in ATFTs associated with non-simply laced Lie algebras. These IQFTs have two different classical limits. Namely, the weak and strong coupling limits correspond to the dual pairs of affine Toda theories. As a result, the mass ratios in these IQFTs depend on the coupling constant and flow from the classical values characteristic for Lie algebra $`G`$<sup>1</sup><sup>1</sup>1Throughout the paper, we denote an untwisted algebra as $`G`$, while $`G^{}`$ refers to a twisted one. to the same values for the dual algebra $`G^{}`$ . The number of particles in ATFTs is equal to the rank $`r`$ of $`G`$. For large $`r`$ the numerical analysis of TBA equations, especially in the UV region, becomes rather complicated. The analytical approach to the TBA equations does not give, at present, the regular UV expansion. So, it is useful to have the full logarithmic expansion for $`c_{\mathrm{eff}}(R)`$ following from CFT data. The agreement of this expansion with the TBA results confirms the $`S`$-matrix as well as the relations between the parameters of the action and masses of particles in non-simply laced ATFT.
The remarkable feature of ATFT is that effective central charge calculated from the CFT data with subtracted bulk free energy term (like in TBA approach) gives a good agreement with the TBA results even outside the UV region (at $`R𝒪(1)`$). This “experimental” fact still needs the explanation.
The rest of the paper is organized as follows. After introduction of some basic notations we give the exact relations between the parameters of the action and masses of particles in non-simply laced ATFTs. Then following the procedure of ref. , we obtain the reflection amplitudes and quantization conditions for the wave function, describing the vacuum zero mode dynamics. Using these results we calculate the UV asymptotics of the effective central charges for ATFTs and compare these asymptotics with numerical data following from TBA equations. We omit here the details, which can be found in ref. , devoted to the analysis of UV asymptotics in simply laced ATFTs.
## 2 Mass-$`𝝁`$ Relations and Reflection Amplitudes
The ATFTs corresponding to Lie algebra $`G`$ is described by the action
$$𝒜=d^2x\left[\frac{1}{8\pi }(_\mu 𝝋)^2+\underset{i=1}{\overset{r}{}}\mu _ie^{b𝐞_i𝝋}+\mu _0e^{b𝐞_0𝝋}\right],$$
(1)
where $`𝐞_i,i=1,\mathrm{},r`$ are the simple roots of the Lie algebra $`G`$ of rank $`r`$ and $`𝐞_0`$ is a maximal root, satisfying the relation:
$$\underset{i=0}{\overset{r}{}}n_i𝐞_i=0,n_0=1.$$
(2)
Non-simply laced ATFTs have standard simple roots with $`𝐞_i^2=2`$ and nonstandard simple roots with $`𝐞_i^2\xi ^2(2)`$. We choose the corresponding parameters $`\mu _i`$ as $`\mu `$ (for standard roots) and $`\mu ^{}`$ (for nonstandard ones) respectively<sup>2</sup><sup>2</sup>2We choose the convention that the length squared of the long roots are four for $`C_r^{(1)}`$ and two for the other untwisted algebras..
In the case of non-simply laced ATFTs, the exact mass ratios are different from the classical ones and get quantum corrections . To describe the spectrum it is convenient to introduce the notations:
$$B=\frac{b^2}{1+b^2},H=\frac{h+b^2h^{}}{1+b^2},$$
(3)
where $`h`$ and $`h^{}`$ are Coxeter and dual Coxeter numbers of the algebra. Then the spectrum of ATFTs can be expressed in terms of one mass parameter $`\overline{m}`$ as:
$`B_r^{(1)}:`$ $`M_r=\overline{m},M_a=2\overline{m}\mathrm{sin}(\pi a/H),a=1,2,\mathrm{},r1`$
$`C_r^{(1)}:`$ $`M_a=2\overline{m}\mathrm{sin}(\pi a/H),a=1,2,\mathrm{},r`$
$`G_2^{(1)}:`$ $`M_1=\overline{m},M_2=2\overline{m}\mathrm{cos}(\pi (1/31/H))`$
$`F_4^{(1)}:`$ $`M_1=\overline{m},M_2=2\overline{m}\mathrm{cos}(\pi (1/31/H)),`$ (4)
$`M_3=2\overline{m}\mathrm{cos}(\pi (1/61/H)),M_4=2M_2\mathrm{cos}(\pi /H).`$
The relation between the parameter $`\overline{m}`$ in the above spectra and the parameters $`\mu _i`$ in the action (1) can be obtained by Bethe Ansatz method (see for example ). The corresponding analysis gives:
$$\underset{i=0}{\overset{r}{}}[\pi \mu _i\gamma (1+𝐞_i^2b^2/2)]^{n_i}=\left[\frac{\overline{m}k(G)}{2}\mathrm{\Gamma }\left(\frac{1B}{H}\right)\mathrm{\Gamma }\left(1+\frac{B}{H}\right)\right]^{2H(1+b^2)},$$
(5)
where, as usual $`\gamma (x)=\mathrm{\Gamma }(x)/\mathrm{\Gamma }(1x)`$, and $`k(G)`$ is a function depending on the algebra:
$`k(B_r^{(1)})={\displaystyle \frac{2^{2/H}}{\mathrm{\Gamma }(1/H)}},`$ $`k(C_r^{(1)})={\displaystyle \frac{2^{2B/H}}{\mathrm{\Gamma }(1/H)}},`$
$`k(G_2^{(1)})={\displaystyle \frac{\mathrm{\Gamma }(2/3)}{2\mathrm{\Gamma }(1/2)\mathrm{\Gamma }(1/6+1/H)}},`$ $`k(F_4^{(1)})={\displaystyle \frac{\mathrm{\Gamma }(2/3)}{2\mathrm{\Gamma }(1/2)\mathrm{\Gamma }(1/6+1/H)}}.`$ (6)
The similar relations for the dual ATFTs can be easily obtained from Eqs.(5, 2) if we use the duality relations for the parameters $`\mu _i`$ and $`\mu _i^{}`$ corresponding to the dual pairs of ATFTs:
$$\pi \mu _i\gamma \left(\frac{𝐞_i^2b^2}{2}\right)=\left(\pi \mu _i^{}\gamma \left(\frac{2}{𝐞_i^2b^2}\right)\right)^{𝐞_i^2b^2/2}$$
(7)
The ATFTs can be considered as perturbed CFTs. Without the last term with the zeroth root $`𝐞_0`$, the action in Eq.(1) describes the non-affine Toda theory (NATT), which is conformal. To describe the generator of conformal symmetry we introduce the complex coordinates $`z=x_1+ix_2`$ and $`\overline{z}=x_1ix_2`$ and vector:
$$𝐐=b𝝆+\frac{1}{b}𝝆^{},𝝆=\frac{1}{2}\underset{𝜶>0}{}𝜶,𝝆^{}=\frac{1}{2}\underset{𝜶>0}{}𝜶^{},$$
(8)
where the sum in definition of Weyl vector $`𝝆`$ ($`𝝆^{}`$) runs over all positive roots $`𝜶`$ (co-roots $`𝜶^{}`$) of $`G`$.
The holomorphic stress-energy tensor
$$T(z)=\frac{1}{2}(_z𝝋)^2+𝐐_z^2𝝋$$
(9)
ensures the local conformal invariance of the NATT with the central charge $`c=r+12𝐐^2`$.
Besides the conformal invariance the NATT possesses extended symmetry generated by $`W(G)`$-algebra. The full chiral $`W(G)`$-algebra contains $`r`$ holomorphic fields $`W_j(z)`$ ($`W_2(z)=T(z)`$) with spins $`j`$ which follows the exponents of Lie algebra $`G`$. The primary fields $`\mathrm{\Phi }_w`$ of $`W(G)`$ algebra are classified by $`r`$ eigenvalues $`w_j,j=1,\mathrm{},r`$ of the operator $`W_{j,0}`$ (the zeroth Fourier component of the current $`W_j(z)`$):
$$W_{j,0}\mathrm{\Phi }_w=w_j\mathrm{\Phi }_w,W_{j,n}\mathrm{\Phi }_w=0,n>0.$$
(10)
The exponential fields
$$V_𝒂(x)=e^{(𝐐+𝒂)𝝋(x)}$$
(11)
are spinless conformal primary fields with dimensions $`\mathrm{\Delta }(𝒂)=w_2(𝒂)=(𝐐^2𝒂^2)/2`$. The fields $`V_𝒂`$ are also primary with respect to all chiral algebra $`W(G)`$ with the eigenvalues $`w_j`$ depending on $`𝒂`$. The functions $`w_j(𝒂)`$, which define the representation of $`W(G)`$-algebra possess the symmetry with respect to the Weyl group $`𝒲`$ of Lie algebra $`G`$ , i.e. $`w_j(\widehat{s}𝒂)=w_j(𝒂)`$; for any $`\widehat{s}𝒲`$. It means that the fields $`V_{\widehat{s}𝒂}`$ for different $`\widehat{s}𝒲`$ are reflection images of each other and are related by the linear transformation:
$$V_𝒂(x)=R_{\widehat{s}}(𝒂)V_{\widehat{s}𝒂}(x)$$
(12)
where $`R_{\widehat{s}}(𝒂)`$ is the “reflection amplitude”.
This function plays an important role in the analysis of perturbed CFTs. It can be calculated by the CFT methods (exactly in the same way as it was done for the simply laced NATTs in ) and has the form:
$$R_{\widehat{s}}(𝒂)=\frac{A_{\widehat{s}𝒂}}{A_𝒂}$$
(13)
where
$$A_𝒂=\underset{i=1}{\overset{r}{}}[\pi \mu _i\gamma (𝐞_i^2b^2/2)]^{𝝎_i^{}𝒂/b}\underset{𝜶>0}{}\mathrm{\Gamma }(1a_𝜶^{}/b)\mathrm{\Gamma }(1a_𝜶b),$$
(14)
here $`a_𝜶=𝒂𝜶`$, $`a_𝜶^{}=𝒂𝜶^{}`$ and vectors $`𝝎_i^{}`$ are the co-weights of $`G`$, satisfying the condition $`𝝎_i^{}𝐞_j=\delta _{ij}`$
In following we will be interested in the values of function $`A(𝐏)=A_{i𝐏}`$ . We note that in the semiclassical limit ($`b0`$ with $`𝐏/b`$ fixed) the functions $`A(\widehat{s}𝐏)`$ coincide with the amplitudes describing the asymptotics of the wave function of quantum mechanical non-affine Toda chain (16) (see for example ).
## 3 Quantization Condition and UV expansion
Function $`A(𝐏)`$ plays an important role in study of quantum mechanical problem for zero modes
$$𝝋_0=_0^{2\pi }𝝋(x)\frac{dx_1}{2\pi },$$
(15)
of the fields $`𝝋(x)`$ defined on an infinite cylinder of circumference $`2\pi `$ with coordinate $`x_2`$ along the cylinder playing the role of imaginary time. In the semiclassical limit $`b0`$, where one can neglect the oscillator modes of $`𝝋(x)`$, the Schrödinger equation governing the zero-mode dynamics is given by:
$$\left[\frac{r}{12}_{𝝋_0}^2+\underset{i=1}{\overset{r}{}}2\pi \mu _ie^{b𝐞_i𝝋_0}\right]\mathrm{\Psi }_𝐏(𝝋_0)=E_0\mathrm{\Psi }_𝐏(𝝋_0)$$
(16)
with the energy
$$E_0=\frac{r}{12}+𝐏^2.$$
(17)
where the momentum $`𝐏`$ is a real vector. The full quantum effect can be implemented simply by introducing the exact reflection amplitudes which take into account also non-zero-mode contributions .
The wave function $`\mathrm{\Psi }_𝐏(𝝋_0)`$ in the asymptotic region (Weyl chamber) can be found by using the same arguments as was given in for simply laced NATTs. The only difference is that there are now two kinds of roots with different lengths. Namely, each exponential term $`\mu _ie^{b𝐞_i𝝋_0}`$ in the Hamiltonian can be considered as a potential wall normal to the $`𝐞_i`$ direction. An incident wave is reflected by this wall to the wave with the Weyl-reflected momentum. The phase change corresponding to this process should be the same as in Liouville field theory. By considering the reflections from all potential walls, we find that the wave function $`\mathrm{\Psi }_𝐏(𝝋_0)`$ can be written as a superposition of plane waves with the momenta forming the orbit of the Weyl group $`𝒲`$ of Lie algebra $`G`$,
$$\mathrm{\Psi }_𝐏(𝝋_0)\underset{\widehat{s}𝒲}{}A(\widehat{s}𝐏)e^{i\widehat{s}𝐏𝝋_0},$$
(18)
where
$$A(𝐏)=\underset{i=1}{\overset{r}{}}[\pi \mu _i\gamma (𝐞_i^2b^2/2)]^{i𝝎_i^{}𝐏/b}\underset{𝜶>0}{}\mathrm{\Gamma }(1iP_𝜶b)\mathrm{\Gamma }(1iP_𝜶^{}/b),$$
(19)
For the Weyl element $`\widehat{s}_i`$, associated with the simple root $`𝐞_i`$, the ratio $`A(\widehat{s}_i𝐏)/A(𝐏)`$ should be given by the reflection amplitude $`S_L(𝐞_i,𝐏)`$ of the Liouville field theory
$`{\displaystyle \frac{A(\widehat{s}_i𝐏)}{A(𝐏)}}`$ $`=`$ $`S_L(𝐞_i,𝐏)`$ (20)
$`=`$ $`[\pi \mu _i\gamma (𝐞_i^2b^2/2)]^{i𝐏𝐞_i^{}/b}{\displaystyle \frac{\mathrm{\Gamma }(1+i𝐏𝐞_ib)\mathrm{\Gamma }(1+i𝐏𝐞_i^{}/b)}{\mathrm{\Gamma }(1i𝐏𝐞_ib)\mathrm{\Gamma }(1i𝐏𝐞_i^{}/b)}}.`$
One can easily check that function $`A(𝐏)`$ satisfies this functional equation. With this function one can proceed to obtain the scaling functions in the UV region of the ATFTs defined on a cylinder with circumference $`R0`$. The additional term in the ATFT Lagrangian corresponding to the zeroth root $`𝐞_0`$ introduces new potential wall in that direction. With this term the Weyl chamber is now closed and the momentum $`𝐏`$ of the wave function should be quantized. It depends on the size of the enclosed region, which is proportional to $`\mathrm{log}(1/R)`$. This quantized momentum $`𝐏(R)`$ defines the scaling function $`c_{\mathrm{eff}}`$ in the UV region by Eq.(17).
It is convenient to rescale back the size of the system from $`R`$ to $`2\pi `$. This leads to the following rescaling of the parameters $`\mu _i`$ in the action (1):
$$\mu _i\nu _i=\mu _i\left(\frac{R}{2\pi }\right)^{2+b^2𝐞_i^2},$$
(21)
In the UV limit the size of enclosed region is rather big and we can neglect the subtleties of interaction (which give only exponential corrections) taking into account only the phase shifts coming from the reflections of the waves by the potential walls. Since the additional potential term is not different from the others, the amplitude $`A(\widehat{s}𝐏)`$ with the momentum $`\widehat{s}𝐏`$ (where $`\widehat{s}`$ is an arbitrary element of Weyl group) has to satisfy also the reflection relation (20) with respect to the zeroth root $`𝐞_0`$
$$\frac{A(\widehat{s}_0\widehat{s}𝐏)}{A(\widehat{s}𝐏)}=S_L(𝐞_0,\widehat{s}𝐏).$$
(22)
Inserting Eqs.(19) and (20) into Eq.(22), we obtain the condition for $`𝐏`$. After some transformations (see ref. for details), it can be written in the form:
$$\left[\underset{i=0}{\overset{r}{}}\left(\pi \nu _i\gamma (𝐞_i^2b^2/2)\right)^{n_i}\right]^{i𝐏\widehat{s}𝐞_0^{}/b}\underset{𝜶>0}{}\left[\frac{𝒢(𝜶,𝐏)}{𝒢(𝜶,𝐏)}\right]^{𝜶\widehat{s}𝐞_0^{}}=1,$$
(23)
where $`\nu _i`$ are defined by Eq.(21) and
$$𝒢(𝜶,𝐏)=\mathrm{\Gamma }(1iP_𝜶b)\mathrm{\Gamma }(1iP_𝜶^{}/b).$$
For the lowest energy state, Eq.(23) reduces to the following equation:
$$L𝐏=2\pi 𝝆\underset{𝜶>0}{}𝜶\delta (𝜶,𝐏),$$
(24)
where
$$L=\frac{2}{b}(h+b^2h^{})\mathrm{ln}\frac{R}{2\pi }\frac{1}{b}\mathrm{ln}\left[\underset{i=0}{\overset{r}{}}\left(\pi \mu _i\gamma (𝐞_i^2b^2/2)\right)^{n_i}\right],$$
(25)
and
$$\delta (𝜶,𝐏)=i\mathrm{log}\frac{\mathrm{\Gamma }(1+iP_𝜶b)\mathrm{\Gamma }(1+iP_𝜶^{}/b)}{\mathrm{\Gamma }(1iP_𝜶b)\mathrm{\Gamma }(1iP_𝜶^{}/b)}.$$
(26)
This is the quantization condition for the momentum $`𝐏`$ in the UV region $`R0`$. The ground state energy of the system on the circle of size $`R`$ is then given by
$$E(R)=\frac{\pi c_{\mathrm{eff}}}{6R}\mathrm{with}c_{\mathrm{eff}}=r12𝐏^2$$
(27)
where $`𝐏`$ is the solution of Eq.(24).
In the UV region we can solve Eq.(24) perturbatively by expanding $`\delta (𝜶,𝐏)`$ in powers of $`P_𝜶`$,
$$\delta (𝜶,𝐏)=\delta _1(𝜶,b)P_𝜶+\delta _3(𝜶,b)P_𝜶^3+\delta _5(𝜶,b)P_𝜶^5\mathrm{},$$
(28)
where the coefficients $`\delta _1(𝜶,b)`$ and $`\delta _s(𝜶,b)`$, $`s=3,5`$ are:
$$\delta _1(𝜶,b)=2\gamma _E\left(b+\frac{2}{𝜶^2b}\right),\delta _s(𝜶,b)=()^{\frac{s3}{2}}\frac{2}{s}\zeta (s)\left(b^s+\left(\frac{2}{𝜶^2b}\right)^s\right).$$
(29)
Using the relations: $`_{𝜶>0}(𝜶)^a(𝜶)^b=h^{}\delta ^{ab},`$ and $`_{𝜶>0}(𝜶)^a(𝜶^{})^b=h\delta ^{ab},`$ we obtain that:
$$l𝐏=2\pi 𝝆\underset{𝜶>0}{}\delta _3(𝜶,b)𝜶P_𝜶^3\underset{𝜶>0}{}\delta _5(𝜶,b)𝜶P_𝜶^5\mathrm{},$$
with
$$l=L2\gamma _E(bh^{}+h/b)LL_0.$$
(30)
The above equation can be solved iteratively in powers of $`1/l`$. Inserting the solution into Eq.(27), we find:
$`c_{\mathrm{eff}}`$ $`=`$ $`rr(h+1)h^{}\left({\displaystyle \frac{2\pi }{l}}\right)^2+{\displaystyle \frac{8}{\pi }}\zeta (3)[C_4(G^{})b^3+C_4(G)/b^3]\left({\displaystyle \frac{2\pi }{l}}\right)^5`$ (31)
$`{\displaystyle \frac{24}{5\pi }}\zeta (5)[C_6(G^{})b^5+C_6(G)/b^5]\left({\displaystyle \frac{2\pi }{l}}\right)^7+𝒪(l^8),`$
where the coefficients $`C(G)`$ are defined as:
$`C_4(G)={\displaystyle \underset{𝜶>0}{}}\rho _𝜶\rho _𝜶^{}^3,C_4(G^{})={\displaystyle \underset{𝜶>0}{}}\rho _𝜶^4,`$
$`C_6(G)={\displaystyle \underset{𝜶>0}{}}\rho _𝜶\rho _𝜶^{}^5,C_6(G^{})={\displaystyle \underset{𝜶>0}{}}\rho _𝜶^6.`$
For simply laced algebras, these coefficients were calculated in and have the values:
$`C_4(A_{n1}^{(1)})`$ $`=`$ $`{\displaystyle \frac{1}{60}}n^2(n^21)(2n^23),`$
$`C_6(A_{n1}^{(1)})`$ $`=`$ $`{\displaystyle \frac{1}{168}}n^2(n^21)(n^22)(3n^25),`$
$`C_4(D_n^{(1)})`$ $`=`$ $`{\displaystyle \frac{1}{30}}(16n^345n^2+27n+8)n(n1)(2n1),`$
$`C_6(D_n^{(1)})`$ $`=`$ $`{\displaystyle \frac{1}{42}}(48n^5213n^4+262n^3+6n^2101n32)n(n1)(2n1).`$ (32)
For the non-simply laced algebras $`B_n^{(1)}`$ and $`C_n^{(1)}`$, we can express the results through these values. Namely, we find:
$`C_i(B_n^{(1)})={\displaystyle \frac{1}{2}}C_i(A_{2n1}^{(1)}),`$ $`C_i(B_n^{(1)})=C_i(D_{n+1/2}^{(1)}),`$ (33)
$`C_i(C_n^{(1)})=C_i(D_{n+1}^{(1)}),`$ $`C_i(C_n^{(1)})=C_i(D_n^{(1)}),(i=4,6).`$
For exceptional algebras $`G_2^{(1)}`$ and $`F_4^{(1)}`$, we obtain:
$`C_4(G_2^{(1)})={\displaystyle \frac{1}{3}}C_4(D_4^{(1)})=392,`$ $`C_4(G_2^{(1)})={\displaystyle \frac{980}{9}},`$ (34)
$`C_6(G_2^{(1)})={\displaystyle \frac{1}{3}}C_6(D_4^{(1)})=7386,`$ $`C_6(G_2^{(1)})={\displaystyle \frac{199516}{243}},`$
$`C_4(F_4^{(1)})={\displaystyle \frac{1}{2}}C_4(E_6^{(1)})=27378,`$ $`C_4(F_4^{(1)})={\displaystyle \frac{22815}{2}},`$
$`C_6(F_4^{(1)})={\displaystyle \frac{1}{2}}C_6(E_6^{(1)})=2203578,`$ $`C_6(F_4^{(1)})={\displaystyle \frac{4052763}{8}}.`$
We note that above equations relating coefficients $`C_i(G)`$ for different Lie algebras follow from the similar exact relations between the ground state energies $`e(G)`$ of quantum affine Toda chains associated with these Lie algebras. These exact relations are valid if the parameters $`\mu `$, $`\mu ^{}`$ for non-simply laced Lie algebras and corresponding parameter $`\mu _{sl}`$ for simply laced ones satisfy the condition: $`\mu ^{hz}(2\mu ^{}/\xi ^2)^z=\mu _{sl}^h`$, where $`z=\frac{2(hh^{})}{2\xi ^2}`$.
## 4 Comparison with TBA results
The effective central charge calculated above from the CFT data (reflection amplitudes) can be compared with the same function determined from numerical solution of the TBA equations for ATFTs. Namely:
$$c_{\mathrm{eff}}^{(\mathrm{TBA})}(R)=\underset{i=1}{\overset{r}{}}\frac{3Rm_i}{\pi ^2}\mathrm{cosh}\theta \mathrm{log}\left(1+e^{ϵ_i(\theta ,R)}\right)d\theta .$$
(35)
where functions $`ϵ_i(\theta ,R)`$ ($`i=1,\mathrm{},r`$) satisfy the system of $`r`$ coupled integral equations:
$$m_iR\mathrm{cosh}\theta =ϵ_i(\theta ,R)+\underset{j=1}{\overset{r}{}}\phi _{ij}(\theta \theta ^{})\mathrm{log}\left(1+e^{ϵ_i(\theta ^{},R)}\right)\frac{d\theta ^{}}{2\pi },$$
(36)
with the kernels $`\phi _{ij}`$, equal to the logarithmic derivatives of the $`S`$-matrices $`S_{ij}(\theta )`$ of ATFTs, conjectured in .
The function $`E^{(TBA)}(R)`$ defined from the TBA equations differs from the ground state energy $`E(R)`$ of the system on the circle of size $`R`$ by the bulk term: $`E^{(TBA)}(R)=E(R)fR`$, where $`f`$ is a specific bulk free energy . To compare the same functions we should subtract this term from the function $`E(R)`$ defined by Eq.(31) i.e.
$$c_{\mathrm{eff}}^{(\mathrm{TBA})}(R)=c_{\mathrm{eff}}^{(\mathrm{RA})}(R)+\frac{6R^2}{\pi }f(G).$$
(37)
The specific bulk free energy $`f(G)`$ for non-simply laced ATFTs can be calculated by Bethe Ansatz method with the result:
$`f(G)`$ $`=`$ $`{\displaystyle \frac{\overline{m}^2\mathrm{sin}(\pi /H)}{8\mathrm{sin}(\pi B/H)\mathrm{sin}(\pi (1B)/H)}},G=B_r^{(1)},C_r^{(1)},`$
$`f(G)`$ $`=`$ $`{\displaystyle \frac{\overline{m}^2\mathrm{cos}(\pi (1/31/H))}{16\mathrm{cos}(\pi /6)\mathrm{sin}(\pi B/H)\mathrm{sin}(\pi (1B)/H)}},G=G_2^{(1)},F_4^{(1)}.`$ (38)
The contribution of bulk term $`f(G)`$ becomes quite essential at $`R𝒪(1)`$.
The TBA Eqs.(36) were solved numerically for non-simply laced algebras, $`B_2^{(1)}`$ ($`=C_2^{(1)}`$), $`B_3^{(1)}`$, $`B_4^{(1)}`$, $`C_3^{(1)}`$, $`C_4^{(1)}`$, $`G_2^{(1)}`$ and $`F_4^{(1)}`$. The effective central charge $`c_{\mathrm{eff}}^{(\mathrm{TBA})}(R)`$ was then computed from Eq.(35) for many different values of parameter $`\overline{m}R`$. After taking into account the bulk term, the numerical solution for $`c_{\mathrm{eff}}^{(\mathrm{TBA})}(R)`$ was fitted with the expansion (31) (neglecting higher order terms in $`1/l`$):
$$c_{\mathrm{eff}}^{(\mathrm{RA})}(R)=rr(h+1)h^{}\left(\frac{2\pi }{l}\right)^2+c_5\left(\frac{2\pi }{l}\right)^5+c_7\left(\frac{2\pi }{l}\right)^7.$$
(39)
with fitting parameters $`L_0`$, $`c_5`$ and $`c_7`$, where parameter $`L_0`$ is defined by the Eq.(30). The exact values of these parameters can be easily identified from Eqs.(30) and (31). To compare the expansion (39) with TBA results we use the relations (5) between parameters $`\mu _i`$ in the action and the parameter $`\overline{m}`$ characterizing the spectrum of particles. It gives the following expression for function $`L(R)`$ in Eqs.(30):
$$L=\frac{2}{b}(h+b^2h^{})\mathrm{ln}\left[\frac{\overline{m}R}{4\pi }k(G)\mathrm{\Gamma }\left(\frac{1B}{H}\right)\mathrm{\Gamma }\left(1+\frac{B}{H}\right)\right]+\frac{2}{b}\mathrm{ln}(b^{2h}(\xi ^2/2)^z).$$
(40)
Tables 1–3 show the values of parameters $`L_0`$, $`c_5`$ and $`c_7`$ obtained numerically from TBA equations (denoted with the superscript (TBA)) and those obtained analytically (Eqs.(30) and (31)) from reflection amplitudes (denoted with the superscript (RA)) for $`C_2^{(1)}`$, $`C_3^{(1)}`$, $`C_4^{(1)}`$, $`B_3^{(1)}`$, $`B_4^{(1)}`$, $`G_2^{(1)}`$ and $`F_4^{(1)}`$ ATFTs with different values of the parameter $`B`$. We see that both data are in excellent agreement. (Relatively poor accuracy for $`c_7`$ is mainly due to the limitation of numerical accuracy and the influence of higher order term ($`𝒪(l^8)`$) in the expansion (39).) This agreement supports the approach based on the reflection amplitudes, $`\mu `$-$`\overline{m}`$ relations and quantization conditions as well as the $`S`$-matrices for non-simply laced ATFTs.
In Fig.1, we plot the functions $`c_{\mathrm{eff}}^{(\mathrm{TBA})}(R)`$ and $`c_{\mathrm{eff}}^{(\mathrm{RA})}(R)`$ for different ATFTs setting $`\overline{m}=1`$. The first function is computed numerically from TBA equations. The second one is calculated using Eqs.(24) and (27), based on the reflection amplitudes, with taking into account the bulk free energy term according to Eq.(37). For all models, the two curves are almost identical without essential difference in the graphs even at $`R𝒪(1)`$. This good agreement outside the UV region looks not to be accidental. However, at present, we have no satisfactory explanation of this interesting phenomena in ATFTs.
## 5 Concluding remarks
In the main part of this paper we considered the UV asymptotics of the effective central charges in ATFTs. The most important CFT data, which we used for this analysis were the reflection amplitudes (13) of NATTs. It was mentioned in Introduction, that these functions play also a crucial role in the calculation of the one point functions in perturbed CFT. The one point functions of the exponential fields in ATFTs:
$$𝒯(𝒂)=<\mathrm{exp}𝒂𝝋>$$
(41)
can be reconstructed from from the same reflection amplitudes. It follows from the results of the paper that functions (41) satisfy the functional equations similar to the relations (12) for the vertex operators. These equations together with analyticity and symmetry conditions fix one point functions in perturbed CFTs. One can find the solution of these functional equations with proper analyticity properties and respecting all symmetries of extended Dynkin diagram of Lie algebra $`G`$. This solution is a natural generalization to the non-simply laced case of the one point function for $`ADE`$ series of ATFTs calculated in and can be written in the form:
$`𝒯(𝒂)`$ $`=`$ $`\left[{\displaystyle \frac{\overline{m}k(G)}{2}}\mathrm{\Gamma }\left({\displaystyle \frac{1B}{H}}\right)\mathrm{\Gamma }\left(1+{\displaystyle \frac{B}{H}}\right)\right]^{2𝐐𝒂𝒂^2}{\displaystyle \underset{i=1}{\overset{r}{}}}[\pi \mu _i\gamma (1+𝐞_i^2b^2/2)]^{𝝎_i^{}𝒂/b}`$ (42)
$`\times \mathrm{exp}{\displaystyle \frac{dt}{t}\left[𝒂^2e^{2t}(𝒂,t)\right]},`$
where
$$(𝒂,t)=\underset{𝜶>0}{}\frac{\mathrm{sinh}(a_𝜶bt)\mathrm{sinh}((ba_𝜶2bQ_𝜶+(1+b^2)H)t)\mathrm{sinh}((b^2𝜶^2/2+1)t)}{\mathrm{sinh}t\mathrm{sinh}(b^2𝜶^2t/2)\mathrm{sinh}((1+b^2)Ht)}.$$
(43)
The one point function $`𝒯(𝒂)`$ can be used for the analysis of ATFTs. In particular, it contains the information about the bulk free energy $`f(G)`$, which was calculated independently by Bethe Ansatz method. One can easily derive from Eqs.(1) and (5) that:
$$\frac{n_if(G)}{H(1+b^2)}=\mu _i𝒯(b𝐞_i)$$
(44)
Using Eq.(42) for function $`𝒯(𝒂)`$ one finds:
$`{\displaystyle \frac{4\pi \gamma (1+𝐞_i^2b^2/2)n_if(G)}{H(1+b^2)(\overline{m}k(G))^2}}`$ $`=`$ $`\left[\mathrm{\Gamma }\left({\displaystyle \frac{1B}{H}}\right)\mathrm{\Gamma }\left(1+{\displaystyle \frac{B}{H}}\right)\right]^2`$ (45)
$`\times \mathrm{exp}{\displaystyle \frac{dt}{t}\left[(b𝐞_i)^2e^{2t}(b𝐞_i,t)\right]}`$
The integral in the exponent can be calculated and results coincides with Eq.(4). This gives the nonperturbative test to the one point function $`𝒯(𝒂)`$. In particular, taking the limit $`b0`$ in Eq.(45) (and the dual limit) one can derive the amusing relations for gamma-functions associated with Lie algebras $`G`$. It is convenient to introduce the integers $`n_i^{}=n_i𝐞_i^2/2`$ $`i=0,\mathrm{},r`$. Then these relations can be written as:
$$\underset{𝜶>0}{}\left(\gamma (𝜶𝝆^{}/h)\right)^{𝐞_i𝜶^{}}=n_i^{}\left(\underset{i=0}{\overset{r}{}}(n_i^{})^{n_i}\right)^{1/h},$$
(46)
and
$$\underset{𝜶>0}{}\left(\gamma (𝜶𝝆/h^{})\right)^{𝐞_i^{}𝜶}=n_i^{}𝐞_i^2/2\left(\underset{i=0}{\overset{r}{}}(n_i^{}𝐞_i^2/2)^{n_i^{}}\right)^{1/h^{}}.$$
(47)
More detailed consideration of one point functions in ATFTs we suppose to give in another publication.
## Acknowledgement
We thank F. Smirnov and Al. Zamolodchikov for valuable discussions. PB gratefully thanks for the hospitality of KIAS where part of this work was done. This work is supported in part by MOST 98-N6-01-01-A-05 (CA), Korea Research Foundation KRF-99-015-DI0021 (CA) 1998-015-D00071 (CR), and KOSEF 1999-2-112-001-5(CA,CR). PB’s work is supported in part by the EU under contract ERBFMRX. VF’s work is supported in part by the EU under contract ERBFMRX CT960012. |
warning/0002/hep-th0002146.html | ar5iv | text | # A New Class of Supersymmetric Orientifolds with D-branes at Angles
## 1 Introduction
Orientifolds of type II string theories or, equivalently, compactifications of the type I open string theory have considerably enlarged our view on consistent vacua of string theory. The plethora of Calabi-Yau compactifications known for the heterotic string carries over to type I via the ten dimensional S-duality . While in ten dimensions it involves an inversion of the string coupling, this is no longer necessary in six and four dimensions. However, in the type I setting one can perturbatively study backgrounds which are non-perturbative from the heterotic viewpoint, as they they involve lower dimensional D-branes. Besides all those recent ideas arising in this context, like lower string scale unification, type I vacua still leave a variety of open questions and puzzles.
First, employing the description of type I compactifications as orientifolds of type II, one chooses an exactly solvable point in the moduli space of the Calabi-Yau manifold and performs the tadpole cancellation computation. In all known cases these exactly solvable points are given by toroidal orbifolds or Gepner models. In general it is not clear how to deform away from these special points in moduli space. Second, at first sight surprisingly it appeared that some of the four dimensional standard orbifolds do not allow cancellation of all massless tadpoles . This perturbative inconsistency has been conjectured to be cured due to non-perturbative states in . However, not all cases for which perturbative inconsistencies arise have been resolved in this fashion.
Finally, the fate of the discrete NS-NS antisymmetric B-field modulus in the blown-up version and its translation to the heterotic side have not been explored in detail. Related to the presence of a B-field background is the occurance of extra multiplicities for certain open string sectors, which are required in order to achieve tadpole cancellation and anomaly freedom.
Some of these issues may possibly be tackled in a new class of orientifold models which we have presented in <sup>4</sup><sup>4</sup>4Some of the examples have been recovered in a completely different approach in . and which has also been analyzed in . They combine the world-sheet parity transformation with a reflection $``$ of half of the internal coordinates, which can be written as a complex conjugation on any two dimensional torus. Note, that this combination still preserves half of the supersymmetry. The new features arising in this class of orientifolds are that the operation $`\mathrm{\Omega }`$ does not exchange twisted sectors and that the internal B-field becomes a continuous modulus (whereas the off-diagonal entries of the metric get discretized). The extra multiplicities for certain open string sectors arise very naturally as intersection numbers of D-branes. We find solutions for any of the orbifold groups we have investigated. Even more surprisingly, we get more than just one orientifold per orbifold group, distinguished by the particular choice of the complex structures on the tori. From the phenomenological point of view the four dimensional models are of little interest, as they generically show non-chiral matter content.
In section 2 we describe the general features of $`\mathrm{\Omega }`$ orientifolds. In section 3 we investigate the $`_3`$ example as the technically most simple one in some detail. It already exhibits all conceptually new aspects of $`\mathrm{\Omega }`$ orientifolds. The section 4 finally summarizes results of our computations for the $`_4`$ and the two $`_6`$ examples.
## 2 $`\mathrm{\Omega }`$ orientifolds
The orientifolds with $`\mathrm{\Omega }`$ projection combine the ordinary world sheet parity transformation $`\mathrm{\Omega }`$ with a conjugation $``$ of the complex coordinates of the compactification torus $`T^4`$ or $`T^6`$, for six- or four dimensional models respectively. The tori are taken to factorize into two dimensional tori in all cases, with coordinates given by
$`X_ix_{102i}+ix_{112i},i=1,2,3.`$ (1)
The action of $``$ then reads
$`:X_i\overline{X}_i.`$ (2)
The complex structure of each two dimensional torus is further restricted by the requirement of the orbifold group acting crystallographically. We have chosen (up to rescaling) the $`SU(3)`$ or $`SU(2)^2`$ root lattices for simplicity. The relative choice of the complex structures on the various $`T^2`$ tori is of crucial importance for the gauge group and the spectra, which leads to a variety of distinct models. Since $``$ can also be considered as a reflection of one half of the compact real coordinates, $`\mathrm{\Omega }`$ is a symmetry of type IIA in four and of type IIB in six dimensions. This operation is then accompanied by one of the well known cyclic $`_N`$ orbifold groups preserving $`𝒩=1`$ supersymmetry. We denote the groups by $`_N=\{1,\mathrm{\Theta },\mathrm{},\mathrm{\Theta }^{N1}\}`$. The operation of the generator is diagonal in the complex basis:
$`\mathrm{\Theta }:X_i\mathrm{exp}\left(2\pi iv_i\right)X_i`$ (3)
and with opposite phases on the conjugate variables. Also the fermionic coordinates are being complexified in order to diagonalize $`\mathrm{\Theta }`$. The operation on the various ground states reads
$`:|s_0,s_1,s_2,s_3`$ $``$ $`|s_0,s_1,s_2,s_3,`$ (4)
$`\mathrm{\Theta }:|s_0,s_1,s_2,s_3`$ $``$ $`\mathrm{exp}\left(2\pi i\stackrel{}{v}\stackrel{}{s}\right)|s_0,s_1,s_2,s_3.`$
If in some twisted sector the ground state is only a spinor of a subgroup of $`SO(8)`$, the respective $`s_i`$ are set to zero formally. A further subtlety arises with the GSO projection which is ambiguous in the twisted closed string and the open string sectors a priori. It can be fixed by imposing the $`\mathrm{\Omega }`$ symmetry of the spectrum and requiring supersymmetry. The entire orientifold group is given by $`_N\mathrm{\Omega }_N`$ and these models cannot be mapped by T-duality to standard orientifolds $`_N\mathrm{\Omega }_N`$. Instead they are dual to asymmetric orientifolds $`\widehat{}_N\mathrm{\Omega }\widehat{}_N`$ with standard $`\mathrm{\Omega }`$ projection. The way how the T-duality acts on the various moduli of the models has been investigated in .
Six dimensional $`_N`$ orbifolds preserving $`𝒩=1`$ supersymmetry are easily defined by $`v_1=1/N,`$ $`v_2=1/N`$. Orbifolds featuring $`𝒩=1`$ supersymmetry in four dimensions have been classified in and we display their action on the complex basis in terms of the $`v_i`$ in table 1.
We will discuss the $`_3`$ example in greater detail and will only briefly state the results for the $`_4`$, $`_6`$ and $`_6^{}`$ cases. For all the other orbifold groups we do neither expect any conceptual novelties, nor will they produce more realistic models in terms of gauge groups or massless spectra, while surely demanding a large amount of tedious computation. In contrast to our results for the $`\mathrm{\Omega }`$ orientifold it has been established for the ordinary four dimensional $`_4`$ case that there exists no solution to the tadpole condition. This is also the case for the standard $`_8`$, $`_8^{}`$ and $`_{12}`$ orbifolds, which we believe to have solutions as $`\mathrm{\Omega }`$ orientifolds, as well.
## 3 The $`_3`$ example
In the following we describe the computations relevant for the $`_3`$ orientifolds in $`d=4`$ and $`d=6`$ with $`\mathrm{\Omega }`$ projection in detail. The two kinds of models are technically very similar but a number of subtle distinctions need to be made with respect to the dimensionality. We shall mostly display the formulas for the six dimensional case and explain the changes in four dimensions separately. The generator $`\mathrm{\Theta }`$ of the orbifold group acts diagonally by
$`\mathrm{\Theta }:X_{1,2}\mathrm{exp}\left(\pm 2\pi i/3\right)X_{1,2}`$ (5)
in six dimensions or by
$`\mathrm{\Theta }:X_{1,2}`$ $``$ $`\mathrm{exp}\left(2\pi i/3\right)X_{1,2},`$
$`\mathrm{\Theta }:X_3`$ $``$ $`\mathrm{exp}\left(4\pi i/3\right)X_3`$ (6)
in four dimensions. The torus we employ is defined by the root lattice of $`SU(3)`$ with basis vectors of length $`\sqrt{2}`$ for any of the complex directions, being depicted in Figure 1. It also displays the fixed points of $`\mathrm{\Theta }`$ by circles.
The action of $``$ is chosen to be a reflection of the vertical axis, leaving two inequivalent orientations of the lattice, denoted by A and B. Both lattices allow a crystallographic reflection, but lead to different results for the invariant Kaluza-Klein (KK) momenta and winding (W) states as well as for the number of $`\mathrm{\Omega }`$ invariant fixed points of $`\mathrm{\Theta }`$.
### 3.1 Closed string sector: The Klein bottle
In the closed string sector we have twisted sectors for all powers of $`\mathrm{\Theta }`$. The Fourier decomposition of the coordinate fields carries modings shifted by $`kv_i`$ in the $`k`$-twisted sector. An important feature is that $`\mathrm{\Omega }`$ as well as $``$ exchange twisted sectors implying that the twist fields are left invariant by the combined action. We then get the $`\mathrm{\Omega }`$ invariant oscillator excitations
$`\left(\mathrm{\Omega }\right)\alpha _{n+kv_i}\stackrel{~}{\overline{\alpha }}_{n+kv_i}\left(\mathrm{\Omega }\right)^1`$ $`=`$ $`\alpha _{n+kv_i}\stackrel{~}{\overline{\alpha }}_{n+kv_i},`$
$`\left(\mathrm{\Omega }\right)\overline{\alpha }_{nkv_i}\stackrel{~}{\alpha }_{nkv_i}\left(\mathrm{\Omega }\right)^1`$ $`=`$ $`\overline{\alpha }_{nkv_i}\stackrel{~}{\alpha }_{nkv_i}.`$
The $`\alpha `$ denotes any kind of ladder operator in the $`k`$-twisted sector with $`n`$ or $`n+1/2`$. The states are also invariant under the action of $`\mathrm{\Theta }`$:
$`\mathrm{\Theta }\alpha _{n+kv_i}\stackrel{~}{\overline{\alpha }}_{n+kv_i}\mathrm{\Theta }^1`$ $`=`$ $`\alpha _{n+kv_i}\stackrel{~}{\overline{\alpha }}_{n+kv_i},`$
$`\mathrm{\Theta }\overline{\alpha }_{nkv_i}\stackrel{~}{\alpha }_{nkv_i}\mathrm{\Theta }^1`$ $`=`$ $`\overline{\alpha }_{nkv_i}\stackrel{~}{\alpha }_{nkv_i}.`$
The left- and right-moving momenta are given by
$`p_\mathrm{L}`$ $`=`$ $`{\displaystyle \frac{1}{i\sqrt{\alpha ^{}U_2T_2}}}\left(m_1Um_2\overline{T}\left(n_1+Un_2\right)\right),`$
$`p_\mathrm{R}`$ $`=`$ $`{\displaystyle \frac{1}{i\sqrt{\alpha ^{}U_2T_2}}}\left(m_1Um_2T\left(n_1+Un_2\right)\right),`$
where $`U`$ and $`T`$ are the complex structure and the complexified Kähler modulus of the $`T^2`$. They are given in terms of the unit vectors spanning the lattice:
$`U={\displaystyle \frac{e_2}{e_1}}`$ , $`T=iV,`$ (7)
$`e_1^𝐀=\sqrt{2}`$ , $`e_2^𝐀={\displaystyle \frac{1}{\sqrt{2}}}+i\sqrt{{\displaystyle \frac{3}{2}}},`$
$`e_1^𝐁=\sqrt{2}`$ , $`e_2^𝐁={\displaystyle \frac{1}{\sqrt{2}}}+i{\displaystyle \frac{1}{\sqrt{6}}},`$
where we have adopted the choice of Figure 1. This allows to determine the KK+W states invariant under $`\mathrm{\Omega }\mathrm{\Theta }^k`$. In fact, the identity
$`\mathrm{\Theta }^{1/2}\left(\mathrm{\Omega }\mathrm{\Theta }^k\right)\mathrm{\Theta }^{1/2}=\left(\mathrm{\Omega }\mathrm{\Theta }^k\right)\mathrm{\Theta }`$ (8)
and the $`\mathrm{\Theta }^{1/2}`$ symmetry of the lattice implies that their contribution to the partition function is independent of $`k`$.
Using the above mentioned ingredients, the computation of the loop channel Klein bottle amplitude becomes a straightforward exercise:
$`𝒦=`$ $`2^{d/2}c{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dt}{t^{d/2+1}}}\text{Tr}_{\text{U+T}}({\displaystyle \frac{\mathrm{\Omega }}{2}}`$
$`{\displaystyle \frac{\left(1+\mathrm{\Theta }+\mathrm{\Theta }^2\right)}{3}}{\displaystyle \frac{\left(1+(1)^F\right)}{2}}e^{2\pi t\left(L_0+\overline{L}_0\right)}).`$
We define $`cV_d/\left(8\pi ^2\alpha ^{}\right)^{d/2}`$ and the momentum integration in the non-compact space-time has already been performed. While in the loop channel there are contributions in all twisted sectors, the relation
$`\mathrm{\Omega }\mathrm{\Theta }^k=\mathrm{\Theta }^k\mathrm{\Omega }`$ (10)
implies that only untwisted closed string states propagate in the tree channel Klein bottle amplitude. More precisely, the modular transformation maps the contributions of the $`k`$-twisted sector to the untwisted sector with the term $`\mathrm{\Theta }^k`$ of the orbifold projector $`(1+\mathrm{\Theta }+\mathrm{}+\mathrm{\Theta }^{N1})/N`$ inserted. In order for the corresponding boundary states to be invariant under $`\mathrm{\Theta }`$, it is necessary that the modular transformation produces the complete projector in the tree channel. It turns out that the latter condition is not always met automatically. We consider the completion of the projector as the guiding principle, which determines how to choose the relative orientations of the $`T^2`$ lattices, say of A or B type. Whenever this loop-tree channel equivalence is satisfied we find consistent models. In case it is not, one is actually missing states in the tree channel to complete the projector perturbatively, a circumstance that may hint towards the existence of non-perturbative states, that cure the discrepancy .
Luckily, for the $`_3`$ orbifold this condition does not impose any constraints, while it does so for all the other orbifold groups. Thus we have three inequivalent choices of the lattice (their complex structures) in six dimensions:
$`\mathrm{𝐀𝐀},\mathrm{𝐀𝐁},\mathrm{𝐁𝐁}`$ (11)
and four in four dimensions:
$`\mathrm{𝐀𝐀𝐀},\mathrm{𝐀𝐀𝐁},\mathrm{𝐀𝐁𝐁},\mathrm{𝐁𝐁𝐁}.`$ (12)
Another crucial input is that in the loop channel amplitude a $`k`$-twisted sector contribution has to be weighted by the number of fixed points of $`\mathrm{\Theta }^k`$, which are as well invariant under $``$. This factor differs for the given choices of the lattice orientations. For instance for the six dimensional AA case the loop channel Klein bottle amplitude reads
$`𝒦`$ $`=`$ $`2c(11){\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dt}{t^4}}({\displaystyle \frac{\vartheta \left[\genfrac{}{}{0pt}{}{0}{1/2}\right]^4}{\eta ^{12}}}`$ (13)
$`\left({\displaystyle \underset{m}{}}e^{4\pi tm^2/r^2}\right)^2\left({\displaystyle \underset{n}{}}e^{3\pi tn^2r^2}\right)^2`$
$`+{\displaystyle \frac{\vartheta \left[\genfrac{}{}{0pt}{}{0}{1/2}\right]^2}{\eta ^6}}{\displaystyle \frac{\vartheta \left[\genfrac{}{}{0pt}{}{1/3}{1/2}\right]\vartheta \left[\genfrac{}{}{0pt}{}{1/3}{1/2}\right]}{\vartheta \left[\genfrac{}{}{0pt}{}{1/6}{1/2}\right]\vartheta \left[\genfrac{}{}{0pt}{}{1/6}{1/2}\right]}}`$
$`+{\displaystyle \frac{\vartheta \left[\genfrac{}{}{0pt}{}{0}{1/2}\right]^2}{\eta ^6}}{\displaystyle \frac{\vartheta \left[\genfrac{}{}{0pt}{}{1/3}{1/2}\right]\vartheta \left[\genfrac{}{}{0pt}{}{1/3}{1/2}\right]}{\vartheta \left[\genfrac{}{}{0pt}{}{1/6}{1/2}\right]\vartheta \left[\genfrac{}{}{0pt}{}{1/6}{1/2}\right]}}),`$
the argument being $`\mathrm{exp}\left(4\pi t\right)`$. After a modular transformation one gets the tree channel amplitude
$`\stackrel{~}{𝒦}`$ $`=`$ $`c{\displaystyle \frac{32}{3}}(11){\displaystyle _0^{\mathrm{}}}dl({\displaystyle \frac{\vartheta \left[\genfrac{}{}{0pt}{}{1/2}{0}\right]^4}{\eta ^{12}}}`$ (14)
$`\left({\displaystyle \underset{m}{}}e^{\pi lm^2r^2}\right)^2\left({\displaystyle \underset{n}{}}e^{4\pi ln^2/(3r^2)}\right)^2`$
$`+3{\displaystyle \frac{\vartheta \left[\genfrac{}{}{0pt}{}{1/2}{0}\right]^2}{\eta ^6}}{\displaystyle \frac{\vartheta \left[\genfrac{}{}{0pt}{}{1/2}{1/3}\right]\vartheta \left[\genfrac{}{}{0pt}{}{1/2}{1/3}\right]}{\vartheta \left[\genfrac{}{}{0pt}{}{1/2}{1/6}\right]\vartheta \left[\genfrac{}{}{0pt}{}{1/2}{1/6}\right]}}`$
$`+3{\displaystyle \frac{\vartheta \left[\genfrac{}{}{0pt}{}{1/2}{0}\right]^2}{\eta ^6}}{\displaystyle \frac{\vartheta \left[\genfrac{}{}{0pt}{}{1/2}{1/3}\right]\vartheta \left[\genfrac{}{}{0pt}{}{1/2}{1/3}\right]}{\vartheta \left[\genfrac{}{}{0pt}{}{1/2}{1/6}\right]\vartheta \left[\genfrac{}{}{0pt}{}{1/2}{1/6}\right]}}).`$
The relative factors $`\left(2\mathrm{sin}(\pi /3)\right)^2=3`$ between the different contributions to the tree channel amplitude come out just as expected to complete the projector. Switching to the B lattice in any complex direction results in an overall factor of 3, while changing the lattice contributions and the number of $``$ invariant fixed points of $`\mathrm{\Theta }`$.
For computing the massless closed string spectra one also needs to distinguish the various lattice choices very carefully. The untwisted $`\mathrm{\Theta }`$ invariant states simply have to be symmetrized and antisymmetrized under $`\mathrm{\Omega }`$ in the NSNS and the RR sector respectively. This contribution to the spectrum is generic and together with the spectrum of the twisted sectors has been summarized in table 2, where we left out the graviton and dilaton multiplets.
The spectrum coming from twisted sectors differs for fixed points invariant under $`\mathrm{\Omega }`$ and for those fixed points, that form pairs under $`\mathrm{\Omega }`$. The former fixed points contribute two hypermultiplets in six and a chiral multiplet in four dimensions. The latter fixed points contribute two hypermultiplets in addition to two tensor multiplets in six dimensions and a chiral plus a vector multiplet in four dimensions. The sum of the neutral multiplets from the closed string sector finally matches the sum of the Hodge numbers $`h^{1,1}+h^{1,2}`$ in all cases.
### 3.2 Open string sector: The annulus and the Möbius strip
To cancel the tadpoles from the Klein bottle one has to introduce D($`d/2+5`$)-branes into the background, which intersect each other on the tori at non-trivial angles. They are extended in one dimension on each two dimensional torus and need to be located in a $`\mathrm{\Theta }`$ and $``$ invariant way. Recalling that open string coordinates with boundary conditions
$`\mathrm{}\left({\displaystyle \frac{}{\sigma }}X^i\right)=0,\mathrm{}\left({\displaystyle \frac{}{\tau }}X^i\right)=0`$
at $`\sigma =0`$ and
$`\mathrm{}\left(e^{\pi iw_i}{\displaystyle \frac{}{\sigma }}X^i\right)=0,\mathrm{}\left(e^{\pi iw_i}{\displaystyle \frac{}{\tau }}X^i\right)=0`$
at $`\sigma =\pi `$ have a Fourier decomposition with modings shifted by $`w_i`$, we are led to consider arrays of 3 kinds of D-branes at relative angle $`\pi /3`$. On each torus $`T^2`$ one of them is located in the fixed plane of the reflection $``$ and the remaining ones are obtained by successively applying the rotation $`\mathrm{\Theta }^{1/2}`$. Such a configuration is shown in Figure 2. The type of lattice is being distinguished, as the number of intersection points of any two kinds of branes differs.
One now needs to compute the annulus
$`𝒜=c`$ $`{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dt}{t^{d/2+1}}}\text{Tr}_{\text{open}}({\displaystyle \frac{1}{2}}{\displaystyle \frac{1+\mathrm{\Theta }+\mathrm{\Theta }^2}{3}}`$ (15)
$`{\displaystyle \frac{1+(1)^F}{2}}e^{2\pi tL_0})`$
and Möbius strip amplitudes
$`=c`$ $`{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dt}{t^{d/2+1}}}\text{Tr}_{\text{open}}({\displaystyle \frac{\mathrm{\Omega }}{2}}{\displaystyle \frac{1+\mathrm{\Theta }+\mathrm{\Theta }^2}{3}}`$ (16)
$`{\displaystyle \frac{1+(1)^F}{2}}e^{2\pi tL_0}).`$
There are only non-vanishing contributions in the traces, if the operator in the trace leaves the two D-branes $`(\mathrm{D}_i,\mathrm{D}_{i+n})`$, the string stretches between, and its orientation invariant. If the moding of the fields in the $`(\mathrm{D}_i,\mathrm{D}_{i+n})`$ open string sector is identical to the moding of the fields in the $`\mathrm{\Theta }^k`$ twisted closed string sector, we call such a sector the “$`\mathrm{\Theta }^k`$ twisted” open string sector.
In order to compute the zero-mode contributions to the $`(\mathrm{D}_i,\mathrm{D}_i)`$ amplitudes, one needs to inspect invariant KK+W states. The result is sensitive to the type of lattice. The additional requirement, that the Möbius strip contributions have to be $`\mathrm{\Omega }`$ invariant, induces a doubling of winding numbers as compared to the annulus. A very interesting point to notice is an extra multiplicity of some twisted open string sectors. As can be seen in Figure 2, the different kinds of branes intersect only in the origin of the A type unit cell, while they intersect in all three fixed points of $`\mathrm{\Theta }`$ in the B type unit cell. Thus the corresponding twisted sectors have to be weighted with extra multiplicities, leading to the same extra multiplicity for the number of states at each mass level. This turns out to be necessary to obtain anomaly free massless spectra. Note, that this issue has been related via T-duality to extra factors for open string sectors in ordinary $`\mathrm{\Omega }`$ orientifolds with background NS-NS B-field .
We present the tree channel annulus and Möbius strip amplitude in six dimensions for the AA choice of the lattice. Switching to a B type lattice results in an overall factor of 3 for all amplitudes and has an effect on the massless spectrum. The annulus reads
$`\stackrel{~}{𝒜}`$ $`=`$ $`c{\displaystyle \frac{\mathrm{M}^2}{6}}(11){\displaystyle _0^{\mathrm{}}}dl({\displaystyle \frac{\vartheta \left[\genfrac{}{}{0pt}{}{1/2}{0}\right]^4}{\eta ^{12}}}`$ (17)
$`\left({\displaystyle \underset{m}{}}e^{\pi lm^2r^2}\right)^2\left({\displaystyle \underset{n}{}}e^{4\pi ln^2/(3r^2)}\right)^2`$
$`+3{\displaystyle \frac{\vartheta \left[\genfrac{}{}{0pt}{}{1/2}{0}\right]^2}{\eta ^6}}{\displaystyle \frac{\vartheta \left[\genfrac{}{}{0pt}{}{1/2}{1/3}\right]\vartheta \left[\genfrac{}{}{0pt}{}{1/2}{1/3}\right]}{\vartheta \left[\genfrac{}{}{0pt}{}{1/2}{1/6}\right]\vartheta \left[\genfrac{}{}{0pt}{}{1/2}{1/6}\right]}}`$
$`+3{\displaystyle \frac{\vartheta \left[\genfrac{}{}{0pt}{}{1/2}{0}\right]^2}{\eta ^6}}{\displaystyle \frac{\vartheta \left[\genfrac{}{}{0pt}{}{1/2}{1/3}\right]\vartheta \left[\genfrac{}{}{0pt}{}{1/2}{1/3}\right]}{\vartheta \left[\genfrac{}{}{0pt}{}{1/2}{1/6}\right]\vartheta \left[\genfrac{}{}{0pt}{}{1/2}{1/6}\right]}}).`$
with argument $`\mathrm{exp}\left(4\pi l\right)`$ and the Möbius strip gives
$`\stackrel{~}{}`$ $`=`$ $`c{\displaystyle \frac{8\mathrm{M}}{3}}(11){\displaystyle _0^{\mathrm{}}}dl({\displaystyle \frac{\vartheta \left[\genfrac{}{}{0pt}{}{1/2}{0}\right]^4\vartheta \left[\genfrac{}{}{0pt}{}{0}{1/2}\right]^4}{\eta ^{12}\vartheta \left[\genfrac{}{}{0pt}{}{0}{0}\right]^4}}`$
$`\left({\displaystyle \underset{m}{}}e^{4\pi lm^2r^2}\right)^2\left({\displaystyle \underset{n}{}}e^{4\pi ln^2/(3r^2)}\right)^2`$
$`+3{\displaystyle \frac{\vartheta \left[\genfrac{}{}{0pt}{}{1/2}{0}\right]^2\vartheta \left[\genfrac{}{}{0pt}{}{0}{1/2}\right]^2}{\eta ^6\vartheta \left[\genfrac{}{}{0pt}{}{0}{0}\right]^2}}`$
$`({\displaystyle \frac{\vartheta \left[\genfrac{}{}{0pt}{}{1/2}{1/3}\right]\vartheta \left[\genfrac{}{}{0pt}{}{0}{1/6}\right]\vartheta \left[\genfrac{}{}{0pt}{}{1/2}{1/3}\right]\vartheta \left[\genfrac{}{}{0pt}{}{0}{1/6}\right]}{\vartheta \left[\genfrac{}{}{0pt}{}{0}{1/3}\right]\vartheta \left[\genfrac{}{}{0pt}{}{1/2}{1/6}\right]\vartheta \left[\genfrac{}{}{0pt}{}{0}{1/3}\right]\vartheta \left[\genfrac{}{}{0pt}{}{1/2}{1/6}\right]}}`$
$`+{\displaystyle \frac{\vartheta \left[\genfrac{}{}{0pt}{}{1/2}{1/3}\right]\vartheta \left[\genfrac{}{}{0pt}{}{0}{1/6}\right]\vartheta \left[\genfrac{}{}{0pt}{}{1/2}{1/3}\right]\vartheta \left[\genfrac{}{}{0pt}{}{0}{1/6}\right]}{\vartheta \left[\genfrac{}{}{0pt}{}{0}{1/3}\right]\vartheta \left[\genfrac{}{}{0pt}{}{1/2}{1/6}\right]\vartheta \left[\genfrac{}{}{0pt}{}{0}{1/3}\right]\vartheta \left[\genfrac{}{}{0pt}{}{1/2}{1/6}\right]}})).`$
with argument $`\mathrm{exp}\left(8\pi l\right)`$. One obtains only a single overall tadpole cancellation condition, which reads
$`{\displaystyle \frac{1}{6}}\left(\mathrm{M}^216\mathrm{M}+64\right)={\displaystyle \frac{1}{6}}\left(\mathrm{M}8\right)^2=0.`$
It implies the presence of $`\mathrm{M}=8`$ D7-branes of each kind. The analogous condition in four dimensions requires $`\mathrm{M}=4`$ D6-branes. The tadpoles are related to contributions of 8- and 7-forms respectively. Thus, we get an $`SO(8)`$ gauge group in six dimensions and $`SO(4)`$ in four. The rank reduction with respect to the $`SO(32)`$ in ten dimensions can be related to the rank reduction in ordinary $`\mathrm{\Omega }`$ orientifolds with background B-field via T-duality. As was pointed out in it is possible via discrete Wilson-lines to get the gauge groups $`Sp(8)`$ and $`Sp(4)`$, respectively.
The massless spectrum in the untwisted open string sector is again generic. In six dimensions there is a hypermultiplet, in four dimensions three chiral multiplets in addition to the gauge vectors, all in the antisymmetric representation. In the twisted sectors one needs to count the extra multiplicity of states deriving from the number of intersection points, each giving rise to a hypermultiplet in the antisymmetric in six and a chiral multiplet in the symmetric representation in four dimensions. The open string spectra are collected in table 3.
All the six dimensional models satisfy the cancellation condition of the irreducible anomaly
$`n_\mathrm{H}n_\mathrm{V}+29n_\mathrm{T}=273`$ (19)
while the four dimensional models are free of any gauge anomaly anyway.
## 4 Results for $`_4`$, $`_6`$ and $`_6^{}`$
In this section we shall only briefly explain some further properties of the other models which we have computed so far. The lattice and brane configurations, the explicit formulas for all the amplitudes to be calculated as well as all their massless spectra have been presented in .
The world-sheet consistency condition is more restrictive than for the $`_3`$, where all combinations of lattices were allowed. For all other examples we find that the relative orientation is fixed between two of the complex tori, while, in four dimensions, the third can be chosen freely of A or B type. With this freedom we obtain another set of eight inequivalent models from the three orbifold groups.
In the open string sector one again has to introduce sets of branes at $`\pi v_i`$ relative angles, where their respective number has to be determined by the tadpole cancellation conditions. Since there always exist two types of branes which are not mapped upon each other by the orbifold, the gauge group will have a product structure with two identical factors. The spectra, however, need not be symmetric under exchanging the two factors. As all the orbifold groups contain an element $`\mathrm{\Theta }^{N/2}`$ of order two, there is also a contribution to the $`_2`$ twisted sector in the tree channel. This generically imposes the second tadpole cancellation condition
$`\text{tr}\left(\gamma _{N/2}^{(i)}\right)=0`$ (20)
for the action of $`\mathrm{\Theta }^{N/2}`$ on the Chan-Paton factors of each type $`i`$ of branes. The two conditions together then exactly reproduce the situation of the standard six dimensional $`_2`$ orientifold discussed by Bianchi and Sagnotti in and later by Gimon and Polchinski in . Their solution implies that the orthogonal $`SO(\mathrm{M})`$ gauge group is broken to its unitary $`U(\mathrm{M}/2)`$ subgroup. For instance, in four dimensions for the $`_4`$ orbifold we find two possible solutions $`U(8)^2`$ and $`U(4)^2`$, while for the various $`_6`$ models we always get $`U(2)^2`$. This is in accord with the rank reduction expected from T-duality.
The determination of the massless spectra is in general more involved than for the $`_3`$. In the twisted sectors of the closed string spectrum one has to keep track of the transformation properties of the various fixed points under the orbifold group generator $`\mathrm{\Theta }`$ and $`\mathrm{\Omega }`$ in order to find the correct symmetrization prescription. Also one needs to distinguish with respect to the action of $`\mathrm{\Theta }^{N/2}`$. In the open string spectrum the same is true for the actions on the intersection points, which do not simply provide all the same states anymore. In four dimensions there are certain subtle phase factors to be regarded, which sometimes require extra signs in the loop channel amplitude in order to complete the tree channel projector. Finally, for the most complicated case, the $`_6^{}`$, one even needs to distinguish different types of branes with different effective gauge theories on them. The results display spectra which are anomaly free in six dimensions and non-chiral in four.
###### Acknowledgments.
We would like to thank C. Angelantonj, who participated in closely related work, A. Kumar, who was also involved in an early stage, as well as M. Gaberdiel and D. Lüst for encouraging discussions and helpful remarks. |
warning/0002/hep-lat0002026.html | ar5iv | text | # Chiral gauged fermions on a lattice
## I Introduction
The problems of chiral gauge theories on a lattice concerning the vectorlike (doubling) phenomenon, chiral gauge symmetries and anomalies have been still bothering theoretical physicists, since the “no-go” theorem of Nielsen and Ninomiya was demonstrated in 1981. However, many progresses in understanding have been made for approaching to the solution of the problem-. These approaches can be very briefly categorized into two major classes, one is the modeling based on appropriately introducing local interactions, another is the modeling, without interactions, instead based on delicately manipulating chiral fermions and gauge fields on the lattice. The latter has been quickly developed as the most exciting approaches with the strong theoretical ground. While as for the modeling with interactions, it has become a general belief that such modeling cannot lead to any chiral gauge theories on the lattice. Nevertheless, following the studies of the multifermion coupling model we proposed in refs., we are attempt to discuss the aspects of the coupling between gauge field and fermion fields, the gauge anomaly, fermion-flavour singlet anomaly (fermion-number violation) and the elimination of residual gauge symmetry breakings in this paper.
The main idea for modeling with interactions , either multifermion couplings (see examples ) or scale-fermion couplings (see example ) is analogous, although the details of interactions can be very different. It is to use strong local interactions to gauge-invariantly decouple extra chiral fermion species and in the meantime obtain the correct gauge anomalies and fermion-number violation in the continuum limit. The Eichten-Preskill and Smit-Swift model were two of such models. The studies of these two models show a broken phase with a hard spontaneous symmetry breaking of Nambu-Jona-Lasinio type at the lattice scale ($`O(1/a)`$). This broken phase separates two chiral gauge symmetric phases in the phase diagram in terms of couplings. This hard spontaneous symmetry breaking enforces the masses of intermediate gauge bosons to be at the lattice scale ($`O(1/a)`$). On the other hand, even within the chiral gauge symmetric phase, due to the exact locality of interactions introduced and the argument of anomalies cancelation, the gauge symmetric spectrum of fermion zero modes must be vectorlike analogously to the context of the “no-go” theorem. These studies lead to the general belief that the chiral gauge symmetry realized by a chiral fermion spectrum in the low-energy cannot be maintained on the lattice. Even though we do not share the view of this general belief, these studies tell us that in such models with interactions, the hard spontaneous symmetry breaking at the lattice scale absolutely cannot be tolerated in order to have a loophole to achieve a correct chiral gauge theory on the lattice. In the model of multifermion couplings with an anomaly-free fermion content, for many years, we have been trying to find the appropriate multifermion couplings that can avoid the hard spontaneous symmetry breaking within a scaling region in the phase space, and in this scaling region composite particles with “wrong” chirality can dissolve to their cuts, to which the “no-go” theorem does not apply, so as to preserve exact chiral gauge symmetries with the chiral fermion spectrum in the continuum limit. Due to the complex dynamics of strong interactions of such models that belong to various different universality classes, it is difficult to demonstrate a model with peculiar interactions to work, but it even much more difficult to prove a general “no-go” theorem for the failure of any models so constructed. The recent development of lattice chiral gauge theories shows that the chiral gauge symmetries can be preserved on the lattice, and this shines a light on this great puzzle whether chiral gauge symmetries realized by the chiral fermion spectrum can be achieved.
The essential spirit of the “no-go” theorem of Nielsen and Ninomiya is that, under very general prerequisites, e.g., locality and gauge symmetry, which are usually required in quantum field theories, the paradox concerning chiral gauge symmetries, vectorlike doubling and anomalies are unavoidably entwined. In this paper, within the context of the multifermion coupling model and its scaling region we proposed in refs., where no any hard spontaneous symmetry breaking occurs, we will discuss the chiral fermion spectrum in the low energy, vectorlike fermion spectrum in the high energy and their coupling to the gauge field to show how chiral gauge symmetry is realized. In addition, we will show the computations of the correct gauge anomaly and fermion-flavour singlet anomaly and discuss the reasons for such anomalies arise in the scaling region of our model. The elimination of residual breakings of the gauge symmetry after concelation of the gauge anomaly is discussed.
The paper is organized as follow. In section II, we give a detail review of the properties of our model in the scaling region, which were obtained in Ref.. In section III, we demonstrate the exact decoupling of the right-handed chiral fermion from the gauge field by the exact shift-symmetry of the right-handed chiral fermion in our model. This plays an extremely important rôle in obtaining the correct fermion-flavour singlet anomaly. In section IV, based on (i) the recursion relations obeyed by the three-point Green functions of the gauge and fermion fields in the strong coupling limit; (ii) the analytical continuation of these Green functions in the energy-momentum space; (iii) the weak-coupling computations of these Green functions, we show the exact chiral gauge symmetry and associated Ward identities are realized by both chiral fermion spectra in low-energy region and vectorlike fermion spectra in high energy regions. In section V, we present an explicit perturbative computation for achieving the correct gauge anomaly and discuss how the correct gauge anomaly arises from a soft chiral symmetry breaking at the scale much smaller than the lattice scale. We discuss the cancelation of the gauge anomalies and elimination of residual chiral gauge symmetry breaking. These are important for our model to have a consistent scenario of preserving chiral gauge symmetries with the correct low-energy spectrum and to have the normal renormalization prescription in perturbative calculations of the small gauge coupling. Some preliminary results of sections IV and V were reported in Phys. Lett. B408 (1997) 299. In section VI, we explicitly derive the fermion-flavour singlet anomaly via the mixing anomaly due to a soft chiral symmetry breaking and we discuss how this fermion-flavour singlet anomaly is obtained from the action (1), which at the lattice scale is global fermion flavour symmetric. In the conluding section, we also present a very general and brief discussion and remark on the possible connection to the recent development of lattice chiral gauge theories based on the Ginsparg-Wilson relation .
## II A plausible scaling region of lattice chiral fermions
We summarize the main points and results of the $`SU_L(2)`$ model of multifermion couplings and its scaling region proposed in the ref.. Note that $`\psi _L^i`$ ($`i=1,2`$) is an $`SU_L(2)`$ gauged doublet, $`\chi _R`$ is an $`SU_L(2)`$ singlet and both are two-component Weyl fermions. $`\chi _R`$ is treated as a “spectator” fermion. $`\psi _L^i`$ and $`\chi _R`$ fields are dimensionful $`[a^{\frac{1}{2}}]`$. We suggested the following action for the chiral fermions $`\psi _L^i`$ and $`\chi _R`$ with the $`SU_L(2)U_R(1)`$ chiral symmetries on the lattice:
$`S`$ $`=`$ $`S_f+S_1+S_2,`$ (1)
$`S_f`$ $`=`$ $`{\displaystyle \frac{1}{2a}}{\displaystyle \underset{x}{}}{\displaystyle \underset{\mu }{}}\left(\overline{\psi }_L^i(x)\gamma _\mu D_{ij}^\mu \psi _L^j(x)+\overline{\chi }_R(x)\gamma _\mu ^\mu \chi _R(x)\right),`$ (2)
$`S_1`$ $`=`$ $`g_1{\displaystyle \underset{x}{}}\overline{\psi }_L^i(x)\chi _R(x)\overline{\chi }_R(x)\psi _L^i(x),`$ (3)
$`S_2`$ $`=`$ $`g_2{\displaystyle \underset{x}{}}\overline{\psi }_L^i(x)\left[\mathrm{\Delta }\chi _R(x)\right]\left[\mathrm{\Delta }\overline{\chi }_R(x)\right]\psi _L^i(x),`$ (4)
where $`S_f`$ is the naive lattice action for chiral fermions, $`a`$ is the lattice spacing fixed as a constant. $`S_1`$ and $`S_2`$ are two external multifermion couplings, where the $`g_1`$ and $`g_2`$ have dimension $`[a^2]`$, and the Wilson factor is given as,
$`\mathrm{\Delta }\chi _R(x)`$ $``$ $`{\displaystyle \underset{\mu }{}}\left[\chi _R(x+\mu )+\chi _R(x\mu )2\chi _R(x)\right],`$ (5)
$`2w(p)`$ $`=`$ $`{\displaystyle _x}e^{ipx}\mathrm{\Delta }(x)={\displaystyle \underset{\mu }{}}\left(1\mathrm{cos}(p_\mu )\right).`$ (6)
Note that all momenta are scaled to be dimensionless, $`p=\stackrel{~}{p}+\pi _A`$ where $`\pi _A`$ runs over fifteen lattice momenta ($`\pi _A0`$).
The action (1) has an exact local $`SU_L(2)`$ chiral gauge symmetry,
$$\underset{\mu }{}\gamma _\mu D^\mu =\underset{\mu }{}(U_\mu (x)\delta _{x,x+\mu }U_\mu ^{}(x)\delta _{x,x\mu }),U_\mu (x)SU_L(2),$$
(7)
which is the gauge symmetry that the continuum theory (the target theory) possesses. The global flavour symmetry $`U_L(1)U_R(1)`$ is not explicitly broken in eq. (1). It has been advocated that there exists a plausible scaling region for chiral fermions in the low-energy limit. This is a peculiar segment in the phase space of the multifermion couplings $`g_1`$ and $`g_2`$,
$$𝒜=\left[g_10,g_2^{c,a}<g_2<g_2^{c,\mathrm{}}\right],a^2g_2^{c,a}=0.124,1g_2^{c,\mathrm{}}<\mathrm{},$$
(8)
where $`g_2^{c,\mathrm{}}`$ is a finite number. The crucial points and results for this scaling region to exist are briefly described in the following.
In the segment $`𝒜`$ (8), the action (1) possesses a global $`\chi _R`$-shift-symmetry , i.e., the action is invariant under the transformation: $`\overline{\chi }_R(x)\overline{\chi }_R(x)+\overline{\delta },\chi _R(x)\chi _R(x)+\delta ,`$ where $`\delta `$ is independent of space-time. The Ward identity corresponding to this $`\chi _R`$-shift-symmetry is given as ($`g_10`$),
$$\frac{1}{2a}\gamma _\mu ^\mu \chi _R^{}(x)+g_2\mathrm{\Delta }\left(\overline{\psi }_L^i(x)\left[\mathrm{\Delta }\chi _R(x)\right]\psi _L^i(x)\right)\frac{\delta \mathrm{\Gamma }}{\delta \overline{\chi }_R^{}(x)}=0.$$
(9)
Note that for studying Ward identities associated to the gauge symmetry and other global symmetries in this segment (8), the “primed” fermion field $`\chi _R^{}(x),\psi _R^{}(x)`$, gauge field $`A_\mu ^{}(x)`$ and the vacuum functional “$`\mathrm{\Gamma }`$” are introduced through the generating functional approach (see eqs.(6-13) in the previous paper).
The important consequences of this Ward identity (9) in segment $`𝒜`$ are:
* the low-energy mode ($`p0`$) of $`\chi _R`$ is a free mode and decoupled for its one particle irreducible (1PI) function:
$$_xe^{ipx}\frac{\delta ^{(2)}\mathrm{\Gamma }}{\delta \chi _R^{}(x)\delta \overline{\chi }_R^{}(0)}=\frac{i}{a}\gamma _\mu \mathrm{sin}(p^\mu );$$
(10)
* no hard spontaneous chiral symmetry breaking at the lattice scale ($`O(1/a)`$) occurs <sup>*</sup><sup>*</sup>*The soft symmetry breaking $`\stackrel{~}{v}`$ at the electroweak scale for the low-energy modes ($`p0`$) is allowed and can be achieved by tuning the multifermion couplings $`g_1`$ and $`g_2`$., as its the self-energy (1PI) function $`\mathrm{\Sigma }^i(p)`$ vanishes both at $`p=0`$ (see eqs.(30) and (31) in ref.):
$$_xe^{ipx}\frac{\delta ^{(2)}\mathrm{\Gamma }}{\delta \psi _L^i(x)\delta \overline{\chi }_R^{}(0)}=\frac{1}{2}\mathrm{\Sigma }^i(p)=0p=0,$$
(11)
and at $`p0`$ (see eq.(104) in ):
$$\mathrm{\Sigma }(p)=0p0$$
(12)
which is demonstrated by the strong coupling expansion in the segment $`𝒜`$.
On the other hand, for the strong coupling $`g_21`$ in the segment $`𝒜`$, the following three-fermion-states comprising of the elementary fields $`\psi _L^i`$ and $`\chi _R`$ in (1) are bound:
$$\mathrm{\Psi }_R^i=\frac{1}{2a}(\overline{\chi }_R\psi _L^i)\chi _R;\mathrm{\Psi }_L^n=\frac{1}{2a}(\overline{\psi }_L^i\chi _R)\psi _L^i.$$
(13)
These fermion bound states are Weyl fermions with the “wrong” chiralities in contrast with the “right” chiralities possessed by the elementary fields $`\psi _L^i`$ and $`\chi _R`$. By using the strong-coupling expansion in powers of $`1/g_2`$, we compute the following two-point Green functions with insertions of appropriate composite operators, for instance, for the charged sector,
$`S_{LL}^{ij}(x)`$ $`=`$ $`\psi _L^i(0)\overline{\psi }_L^j(x),S_{RL}^{ij}(x)=\psi _L^i(0)\overline{\mathrm{\Psi }}_R^j(x),`$ (14)
$`S_{RR}^{ij}(x)`$ $`=`$ $`\mathrm{\Psi }_R^i(0)\overline{\mathrm{\Psi }}_R^j(x),S_{LR}^{ij}(x)=\mathrm{\Psi }_L^i(0)\overline{\psi }_R^j(x).`$ (15)
In the lowest nontrivial order, we obtain following recursion relations ,
$`S_{LL}^{ij}(x)`$ $`=`$ $`{\displaystyle \frac{1}{g_2\mathrm{\Delta }^2(x)}}\left({\displaystyle \frac{1}{2a}}\right)^2{\displaystyle \underset{\mu }{\overset{}{}}}S_{RL}^{ij}(x+\mu )\gamma _\mu ,`$ (16)
$`S_{RL}^{ij}(x)`$ $`=`$ $`\left({\displaystyle \frac{1}{2a}}\right)\left({\displaystyle \frac{\delta (x)\delta _{ij}}{2g_2\mathrm{\Delta }^2(x)}}+{\displaystyle \frac{1}{g_2\mathrm{\Delta }^2(x)}}\left({\displaystyle \frac{1}{2a}}\right){\displaystyle \underset{\mu }{\overset{}{}}}S_{LL}^{ij}(x+\mu )\gamma _\mu \right),`$ (17)
$`S_{RR}^{ij}(x)`$ $`=`$ $`\left({\displaystyle \frac{1}{2a}}\right)^2{\displaystyle \frac{1}{g_2\mathrm{\Delta }^2(x)}}{\displaystyle \underset{\mu }{\overset{}{}}}\gamma _\mu S_{RL}^{ij}(x+\mu ).`$ (18)
The Fourier transformations of these recursion equations for $`p0`$ and $`\mathrm{\Delta }^2(p)=4w^2(p)0`$ lead to,
$`S_{LL}^{ij}(p)`$ $`=`$ $`{\displaystyle _x}e^{ipx}\psi _L^i(0)\overline{\psi }_L^j(x)P_L{\displaystyle \frac{\delta _{ij}\frac{i}{a}\underset{\mu }{}\mathrm{sin}p^\mu \gamma _\mu }{\frac{1}{a^2}_\mu \mathrm{sin}^2p_\mu +M^2(p)}}P_R,`$ (19)
$`S_{RL}^{ij}(p)`$ $`=`$ $`{\displaystyle _x}e^{ipx}\mathrm{\Psi }_R^i(0)\overline{\psi }_L^j(x)P_L{\displaystyle \frac{\delta _{ij}M(p)}{\frac{1}{a^2}_\mu \mathrm{sin}^2p_\mu +M^2(p)}}P_L,`$ (20)
$`S_{RR}^{ij}(p)`$ $`=`$ $`{\displaystyle _x}e^{ipx}\mathrm{\Psi }_R^i(0)\overline{\mathrm{\Psi }}_R^j(x)P_R{\displaystyle \frac{\delta _{ij}\frac{i}{a}\underset{\mu }{}\mathrm{sin}p^\mu \gamma _\mu }{\frac{1}{a^2}_\mu \mathrm{sin}^2p_\mu +M^2(p)}}P_L.`$ (21)
$`M(p)`$ $`=`$ $`8ag_2w^2(p).`$ (22)
However, these two-point Green functions (20-21) show only structure of poles, but not residues of these poles, which are related to the renormalization of Green functions with insertions of composite fermion operators $`\mathrm{\Psi }_R^i(x)`$.
The description of the renormalization of n-point 1PI functions with insertions of composite operators shows that the renormalized n-point 1PI functions $`\mathrm{\Gamma }_{ren}^{(n)}`$ with single and two insertions of composite operators are given by,
$`\mathrm{\Gamma }_{ren}^{(n)}(p_1,q_1,q_2,\mathrm{},q_n)`$ $`=`$ $`Z\mathrm{\Gamma }_{reg}^{(n)}(p_1,q_1,q_2,\mathrm{},q_n),`$ (23)
$`\mathrm{\Gamma }_{ren}^{(n)}(p_1,p_2,q_1,q_2,\mathrm{},q_n)`$ $`=`$ $`Z^2\mathrm{\Gamma }_{reg}^{(n)}(p_1,p_2,q_1,q_2,\mathrm{},q_n),`$ (24)
where $`\mathrm{\Gamma }_{reg}^{(n)}`$ are the regularized n-point 1PI functions, $`p_1`$ and $`p_2`$ stand for the momenta entering the composite operators. The renormalization constant $`Z`$ is the generalized “wave-function renormalization” or “form factor” of composite operators. It is worthwhile to stress that $`Z`$ is a finite positive constant and the wave-function renormalization of composite fields is the exactly same as the wave-function renormalization of elementary fields. In fact, composite particles are indistinguishable from elementary particles in this case.
To be consistent with this description of the renormalization of n-point 1PI functions with insertions of composite operators, the two-point Green functions (20-21) should be modified as follow,
$`S_{RR}^{ij}(p)`$ $`=`$ $`{\displaystyle _x}e^{ipx}\mathrm{\Psi }_R^i(0)\overline{\mathrm{\Psi }}_R^j(x)\delta _{ij}{\displaystyle \frac{Z_R(p)\frac{i}{a}\underset{\mu }{}\mathrm{sin}p^\mu \gamma _\mu Z_R(p)}{\frac{1}{a^2}_\mu \mathrm{sin}^2p_\mu +M^2(p)}}P_L;`$ (25)
$`S_{RL}^{ij}(p)`$ $`=`$ $`{\displaystyle _x}e^{ipx}\mathrm{\Psi }_R^i(0)\overline{\psi }_L^j(x)\delta _{ij}Z_R(p){\displaystyle \frac{M(p)}{\frac{1}{a^2}\mathrm{sin}^2p+M^2(p)}}P_R.`$ (26)
These regularized two-point Green functions with one and two insertions of composite fermion operators identify not only their poles, but also the corresponding residues. The residues $`Z_R(p)`$ (25-26), i.e., the generalized form factors of the composite three-fermion-states (13), are given by one-particle irreducible (1PI) functions,
$$Z_R(p)=_xe^{ipx}\frac{\delta ^{(2)}\mathrm{\Gamma }}{\delta \mathrm{\Psi }_R^i(x)\delta \overline{\psi }_L^j(0)}aM(p),Z_L(p)=_xe^{ipx}\frac{\delta ^{(2)}\mathrm{\Gamma }}{\delta \mathrm{\Psi }_L^n(x)\delta \overline{\chi }_R^{}(0)}=aM(p).$$
(27)
The “primed fields” $`\mathrm{\Psi }_L^n(x)`$ and $`\mathrm{\Psi }_R^i(x)`$ of three-fermion-states are defined by eqs.(41) and (42) in ref. through the generating functional approach and $`Z_R(p)`$ is obtained by the strong coupling expansion for $`p\pi _A0`$. In eq.(27), $`Z_L(p)`$ is the renormalization constant for the three-fermion-state $`\mathrm{\Psi }_L^n(x)`$ in the neutral sector, which is exactly resulted from the Ward identity (9). These residues $`Z_{L,R}(p\pi _A)`$ (generalized form factors of three-fermion-states $`\mathrm{\Psi }_R^i(x),\mathrm{\Psi }_L^n(x)`$) are different positive (non-zero) constants with respect to each doubler ($`p\pi _A`$).
We make a wave-function renormalization of three-fermion-states with respect to each doubler $`p=\pi _A0`$,
$$\mathrm{\Psi }_R^i|_{ren}=Z_R^1(p)\mathrm{\Psi }_R^i;\mathrm{\Psi }_L^n|_{ren}=Z_L^1(p)\mathrm{\Psi }_L^n.$$
(28)
As a result for $`p=\pi _A0`$, the neutral $`\mathrm{\Psi }_n`$ and charged $`\mathrm{\Psi }_c^i`$ Dirac fermions are formed,
$$\mathrm{\Psi }_c^i=(\psi _L^i,\mathrm{\Psi }_R^i|_{ren});\mathrm{\Psi }_n=(\mathrm{\Psi }_L^n|_{ren},\chi _R),$$
(29)
whose propagators are
$`S_c^{ij}(p)={\displaystyle _x}e^{ipx}\mathrm{\Psi }_c^i(0)\overline{\mathrm{\Psi }}_c^j(x)`$ $``$ $`\delta _{ij}{\displaystyle \frac{\frac{i}{a}\gamma _\mu \mathrm{sin}(p)^\mu +M(p)}{\frac{1}{a^2}\mathrm{sin}^2p+M^2(p)}};`$ (30)
$`S_n(p)={\displaystyle _x}e^{ipx}\mathrm{\Psi }_n(0)\overline{\mathrm{\Psi }}_n(x)`$ $``$ $`{\displaystyle \frac{\frac{i}{a}\gamma _\mu \mathrm{sin}(p)^\mu +M(p)}{\frac{1}{a^2}\mathrm{sin}^2p+M^2(p)}}.`$ (31)
These show that all doublers $`(p=\pi _A)`$ are decoupled as very massive Dirac fermions consistently with the $`SU_L(2)U_R(1)`$ chiral symmetries, since the three-fermion-states (13) carry the appropriate quantum numbers of the chiral groups that accommodate $`\psi _L^i`$ and $`\chi _R`$.
Due to the locality of action (1), all Green functions and 1PI functions must be analytically continuous functions in energy-momentum space, provided the dynamics is fixed by given multifermion couplings $`g_1`$ and $`g_2`$. Although two-point Green functions (19) and (25-26) are obtained from the strong coupling expansion $`g_21`$ for the momentum $`p=\pi _A0`$, we can make analytical continuation of these two-point Green functions from $`p=\pi _A`$ to $`p=0`$ in the energy-momentum space, provided no hard spontaneous symmetry breaking takes place (11,12) as the effective multifermion coupling $`g_2(p)`$ is reduced. Since the residues $`Z_{L,R}(p)`$ positively vanish ($`Z_{L,R}(p)O(p^8)`$) in the low-energy limit $`p0`$, eq.(25) has no simple pole at $`p0`$ and mixing (26) vanishes as well. For $`Z_{L,R}(0)=0`$, we are not allowed to make the wave-function renormalization (28) with respect to $`p0`$ and obtain Dirac fermions (30,31) at $`p0`$. While, the propagator (19) of the elementary field $`\psi _L^i`$ has a simple pole at $`p0`$ behaving as a charged chiral particle $`\psi _L^i(x)`$ consistently with the $`SU_L(2)`$ symmetry ($`\stackrel{~}{p}^\mu `$ is the dimensionful continuum momentum),
$$S_{LL}^1(\stackrel{~}{p})^{ij}=i\gamma _\mu \stackrel{~}{p}^\mu \stackrel{~}{Z}_2\delta _{ij}P_L;S_{RR}^1(\stackrel{~}{p})=i\gamma _\mu \stackrel{~}{p}^\mu P_R,$$
(32)
where $`\stackrel{~}{Z}_2`$ is the finite wave-function renormalization constant of the elementary interpolating field $`\psi _L^i(x)`$ and $`S_{RR}^1(\stackrel{~}{p})`$ (10) is the inverse propagator of $`\chi _R(x)`$ at $`\stackrel{~}{p}0`$. All these properties discussed were also obtained by the weak coupling expansion in ref..
The vanishing of $`Z_{L,R}(p)`$ and eqs.(25\- 26) at $`p0`$ indicates that three-fermion-states $`\mathrm{\Psi }_R^i`$ and $`\mathrm{\Psi }_L^n`$ dissolve into the virtual states of three individual chiral fermions $`\psi _L^i(x)`$ and $`\chi _R(x)`$ with the fixed total momentum $`p`$ and a continuous energy spectrum. The “no-go” theorem is entirely inapplicable to these virtual states. We call these virtual states three-fermion-cuts $`𝒞[\mathrm{\Psi }_L^n(x)]`$ and $`𝒞[\mathrm{\Psi }_R^i(x)]`$. These two virtual states carry the exactly same quantum numbers and total momentum $`p`$ as that of the corresponding three-fermion-states. Thus, chiral gauge symmetries are preserved in such a dissolving phenomenon. The energy-threshold $`ϵ`$ of such a dissolving, where the energy-gap between three-fermion-states and their corresponding virtual states goes to zero, must locate at ($`\stackrel{~}{v}\frac{\pi }{a}`$ being a soft spontaneous symmetry breaking at the electroweak scale)
$$\stackrel{~}{v}<ϵ<\frac{\pi }{a},$$
(33)
whose value depends on the values of the multifermion couplings $`g_1`$ and $`g_2`$ in the scaling region $`𝒜`$ (8).
As results, the spectrum of the model (1) in the scaling region $`𝒜`$ (8) is the following. It consists of 15 copies of $`SU_L(2)`$-charged Dirac doublers eq.(30) and 15 copies of $`S_LU(2)`$-neutral Dirac doublers eq.(31). They are very massive and decoupled. Beside, the low energy spectrum contains the massless normal modes eqs.(32) for $`p=\stackrel{~}{p}0`$.
Before ending this section, we emphasize that our scenario in the scaling region is resulted from the dynamics of the special dimension-10 operator $`S_2`$ in the action (1). It is worthwhile to mention that we can replace the Wilson operator $`\mathrm{\Delta }`$ in $`S_2`$ by $`\mathrm{\Delta }_\mu `$ defined as
$`\mathrm{\Delta }_\mu \chi _R(x)`$ $``$ $`\chi _R(x+\mu )\chi _R(x\mu ),`$ (34)
$`2w_\mu (p)`$ $`=`$ $`{\displaystyle _x}e^{ipx}\mathrm{\Delta }_\mu (x)=\mathrm{sin}(p_\mu ),`$ (35)
and the multifermion coupling $`S_2`$ in the action (1) is substituted by
$$\stackrel{~}{S}_2=g_2\underset{x,\mu }{}\overline{\psi }_L^i(x)\left[\mathrm{\Delta }_\mu \chi _R(x)\right]\left[\mathrm{\Delta }_\mu \overline{\chi }_R(x)\right]\psi _L^i(x),$$
(36)
which is a dimension-8 operator. Although this action $`\stackrel{~}{S}_2`$ has the exactly same $`SU_L(2)`$ chiral gauge symmetry, $`SU_L(2)U_R(1)`$ global chiral symmetries and $`\chi _R`$-shift-symmetry, one can check such a dimension-8 operator does not process the proper properties presented in this section. Another example is the multifermion coupling $`S_1`$ with dimension-6 operators that suffers from a hard spontaneous symmetry breaking as that in ref.. These two examples tell us that in principle, chiral gauge symmetries can be preserved by multifermion couplings, but in practice, multifermion couplings do not definitely have desired dynamics, which are: (i) non hard spontaneous symmetry breaking; (ii) the dissolving of three-fermion-state to the three-fermion-cut in the low-energy; (iii) doublers gauge-invariantly decoupled as massive particles and (iv) spectator fermions decoupled as free fermions, to have the scaling region for an asymptotic chiral gauge theories in the continuum limit.
## III The gauge coupling vertices in the neutral sector
In the computations by the strong coupling expansion to obtain the fermion spectrum discussed in section 2, the $`SU_L(2)`$ gauge interaction is neglected as its perturbative nature with respect to the strong multifermion coupling $`g_21`$. In the following two sections, we turn on the gauge field as a dynamical field to examine (i) whether the spectator field $`\chi _R`$ decouples from the gauge field; (ii) whether charged chiral fermion couples to the gauge field in chiral manner, as required by an asymptotic chiral gauge theory in the continuum limit.
As the $`SU_L(2)`$ chiral gauge coupling $`g`$ is turned on and the action (1) is $`SU_L(2)`$-chirally gauged, the properties of the scaling region $`𝒜`$ (8) should not be greatly altered for the reasons that the $`SU_L(2)`$-chiral gauge interaction does not spoil the Ward identity (9) of the $`\chi _R`$-shift-symmetry and the spectrum is gauge symmetric in the scaling region $`𝒜`$ (8).
Now we consider all possible interacting vertices (one particle irreducible (1PI) functions) involving truncated external gauge field $`A_\mu ^{}`$ and fermion fields $`\chi _R^{},\mathrm{\Psi }_L^n`$ in the neutral sector. Based on the Ward identity (9) of the $`\chi _R`$-shift-symmetry, we take functional derivatives with respect of the external gauge field $`A_\mu ^{}`$, fermion fields $`\chi _R^{}`$ and $`\mathrm{\Psi }_L^n`$, to obtain the following Ward identities,
$$\frac{\delta ^{(2)}\mathrm{\Gamma }}{\delta A_\mu ^{}\delta \overline{\chi }_R^{}}=\frac{\delta ^{(3)}\mathrm{\Gamma }}{\delta A_\mu ^{}\delta \chi _R^{}\delta \overline{\psi }_R^{}}=\frac{\delta ^{(3)}\mathrm{\Gamma }}{\delta A_\mu ^{}\delta \mathrm{\Psi }_L^n\delta \overline{\chi }_R^{}}=\mathrm{}=0.$$
(37)
As a consequence of these Ward identities and identical vanishing of interacting 1PI functions containing external gauge field $`A_\mu ^{}`$, “spectator” fermions $`\chi _R^{}(x)`$ and neutral three-fermion states field $`\mathrm{\Psi }_L^n(x)`$, we prove no any interactions between the gauge field $`A_\mu `$, the “spectator” fermion $`\chi _R`$ and the neutral three-fermion states $`\mathrm{\Psi }_L^n(x)`$. It should be emphasized that the decoupling of the neutral sector ($`\chi _R,\mathrm{\Psi }_L^n`$) from the gauge field $`A_\mu `$ is valid not only for perturbative gauge-interaction but also non-perturbative gauge-interaction, since the action is exactly $`SU_L(2)`$-chiral-gauge symmetric and $`\chi _R`$-shift-symmetric for any values of the gauge coupling and multifermion couplings.
As a result, The right-handed current $`j_R^\mu =i\overline{\chi }_R\gamma ^\mu \chi _R`$ is exactly conserved i.e., $`_\mu j^\mu =0`$ and no anomalous contribution is expected from the topological gauge configuration. This is a very important feature, which we will see later in section 5 and 6, for obtaining the gauge anomaly and the fermion-flavour singlet anomaly relating to the $`U_L(1)`$-symmetry (the number of the fermion field $`\psi _L^i`$) violation.
## IV Gauge coupling vertices in the charged sector
In this section we turn to directly compute three-point interacting vertices of gauge field $`A_\nu ^a(x)`$, elementary and composite fermion fields $`\psi _L^i(x)`$ and $`\mathrm{\Psi }_R^i(x)`$, in order to show the gauge-fermion coupling is chiral in the low-energy region, while vector-like in the high-energy region.
Setting 1PI vertex functions of the gauge-fermion coupling to be $`\mathrm{\Lambda }_\mu ^a(p,p^{})`$ and $`q=p^{}+p`$, where $`p^{}`$ and $`p`$ are two external momenta of fermion fields and $`q`$ is the momentum of the gauge field, we can write the following three-point Green functions in the momentum space:
$`{\displaystyle _{x_1xy}}e^{i(p^{}x+px_1qy)}\psi _L(x_1)\overline{\psi }_L(x)A_\nu ^a(y)`$ $`=`$ $`G_{\nu \mu }^{ab}(q)S_{LL}(p)\mathrm{\Lambda }_{\mu LL}^b(p,p^{})S_{LL}(p^{});`$ (38)
$`{\displaystyle _{x_1xy}}e^{i(p^{}x+px_1qy)}\psi _L(x_1)\overline{\mathrm{\Psi }}_R(x)A_\nu ^a(y)`$ $`=`$ $`G_{\nu \mu }^{ab}(q)S_{LL}(p)\mathrm{\Lambda }_{\mu LR}^b(p,p^{})S_{RR}(p^{});`$ (39)
$`{\displaystyle _{x_1xy}}e^{i(p^{}x+px_1qy)}\mathrm{\Psi }_R(x_1)\overline{\mathrm{\Psi }}_L(x)A_\nu ^a(y)`$ $`=`$ $`G_{\nu \mu }^{ab}(q)S_{RR}(p)\mathrm{\Lambda }_{\mu RL}^b(p,p^{})S_{LL}(p^{});`$ (40)
$`{\displaystyle _{x_1xy}}e^{i(p^{}x+px_1qy)}\mathrm{\Psi }_R(x_1)\overline{\mathrm{\Psi }}_R(x)A_\nu ^a(y)`$ $`=`$ $`G_{\nu \mu }^{ab}(q)S_{RR}(p)\mathrm{\Lambda }_{\mu RR}^b(p,p^{})S_{RR}(p^{}),`$ (41)
where $`G_{\nu \mu }^{ab}(q)`$ is the propagator of gauge field $`A_\nu ^a(x)`$; the $`S_{LL}(p)`$ and $`S_{RR}(p)`$ are the propagators of the elementary and composite chiral fermions $`\psi _L(x)`$ and $`\mathrm{\Psi }_R(x)`$, we omit henceforth the $`SU_L(2)`$ indices $`i`$ and $`j`$. We try to compute the 1PI vertex functions $`\mathrm{\Lambda }_{\mu LL}^b(p,p^{}),\mathrm{\Lambda }_{\mu RR}^b(p,p^{})`$, $`\mathrm{\Lambda }_{\mu LR}^b(p,p^{})`$ and $`\mathrm{\Lambda }_{\mu RL}^b(p,p^{})`$ in eqs.(38-41).
Using the small gauge coupling ($`g`$) expansion, one can find the 1PI vertex functions $`\mathrm{\Lambda }_{\mu LL}^b(p,p^{})`$ by calculating
$`\psi _L(x_1)\overline{\psi }_L(x)A_\mu ^a(y)=i{\displaystyle \frac{g}{2}}\left({\displaystyle \frac{\tau ^a}{2}}\right)\psi _L(x_1)\overline{\psi }_L(x)\gamma _\rho `$ (42)
$`{\displaystyle _z}\left[\psi _L(z+\rho )\overline{\psi }_L(x)A_\rho ^b(z+{\displaystyle \frac{\rho }{2}})A_\mu ^a(y)+\psi _L(z\rho )\overline{\psi }_L(x)A_\rho ^b(z{\displaystyle \frac{\rho }{2}})A_\mu ^a(y)\right],`$ (43)
and obtains
$`\mathrm{\Lambda }_{\mu LL}^{(1)}(p,p^{})`$ $`=`$ $`ig\left({\displaystyle \frac{\tau ^a}{2}}\right)\mathrm{cos}\left({\displaystyle \frac{p+p^{}}{2}}\right)_\mu \gamma _\mu P_L,`$ (44)
$`\mathrm{\Lambda }_{\mu \nu LL}^{(2)}(p,p^{})`$ $`=`$ $`i{\displaystyle \frac{g^2}{2}}\left({\displaystyle \frac{\tau ^a\tau ^b}{4}}\right)\mathrm{sin}\left({\displaystyle \frac{p+p^{}}{2}}\right)_\mu \gamma _\mu \delta _{\mu \nu }P_L,`$ (45)
$`\mathrm{},`$ (46)
in terms of the powers of the small gauge coupling $`(g)`$.
Then, we try to find the relationships between the 1PI vertex function $`\mathrm{\Lambda }_{\mu LL}^b(p,p^{})`$ and other three 1PI vertex functions $`\mathrm{\Lambda }_{\mu RR}^b(p,p^{}),\mathrm{\Lambda }_{\mu LR}^b(p,p^{}),\mathrm{\Lambda }_{\mu RL}^b(p,p^{})`$ in eqs.(39-41). By the strong coupling expansion in powers of $`1/g_2`$ for $`p,p^{}\pi _A0`$, analogously to recursion relations (16-18), we obtain the following recursion relations at the nontrivial order,
$`\psi _L(x_1)\overline{\psi }_L(x)A_\nu ^a(y)`$ $`=`$ $`{\displaystyle \frac{1}{g_2\mathrm{\Delta }^2(x)}}\left({\displaystyle \frac{1}{2a}}\right)^2{\displaystyle \underset{\rho }{\overset{}{}}}\psi _L(x_1)\overline{\mathrm{\Psi }}_R(x+\rho )A_\nu ^a(y)\gamma _\rho ,`$ (47)
$`\psi _L(x_1)\overline{\psi }_L(x)A_\nu ^a(y)`$ $`=`$ $`{\displaystyle \frac{1}{g_2\mathrm{\Delta }^2(x_1)}}\left({\displaystyle \frac{1}{2a}}\right)^2{\displaystyle \underset{\rho }{\overset{}{}}}\gamma _\rho \mathrm{\Psi }_R(x_1+\rho )\overline{\psi }_L(x)A_\nu ^a(y),`$ (48)
$`\mathrm{\Psi }_R(x_1)\overline{\mathrm{\Psi }}_R(x)A_\nu ^a(y)`$ $`=`$ $`{\displaystyle \frac{1}{g_2\mathrm{\Delta }^2(x)}}\left({\displaystyle \frac{1}{2a}}\right)^2{\displaystyle \underset{\rho }{\overset{}{}}}\gamma _\rho \psi _L(x_1)\overline{\mathrm{\Psi }}_R(x+\rho )A_\nu ^a(y).`$ (49)
We make the Fourier transform in both sides of the above recursion relations and obtain for $`p,p^{}\pi _A0`$,
$`S_{LL}(p)\mathrm{\Lambda }_{\mu LL}^a(p,p^{})S_{LL}(p^{})`$ $`=`$ $`{\displaystyle \frac{i}{aM(p^{})}}S_{LL}(p)\mathrm{\Lambda }_{\mu LR}^a(p,p^{})S_{RR}(p^{}){\displaystyle \underset{\rho }{}}\mathrm{sin}p_\rho ^{}\gamma ^\rho ,`$ (50)
$`S_{LL}(p)\mathrm{\Lambda }_{\mu LL}^a(p,p^{})S_{LL}(p^{})`$ $`=`$ $`{\displaystyle \frac{i}{aM(p)}}{\displaystyle \underset{\rho }{}}\mathrm{sin}p_\rho \gamma ^\rho S_{RR}(p)\mathrm{\Lambda }_{\mu RL}^a(p,p^{})S_{LL}(p^{}),`$ (51)
$`S_{RR}(p)\mathrm{\Lambda }_{\mu RR}^a(p,p^{})S_{RR}(p^{})`$ $`=`$ $`{\displaystyle \frac{i}{aM(p^{})}}{\displaystyle \underset{\rho }{}}\mathrm{sin}p_\rho ^{}\gamma ^\rho S_{LL}(p)\mathrm{\Lambda }_{\mu LR}^a(p,p^{})S_{RR}(p^{}),`$ (52)
where the propagator of the gauge boson $`G_{\nu \mu }^{ab}(q)`$ is eliminated from the both sides of equations.
Using these recursion relations (50-52) and the propagators $`S_{LL}(p),S_{RR}(p)`$ (19,21) obtained by the same strong coupling expansion $`1/g_2`$ for $`p\pi _A0`$, we can compute the 1PI vertex functions $`\mathrm{\Lambda }_{\mu RL}^a(p,p^{})`$, $`\mathrm{\Lambda }_{\mu LR}^a(p,p^{})`$ and $`\mathrm{\Lambda }_{\mu RR}^a(p,p^{})`$ eqs.(39-41) in terms of the 1PI vertex function $`\mathrm{\Lambda }_{\mu LL}^a(p,p^{})`$),
$`M(p^{})\mathrm{\Lambda }_{\mu LL}^a(p,p^{})`$ $`=`$ $`\mathrm{\Lambda }_{\mu LR}^a(p,p^{})\left({\displaystyle \frac{i}{a}}\right){\displaystyle \underset{\rho }{}}\mathrm{sin}p_\rho ^{}\gamma ^\rho ,`$ (53)
$`M(p)\mathrm{\Lambda }_{\mu LL}^a(p,p^{})`$ $`=`$ $`\left({\displaystyle \frac{i}{a}}\right){\displaystyle \underset{\rho }{}}\mathrm{sin}p_\rho \gamma ^\rho \mathrm{\Lambda }_{\mu RL}^a(p,p^{}),`$ (54)
$`M(p^{})\mathrm{\Lambda }_{\mu RR}^a(p,p^{})`$ $`=`$ $`\left({\displaystyle \frac{i}{a}}\right){\displaystyle \underset{\rho }{}}\mathrm{sin}p_\rho ^{}\gamma ^\rho \mathrm{\Lambda }_{\mu LR}^a(p,p^{}).`$ (55)
These relationships are independent of the strength of the gauge coupling $`g`$. Taking $`\mathrm{\Lambda }_{\mu LL}^a(p,p^{})`$ to be eq.(44) at the leading order, we obtain perturbative results
$`\mathrm{\Lambda }_{\mu RR}^{(1)}(p,p^{})`$ $`=`$ $`ig\left({\displaystyle \frac{\tau ^a}{2}}\right)\mathrm{cos}\left({\displaystyle \frac{p+p^{}}{2}}\right)_\mu \gamma _\mu P_R,`$ (56)
$`\mathrm{\Lambda }_{\mu LR}^{(1)}(p,p^{})\left({\displaystyle \frac{i}{a}}\right)\mathrm{sin}p_\mu ^{}`$ $`=`$ $`{\displaystyle \frac{1}{2}}M(p^{})ig\left({\displaystyle \frac{\tau ^a}{2}}\right)\mathrm{cos}\left({\displaystyle \frac{p+p^{}}{2}}\right)_\mu ,`$ (57)
$`\left({\displaystyle \frac{i}{a}}\right)\mathrm{sin}p_\mu \mathrm{\Lambda }_{\mu RL}^{(1)}(p,p^{})`$ $`=`$ $`{\displaystyle \frac{1}{2}}M(p)ig\left({\displaystyle \frac{\tau ^a}{2}}\right)\mathrm{cos}\left({\displaystyle \frac{p+p^{}}{2}}\right)_\mu .`$ (58)
According to the renormalization of 1PI functions with composite operator insertions (24), we have the renormalized 1PI vertex functions for $`p,p^{}0`$ (see Figure 5 in Ref.),
$`\mathrm{\Lambda }_{\mu RR}^a(p,p^{})|_{ren}`$ $`=`$ $`Z_R(p)Z_R(p^{})\mathrm{\Lambda }_{\mu RR}^a(p,p^{}),`$ (59)
$`\mathrm{\Lambda }_{\mu RL}^a(p,p^{})|_{ren}`$ $`=`$ $`Z_R(p)\mathrm{\Lambda }_{\mu RL}^a(p,p^{}),`$ (60)
$`\mathrm{\Lambda }_{\mu LR}^a(p,p^{})|_{ren}`$ $`=`$ $`Z_R(p^{})\mathrm{\Lambda }_{\mu LR}^a(p,p^{}),`$ (61)
where $`\mathrm{\Lambda }_{\mu RR}^a(p,p^{}),\mathrm{\Lambda }_{\mu LR}^a(p,p^{})`$ and $`\mathrm{\Lambda }_{\mu RL}^a(p,p^{})`$ are the regularized 1PI vertex functions obtained in eqs.(53-55) in terms of $`\mathrm{\Lambda }_{\mu LL}(p,p^{})`$. For $`p,p^{}=\pi _A0`$ with respect to each doublers, $`Z_R(p)`$ and $`Z_R(p^{})0`$, and we can make the renormalization (28) of the composite fermion-operator (field) $`\mathrm{\Psi }_R(x)`$ so that,
$`\mathrm{\Lambda }_{\mu RR}^a(p,p^{},\mathrm{\Psi }_R|_{ren})`$ $`=`$ $`\mathrm{\Lambda }_{\mu RR}^a(p,p^{}),`$ (62)
$`\mathrm{\Lambda }_{\mu RL}^a(p,p^{},\mathrm{\Psi }_R|_{ren})`$ $`=`$ $`\mathrm{\Lambda }_{\mu RL}^a(p,p^{}),`$ (63)
$`\mathrm{\Lambda }_{\mu LR}^a(p,p^{},\mathrm{\Psi }_R|_{ren})`$ $`=`$ $`\mathrm{\Lambda }_{\mu LR}^a(p,p^{}),`$ (64)
where $`\mathrm{\Lambda }_{\mu RR}^a(p,p^{},\mathrm{\Psi }_R|_{ren}),\mathrm{\Lambda }_{\mu LR}^a(p,p^{},\mathrm{\Psi }_R|_{ren})`$ and $`\mathrm{\Lambda }_{\mu RL}^a(p,p^{},\mathrm{\Psi }_R|_{ren})`$ are the renormalized 1PI vertex functions in terms of the renormalized composite fermion-operator(field) $`\mathrm{\Psi }_R|_{ren}`$ (28). Correspondingly, the charged Dirac fermion $`\mathrm{\Psi }_c=(\psi _L,\mathrm{\Psi }_R|_{ren})`$ and its propagator $`S_c(p)`$ in terms of the renormalized composite fermion-operator(field) $`\mathrm{\Psi }_R|_{ren}`$ (28) is given by eq.(30) in section 2.
The interacting 1PI vertex function between the gauge field $`A_\mu ^a(x)`$ and the charged Dirac fermion $`\mathrm{\Psi }_c(x)`$ is related to the following three-point Green functions,
$`\mathrm{\Psi }_c(x_1)\overline{\mathrm{\Psi }}_c(x)A_\nu ^a(y)`$ $`=`$ $`\psi _L(x_1)\overline{\psi }_L(x)A_\nu ^a(y)+\psi _L(x_1)\overline{\mathrm{\Psi }}_R(x)|_{ren}A_\nu ^a(y)`$ (65)
$`+`$ $`\mathrm{\Psi }_R(x_1)|_{ren}\overline{\psi }_L(x)A_\nu ^a(y)+\mathrm{\Psi }_R(x_1)|_{ren}\overline{\mathrm{\Psi }}_R(x)|_{ren}A_\nu ^a(y),`$ (66)
and
$$_{x_1xy}e^{i(p^{}x+px_1qy)}\mathrm{\Psi }_c(x_1)\overline{\mathrm{\Psi }}_c(x)A_\nu ^a(y)=G_{\nu \mu }^{ab}(q)S_c(p)\mathrm{\Lambda }_{\mu c}^b(p,p^{},\mathrm{\Psi }_R|_{ren})S_c(p^{}),$$
(67)
where $`\mathrm{\Lambda }_{\mu c}^b(p,p^{},\mathrm{\Psi }_R|_{ren})`$ is the renormalized 1PI vertex function of the gauge field $`A_\mu ^a(x)`$ and the charged Dirac fermion $`\mathrm{\Psi }_c=(\psi _L,\mathrm{\Psi }_R|_{ren})`$.
From eqs.(38-41) and eqs.(66-67), the 1PI vertex function $`\mathrm{\Lambda }_{\mu c}(p,p^{},\mathrm{\Psi }_R|_{ren})`$ is given by
$`\mathrm{\Lambda }_{\mu c}(p,p^{},\mathrm{\Psi }_R|_{ren})^{(1)}`$ $`=`$ $`\mathrm{\Lambda }_{\mu LL}(p,p^{})^{(1)}+\mathrm{\Lambda }_{\mu LR}(p,p^{},\mathrm{\Psi }_R|_{ren})^{(1)}+\mathrm{\Lambda }_{\mu RL}(p,p^{},\mathrm{\Psi }_R|_{ren})^{(1)}`$ (68)
$`+`$ $`\mathrm{\Lambda }_{\mu RR}(p,p^{},\mathrm{\Psi }_R|_{ren})^{(1)},`$ (69)
up to the first order of the perturbative gauge coupling $`g`$. One can check that the 1PI vertex functions (56-58) and the renormalization for these vertex function (59-61) with respect to each doublers $`p^{},p=\pi _A0`$ precisely obey the following extremely important Ward identity of the exact $`SU_L(2)`$ chiral gauge symmetry at the lattice scale $`\pi /a`$,
$$\left(\frac{i}{a}\right)(\mathrm{sin}p_\mu \mathrm{sin}p_\mu ^{})\mathrm{\Lambda }_{\mu c}^{(1)}(p,p^{},\mathrm{\Psi }_R|_{ren})=S_c^1(p)S_c^1(p^{}).$$
(70)
where the gauge coupling $`g`$ and gauge group generator $`\tau _a/2`$ are eliminated from the vertex function $`\mathrm{\Lambda }_{\mu c}(p,p^{},\mathrm{\Psi }_R|_{ren})`$. This shows that the exact $`SU_L(2)`$ chiral gauge symmetry is realized by the vector-like and massive spectrum of Dirac modes $`p^{},p=\pi _A0`$ (29). Such results are expected, since we are in the symmetric phase ($`1g_2<\mathrm{}`$). These calculations can be straightforwardly generalized to higher orders of the perturbative expansion in powers of the gauge coupling $`g`$.
So far, by the strong coupling expansion for $`p^{},p=\pi _A0`$, we have computed the 1PI regularized vertex functions $`\mathrm{\Lambda }_{\mu RR}^b(p,p^{}),\mathrm{\Lambda }_{\mu LR}^b(p,p^{})`$ and $`\mathrm{\Lambda }_{\mu RL}^b(p,p^{})`$ (53-55) in terms of $`\mathrm{\Lambda }_{\mu LL}^b(p,p^{})`$ (38,44), and the corresponding 1PI renormalized vertex functions (59)-(61). Clearly the computations by the strong coupling expansion are broken down for $`p^{},p0`$. However, these 1PI renormalized vertex functions (62-64) are analytical continuous functions in the energy-momentum plane for the locality of the action (1). We can make an analytical continuation of the 1PI renormalized vertex functions (62)-(64) from $`p^{},p=\pi _A0`$ to $`p^{},p0`$, provided no hard spontaneous symmetry breaking takes place (11,12). Eqs.(59-61) show that the renormalized 1PI vertex functions $`\mathrm{\Lambda }_{\mu RR}^a(p,p^{})|_{ren},\mathrm{\Lambda }_{\mu LR}^a(p,p^{})|_{ren}`$ and $`\mathrm{\Lambda }_{\mu RL}^a(p,p^{})|_{ren}`$ vanish as $`Z_R(p)`$ and $`Z_R(p^{})`$ vanish for $`p^{},p0`$. Furthermore, due to $`Z_R(p)`$ vanishes for $`p0`$, we are not allowed to make the renormalization of the composite fermion-operator(field) $`\mathrm{\Psi }_R(x)`$ (28) to obtain eqs.(62-64) and (69). As a result, the interacting 1PI vertex function between the gauge field $`A_\nu ^a(x)`$ and charged fermion turns out to be $`\mathrm{\Lambda }_{\mu LL}(p,p^{})`$ for $`p^{},p0`$ and the Ward identity (70) is reduced to its counterpart of the continuum theory,
$$\left(\frac{i}{a}\right)(\mathrm{sin}p_\mu \mathrm{sin}p_\mu ^{})\mathrm{\Lambda }_{\mu LL}^{(1)}(p,p^{})=S_{LL}^1(p)S_{LL}^1(p^{}),p^{},p0$$
(71)
where the propagator $`S_{LL}(p)`$ of elementary chiral fermion $`\psi _L^i(x)`$ is given by eq.(19). We recall again that two-point functions (25) and (26) vanish as $`p`$ goes to zero. This Ward identity (71) is consistent with the $`SU_L(2)`$ chiral gauge symmetry realized by the chiral spectrum.
The dissolving of the composite three-fermion-state $`\mathrm{\Psi }_R(x)`$ into its three-fermion-cut $`𝒞[\mathrm{\Psi }_R(x)]`$ at the energy-threshold (33) described in the previous section, results in the vanishing of the 1PI vertex functions $`\mathrm{\Lambda }_{\mu RR}^a(p,p^{})|_{ren},\mathrm{\Lambda }_{\mu LR}^a(p,p^{})|_{ren}`$ and $`\mathrm{\Lambda }_{\mu RL}^a(p,p^{})|_{ren}`$in eqs.(59-61). This implies the decoupling between the gauge field $`A_\mu ^a(x)`$ and the composite three-fermion-state $`\mathrm{\Psi }_R(x)`$ (13) in the low-energy.
All discussions in this section are on the basis of (i) the 1PI vertex functions (59,61) computed by the strong coupling expansion at $`p^{},p=\pi _A0`$ and (ii) the 1PI vertex functions (59-61) analytically extrapolated to $`p^{},p0`$ by the analytical continuation property of these 1PI vertex functions in the energy-momentum plane in the segment $`𝒜`$ (8). However, it is worthwhile on the other hand to directly compute these 1PI vertex functions in the region of $`p^{},p0`$ by the weak coupling expansion.
By the weak coupling expansion for $`p^{},p0`$, we can directly compute these 1PI vertex functions (59,61) to see whether the consistent results can be reached. For the small gauge coupling, analogous to eq.(43), the three-point Green functions with the insertions of composite three-fermion-operator(field) are given by,
$`\mathrm{\Psi }_R(x_1)\overline{\mathrm{\Psi }}_R(x)A_\mu ^a(y)=i{\displaystyle \frac{g}{2}}\left({\displaystyle \frac{\tau ^a}{2}}\right)\mathrm{\Psi }_R(x_1)\overline{\mathrm{\Psi }}_R(x)\gamma _\rho `$ (72)
$`{\displaystyle _z}\left[\mathrm{\Psi }_R(z+\rho )\overline{\mathrm{\Psi }}_R(x)A_\rho ^b(z+{\displaystyle \frac{\rho }{2}})A_\mu ^a(y)+\mathrm{\Psi }_R(z\rho )\overline{\mathrm{\Psi }}_R(x)A_\rho ^b(z{\displaystyle \frac{\rho }{2}})A_\mu ^a(y)\right],`$ (73)
and
$`\psi _L(x_1)\overline{\mathrm{\Psi }}_R(x)A_\mu ^a(y)=i{\displaystyle \frac{g}{2}}\left({\displaystyle \frac{\tau ^a}{2}}\right)\psi _L(x_1)\overline{\mathrm{\Psi }}_R(x)\gamma _\rho `$ (74)
$`{\displaystyle _z}\left[\psi _L(z+\rho )\overline{\mathrm{\Psi }}_R(x)A_\rho ^b(z+{\displaystyle \frac{\rho }{2}})A_\mu ^a(y)+\psi _L(z\rho )\overline{\mathrm{\Psi }}_R(x)A_\rho ^b(z{\displaystyle \frac{\rho }{2}})A_\mu ^a(y)\right].`$ (75)
Actually, in the section 6 of ref., we made the calculations of two-point Green functions $`\mathrm{\Psi }_R(x_1)\overline{\mathrm{\Psi }}_R(x)`$ and $`\psi _L(x_1)\overline{\mathrm{\Psi }}_R(x)`$ by the weak coupling expansion for the small momenta $`p0`$ in the segment $`𝒜`$ (8). The results show that these two-point Green functions in the momentum space,
$$S_{RR}(p)=_xe^{ipx}\mathrm{\Psi }_R(0)\overline{\mathrm{\Psi }}_R(x),S_{LR}(p)=_xe^{ipx}\psi _L(0)\overline{\mathrm{\Psi }}_R(x)$$
(76)
have no simple poles at $`p0`$ contrasting with that in eq.(20,21), instead are regular and vanish as $`p0`$. As a direct consequence, the 1PI vertex functions $`\mathrm{\Lambda }_{\mu RR}^a(p,p^{})|_{ren}`$ and $`\mathrm{\Lambda }_{\mu LR}^a(p,p^{})|_{ren}`$ relating to the three-point Green functions (73,75) vanish as their external momenta $`p`$ and $`p^{}`$ vanish. The Ward identity (70) is reduced to (71) for the small external momenta $`p,p^{}0`$. The results are the same as that obtained by the strong coupling expansion and the analytical continuation of these 1PI vertex functions in the momentum space.
In conclusion, based on the computations and discussions of the spectrum and 1PI gauge-coupling vertices of the model (1) in the segment $`𝒜`$ (8), we demonstrate the following scenario:
* fifteen non-degenerate massive Dirac fermions (30) in the high-energy $`p\pi _A0`$ vectorially coupling to the gauge field consistently with the $`SU_L(2)`$ chiral gauge symmetry;
* one massless chiral fermion $`\psi _L^i`$ (30) in the low-energy $`p0`$ chirally coupling to the $`SU_L(2)`$ gauge field;
* the neutral fermion sector of Dirac doublers (31) and the spectator $`\chi _R`$ entirely decoupling from the gauge field.
The $`SU_L(2)`$ chiral gauge symmetry and $`SU_L(2)U_R(1)`$ global chiral symmetry are exactly preserved at the lattice scale for not only perturbative but also non-perturbative gauge field configurations. It is no doubt that we need to have numerical simulations to show if such a scenario is indeed realized in the scaling region $`𝒜`$ of the action (1) proposed.
## V The vacuum functional and gauge anomaly
The Ward identities (70) and (71) play an extremely important rôle in a guarantee that the gauge perturbation theory in the scaling region $`𝒜`$ (8) is gauge symmetric. To all orders of the gauge coupling perturbation theory, gauge boson masses vanish and the gauge boson propagator is gauge-invariantly transverse. The gauge perturbation theory can be described by the normal renormalization prescription as that of normal vector-like gauge theories, as well as that of gauge theories with the soft spontaneous symmetry breaking like the Standard Model in the continuum. In fact, due to the manifest $`SU_L(2)`$ chiral gauge symmetry and corresponding Ward identities that are respected by the spectra (vector-spectrum for $`p\pi _A0`$ and chiral-spectrum for $`p0`$) in this scaling region, we should then apply the Rome approach (which is based on the conventional wisdom of quantum field theories) to the perturbation theory in the small gauge coupling. It is expected that the Rome approach would work in the same way but all gauge-variant (hard gauge-symmetry-breaking) counterterms are prohibited.
Provided the scenario of the gauge coupling vertex and spectrum given in above sections, we find that the gauge field not only chirally couples to the massless chiral fermion of the $`\psi _L^i`$ in the low-energy region, but also vectorially couples to the massive doublers of Dirac fermion $`\mathrm{\Psi }_c^i`$ in the high-energy regime. In this section, we perturbatively compute the vacuum functional in the small gauge coupling ($`g`$), and discuss the gauge anomaly and the renormalization prescription of the gauge-symmetric perturbation theory in this scenario.
We consider the following $`n`$-point 1PI functional:
$$\mathrm{\Gamma }_{\{\mu \}}^{(n)}=\frac{\delta ^{(n)}\mathrm{\Gamma }(A^{})}{\delta A_{\mu _1}^{}(x_1)\mathrm{}\delta A_{\mu _j}^{}(x_j)\mathrm{}\delta A_{\mu _n}^{}(x_n)},$$
(77)
where $`j=1\mathrm{}n,(n2)`$ and $`\mathrm{\Gamma }(A^{})`$ is the vacuum functional of the external gauge field $`A^{}`$. The perturbative computations of the 1PI vertex functions $`\mathrm{\Gamma }_{\{\mu \}}^{(n)}`$ can be straightforwardly performed by adopting the method presented in ref. for the lattice QCD. Dividing the integration of internal momenta (internal fermion loop) into 16 hypercubes where 16 modes live, we have 16 contributions to the truncated n-point 1PI functional. The region for the chiral fermion modes of continuum limit is defined as
$$\mathrm{\Omega }=[0,ϵ]^4,p<ϵ\frac{\pi }{2},p0,$$
(78)
where the $`ϵ`$ is the energy-threshold (33) where the three-fermion-state $`\mathrm{\Psi }_R(x)`$ dissolve to the three-fermion-cut $`𝒞[\mathrm{\Psi }_R(x)]`$.
As a first example, we deal with the vacuum polarization (in the following we refer $`p`$ to the external momentum of gauge bosons),
$$\mathrm{\Pi }_{\mu \nu }(p)=\underset{i=1}{\overset{16}{}}\mathrm{\Pi }_{\mu \nu }^i(p)=\mathrm{\Pi }_{\mu \nu }^c(p)+\mathrm{\Pi }_{\mu \nu }^d(p),\mathrm{\Pi }_{\mu \nu }^d(p)=\underset{i=2}{\overset{16}{}}\mathrm{\Pi }_{\mu \nu }^i(p),$$
(79)
where $`\mathrm{\Pi }_{\mu \nu }^d(p)`$ is doublers’ contributions and $`\mathrm{\Pi }_{\mu \nu }^c(p)`$ the contribution from the massless chiral mode in the region (78). As for the contributions $`\mathrm{\Pi }_{\mu \nu }^d(p)`$ from the 15 doublers $`(i=2,\mathrm{},16)`$, we make a Taylor expansion in terms of external physical momenta $`p=\stackrel{~}{p}`$ and the following equation is mutatis mutandis valid ,
$`\mathrm{\Pi }_{\mu \nu }^d(p)`$ $`=`$ $`\mathrm{\Pi }_{\mu \nu }^{}(0)+\mathrm{\Pi }_{\mu \nu }^{d(2)}(p)(\delta _{\mu \nu }p^2p_\mu p_\nu )`$ (80)
$`+`$ $`{\displaystyle \underset{i=2}{\overset{16}{}}}\left(1p_\rho |_{}_\rho {\displaystyle \frac{1}{2}}p_\rho p_\sigma |_{}_\rho _\sigma \right)\mathrm{\Pi }_{\mu \nu }^{con}(p,m_i),`$ (81)
where $`|_{}f(p)=f(0)`$ and $`m^i`$ are doubler masses. The first and second terms are specific for the lattice regularization. Since the 15 doublers are gauged as an $`SU_L(2)`$ vector-like gauge theory with propagator (30) and interacting vertex (69), the Ward identity (70) associated with this vectorlike gauge symmetry results in the vanishing of the first divergent term $`\mathrm{\Pi }_{\mu \nu }^{}(0)`$ and the gauge invariance of the second finite contact term in eq.(81). We recall that in Roma approach, this was achieved by adding gauge-variant counterterms at the lattice scale to enforce the valid of Ward identities. The third term in eq.(81) corresponds to the relativistic contribution of the 15 doublers. The $`\mathrm{\Pi }_{\mu \nu }^{con}(p,m_i)`$ is logarithmically divergent and evaluated in some continuum regularization for vectorlike gauge theories. For doubler masses $`m_i`$ of $`O(a^1)`$, the third term in eq.(81) is just finite and gauge-invariant contributions.
We turn to the contribution $`\mathrm{\Pi }_{\mu \nu }^c(p)`$ from the massless chiral mode that is in the first hypercube $`\mathrm{\Omega }=[ϵ,ϵ]^4`$ (78). We can use some continuum regularization to calculate this contribution,
$$\mathrm{\Pi }_{\mu \nu }^c(p)=\mathrm{\Pi }_{\mu \nu }^{c(2)}(p)(\delta _{\mu \nu }p^2p_\mu p_\nu ),$$
(82)
up to some finite local counterterms that are subtracted away in the normal renormalization prescription. The spectrum eq.(32) and gauge-coupling vertex eq.(44) with respect to the chiral mode are $`SU_L(2)`$ chiral-gauge symmetric. The Ward identity (71) associated with this chiral gauge symmetry makes eq.(82) to be gauge invariant and the gauge boson mass is zero to all orders of perturbative calculations. The $`ϵ`$-dependence (logarithmical divergence $`\mathrm{}nϵ`$) in eq.(82) has to be exactly canceled out from those contributions (81) from doublers, because the continuity of 1PI vertex functions in the momentum space. In summary, the total vacuum polarization $`\mathrm{\Pi }_{\mu \nu }(p)`$ contains two parts: (i) the vacuum polarization of the chiral mode $`\psi _L^i`$ in some continuum regularization; (ii) gauge invariant finite terms stemming from doublers’ contributions. The second part is the same as the perturbative lattice QCD, and can be subtracted away in the normal renormalization prescription.
The second example is the 1PI vertex functions $`\mathrm{\Gamma }_{\{\mu \}}^{(n)}(\{p\})(n4)`$,
$`\mathrm{\Gamma }_{\{\mu \}}^{(n)}(\{p\})`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{16}{}}}\mathrm{\Gamma }_{\{\mu \}}^{(n)i}(\{p\},m_i)n4,`$ (83)
$`\{p\}`$ $`=`$ $`p_1,p_2,\mathrm{}`$ (84)
$`\{\mu \}`$ $`=`$ $`\mu _1,\mu _2,\mathrm{},`$ (85)
where the internal momentum integral is analogously divided into the contributions from sixteen sub-regions of the Brillouin zone where sixteen modes live. Based on the gauge invariance and power counting, one concludes that up to some gauge invariant finite terms, the $`\mathrm{\Gamma }_{\{\mu \}}^{(n)}(\{p\})(n4)`$ (83) contain the 15 continuum expressions for 15 massive ($`m_i`$) Dirac doublers and one for the massless Weyl mode. The 15 doubler contributions vanish for $`m_iO(a^1)`$. The $`n`$-point 1PI vertex functions (83) end up with their continuum counterpart for the Weyl fermion and some gauge invariant finite terms. These finite gauge invariant terms come from doublers’ contributions are similar to those in the lattice QCD, and can be subtracted away in the normal renormalization prescription.
The most important contribution to the vacuum functional is the triangle graph $`\mathrm{\Gamma }_{\mu \nu \alpha }(p,q)`$. Again, dividing the integration of the internal momenta into 16 hypercubes, one obtains
$`\mathrm{\Gamma }_{\mu \nu \alpha }(p,q)`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{16}{}}}\mathrm{\Gamma }_{\mu \nu \alpha }^i(p,q)`$ (86)
$`\mathrm{\Gamma }_{\mu \nu \alpha }^i(p,q)`$ $`=`$ $`\mathrm{\Gamma }_{\mu \nu \alpha }^{i()}(0)+p_\rho \mathrm{\Gamma }_{\mu \nu \alpha ,\rho }^{i(1)}(0)+q_\rho \mathrm{\Gamma }_{\mu \nu \alpha ,\rho }^{i(1)}(0)`$ (87)
$`+`$ $`(1|_{}p_\rho |_{}_\rho q_\rho |_{}_\rho )\mathrm{\Gamma }_{\mu \nu \alpha }^{con}(p,q,m_i),`$ (88)
where $`\mathrm{\Gamma }_{\mu \nu \alpha }^{con}(p,q,m_i)`$ can be evaluated in some continuum regularizations. As for the 15 contributions of Dirac doublers ($`i=2\mathrm{}15`$), the first three terms in eq.(88) are zero for the vector-like Ward identity (70). The non-vanishing contributions are the same as the 15 copies of the $`SU(2)`$ vector-like gauge theory of massive Dirac fermions. These contributions are gauge-invariant and finite (as $`m_iO(a^1)`$), thus, disassociate from the gauge anomaly.
The non-trivial contribution from the chiral mode in the hypercube $`\mathrm{\Omega }=[ϵ,ϵ]^4`$ is given by
$`\mathrm{\Gamma }_{\mu \nu \alpha }^{i=1}(p,q)`$ $`=`$ $`{\displaystyle _\mathrm{\Omega }}{\displaystyle \frac{d^4k}{(2\pi )^4}}\text{tr}\left[S(k+{\displaystyle \frac{p}{2}})\mathrm{\Gamma }_\mu (k)S(k{\displaystyle \frac{p}{2}})\mathrm{\Gamma }_\nu (k{\displaystyle \frac{p+q}{2}})S(k{\displaystyle \frac{p}{2}}q)\mathrm{\Gamma }_\alpha (k{\displaystyle \frac{q}{2}})\right]`$ (90)
$`+(\nu \alpha ),`$
where the propagator $`S(p)`$ and vertex $`\mathrm{\Gamma }_\mu (k)`$ are given by eqs.(32,44). Other contributions containing anomalous vertices $`(\psi \overline{\psi }AA,\psi \overline{\psi }AAA)`$ vanish within the hypercube $`\mathrm{\Omega }=[ϵ,ϵ]^4`$. We evaluate eq.(90) by the Pauli Villars regularizationThe number of Pauli-Villars massive fermionic regulators is not infinite, which is different from ref.. in the continuum, which certainly violates chiral gauge symmetry and is linearly divergent at the scale of $`O(ϵ)`$. As a result, modulo possible finite local counterterms, we obtain the consistent gauge anomaly for the non-abelian chiral gauge theories as the continuum one:
$$\delta _g\mathrm{\Gamma }(A^{})=\frac{ig^2}{24\pi ^2}d^4xϵ^{\alpha \beta \mu \nu }\text{tr}\theta _a(x)\tau _a_\nu \left[A_\alpha (x)\left(_\beta A_\mu +\frac{ig}{2}A_\beta (x)A_\mu (x)\right)\right],$$
(91)
where the gauge field $`A_\mu =\frac{\tau ^a}{2}A_\mu ^a`$. The $`SU_L(2)`$ chiral gauge theory is anomaly-free for $`\text{tr}(\tau ^a,\{\tau ^b,\tau ^c\})=0`$, as if an appropriate anomaly-free fermion content in the group space. The gauge current,
$$J_\mu ^a=i\overline{\psi }_L\gamma ^\mu \frac{\tau ^a}{2}\psi _L=\frac{\delta \mathrm{\Gamma }(A)}{\delta A_\mu ^a(x)}^\mu J_\mu ^a=0$$
(92)
is covariantly conserved and gauge invariant. It must be pointed out and emphasized that in the hypercube $`\mathrm{\Omega }=[ϵ,ϵ]^4`$ we actually adopt a continuum (Pauli-Villars) regularization, which explicitly violates chiral gauge symmetries at the scale of $`O(ϵ)`$, which is much smaller than the lattice scale $`O(\pi /a)`$, for evaluating these anomalous terms like (91). The non-renormalization theorem of gauge anomaly guarantees that the resulted gauge anomaly is independent of any explicit breaking of chiral gauge symmetries at the scale of $`O(ϵ)`$ and free from hight-order corrections.
It seems surprising and impossible that we start from a gauge symmetric action (1) and fermion-field measure at the lattice scale, we end up with the correct form of the gauge anomaly (91). Because, one normally claim that the anomaly has to come from the explicit breaking of the chiral gauge symmetry in a regularized action (e.g., a Wilson term) at the lattice scale. This statement is indeed correct if regularized actions are exactly local and bilinear in fermion fields, since this is nothing but what the “no-go” theorem asserts. However, this is not correct in general. The general statement should be that the most essential and intrinsic raison d’être of producing the correct gauge anomaly in the lattice regularization is “decoupling doublers” rather than “explicitly breaking of chiral gauge symmetries at the lattice scale”.
In order to clarify and understant this general statement, let us first briefly review the relationship between doublers and chiral gauge anomalies in the lattice regularization. A most subtle property of the naive lattice regularization of chiral gauge theories is the appearance of 16 fermion zero modes. The gauge anomalies produced by these 16 fermion zero modes cancel each other. As a result, chiral gauge symmetries are exactly preserved at the lattice scale not only in the naive lattice action but also lattice fermion-field measure. This lattice fermion-field measure relates to the finite number of fermion-states (up to the lattice scale) of the vacuum of chiral gauge theories regularized by the lattice regularization for the finite lattice spacing $`a`$. While the anomalous currents of massless chiral fermions, that carry the fermion-states with definite chiral charges, flow into or out from the lattice regularized vacuum. Since the number of fermion-states of the lattice regularized vacuum is finite and these fermion-states are fully occupied, the total net chiral charges carried by anomalous currents of massless chiral fermions must zero in the lattice regularized vacuum. Otherwise the total number of fermion-states of the lattice regularized vacuum would not be finite. This is the reason for the occurrence of 16-modes in 4-dimension, the cancelation of chiral gauge anomalies produced by each chiral mode with definite axial charge $`Q_5`$, and the chiral gauge symmetry is perfectly preserved. This subtle property is in fact an important merit of the lattice regularization. Contrastively, the number of fermion-states of the vacuum in other continuum regularization schemes is not exactly finite. The infinite number of hight-energy fermion-states of the vacuum in the continuum is only exponentially suppressed by explicit chiral-symmetry-violating terms in those continuum regularization schemes. In order to maintain this merit of the lattice regularization, doublers should be decoupled in a gauge invariant way rather than an explicit gauge variant way.
In order to obtain the correct anomaly (91) in the lattice regularization, we obviously need to decouple extra doublers. If we adopt a local bilinear action to decouple doublers, we must either explicitly break chiral gauge symmetries at the lattice scale $`O(1/a)`$ or give up the exactly locality as required by the “no-go” theorem. However, on the other hand, we run into the dilemma that the gauge anomaly (91), obtained from an explicity breaking, is independent of any explicitly breaking parameters(scale) (e.g., the Wilson parameter $`r`$($`r/a`$)), which is consistent with the non-renormalization theorem of gauge anomalies. In this sense, the explicity breaking at the lattice scale leading to the gauge anomaly is just a superficial artifact in bilinear fermion actions.
If we give up the bilinearity of regularized actions in fermion fields and turn to our model and scenario with the exact chiral-gauge symmetry, the 15 doublers are decoupled as massive Dirac fermions that are vectorlike-gauge symmetric (70). Thus, they decouple from the gauge anomaly. Only the gauge anomaly associated with the normal (chiral) mode of the $`\psi _L^i`$ is left and is the same as the continuum one, provided the right-handed three-fermion state $`\mathrm{\Psi }_R^i`$ dissolves to the three-fermion-cut $`𝒞[\mathrm{\Psi }_R^i]`$ in the low-energy scale $`ϵ`$ (33). It would be otherwise that the massless right-handed three-fermion state $`\mathrm{\Psi }_R^i`$ gives rise to a gauge anomaly exactly eliminating the gauge anomaly (91) associated to the massless left-handed fermion $`\psi _L^i(x)`$. We still need to understand how the continuous states (virtual states) of the three-fermion cutes at the scale $`ϵ`$ fill up the lattice regularized vacuum.
However, to be consistent with the manifest chiral gauge symmetry of the regularized theory (action (1)) and fermionic measure, the gauge anomaly (91) must be canceled within the fermion content of the theory. Otherwise, the vectorlike spectrum of fermion zero modes must appear, either doublers do not decouple or the right-handed three-fermion-state $`\mathrm{\Psi }_R^i`$ does not dissolve into its cut and becomes a massless right-handed particle in the low-energy. From this point of view, we see the anomaly-cancelation by the fermion content is a necessary conditionIt is obvious not a sufficient condition. for this scenario to work, in particular, for the Ward identities (70) and (71) to be valid.
Before ending this section, we would like in particularly to discuss the residual breakings $`R(a)`$ of chiral gauge symmetries after the gauge anomaly (91) is canceled $`\delta _g\mathrm{\Gamma }(A)=0`$,
$$\delta \mathrm{\Gamma }(A)=\delta _g\mathrm{\Gamma }(A)+R(a).$$
(93)
If the gauge anomaly is induced by explicit breakings of the chiral gauge symmetry at the lattice scale, there must be the residual breakings $`R(a)`$ of the gauge symmetry at the lattice scale. Normally, given an explicity breaking of the chiral gauge symmetry at the lattice scale, we can perturbatively compute the gauge anomaly (91) and residual breakings of gauge symmetry by a small and smooth background of the external classical gauge field. Here by a smooth and small background, we indicate that the correlations of the gauge field are much larger than the lattice spacing $`O(a)`$ and the fluctuations of the gauge field are much smaller than the lattice scale $`O(1/a)`$. In such a background of the external gauge field, except the non-local gauge anomaly that is eliminated, the residual breakings of gauge symmetry are local and high dimension irrelevant operators at the lattice scale. However, these residual breakings of gauge symmetry could turn out to be relevant for the large and no smooth fluctuations of the longitudinal gauge field at the order of the lattice scale, and attempted lattice chiral gauge theories are jeopardized by breaking chiral gauge symmetries.
While, in our scenario, no smooth ($`O(a)`$) and the non-perturbatively large fluctuations ($`O(1/a)`$) of the longitudinal gauge field at the order of the lattice scale are fully under controlled by the chiral gauge symmetric and vectorlike spectra of 15 non-degenerated massive Dirac fermions. In the background of the external gauge field (the transverse and longitudinal components) with small fluctuations ($`O(ϵ)`$) and smooth correlations ($`O(1/ϵ)`$) at the order of the scale $`ϵ`$ (33), we compute the vacuum functional and obtain the gauge anomaly. The residual breakings $`R(ϵ)`$ of gauge symmetry, which comes out as the companions of the genuine gauge anomaly, are due to the explicity chiral symmetry breaking introduced by a continuum regularization scheme at the scale $`ϵ`$ (33), rather than at the lattice scale. Since the non-local gauge anomaly is canceled, chiral gauge symmetry is exact in the scaling region $`𝒜`$ (8) and the vacuum functional is determined up to local counterterms, we can subtract these residual breakings $`R(ϵ)`$ of gauge symmetry away by adding appropriate local counterterms, in order to achieve an asymptotically chiral gauge field theory in the continuum limit. This is the same as the procedures in the normal renormalization prescription of quantum field theories in the continuum.
## VI The anomalous fermion-flavour singlet current
The non-conservation of fermion numbers is an important feature of the Standard Model. A successful regularization of chiral gauge theories should give this feature in the continuum limit. In the Eichten-Preskills model of multifermion couplings for the $`SU(5)`$ and $`SO(10)`$ theories, inspired by the origination of the axial anomaly in the lattice QCD, it was suggested that the anomalous global current should be originated from the explicit breaking of the global symmetry at the lattice scale. In the context of the standard model, an elegant four-fermion interacting vertex of explicitly violating fermion number (B-L) was introduced. It is expected that the correct anomalous fermion-flavour singlet currents and violating fermion numbers should consistently be obtained. Nevertheless, we need to do explicit calculations to obtain the anomalous fermion-flavour singlet current in the standard model and $`SU(5),SO(10)`$ unification theories.
In this section, within the scenario presented in previous sections, we show that this fermion-flavour singlet anomaly can be consistently obtained from the explicit chiral gauge and fermion-flavour symmetric action (1) and fermion-field measure at the lattice scale. We first present the explicit calculations to obtain the correct fermion-flavour singlet anomaly and then discuss the consistency and reason for achieving this anomaly in our scenario.
Our action (1) processes the $`U_L(1)`$ and $`U_R(1)`$ global chiral symmetries. At the lattice scale, the action is invariant under the following global transformations:
$$\psi _L^ie^{i\theta _L}\psi _L^i\chi _Re^{i\theta _R}\chi _R,$$
(94)
where $`\theta _{L,R}`$ are the $`U_{L,R}(1)`$ chiral phases. These global symmetries lead to the conservation of the singlet chiral fermion currents,
$`_\mu j_L^\mu (x)`$ $`=`$ $`0,j_L^\mu =i\overline{\psi }_L^i\gamma ^\mu \psi _L^i`$ (95)
$`_\mu j_R^\mu (x)`$ $`=`$ $`0,j_R^\mu =i\overline{\chi }_R\gamma ^\mu \chi _R+O(a^2),`$ (96)
where $`j_{L,R}^\mu `$ are Noether currents and $`_\mu `$ is the derivative on the lattice. Eqs.(95,96) correspond to the conservation of the fermion numbers of $`\psi _L^i`$ and $`\chi _R`$. However, as we know, eq.(95) should be anomalous.
In order to see whether the conservations of the currents are violated when the external chiral gauge field is coupled to chiral fermions, we consider the source currents $`j_L^\mu (x)`$ and $`j_R^\mu (x)`$ defined as
$`j_L^\mu (x)`$ $`=`$ $`{\displaystyle \frac{\delta _L\mathrm{\Gamma }(A^{})}{\delta V_\mu ^L(x)}};\delta V_\mu ^L(x)=_\mu \theta _L(x);`$ (97)
$`j_R^\mu (x)`$ $`=`$ $`{\displaystyle \frac{\delta _R\mathrm{\Gamma }(A^{})}{\delta V_\mu ^R(x)}};\delta V_\mu ^R(x)=_\mu \theta _R(x),`$ (98)
where $`\mathrm{\Gamma }(A^{})`$ is the vacuum functional. Under the variations $`\delta _L`$ and $`\delta _R`$ of the $`U_{L,R}(1)`$-phases $`\theta _L(x)`$ and $`\theta _R(x)`$, the vacuum functional $`\mathrm{\Gamma }`$ is transformed up to $`O(\theta _L)`$ and $`O(\theta _R)`$,
$`\delta _L\mathrm{\Gamma }`$ $`=`$ $`{\displaystyle d^4x\delta V_\mu ^L(x)j_L^\mu (x)}={\displaystyle d^4x\theta _L(x)_\mu j_L^\mu (x)},`$ (99)
$`\delta _R\mathrm{\Gamma }`$ $`=`$ $`{\displaystyle d^4x\delta V_\mu ^R(x)j_R^\mu (x)}={\displaystyle d^4x\theta _R(x)_\mu j_R^\mu (x)},`$ (100)
where
$`\delta _R\mathrm{\Gamma }`$ $`=`$ $`\mathrm{\Gamma }(A_\mu ^{}+\delta V_\mu ^R(x))\mathrm{\Gamma }(A_\mu ^{}),`$ (101)
$`\delta _L\mathrm{\Gamma }`$ $`=`$ $`\mathrm{\Gamma }(A_\mu ^{}+\delta V_\mu ^L(x))\mathrm{\Gamma }(A_\mu ^{}).`$ (102)
In our action (1), the $`\chi _R`$ does not directly couple to the external gauge field and this decoupling strictly holds due to the Ward identity (37). Thus, the fact that the $`j_R^\mu (x)`$ defined formally in eq.(98) does not couple to the external gauge field $`A_\mu `$, i.e., in eq.(101)
$$\delta _R\mathrm{\Gamma }(A^{})=0,$$
(103)
leads to two direct consequences. One is the exact conservation of the source current $`j_R^\mu (x)`$
$$_\mu j_R^\mu (x)=0,$$
(104)
from (100). Another the exact gauge invariant source current $`j_R^\mu (x)`$,
$$\delta _gj_R^\mu (x)=\delta _g\frac{\delta _R\mathrm{\Gamma }(A^{})}{\delta V_\mu ^R(x)}=0,$$
(105)
where $`\delta _g`$ is a gauge variation. The facts that the Ward identity of the $`\chi _R`$-shift-symmetry is not violated by the gauge interaction in the action (1) and the gauge field is completely decoupled from the neutral sector (37), are extremely crucial to the conservation of the right-handed fermion numbers (104). The same reason leads to the exact conservation of the neutral and composite left-handed current $`J_L^{\mu ,n}=i\overline{\mathrm{\Psi }}_L^n\gamma _\mu \mathrm{\Psi }_L^n`$, i.e., $`_\mu J_L^{\mu ,n}=0`$.
As discussed in section 5, under a gauge variation $`\delta _g`$ , in general we have,
$$\delta _g\mathrm{\Gamma }_{\mu \nu \alpha }(p,q)0,\delta _g\mathrm{\Gamma }_{\{\mu \}}^{(n)}=0(n3),$$
(106)
where the vacuum functional $`\mathrm{\Gamma }(A^{})`$ is just the same as the continuum counterpart up to some gauge-invariant finite terms. In the anomaly-free $`SU_L(2)`$ case, i.e., $`\delta _g\mathrm{\Gamma }(A)=0`$, one may conclude that the source current $`j_L^\mu (x)`$ defined in eq.(97) is gauge invariant,
$$\delta _gj_L^\mu (x)=\delta _g\frac{\delta _L\mathrm{\Gamma }(A)}{\delta V_\mu ^L(x)}=\frac{\delta _L\delta _g\mathrm{\Gamma }(A)}{\delta V_\mu ^L(x)}=0,$$
(107)
and the variation (102) vanishes, i.e., $`\delta _L\mathrm{\Gamma }=0`$, leading to $`_\mu j_L^\mu (x)=0`$ by eq.(99) . However, these are not true. The order of the differentiations $`\delta _g`$ and $`\delta _L`$ can not be exchanged and $`\delta _L\mathrm{\Gamma }0`$ because of the mixing anomaly.
We know that in our action (1), the left-handed variation
$$\delta V_\mu ^L(x)=_\mu \theta _L(x),$$
(108)
can be considered as a commuting $`U_L(1)`$ factor in the $`SU_L(2)`$ chiral gauge group, i.e.,
$$\stackrel{~}{A}_\mu =A_\mu +V_\mu ^L,(A_\mu =\frac{\tau ^a}{2}A_\mu ^a).$$
(109)
Actually, this is an $`SU_L(2)U_L(1)`$ chiral gauge group and there is a mixing anomaly,
$`\delta _L\mathrm{\Gamma }`$ $`=`$ $`C_1{\displaystyle \frac{ig^2}{32\pi ^2}}{\displaystyle d^4x\theta _L\text{tr}\left(F_{\mu \nu }\stackrel{~}{F}^{\mu \nu }\right)},`$ (110)
$`\delta _g\mathrm{\Gamma }`$ $`=`$ $`C_2{\displaystyle \frac{ig}{16\pi ^2}}{\displaystyle d^4x\stackrel{~}{F}_1^{\mu \nu }\text{tr}\left(\theta _g_\mu A_\nu \right)},`$ (111)
where $`\theta _g=\theta _g^a\tau _a`$ is the $`SU_L(2)`$-transformation parameter, $`C_1,C_2`$ are arbitrary constants with $`(C_1+C_2=1)`$, and
$`F^{\mu \nu }`$ $`=`$ $`^\mu A^\nu ^\nu A^\mu ,`$ (112)
$`F_1^{\mu \nu }`$ $`=`$ $`^\mu V^{L\nu }^\nu V^{L\mu }.`$ (113)
The reason is that one of the Pauli matrices $`\tau ^a/2`$ in the triangle graph is replaced by the generator (identity) of the $`U_L(1)`$, i.e., the $`U_L(1)`$ global current, therefore the vanishing of the $`SU_L(2)`$ anomaly for $`\text{tr}(\tau ^a,\{\tau ^b,\tau ^c\})=0`$ is no longer true. Note that in eqs.(110,111), we only consider the triangle diagram $`(n=3)`$, since
$$\delta _L\mathrm{\Gamma }_{\{\mu \}}^{(n)}=0,\delta _g\mathrm{\Gamma }_{\{\mu \}}^{(n)}=0(n3)$$
(114)
for $`\mathrm{\Gamma }_{\{\mu \}}^{(n)}(n3)`$ being gauge-invariant, as we discussed in section 5.
The mixing anomaly (111) has arbitrariness $`C_1`$ and $`C_2`$, which arise because the triangle graphs with one insertion of the $`U_L(1)`$ global current determine the $`\mathrm{\Gamma }_{\mu \nu \alpha }(A^{})`$ up to local counterterms. As the Feynman diagrams determine the vacuum functional $`\mathrm{\Gamma }(A^{})`$ only up to an arbitrary choice of local counterterms, we are allowed to add local counterterms into the vacuum functional
$$\mathrm{\Gamma }^{}(A^{})=\mathrm{\Gamma }(A^{})+\mathrm{\Gamma }_{c.t.}(A^{}),$$
(115)
which is equivalent to the re-definition of the chiral fermion current defined by eq.(97),
$$j_L^{}_{}{}^{}\mu =j_L^{}_{}{}^{}\mu +j_{L,c.t.}^{}_{}{}^{}\mu .$$
(116)
Due to the fact that in our scenario the vacuum functional $`\mathrm{\Gamma }(A^{})`$ we obtained for the $`SU_L(2)`$ case is free from the non-local gauge anomaly and local gauge-symmetry-breaking terms, $`\delta _g\mathrm{\Gamma }`$ (111) must vanish and the arbitrariness in eq.(110,111) can be fixed,
$$C_1=1,C_2=0$$
(117)
by choosing an adequate local counterterm. As a result, the vacuum functional and the left-handed current are gauge invariant,
$$\delta _g\mathrm{\Gamma }^{}(A^{})=0,\delta _gj_L^{}_{}{}^{}\mu =0.$$
(118)
From eqs.(99,110), we obtain
$$\delta _L\mathrm{\Gamma }^{}=\frac{ig^2}{32\pi ^2}d^4x\theta _L\text{tr}\left(F_{\mu \nu }\stackrel{~}{F}^{\mu \nu }\right);_\mu j_L^{}_{}{}^{}\mu =\frac{ig^2}{32\pi ^2}\text{tr}\left(F_{\mu \nu }\stackrel{~}{F}^{\mu \nu }\right).$$
(119)
This is just the desired result, showing the left-handed fermion number is violated by the $`SU(2)`$ instanton effect, which attributes to the topological configuration of the gauge field carrying fermion numbers. By the definition $`j_5^\mu =j_L^{}_{}{}^{}\mu j_R^\mu `$ and $`_\mu j_R^\mu =0`$ (104), we obtain,
$$_\mu j_5^\mu =\frac{ig^2}{32\pi ^2}\text{tr}\left(F_{\mu \nu }\stackrel{~}{F}^{\mu \nu }\right).$$
(120)
We emphasize crucial points for achieving the correct form of fermion-flavour singlet anomaly (119) up to gauge invariant local counterterms in this scenario. The first is the right-handed composite three-fermion-state $`\mathrm{\Psi }_R^i(x)`$ dissolving into its three-fermion-cut $`𝒞\mathrm{\Psi }_R^i(x)`$, in another word, no massless right-handed composite three-fermion-state $`\mathrm{\Psi }_R^i(x)`$ exists in the low energy so that the correct from of the gauge anomaly is obtained (91). The second is the decoupling of the right-handed massless fermion $`\chi _R`$ from the gauge field, which leads to $`_\mu j_R^\mu =0`$. The third is exact chiral gauge symmetry, i.e., no chiral gauge-symmetry-breakings, and the vanishing of the gauge anomaly so as to have the choice (117) up to local gauge invariant counterterms. The fourth is the correct gauge anomaly obtained in the continuum regularization scheme which explicitly breaks the chiral symmetry at the scale $`O(ϵ)`$ being much smaller than the lattice scale. In fact, instead of using the formula of the mixing anomaly (110, 111), we can directly compute the triangle diagrams (90) with the left-handed fermion-flavour singlet current insertions in the continuum Pauli-Villars regularization scheme for the hypercube $`\mathrm{\Omega }=[ϵ,ϵ]^4`$, to achieve the correct fermion-flavour singlet anomaly (119). However, the mixing anomaly let us have a clear connection between the gauge anomaly-free, chiral gauge symmetry and the fermion-flavour singlet anomaly.
We turn to the discussions why we obtain the correct fermion-flavour $`U_L(1)`$-anomaly from the $`U_L(1)`$ symmetric action (1) at the lattice scale. This question arises because of our knowledge of the lattice QCD where the axial current anomaly,
$$_\mu j_5^\mu =\frac{ig^2}{16\pi ^2}\text{tr}\left(F_{\mu \nu }\stackrel{~}{F}^{\mu \nu }\right)$$
(121)
is due to the flavour $`SU_L(3)SU_R(3)`$ asymmetric Wilson term at the lattice scale. Certainly, since we start from a global symmetric action and no spontaneous symmetry breaking occurs, the relevant operators in the scaling region must have the same symmetries as they are at the lattice scale. No way to generate chirally asymmetric relevant operators by its own dynamics to produce the fermion-flavour singlet anomaly. All what we have done for producing the fermion-flavour singlet anomaly is introducing an explicit and soft chiral symmetry breaking at the scale $`O(ϵ)`$ consistently with the dissolving scale of the right-handed three-fermion state $`\mathrm{\Psi }_R^i`$. The reason that we obtain the correct fermion-flavour singlet anomaly is very analogous to that we obtain the correct gauge anomaly discussed in the end of section 5.
A priori, we have no any dynamical reasons to expect that the non-conservation of fermion number must be due to the explicit chiral symmetry breaking of the global $`U_L(1)`$ symmetry at the lattice scale. Further, the fact that the resulted fermion-flavour singlet current anomaly is independent of the explicit-symmetry breaking parameters implies that the explicit-symmetry-breaking is not necessary at the lattice scale, and can be at other scales, for instance $`ϵ`$ (33), much smaller than the lattice scale. This is also consistent with the non-renormalization theorem of anomalies. However, in order to achieve the fermion-flavour singlet anomaly by the explicit chiral symmetry breaking at a much softer scale than the lattice scale, we have to pay price as required by the “no-go” theorem of Nielsen-Ninomiya, either relaxing the exact locality of the theory, which is fundamental property of quantum field theories, or giving up the bilinear construction of fermionic action and running to multifermion interactions at the lattice scale. The appearance of two scales in our scenario, the lattice scale and the dissolving scale $`ϵ`$ (33), is reminiscence of the two cutoff approach to chiral gauge theories on the lattice.
In fact, the fermion-flavour singlet anomaly (119) must disappear as the external gauge field is turned off. The conservations of the fermion-flavour singlet currents (104,119) must be related to the conservations of Noether currents (95,96) coming from the explicit global $`U_L(1)`$ and $`U_R(1)`$ symmetries of the action (1) in the absence of gauge field at the lattice scale(tree-level). Otherwise, we would run to the dilemma that when the external gauge field (or the topological configuration of the gauge field) is turned off and the fermion-flavour singlet current (119) becomes conserved, however, on the other hand, the corresponding Noether current (95) is not conserved due to the explicit $`U_L(1)`$ chiral-symmetry-breaking introduced at the lattice scale (tree level).
In our scenario, we still need to have more intuitive understanding of how this anomalous fermion number, contributed by the presence of the topological configuration of the gauge field, flows into and emerges out from the regularized vacuum and how the continuous states (virtual states) of the three-fermion cutes at the scale $`ϵ`$ relate to the anomalous fermion number flows.
The multifermion couplings we considered in this paper is particular one in the context of the standard model. Taking the doublet $`\psi _L^i`$ and spectator $`\chi _R`$ as the $`SU_L(2)`$ doublet and right-handed neutrino in the leptonic sector of the standard model, we have the gauge symmetric multifermion couplings in the action (1). However, in the fermion content of the standard model, we certainly have the possibilities of fermion-number violation, but chiral gauge symmetric multifermion couplings at the lattice scale. The nice examples are given in refs.. If a scaling region with the desired gauge symmetric spectrum in the low-energy can be achieved, these fermion-number violating multifermion couplings would turn out to be relevant operators in such a scaling region to give the $`BL`$ fermion number violation. It could be the Nature choice.
## VII Conclusions and discussions
In this paper, we present the analysis and scenario of the chiral fermion spectrum and 1PI interacting vertices between fermions and gauge field in the low-energy scaling region (8) of the model (1) proposed in ref. for chiral gauged fermions on the lattice. In addition, we show how the gauge anomaly and fermion-flavour singlet anomaly are correctly produced in such a scenario. This shows that an asymptotic chiral gauge theory in the continuum limit can be realized in the low-energy scaling region (8) of our model (1). We conclude that our model and scenario provide a plausible solution to the long-standing problem of chiral gauged fermion on a lattice. It is very inviting that numerical computations and other techniques to verify our scenario, in particular, the phenomena of no hard spontaneous symmetry breaking at the lattice scale and the intermediate energy-threshold $`ϵ`$ (33) where the three-fermion-state $`\mathrm{\Psi }_R^i(x)`$ dissolve into its corresponding cut $`𝒞[\mathrm{\Psi }_R^i(x)]`$. It is also very necessary to analyze an analogous model in 1+1 dimensions so as to make our conclusion be more convincible. In addition, the consistency of our model in the scaling region regarding the $`SU_L(2)`$ global Witten anomaly is still open question. These are subject to the future work.
In this section, we wish to make a very general and brief discussion on any possible relationships between the multifermion coupling and bilinear fermion coupling approaches for anomaly-free chiral gauge theories on the lattice. In both bilinear fermion and multifermion coupling models, extra fermionic species must be decoupled and right-handed and left-handed fermionic species must couple in some ways to have anomalies. The couplings between right- and left-handed fermionic species appear either in the action or in the fermionic measure.
The Wilson fermion is exact local (in the range of the lattice spacing) and right- and left-handed fermionic species couple at the lattice scale. Doublers are very massive and decoupled. As required by the “no-go” theorem, the residual breakings of the gauge symmetry are at the lattice scale. These residual breakings of the chiral gauge symmetry can be eliminated by adding and fine-tuning appropriate counterterms so as to enforce the Ward identities associated to exact chiral gauge symmetries at the continuum limit.
Contrasting with the Wilson fermion, the regularization of the fermion sector adopted by the “overlap” and Lüscher approaches use the Ginsparg-Wilson equation that was obtained from the renormalization group equation. Owing to the Dirac operator satisfying the Ginsparg-Wilson equation, the residual breakings of the gauge symmetry are reduced and supposed to be eliminated by either average over gauge configurations or adding local counterterms in order to preserve exact chiral gauge symmetries. In the Lüscher approach for the abelian gauge theory, all residual breakings of the gauge symmetry including the gauge anomaly can be rewritten as a total divergence for its topological nature, thus they are eliminated by redefining the gauge current<sup>§</sup><sup>§</sup>§Private conmmunication with Lüscher. for the finite lattice spacing and without fine-tuning. In order to completely decouple extra fermion species, as required by the “no-go” theorem, the approaches relax the exact locality to the locality whose range extends to a few lattice spacings with an exponential tail.
While in the models of multifermion couplings, the couplings of right-handed and left-handed fermions can be made exactly local and chiral gauge symmetric at the lattice scale. However, as discussed at the end of section 2, we have to find a peculiar multifermion coupling and a scaling region desired for the low energy. The hard spontaneous symmetry breaking is absolutely not tolerated so that residual breakings of gauge symmetries are not at the lattice scale. The strong coupling at the high energy is needed and three-fermion cut must be realized at the low-energy so as to decouple extra fermion species with “wrong” chirality.
Taking our action (1) as an example, we can formally integrate away the spectator field $`\chi _R`$ and obtain the effective Dirac action bilinearly in the fermion field $`\psi _L^i`$. Such an effective Dirac action is obviously not exactly local. It is worthwhile to examine whether such an effective Dirac operator could be the solution to the Ginsparg-Wilson equation in the sense of the renormalization group invariance. Most importantly, we need to show in the scaling region $`𝒜`$ (8) for the low-energy, whether the relevant spectra and operators induced from the multifermion couplings(high dimension operators) at the lattice scale are in the same universal class with the solutions to the Ginsparg-Wilson equation, in the view of the renormalization group invariance.
In fact, the recent successful progress based on the Ginsparg-Wilson equation strongly implies the existence of the scaling region 8) for exactly chiral-gauge symmetric theories in the low-energy, obtained from our model (1). This has been generally believed to be impossible. The studies of appropriate multifermion couplings at the lattice scale and desired scaling region are highly deserved, since they could be Nature’s choice for chiral gauge theories, e.g., the Standard Model, at the high-energy.
## VIII Acknowledgements
I thank to the organizers of the winter workshop at Benasaque Spain (January 2000), where the manuscript of paper was completed. |
warning/0002/math0002244.html | ar5iv | text | # Complementary Algorithms for Tableaux
## Introduction
Schützenberger’s theory of jeu de taquin and the insertion procedure of Schensted have found their way into the standard toolkit of the combinatorial representation theorist. These algorithms were originally developed for tableaux of partition shape, but more recent work uses them to define operations on pairs of skew tableaux. A striking duality between Schensted insertion and Schützenberger’s jeu de taquin was noted by Stembridge in his theory of rational tableaux.
We introduce the involutive tableau operation of complementation to formalise this duality, showing that internal row insertion and the jeu de taquin give complementary operations on pairs of tableaux. This is readily seen using the growth diagram formulation for operations on tableaux introduced by Fomin and rediscovered several times, most notably by van Leeuwen . Column insertion is another tableau algorithm that can be extended to define the operation of internal column insertion between pairs of tableaux. The operation hopscotch which is complementary to internal column insertion is new, although it extends both Stroomer’s algorithm of column-sliding and Tesler’s rightward- and leftward-shift games.
These four operations of jeu de taquin, internal row insertion, internal column insertion, and hopscotch all preserve Knuth-equivalence. A consequence is that hopscotch gives rise to new algorithms to rectify tableaux. We also show that these four operations have the same dual-equivalence theory.
This paper is organised as follows. Section 1 recalls basic definitions and terminology concerning tableaux, the jeu de taquin, and Knuth equivalence. In Section 2 we introduce complementation and determine its behaviour with respect to Knuth equivalence and dual equivalence. Section 3 is devoted to operations on pairs of tableaux given by growth diagrams constructed from local rules, and introduces the notion of complementary operations on pairs of tableaux. In Section 4 we show that the two operations of jeu de taquin and internal row insertion are complementary. In Section 5 we define internal column insertion and its complementary operation hopscotch, and in Section 6, we investigate some features of hopscotch. Finally, in Section 7, we supplement the jeu de taquin, describing four additional algorithms to rectify skew tableaux. The first three are not well-known, while the fourth is new.
Schematically, the four operations of jeu de taquin, internal row insertion, internal column insertion, and hopscotch are linked through complementation and conjugation as follows:
## 1. Tableaux
The tableaux operations of internal row insertion, internal column insertion, and hopscotch require us to extend the usual notion of tableaux. While these operations are defined on a class of more general tableaux, we show in Section 7 how they provide new algorithms for ordinary tableaux. Rather than give a common generalisation, we instead give a definition that will suffice until Section 5, where we will be precise about necessary further extensions to our notion of tableau.
A partition $`\alpha `$ is a weakly decreasing sequence of positive integers $`\alpha _m\mathrm{}\alpha _n`$ indexed by an interval $`m,m+1,\mathrm{},n`$ of integers. The shape of a partition $`\alpha `$ is a left justified array of boxes whose $`i`$th row contains $`\alpha _i`$ boxes. We make no distinction between a partition and its shape. Thus we identify partitions that differ only in their number of trailing zeroes. Partitions are partially ordered $``$ by componentwise comparison of sequences. If $`\alpha \beta `$, then we may form the (skew) shape $`\beta /\alpha `$, consisting of those boxes which are in $`\beta `$ but not in $`\alpha `$. The shape $`\beta /\alpha `$ has inner border $`\alpha `$ and outer border $`\beta `$. We call $`\beta /\alpha `$ a horizontal strip if no column of $`\beta /\alpha `$ contains two or more boxes. That is, $`\beta _m\alpha _m\beta _{m+1}\mathrm{}\beta _n\alpha _n`$. When $`\alpha \beta \delta `$, we say that the shape $`\delta /\beta `$ extends the shape $`\beta /\alpha `$.
A tableau $`T`$ with shape $`\beta /\alpha `$ and entries from the alphabet $`[k]:=\{1,\mathrm{},k\}`$ is a chain $`\alpha =\beta ^0\beta ^1\mathrm{}\beta ^k=\beta `$, where the successive skew shapes $`\beta ^i/\beta ^{i1}`$ are horizontal strips. The content of $`T`$ is the sequence whose $`i`$th component is the number of boxes in the horizontal strip $`\beta ^i/\beta ^{i1}`$. Filling the boxes of $`\beta ^i/\beta ^{i1}`$ with the integer $`i`$ shows this is equivalent to the usual definition of a column-strict tableau. We represent tableaux both as chains and as fillings of a shape with a totally ordered alphabet (in practice, the positive integers or Roman alphabet). For example, the chain and the filling of the shape $`(4,2,1)/(1)`$ below both represent the same tableau.
A tableau is standard if each $`\beta ^i/\beta ^{i1}`$ consists of a single box, so that the chain is saturated. In this case, we say $`\beta ^{i1}\beta ^i`$ is a cover. Many (but not all) of our tableaux algorithms may be computed via the standard renumbering of a tableau $`T`$. The standard renumbering of $`T`$ is the refinement of the chain representing $`T`$ where each horizontal strip is filled in ‘left-to-right’. For example, here is the standard renumbering of the tableau above.
Let $`\alpha \beta \delta `$ be partitions. If $`\alpha \beta `$ is a cover, then $`\beta /\alpha `$ is a single box $`b`$, which we call an inner corner of $`\delta /\beta `$. If instead $`\beta \delta `$, then $`\delta /\beta `$ is a single box $`b`$, called an outer corner of $`\beta /\alpha `$.
A jeu de taquin slide is a specific reversible procedure that, given a tableau $`T`$ and an inner corner $`b`$ of the shape of $`T`$, moves $`b`$ through $`T`$, producing a new tableau together with an outer corner. It proceeds by successively interchanging the empty box $`b`$ with one of its neighbours to the right or below in such a way that at every step the figure has increasing columns and weakly increasing rows. When the box has neighbours both right and below it moves as indicated:
(1.1)
$$\text{}\{\begin{array}{ccc}\text{}\hfill & & a<b\hfill \\ \text{}\hfill & & ba\text{}\hfill \end{array}$$
When the box has only a single such neighbour, it interchanges with that neighbour, and the slide concludes when the box has no neighbours. The reverse of this procedure is also called a jeu de taquin slide.
Schützenberger’s jeu de taquin is the procedure that, given a tableau $`P`$, applies jeu de taquin slides beginning at inner corners of $`P`$, and concludes when there are no more such corners. Schützenberger proves that the result, which we call the rectification of $`P`$, is independent of choices of inner corners.
A tableau $`U`$ extends another tableau $`T`$ if the shape of $`U`$ extends the shape of $`T`$. If $`T`$ is a tableau extended by $`U`$, then the standard renumbering of $`T`$ gives a set of instructions for applying jeu de taquin slides to $`U`$: The last cover of $`T`$ is an inner corner of $`U`$, and after applying the slide beginning with that corner, the next slide begins at the new inner corner given by the next to last cover of $`T`$, and so on. Write $`Q`$ for the resulting tableau. One could consider the sequence of outer corners obtained from this procedure as a tableau $`P`$ or regard $`U`$ as a set of instructions for slides on $`T`$; these both produce the same tableau (up to standard renumbering). If we set $`𝒥(T,U)=(P,Q)`$, then $`𝒥`$ defines an involution, which we call the jeu de taquin, on pairs of tableaux where one extends the other.
Two fundamental equivalence relations among tableaux are Knuth equivalence and dual equivalence. We call two tableaux $`T`$ and $`U`$ Knuth-equivalent if one can be obtained from the other by a sequence of jeu de taquin slides. (This is equivalent to their reading words being Knuth-equivalent in the standard sense .) Two tableaux $`T`$ and $`U`$ with the same shape are dual equivalent if applying the same sequence of jeu de taquin slides to $`T`$ and to $`U`$ always gives tableaux of the same shape. These are related by a result of Haiman : The intersection of any Knuth-equivalence class and any dual-equivalence class is either empty or consists of a unique tableau.
## 2. Complementation
Complementation originated in a combinatorial procedure of Stanley to model the evaluation $`(x_1\mathrm{}x_k)^{\lambda _1}s_\lambda (1/x_1,\mathrm{},1/x_k)`$, where $`s_\lambda (x_1,\mathrm{},x_k)`$ is the Schur polynomial . In combinatorial representation theory, complementation provides a combinatorial model for the procedure of dividing by the determinantal representation of $`GL_k`$. In this context, it was used by Stembridge and Stroomer to develop combinatorial algorithms for studying rational representations of $`GL_k`$. (The classical Robinson-Schensted-Knuth correspondence is used to study polynomial representations of $`GL_k`$.) Reiner and Shimozono studied the commutation of a version of complementation (which they called “Boxcomp”) with other tableaux operations. Given a tableau $`T`$ filled with integers from $`[k]`$, these authors formed a tableau whose columns were obtained from the columns of $`T`$ by complementing each in the set $`[k]`$ and rotating the resulting tableau by $`180^{}`$. Rather than rotate the result, we instead choose to renumber the resulting tableau.
More precisely, let $`T`$ be a tableau with shape $`\beta /\alpha `$ and entries from the alphabet $`[k]`$, and fix $`l\beta _m`$, the initial and largest part of $`\beta `$. Form a new tableau $`T^C`$ with $`l`$ columns as follows. For each $`1jl`$, let $`A`$ be the set-theoretic complement of the entries of column $`j`$ of $`T`$ in the set $`[k]`$, considered in the dual order: $`k<k1<\mathrm{}<1`$. Since it is inconvenient (and possibly ambiguous) to work with tableau whose alphabet has more than one order, we make the replacement $`jk+1j`$, which indicates the same chain of shapes in the usual order on $`[k]`$. The figure below illustrates this two-step process, complementing the tableau $`T`$ on the left with $`k=4`$ and $`l=5`$. In the middle figure, we write the complement of a column of $`T`$ below that column, but in reverse order, and the rightmost figure is the complement $`T^C`$, where we have applied the substitution $`jk+1j`$ to the complemented columns and omitted writing $`T`$.
(2.1)
$$\begin{array}{c}\text{}\end{array}$$
This figure also illustrates the fact that if some columns of $`T`$ are empty (for instance, if $`l>\beta _m`$), then the corresponding columns of $`T^C`$ will consist of the full set $`[k]`$. If we identify tableaux which differ by a vertical shift, we may also write the columns of $`T^C`$ above the corresponding columns of $`T`$. This proves that complementation is involutive, under this identification.
###### Theorem 2.1.
Let $`T`$ be a tableau with shape $`\beta /\alpha `$ and entries from $`[k]`$ and suppose $`l\beta _m`$, the initial part of $`\beta `$. Then $`T^{CC}=T`$.
Complementation depends upon both $`k`$ and $`l`$. Our notation, $`T^C`$, intentionally disregards this dependence. We adopt the convention that we use the same integers $`k`$ and $`l`$ when two or more tableaux are to be complemented. Using an extension of notation, we observe that the shape of $`T^C`$ is $`(l^k,\beta )/\alpha `$. (Here $`l^k`$ denotes the rectangular shape consisting of $`k`$ parts of length $`l`$; context will distinguish this from the superscripts on Greek letters indicating chains of shapes.)
We express complementation in terms of chains in Young’s lattice. Let $`T`$ be a tableau, written as a chain $`\alpha =\beta ^0\beta ^1\mathrm{}\beta ^k=\beta `$ in Young’s lattice, where $`\beta ^i/\beta ^{i1}`$ is the horizontal strip of $`i`$’s in $`T`$. Then $`T^C`$ is the chain $`\beta (l,\beta ^{k1})(l^2,\beta ^{k2})\mathrm{}(l^i,\beta ^{ki})\mathrm{}(l^k,\beta ^0)`$. It is an exercise that this agrees with the definition given above. In particular, it can immediately be seen that the shape is as claimed. Each row of partitions in Figure 1 is one of the two complementary tableaux in (2.1).
We write the second row in reverse order to illustrate the key feature of complementation—that the $`i`$th horizontal strip in the complement contains boxes in exactly the columns complementary to the corresponding ($`k+1i`$)th horizontal strip in the original tableau.
Complementation preserves both Knuth equivalence and dual equivalence.
###### Theorem 2.2.
Suppose $`T`$ and $`U`$ are tableaux with at most $`l`$ columns and entries from the alphabet $`[k]`$. Then
* $`U`$ is Knuth-equivalent to $`T`$ if and only if $`U^C`$ is Knuth-equivalent to $`T^C`$.
* $`U`$ is dual-equivalent to $`T`$ if and only if $`U^C`$ is dual-equivalent to $`T^C`$.
We prove Theorem 2.2 in Section 4.
Remark. Theorem 2.2 shows that complementation commutes with the involution reversal $`TT^e`$ of . To see this, given a tableau $`T`$, let $`T^{}`$ be the tableau obtained by rotating $`T`$ $`180^{}`$ and replacing each entry $`j`$ with $`k+1j`$. Then the reversal $`T^e`$ of $`T`$ is the unique tableau dual-equivalent to $`T`$ (and hence with the same shape as $`T`$) and Knuth-equivalent to $`T^{}`$. For tableaux of partition shape, reversal coincides with Schützenberger’s evacuation procedure (called “promotion” there), but the two procedures differ for general skew tableaux. This extends the result of Reiner and Shimozono (, Theorem 2) that complementation commutes with evacuation.
## 3. Growth diagrams and local rules
Many properties of tableaux algorithms such as symmetry become clear when the algorithms are formulated in terms of growth diagrams governed by local rules. Fomin introduced this approach to the Robinson-Schensted correspondence, it was rediscovered by van Leeuwen , and Roby developed it further. We study tableaux algorithms related via complementation of their growth diagrams.
### 3.1. Growth diagrams
A growth diagram is a rectangular array of partitions where every row and every column is a tableau, with the additional restriction that all tableaux formed by the rows have the same content, as do all tableaux formed by the columns. Specifically, the sequence $`\beta ^0,\beta ^1,\mathrm{},\beta ^k`$ of partitions (read left-to-right) in each row forms a chain $`\beta ^0\beta ^1\mathrm{}\beta ^k`$, with each $`\beta ^i/\beta ^{i1}`$ a horizontal strip, and the number of boxes in $`\beta ^i/\beta ^{i1}`$ does not depend upon which row this chain came from. (We require the same to hold for the sequence of partitions read top-to-bottom in each column.) For example, here is a growth diagram where the horizontal tableaux have content $`(1,2,2)`$ and the vertical tableaux have content $`(2,1,1)`$.
(3.1)
$$\begin{array}{cccc}1\hfill & 2\hfill & 31\hfill & 33\hfill \\ 21\hfill & 22\hfill & 321\hfill & 332\hfill \\ 211\hfill & 221\hfill & 3211\hfill & 3321\hfill \\ 221\hfill & 222\hfill & 3221\hfill & 3322\hfill \end{array}$$
Traditionally the tableaux in a growth diagram are standard. We relax this in order to define complementation of a growth diagram. Given an integer $`l\delta _1`$, where $`\delta `$ is the lower right partition in the growth diagram, we complement the tableaux represented by each row to obtain new tableaux, also written as chains of shapes. These combine together to give a new growth diagram. We then complement the columns of these diagrams to obtain two more growth diagrams. Figure 2 shows the four growth diagrams we obtain from (3.1) by this process. Here, $`k`$ is 3 for both the vertical and horizontal tableaux, and the complementation parameter is $`l=4`$.
The lower right growth diagram may also be obtained from the lower left diagram by complementing rows. Indeed, if we number the rows of the original growth diagram $`1,2,\mathrm{},k`$ and the columns $`1,2,\mathrm{},m`$ and if $`\beta `$ is the partition in position $`(i,j)`$ of the original diagram, then $`(l^{(ki)+(mj)},\beta )`$ is the partition in position $`(ki+1,mj+1)`$ of the lower right growth diagram. By Theorem 2.1, further complementation of the rows or columns yields no new growth diagrams. (Here, we identify tableaux which differ by a vertical shift.)
### 3.2. Local rules
A local rule $``$ is a rule for completing the missing corner of a $`2\times 2`$ growth diagram given 3 partitions that are related by horizontal strips. Suppose we have partitions $`\alpha \beta \delta `$ with $`\beta /\alpha `$ and $`\delta /\gamma `$ horizontal strips. A switching local rule $``$ is a rule for completing the missing lower left corner $`\gamma `$ of a partial growth diagram
$$\begin{array}{cc}\alpha & \beta \\ & \delta \end{array}.$$
Write $`\gamma =(\alpha ,\beta ,\delta )`$. A switching local rule also gives a rule for completing a missing upper right corner of a partial growth diagram, by the obvious symmetry of the two cases. In order to make $``$ reversible, we insist that the rule be symmetrical: $`\gamma =(\alpha ,\beta ,\delta )\beta =(\alpha ,\gamma ,\delta )`$.
Suppose we have three partitions $`\alpha ,\beta `$, and $`\gamma `$ with $`\gamma /\alpha `$ and $`\beta /\alpha `$ horizontal strips, so that
$$\begin{array}{cc}\alpha & \beta \\ \gamma & \end{array}$$
is a (partial) growth diagram. An insertion local rule $``$ is a rule that associates a fourth partition $`\delta `$ to such a partial growth diagram such that
$$\begin{array}{cc}\alpha & \beta \\ \gamma & \delta \end{array}$$
is a growth diagram. We require that $``$ be invariant under vertical shifts of the horizontal strips $`\beta /\alpha `$ and $`\gamma /\alpha `$. By this we mean that if $`l`$ is at least as long as the initial part of $`\delta `$, then the rule $``$ completes the partial growth diagram
$$\begin{array}{cc}(l,\alpha )& (l,\beta )\\ (l,\gamma )& \end{array}$$
with the partition $`(l,\delta )`$. We write $`\delta =(\alpha ;\beta ,\gamma )`$. We further require $``$ to be symmetric in $`\beta `$ and $`\gamma `$.
We would like to define a reverse map $`(\beta ,\gamma ;\delta )`$ for all $`\beta `$, $`\gamma `$, and $`\delta `$ with $`\delta /\beta `$ and $`\delta /\gamma `$ horizontal strips with the property that $`(\alpha ;\beta ,\gamma )=\delta `$ if and only if $`\alpha =(\beta ,\gamma ;\delta )`$. Unfortunately, this is impossible as the following example shows. Let $`\beta =\gamma =0`$ be the empty partition and $`\delta =1`$ be the partition with a single part of size 1. Then there simply does not exist a partition $`\alpha `$ with $`\beta /\alpha `$ and $`\gamma /\alpha `$ a single box.
To circumvent this problem, we call an insertion local rule $``$ reversible if given partitions $`\beta `$, $`\gamma `$, and $`\delta `$ with $`\delta /\beta `$ and $`\delta /\gamma `$ horizontal strips and an integer $`l`$ at least as large as the initial part of $`\delta `$, then there exists a unique partition $`\alpha `$ such that $`(\alpha ;(l,\beta ),(l,\gamma ))=(l,\delta )`$.
A reversible local rule $``$ and an integer $`l`$ together define a local rule for computing the missing upper left corner of a growth diagram
$$\begin{array}{cc}& \beta \\ \gamma & \delta \end{array}$$
when $`l\delta _m`$, the initial part of $`\delta `$. First prepend a single, $`(m1)`$th, part $`l`$ to each of $`\beta ,\gamma `$, and $`\delta `$ and then set $`(\beta ,\gamma ;\delta ):=\alpha `$, where $`\alpha `$ is the unique partition such that $`(\alpha ;(l,\beta ),(l,\gamma ))=(l,\delta )`$. Prepending the part of size $`l`$ does not alter the horizontal strips $`\delta /\beta `$ and $`\delta /\gamma `$, as it only shifts them vertically.
As with complementation, we suppress this parameter $`l`$ in our notation. However, we insist that it coincides with the complementation parameter when combining an insertion local rule with complementation.
### 3.3. Tableaux algorithms
An insertion local rule $``$ determines a bijection on pairs of tableaux which share a common border as follows. Given a pair $`(P,Q)`$ sharing an inner border, write the tableau $`P`$ across the first row of an array and the tableau $`Q`$ down the first column. Then use $``$ to fill in the array and obtain a growth diagram. If $`U`$ is the tableau of the last column in this diagram and $`T`$ the tableau of the last row in this diagram, then $`T`$ and $`U`$ share the same outer border and the pair $`(T,U)`$ is determined from the pair $`(P,Q)`$ by the local rule $``$. When $``$ is reversible and $`P`$ and $`Q`$ share an outer border occupying columns (at most) $`1,\mathrm{},l`$, then we write $`P`$ and $`Q`$ across the last row and column of the array and use the (reverse) local rule to complete the growth diagram, obtaining a pair $`(T,U)`$ which share an inner border. We indicate this by writing $`(P,Q)=(T,U)`$.
Similarly, a switching local rule $``$ determines an involution on pairs of Young tableaux $`(P,Q)`$ where $`P`$ extends $`Q`$ and we write $`(P,Q)=(T,U)`$. These mappings have the following properties.
###### Theorem 3.1.
A switching local rule $``$ determines an involution
$$\left\{\text{Tableaux }P\text{ and }Q\text{ where }Q\text{ extends }P\text{.}\right\}\stackrel{}{\leftarrow -\to }\left\{\text{Tableaux }T\text{ and }U\text{ where }T\text{ extends }U\text{.}\right\}$$
such that if $`(P,Q)=(T,U)`$, then $`P`$ and $`T`$ have the same content, as do $`Q`$ and $`U`$. Also, $`P`$ and $`U`$ have the same inner border and $`Q`$ and $`T`$ have the same outer border.
A reversible insertion local rule $``$ determines a bijection
$$\left\{\text{Tableaux }P\text{ and }Q\text{ which share an inner border and occupy columns }1,\mathrm{},l\text{.}\right\}\stackrel{}{\leftarrow -\to }\left\{\text{Tableaux }T\text{ and }U\text{ which share an outer border and occupy columns }1,\mathrm{},l\text{.}\right\}$$
with $``$ the identity such that if $`(P,Q)=(T,U)`$, then $`P`$ and $`T`$ have the same content, as do $`Q`$ and $`U`$. Also, the outer border of $`P`$ equals the inner border of $`U`$, and the outer border of $`Q`$ equals the inner border of $`T`$.
We combine complementation of growth diagrams with this local rules construction of tableaux algorithms. Given a local rule (and mapping) $``$, define its complement $`^C`$ by
(3.2)
$$^C(P,Q)=(T,U^C)\text{, where }(P,Q^C)=(T,U),$$
when this is well-defined. (Recall our convention that $`k`$ and $`l`$ are fixed when complementing several tableaux.) In other words, complement one tableau, perform the operation determined by $``$, then complement back the appropiate tableau.
One set of conditions on $``$ which ensures this is well-defined is the following.
* $``$ does not increase the number of columns in tableaux. By this we mean that if $`(P,Q)=(T,U)`$ with $`P`$ and $`Q`$ having at most $`l`$ columns, then $`T`$ and $`U`$ have at most $`l`$ columns. (This is automatically satisfied by switching local rules.)
* We have $`(P,Q)=(T,U)`$ if and only if $`(P^C,Q^C)=(T^C,U^C)`$.
Any local rule $``$ satisfying conditions I and II has a complement $`^C`$ which also satisfies conditions I and II.
When an insertion local rule satisfies conditions I and II, the tableaux operation $`^C`$ is given by the switching local rule $`^C(\alpha ,\beta ,\delta )=\gamma `$, where $``$ completes the partial growth diagram
$$\begin{array}{cc}\beta & (l,\alpha )\\ \delta & \end{array}$$
with the partition $`(l,\gamma )`$, for $`l\delta _m`$, the initial part of $`\delta `$. Also, the growth diagrams for $``$ and $`^C`$ fit into an array of four growth diagrams as in Section 3.1, displayed schematically in Figure 3. (Shown here for an insertion local rule $``$.)
### 3.4. Dual equivalence
Let $``$ be a reversible insertion local rule. Given a tableau $`Q`$ sharing a border (inner or outer) with a tableau $`P`$, set $`(T,U)=(P,Q)`$. We call this passage from $`P`$ to $`T`$ (which is determined by $`Q`$) an $``$-move applied to $`P`$. Two tableaux $`P`$ and $`P^{}`$ with the same shape are $``$-dual equivalent if applying the same sequence of $``$-moves to $`P`$ and to $`P^{}`$ gives tableaux of the same shape. When $``$ is a switching local rule, we may similarly define $``$-moves and $``$-dual equivalence. An elementary consequence of these definitions is that if $`Q`$ and $`Q^{}`$ are $``$-dual equivalent, then the $``$-moves determined by $`Q`$ and $`Q^{}`$ are identical.
When a local rule $``$ has a complement $`^C`$, $``$-dual equivalence coincides with $`^C`$-dual equivalence.
###### Theorem 3.2.
Let $``$ be a local rule. If (3.2) defines a complementary mapping $`^C`$, then $``$-dual equivalence coincides with $`^C`$-dual equivalence. If $``$ satisfies condition II, then a tableau $`P`$ is $``$-dual equivalent to $`Q`$ if and only if $`P^C`$ is $``$-dual equivalent to $`Q^C`$.
Proof. Complementing the second coordinates of a sequence of $``$-moves applied to a tableau $`P`$ gives a sequence of $`^C`$-moves applied to $`P`$. Thus $``$-dual equivalence coincides with $`^C`$-dual equivalence.
When $``$ satisfies condition II, complementing both coordinates of a sequence of $``$-moves applied to a tableau $`P`$ gives a sequence of $``$-moves applied to $`P^C`$. Thus $`P`$ is $``$-dual equivalent to $`Q`$ if and only if $`P^C`$ is $``$-dual equivalent to $`Q^C`$.
## 4. Internal row insertion
We apply the formalism of Section 3 to show Schützenberger’s jeu de taquin and internal row insertion (a modification of the procedure introduced in ) are complementary operations on pairs of tableaux. We first define the insertion local rule $``$ for internal row insertion. Suppose $`\alpha \beta `$ and $`\alpha \gamma `$ with $`\beta /\alpha `$ and $`\gamma /\alpha `$ horizontal strips, and prepending an integer if necessary, we assume that the initial ($`m`$th) parts of $`\alpha ,\beta `$, and $`\gamma `$ are equal to the same number $`l`$. Set $`\delta _m=l`$ and for $`i>m`$,
(4.1)
$$\delta _i:=\mathrm{max}\{\beta _i,\gamma _i\}+\mathrm{min}\{\beta _{i1},\gamma _{i1}\}\alpha _{i1}.$$
Define $`(\alpha ;\beta ,\gamma ):=\delta `$.
###### Lemma 4.1.
The rule (4.1) defines a reversible insertion local rule $``$.
Proof. The rule (4.1) is symmetric in $`\beta `$ and $`\gamma `$, and $`\alpha `$ is determined by $`\beta `$, $`\gamma `$, and $`\delta `$, and so it is reversible. If $`\delta `$ is defined by (4.1), then
$$\begin{array}{cc}\alpha & \beta \\ \gamma & \delta \end{array}$$
is a growth diagram. Indeed, as $`\mathrm{max}\{\beta _i,\gamma _i\}\delta _i`$, both $`\delta /\beta `$ and $`\delta /\gamma `$ will be horizontal strips if $`\delta _i\mathrm{min}\{\beta _{i1},\gamma _{i1}\}`$. But this follows since $`\mathrm{max}\{\beta _i,\gamma _i\}\alpha _{i1}`$, as $`\beta /\alpha `$ and $`\gamma /\alpha `$ are horizontal strips.
It remains to show that $`_i\delta _i_i\gamma _i=_i\beta _i_i\alpha _i`$. Let $`n`$ be an index such that $`\alpha _n=\beta _n=\gamma _n=0`$. Since $`\delta _m=l`$, we have
$`{\displaystyle \underset{i}{}}\delta _i`$ $`=`$ $`l+{\displaystyle \underset{i=m+1}{\overset{n}{}}}\mathrm{max}\{\beta _i,\gamma _i\}+\mathrm{min}\{\beta _{i1},\gamma _{i1}\}\alpha _{i1}`$
$`=`$ $`{\displaystyle \underset{i=m}{\overset{n}{}}}\mathrm{max}\{\beta _i,\gamma _i\}+\mathrm{min}\{\beta _i,\gamma _i\}\alpha _{i1}`$
$`=`$ $`{\displaystyle \underset{i}{}}\gamma _i+{\displaystyle \underset{i}{}}\beta _i{\displaystyle \underset{i}{}}\alpha _i,`$
which completes the proof.
We define internal row insertion to be the tableaux operation determined by the insertion local rule $``$, as in Section 3.3. Observe that $``$ satisfies the conditions I and II of Section 3.3, and so it has a complement, $`^C`$. We show that $`^C`$ coincides with the jeu de taquin operation $`𝒥`$ of Section 1, formalising the duality between jeu de taquin slides and row insertion discovered by Stembridge. We later relate these formulations of $`𝒥`$ and $``$ to the traditional descriptions found in .
###### Theorem 4.2.
The tableaux operation $`^C`$ coincides with the jeu de taquin $`𝒥`$.
Figure 4 shows $``$ and $`𝒥`$ applied to complementary pairs of horizontal strips. Our convention is to number one tableau with letters, the other with numbers to help distinguish them; it is only the chain of shapes that matters. When one tableau extends another, we juxtapose them and use a thick line to indicate the boundary between the two.
Proof of Theorem 2.2. For statement (i), suppose we have a tableau $`Q`$ extending a tableau $`P`$. If we complement both the rows and the columns of the growth diagram computing $`𝒥(P,Q)`$, then we obtain a growth diagram computing $`𝒥(P^C,Q^C)`$. Thus
$$(T,U)=𝒥(P,Q)(T^C,U^C)=𝒥(P^C,Q^C),$$
from which the statement (i) follows.
In the notation of this section, assertion (ii) becomes
* $`U`$ is $`𝒥`$-dual equivalent to $`T`$ if and only if $`U^C`$ is $`𝒥`$-dual equivalent to $`T^C`$.
But this follows by Theorem 3.2.
Proof of Theorem 4.2. Given a tableau $`T:\alpha \beta \delta `$ with $`l=\delta _m`$, the initial part of $`\delta `$, we have $`^C(\alpha ,\beta ,\delta )=\gamma `$, where $``$ completes the partial growth diagram
$$\begin{array}{cc}\beta & (l,\alpha )\\ \delta & \end{array}$$
with the partition $`(l,\gamma )`$. Thus for each $`i`$,
$$(l,\gamma )_{i+1}=\mathrm{max}\{\delta _{i+1},(l,\alpha )_{i+1}\}+\mathrm{min}\{\delta _i,(l,\alpha )_i\}\beta _i$$
or
(4.2)
$$\gamma _i=\mathrm{max}\{\delta _{i+1},\alpha _i\}+\mathrm{min}\{\delta _i,\alpha _{i1}\}\beta _i.$$
We describe this in terms of the tableau $`T`$. Suppose $`\beta /\alpha `$ is filled with $`a`$’s and $`\delta /\beta `$ is filled with $`1`$’s. Form a new tableau $`U`$ as follows.
* In each column of $`T`$ that contains both an $`a`$ and a $`1`$, interchange these symbols, and afterwards
* in each row segment containing both $`a`$’s and $`1`$’s left fixed under (i), shift the $`a`$’s to the right and the $`1`$’s to the left.
We claim that $`U`$ is the tableau $`\alpha \gamma \delta `$ with $`\gamma /\alpha `$ filled with $`1`$’s and $`\delta /\gamma `$ filled $`a`$’s. Any tableau $`U^{}:\alpha \gamma ^{}\delta `$ has each column of length 2 filled with a $`1`$ above an $`a`$, as in (i). Observe that $`\beta _i\mathrm{max}\{\delta _{i+1},\alpha _i\}`$ is the number of $`a`$’s in the $`i`$th row of $`T`$ left fixed by (i) and $`\mathrm{min}\{\delta _i,\alpha _{i1}\}\beta _i`$ the number of $`1`$’s left fixed by (i). The claim follows by observing that $`\gamma _i\mathrm{max}\{\delta _{i+1},\alpha _i\}=\mathrm{min}\{\delta _i,\alpha _{i1}\}\beta _i`$.
The theorem follows as rules (i) and (ii) describe the action of the jeu de taquin on horizontal strips as given by James and Kerber \[5, pp. 91-92\], and this suffices to describe the jeu de taquin (see also ).
The local rule $``$ is described in terms of horizontal strips filled with entries $`a`$ and $`1`$ as follows. Given a horizontal strip $`\beta /\alpha `$ filled with $`a`$’s and a horizontal strip $`\gamma /\alpha `$ filled with $`1`$’s, transfer any $`a`$’s and $`1`$’s which occupy the same boxes in row $`i`$ into the next row, beginning with $`\mathrm{max}\{\beta _i,\gamma _i\}`$. The reverse operation is similarly described.
###### Example 4.3.
The first column $`P`$ and first row $`Q`$ of the growth diagram (3.1) are the following tableaux:
$$P=\text{}\mathrm{and}Q=\text{}.$$
We claim that this growth diagram is obtained from $`P`$ and $`Q`$ using $``$. For example, consider the upper left square of (3.1). Set $`l:=3`$, the length of the first (undisplayed) part of the partitions, and set $`\alpha =(3,1,0)`$, $`\beta =(3,2,0)`$, $`\gamma =(3,2,1)`$, and $`\delta =(3,2,2)`$. Note that $`\delta _1=3`$, and
$$\begin{array}{ccccccc}\delta _2& =& 2& =& 2+33& =& \mathrm{max}\{\beta _2,\gamma _2\}+\mathrm{min}\{\beta _1,\gamma _1\}\alpha _1,\hfill \\ \delta _3& =& 2& =& 1+21& =& \mathrm{max}\{\beta _3,\gamma _3\}+\mathrm{min}\{\beta _2,\gamma _2\}\alpha _2,\hfill \end{array}$$
in agreement with (4.1). If we consider the last row and last column of (3.1), we see that $`(P,Q)=(T,U)`$, where
$$T=\text{}\mathrm{and}U=\text{}.$$
The additional 3 growth diagrams in Figure 2 provide further examples of $`𝒥`$ and $``$. Consider the lower left growth diagram of Figure 2.
$$\begin{array}{cccc}221\hfill & 222\hfill & 3221\hfill & 3322\hfill \\ 4211\hfill & 4221\hfill & 43211\hfill & 43321\hfill \\ 4421\hfill & 4422\hfill & 44321\hfill & 44332\hfill \\ 4441\hfill & 4442\hfill & 44431\hfill & 44433\hfill \end{array}$$
Let $`P^C`$ be the tableau of the first column (filled with $`a<b<c`$) and $`Q`$ the tableau of the last row. From this, we see that $`𝒥(P^C,Q)=(T^C,U)`$, i.e.,
$$𝒥\left(\text{}\right)=\text{}.$$
Both the jeu de taquin and the Robinson-Schensted correspondence (of which internal row insertion is a variant) are described in via growth diagrams consisting of standard tableaux as follows.
In Appendix 1 to Chapter 7 in , Fomin gives the following local rule for the jeu de taquin: If $`\alpha \beta \delta `$, then either the interval $`[\alpha ,\delta ]`$ in Young’s lattice contains a fourth partition $`\gamma \beta `$, or else the interval $`[\alpha ,\delta ]`$ is a chain, in which case we set $`\gamma :=\beta `$. One may verify from the definition (4.2) of $`𝒥`$ that $`\gamma =𝒥(\alpha ,\beta ,\delta )`$.
Similarly, suppose that $`\alpha \beta `$ and $`\alpha \gamma `$. In Section 7.13 of , Stanley gives Fomin’s local rule for Robinson-Schensted insertion.
* If $`\beta \gamma `$, then let $`\delta :=\beta \gamma `$, the unique partition covering both $`\beta `$ and $`\gamma `$.
* If $`\beta =\gamma `$, then $`\beta /\alpha `$ is a single box in the $`i`$th row, and we define $`\delta `$ so that $`\delta /\beta `$ is a single box in the $`(i+1)`$st row.
One may verify from the definition (4.1) of $``$ that $`\delta =(\alpha ;\beta ,\gamma )`$. (The other 2 possibilities in of $`\alpha =\beta =\gamma (=\delta )`$ and $`\alpha =\beta =\gamma `$ with a mark in the square—indicating an external insertion—do not occur for us.)
Let $`st(\beta /\alpha )`$ denote the standard renumbering of a horizontal strip $`\beta /\alpha `$, which is now a standard tableau. The familiar fact that the jeu de taquin and Schensted insertion commute with standard renumbering manifests itself here as follows: $`\delta =(\alpha ;\beta ,\gamma )`$ if and only if $`(st(\delta /\gamma ),st(\delta /\beta ))=(st(\beta /\alpha ),st(\gamma /\alpha ))`$, and $`\gamma =𝒥(\alpha ,\beta ,\delta )`$ if and only if $`(st(\delta /\gamma ),st(\gamma /\alpha ))=𝒥(st(\beta /\alpha ),st(\delta /\beta ))`$.
The skew insertion procedure of Sagan and Stanley mixes Schensted insertion with the internal insertion procedure $``$. If tableaux $`T`$ and $`U`$ share a common inner border, then $``$ acts on the pair $`(T,U)`$ in the same way as the forward direction of the procedure of Theorem 6.11 when the matrix word $`\pi `$ is empty.
The difference lies in the reverse procedure, where $`P`$ and $`Q`$ share an outer border. The essence of this difference occurs when $`P`$ and $`Q`$ are single boxes. Suppose $`\beta \delta `$ and $`\gamma \delta `$ (so that $`\delta /\beta `$ and $`\delta /\gamma `$ are single boxes) and suppose we have $`\alpha =(\beta ,\gamma ;\delta )`$. By the definition of $``$, if $`\beta \gamma `$, then $`\alpha `$ is the unique partition covered by both $`\beta `$ and $`\gamma `$. If however $`\beta =\gamma `$ and $`\delta /\beta `$ is a single box in the $`i`$th row, then $`\beta /\alpha `$ is a single box in the ($`i1`$)st row.
This differs from the procedure in only in the case (2) when $`i`$ is the initial row, that is, $`\beta =\gamma `$ and $`\delta /\beta `$ is a single box in the first row. Then Sagan and Stanley set $`\alpha =\beta `$, and bump a number out of their tableaux, forming part of the matrix word $`\pi `$. We avoid this by assuming in effect that our partitions have a previous row of length $`l=\delta _1`$, which is empty in the skew shapes $`\delta /\beta `$ and $`\delta /\gamma `$.
## 5. Internal Column Insertion, Hopscotch, and Stable Tableaux
For standard tableaux (and hence for all tableaux via standard renumbering), internal column insertion is essentially the same as internal row insertion—one simply replaces ‘row’ by ‘column’ in the definitions to obtain internal column insertion. While internal column insertion gives nothing new of itself, when combined with complementation, we do get something interesting. We call this new operation hopscotch. Hopscotch is defined on pairs $`P`$ and $`Q`$, where $`Q`$ extends $`P`$, and one of $`P`$, $`Q`$ is a tableau, while the other is a stable tableau, which is defined in Section 5.3 below.
### 5.1. Local rules for column insertion
We give the following local rules formulation of internal column insertion for standard tableaux, which is the matrix transpose of that for internal row insertion. Given partitions $`\alpha \beta `$ and $`\alpha \gamma `$, define a partition $`\delta `$ by
1. If $`\gamma \beta `$, then $`\delta `$ is their least upper bound, $`\beta \gamma `$.
2. If $`\gamma =\beta `$, and the box $`\beta /\alpha `$ is in the $`j`$th column, then $`\delta /\beta `$ is a box in the ($`j+1`$)st column.
Set $`𝒞(\alpha ;\beta ,\gamma ):=\delta `$.
This gives an algorithm for standard tableaux which can be generalised to arbitrary tableaux by standard renumbering. An explicit rule for operating on horizontal strips is given in Section 5.2.
Consider now the reverse of this procedure. Given partitions $`\beta \delta `$ and $`\gamma \delta `$, we define $`\alpha `$ by
1. If $`\gamma \beta `$, then $`\alpha `$ is their greatest lower bound, $`\beta \gamma `$.
2. If $`\gamma =\beta `$, and the box $`\delta /\beta `$ is in the $`j`$th column, then $`\beta /\alpha `$ is a box in the ($`j1`$)st column.
Set $`𝒞(\beta ,\gamma ;\delta ):=\alpha `$.
The reverse procedure does not work when $`\delta /\beta `$ is a box in the first column. This forces us to generalise our notions of shape and tableaux. The basic idea is to allow the creation of new columns, labeled with non-positive integers, to the left of existing columns when we need them.
Henceforth we define a shape to be a finite weakly decreasing sequence of integers (positive or negative), called parts. This set of shapes with a fixed length $`m`$ forms a poset under componentwise comparison, which extends Young’s lattice of partitions. Using $``$ for this partial order, we define skew shapes $`\beta /\alpha `$ for $`\alpha \beta `$ as before. A horizontal strip is (as before), a skew shape $`\beta /\alpha `$ with at most one box in each column. We define a tableau of shape $`\beta /\alpha `$ with entries from $`[k]`$ to be a chain $`\alpha =\beta ^0\beta ^1\mathrm{}\beta ^k=\beta `$ where each $`\beta ^i/\beta ^{i1}`$ is a horizontal strip. We can convert a tableau $`T`$ (as defined here) into a tableau as defined in , merely by adding a large enough number to each part of the shapes that define $`T`$, in effect shifting $`T`$ horizontally.
### 5.2. Internal column insertion
We extend this rule $`𝒞`$ for standard tableaux to an insertion local rule $`𝒞`$ for tableaux so that the resulting tableaux operation commutes with standard renumbering. Let $`\beta /\alpha `$ and $`\gamma /\alpha `$ be horizontal strips, with $`\alpha `$, $`\beta `$, and $`\gamma `$ shapes. Consider applying the tableaux algorithm $`𝒞`$ to the standard renumberings of the horizontal strips $`\beta /\alpha `$ and $`\gamma /\alpha `$. This proceeds from left to right (from the last row to the first). If there is any overlap between $`\beta /\alpha `$ and $`\gamma /\alpha `$ in the $`i`$th row, then the entries in that row are shifted to the right by the amount of that overlap, displacing entries in the previous row, if necessary. For example, suppose $`\alpha =(6,3,0)`$, $`\beta =(7,6,1)`$, and $`\gamma =(8,4,1)`$. Then we have
$`\stackrel{𝒞}{–⟶}`$ .
and so $`\delta =(10,6,2)`$.
Formally, given shapes $`\alpha `$, $`\beta `$, and $`\gamma `$ with $`\beta /\alpha `$ and $`\gamma /\alpha `$ horizontal strips, we define the shape $`\delta `$ recursively as follows. Suppose that the shapes $`\alpha `$, $`\beta `$, and $`\gamma `$ each have $`m`$ parts. Set $`d_m=0`$, and for $`i=m,m1,\mathrm{},2`$, set
(5.1)
$$\begin{array}{ccc}\hfill \delta _i& :=& \mathrm{min}\{\begin{array}{c}\beta _i+\gamma _i+d_i\alpha _i\hfill \\ \alpha _{i1}\hfill \end{array},\hfill \\ \hfill d_{i1}& :=& \mathrm{max}\{\begin{array}{c}\beta _i+\gamma _i+d_i\alpha _i\alpha _{i1}\hfill \\ 0\hfill \end{array},\text{}\hfill \end{array}$$
and set $`\delta _1:=\beta _1+\gamma _1+d_1\alpha _1`$. We define $`𝒞(\alpha ;\beta ,\gamma ):=\delta `$. Note that $`𝒞`$ does not change the number of rows of shapes. In the above example, $`d_3=d_2=0`$ while $`d_1=1`$. The numbers $`d_i`$’s keep track of boxes ‘bumped’ from row $`i1`$ to row $`i`$.
###### Example 5.1.
We apply $`𝒞`$ to the tableaux $`P`$ and $`Q`$ of the upper left growth diagram in Figure 2 to obtain the following growth diagram.
(5.2)
$$\begin{array}{cccc}1\hfill & 2\hfill & 31\hfill & 33\hfill \\ 21\hfill & 31\hfill & 42\hfill & 53\hfill \\ 211\hfill & 311\hfill & 421\hfill & 531\hfill \\ 221\hfill & 321\hfill & 431\hfill & 541\hfill \end{array}$$
###### Lemma 5.2.
The rule (5.1) defines a local rule $`𝒞`$ which is reversible in that, given $`\beta ,\gamma `$, and $`\delta `$ with $`\delta /\beta `$ and $`\delta /\gamma `$ horizontal strips, there is a unique shape $`\alpha `$ with $`𝒞(\alpha ;\beta ,\gamma )=\delta `$.
Proof. The only condition that needs to be checked is that $`\delta `$, $`\beta `$, and $`\gamma `$ determine $`\alpha `$. But this follows from the corresponding property of $``$ as we may compute $`\alpha `$ from $`\delta `$, $`\beta `$, and $`\gamma `$ by the standard renumbering of $`\delta /\beta `$ and $`\delta /\gamma `$, using the local rule $`𝒞`$ for standard tableaux, which is equivalent to $``$ via conjugation.
We define internal column insertion $`𝒞`$ to be the tableaux operation determined by the local rule $`𝒞`$.
###### Theorem 5.3.
Internal column insertion is a bijection
$$\left\{\text{Tableaux }P\text{ and }Q\text{ which share an inner border.}\right\}\stackrel{𝒞}{\leftarrow -\to }\left\{\text{Tableaux }T\text{ and }U\text{ which share an outer border.}\right\}$$
with $`𝒞𝒞`$ the identity such that if $`𝒞(P,Q)=(T,U)`$, then $`P`$ is Knuth-equivalent to $`T`$ and $`Q`$ is Knuth-equivalent to $`U`$. Also, the outer border of $`P`$ equals the inner border of $`U`$, and outer border of $`Q`$ equals the inner border of $`T`$.
Furthermore, $`𝒞`$-dual equivalence coincides with $`𝒥`$-dual equivalence.
Proof. To show Knuth-equivalence, let $`𝒞(P,Q)=(T,U)`$. Since $`𝒞`$ commutes with standard renumbering, we may assume $`Q`$ is standard. Then $`T`$ is obtained from $`P`$ by a series of internal column insertions. Adapting the arguments of , we see that internal column insertion preserves Knuth equivalence. Thus $`P`$ is Knuth-equivalent to $`T`$. Since $`𝒞`$ is symmetric in its two arguments, $`Q`$ is Knuth-equivalent to $`U`$.
Consider applying a sequence of $`𝒞`$-moves to a tableau $`P`$. Since $`𝒞`$ commutes with standard renumbering, we may assume that each $`𝒞`$-move in that sequence is given by a tableau consisting of a single box. The effect of these moves on the shape of $`P`$ will be the same as the effect on the shape of the standard renumbering $`st(P)`$ of $`P`$. If we take the conjugate (matrix transpose) of tableaux in this sequence of $`𝒞`$-moves applied to $`st(P)`$, we obtain a sequence of $``$-moves applied to the conjugate $`st(P)^t`$ of $`st(P)`$. Since $`𝒥=^C`$, this observation and Theorem 3.2 imply that $`P`$ and $`Q`$ are $`𝒞`$-dual equivalent if and only if $`st(P)^t`$ and $`st(Q)^t`$ are $`𝒥`$-dual-equivalent.
However, a tableau $`P`$ is dual-equivalent to its standard renumbering, as $`𝒥`$ commutes with standard renumbering. Also, two dual-equivalent standard tableaux have dual equivalent conjugates, as the conjugate of a jeu de taquin slide is another jeu de taquin slide, by (1.1). This proves that $`𝒞`$ and $`𝒥`$ have the same dual equivalence classes.
### 5.3. Hopscotch and stable tableaux
We would like to define the new tableaux operation $``$ of hopscotch to be the tableaux operation complementary to internal column insertion $`𝒞`$. That is, if $`P`$ and $`Q`$ are tableaux with $`Q`$ extending $`P`$, then we would set
(5.3)
$$(P,Q):=(T,U^C)\text{where}𝒞(P,Q^C)=(T,U).$$
(Recall that $`U`$ and $`Q^C`$ have the same content.) Unfortunately, this does not work in general. This is because $`𝒞`$ can increase the number of columns in tableaux, violating condition I of Section 3.3. However, we will salvage something from this idea, by extending our notion of “tableau”.
The essential problem comes from the fact that $`𝒞`$ may increase the number of columns of a tableau. Thus, if we use a complementation parameter $`l`$ to form $`Q^C`$, then the tableau $`U`$ may extend beyond the $`l`$th column, and thus forming $`U^C`$ will require a different complementation parameter, larger than $`l`$. This cannot in general be remedied by increasing $`l`$ to some other integer. To solve this problem we build an asymmetry into $``$, requiring that one of $`P`$ or $`Q`$ is a new object called a “stable tableau”, which is the complement of an ordinary tableau with respect to infinitely many columns.
We formalise this idea. A stable shape $`\alpha `$ is a finite weakly decreasing sequence of integers and the symbols $`\mathrm{}`$ and $`\overline{\mathrm{}}`$, where $`\overline{\mathrm{}}<n<\mathrm{}`$ for any integer $`n`$. If we regard a shape as having arbitrarily many trailing $`\overline{\mathrm{}}`$s, then stable shapes form a lattice under componentwise comparison, which contains our earlier poset of shapes. We form the stable skew shape $`\beta /\alpha `$ as before. If the finite parts of $`\alpha `$ and $`\beta `$ occur in the same rows, then we can ignore initial $`\mathrm{}`$’s and trailing $`\overline{\mathrm{}}`$’s and regard $`\beta /\alpha `$ as an ordinary skew shape. In this way, stable skew shapes and ordinary skew shapes may extend each other.
A stable horizontal strip $`\beta /\alpha `$ is a pair of stable shapes $`\alpha \beta `$ where, if $`\alpha _m`$ is the first finite part of $`\alpha `$ and $`\beta _n`$ the last finite part of $`\beta `$, then
$$\alpha _{m1}=\mathrm{}=\beta _m>\alpha _m\beta _{m+1}\mathrm{}\alpha _{n1}\beta _n>\alpha _n=\overline{\mathrm{}}=\beta _{n+1}.$$
If we consider $`\beta /\alpha `$ as a collection of boxes in the plane, then $`\beta /\alpha `$ has at most one box in each column, and it has 2 half-infinite rows of boxes extending from either end. Equivalently, if we let $`\alpha ^S:=(\mathrm{},\alpha )`$, then $`\beta /\alpha `$ is a stable horizontal strip if and only if $`\beta \alpha ^S`$, and, upon removing the $`\mathrm{}`$’s and $`\overline{\mathrm{}}`$’s from $`\beta `$ and $`\alpha ^S`$, $`\alpha ^S/\beta `$ is an ordinary horizontal strip with boxes in the empty columns of $`\beta /\alpha `$.
Let $`\alpha =(\mathrm{},4,3,1,\overline{2},\overline{\mathrm{}},\overline{\mathrm{}})`$ and $`\beta =(\mathrm{},\mathrm{},3,3,\overline{1},\overline{3},\overline{\mathrm{}})`$, where we write $`\overline{a}`$ for a negative integer $`a`$. Figure 5 shows the stable horizontal strip $`\beta /\alpha `$ and its complementary horizontal strip $`\alpha ^S/\beta `$.
A stable tableau $`T`$ with stable shape $`\beta /\alpha `$ and entries from $`[k]`$ is a chain $`\alpha =\beta ^0\beta ^1\mathrm{}\beta ^k=\beta `$ of stable shapes, where the successive stable skew shapes $`\beta ^i/\beta ^{i1}`$ are stable horizontal strips. The appropriate notion of complementation for stable tableaux involves complementing all columns. Thus if $`T:\beta ^0\beta ^1\mathrm{}\beta ^k`$ is a stable tableaux, its stable complement $`T^S`$ is the chain $`\beta ^k(\mathrm{},\beta ^{k1})\mathrm{}(\mathrm{}^k,\beta ^0)`$. Since $`T`$ is a stable tableau, each shape $`\beta ^i`$ has one more part of $`\mathrm{}`$ than its predecessor $`\beta ^{i1}`$ and the same number of finite parts, and so we may remove the same number of parts of $`\mathrm{}`$ from each shape $`(\mathrm{}^{ki},\beta ^i)`$ of $`T^S`$ and obtain an ordinary tableau. Similarly, if we complement an ordinary tableau $`T`$ using the complementation parameter $`\mathrm{}`$, then we obtain its stable complement $`T^S`$, a stable tableau. Figure 6 shows a tableau $`T`$ and its stable complement $`T^S`$ when $`k=5`$.
We summarise the properties of stable complementation.
###### Theorem 5.4.
Let $`T`$ be a tableau of shape $`\beta /\alpha `$ and entries from $`[k]`$. Then $`T^S`$ is a stable tableau with stable shape $`(\mathrm{}^k,\alpha )/\beta `$ and entries from $`[k]`$. Similarly, if $`T`$ is a stable tableau with entries from $`[k]`$, then $`T^S`$ is an ordinary tableau with entries from $`[k]`$. In either case, we have $`T^{SS}=T`$.
(For the last assertion, we identify tableaux that differ only by a vertical shift.)
###### Definition 5.5.
We define hopscotch, $``$, to be the stable complement of internal column insertion. That is, given a tableau $`P`$ and a stable tableau $`Q`$ with either $`P`$ extending $`Q`$ or $`Q`$ extending $`P`$, then we set
$$(P;Q):=(T;U^S)\text{where}𝒞(P,Q^S)=(T,U).$$
Suppose we have $`(P;Q)=(T;U)`$. We call the passage from the tableau $`P`$ to the tableau $`T`$ an $``$-move applied to $`P`$ (which depends upon the stable tableau $`Q`$). Likewise, the passage from $`Q`$ to the stable tableau $`U`$ is also called a $``$-move. As in Section 3.4, this gives rise to the notion of $``$-dual equivalence for tableaux and also for stable tableaux.
###### Theorem 5.6.
The tableaux operation hopscotch gives a bijection
$$\left\{\text{Pairs }(P;Q)\text{ with }P\text{ a tableau, }Q\text{ a stable tableau, and }Q\text{ extending }P\text{.}\right\}\stackrel{}{\leftarrow -\to }\left\{\text{Pairs }(T;U)\text{ with }T\text{ a tableau, }U\text{ a stable tableau, and }T\text{ extending }U\text{.}\right\}$$
with $``$ the identity. If $`(P;Q)=(T;U)`$, then $`P`$ is Knuth-equivalent to $`T`$.
Furthermore, two tableaux $`P`$ and $`T`$ are $``$-dual equivalent if and only if $`P`$ and $`T`$ are $`𝒥`$-dual equivalent, and two stable tableaux $`Q`$ and $`U`$ are $``$-dual equivalent if and only if the tableaux $`Q^S`$ and $`U^S`$ are $`𝒥`$-dual equivalent.
Proof. These properties of hopscotch all follow from the corresponding properties of internal column insertion (Theorem 5.3), the definition of hopscotch, and properties of stable complementation (Theorem 5.4).
## 6. Hopscotch and Tesler’s Shift Games
We give local rules for hopscotch and relate it to Tesler’s rightward- and leftward-shift games.
### 6.1. Local Rules for Hopscotch
We give local rules for $``$ when the ordinary tableau is standard, that is when each horizontal strip in $`P`$ is a single box. As hopscotch commutes with standard renumbering (because internal column insertion does), this enables the computation of $`(P;Q)`$ when $`P`$ is an arbitrary tableau.
###### Definition 6.1.
Let $`\alpha \beta \delta `$ be shapes with $`\beta /\alpha `$ a stable horizontal strip. Let $`c`$ denote the column of $`\delta /\beta `$, and $`H`$ denote the set of columns of $`\beta /\alpha `$. Define the stable shape $`\gamma `$ with $`\alpha \gamma \delta `$ as follows.
1. If $`cH`$, then define $`\gamma `$ so that $`\alpha `$ and $`\gamma `$ also differ in column $`c`$. (So $`\delta /\gamma `$ will have the same set of columns $`H`$.)
2. If $`cH`$, then let $`c^{}`$ denote the smallest element of $`H`$ greater than $`c`$, and define $`\gamma `$ so that $`\gamma `$ and $`\alpha `$ differ in column $`c^{}`$. (So $`\delta /\gamma `$ will have the set of columns $`H^{}:=H\{c^{}\}\backslash \{c\}`$.)
###### Remark 6.2.
In either case of Definition 6.1, if $`\beta `$ and $`\delta `$ differ in row $`r`$, then $`\alpha `$ and $`\gamma `$ will differ in row $`r^{}`$, where $`r^{}`$ is the largest row less than $`r`$ where $`\alpha `$ and $`\beta `$ differ.
The local rule described in Definition 6.1 may sometimes be computed when $`\beta /\alpha `$ is an ordinary horizontal strip—which may be considered to be the truncation to a finite interval of columns of a stable horizontal strip. The $`2\times 6`$ array in Figure 7 is a growth diagram of ordinary shapes, with the second row filled in from right to left using the local rule of Definition 6.1.
###### Theorem 6.3.
Given stable shapes $`\alpha \beta \delta `$ with $`\beta /\alpha `$ a stable horizontal strip, the stable shape $`\gamma `$ defined in Definition 6.1 satisfies $`\gamma =(\alpha ,\beta ,\delta )`$, i.e., Definition 6.1 gives a local rules description of $``$.
Proof. The theorem follows from the description for $`𝒞`$ when one strip is standard. Let $`\alpha `$, $`\beta `$, and $`\gamma `$ be shapes with $`\alpha \gamma `$, $`\alpha \beta `$ with $`\beta /\alpha `$ a horizontal strip. Then the shape $`\delta =𝒞(\alpha ;\beta ,\gamma )`$ may also be defined by
1. If $`\gamma \beta `$, then set $`\delta :=\gamma \beta `$.
2. Otherwise, $`\gamma \beta `$. In this case, $`\gamma /\alpha `$ is a box in column $`c`$. Choose $`\delta `$ covering $`\beta `$ so that the box $`\delta /\beta `$ is in column $`c^{}`$ minimal subject to $`c^{}>c`$ with $`\delta /\alpha `$ is a horizontal strip.
The reverse local rule for $``$ given $`\alpha \beta \delta `$, when the inner horizontal strip consists of a single box, is described by replacing Case 2 in Definition 6.1 by
1. If $`cH`$, then let $`c^{}`$ denote the largest element of $`H`$ less than $`c`$, and define $`\gamma `$ so that $`\gamma `$ and $`\alpha `$ differ in column $`c^{}`$. (So $`\delta /\gamma `$ will be the set of columns $`H^{}:=H\{c^{}\}\backslash \{c\}`$.)
### 6.2. Hopscotch and Tesler’s shift games
In his study of semi-primary lattices, Tesler defines certain leftward- and rightward-shift games applied to standard skew tableaux which model the effect of certain lattice-theoretic procedures. These games give rise to an algorithm to construct a tableau $`T`$ of partition shape from a standard skew tableau $`P`$. Studying the corresponding objects in a semi-primary lattice, he then shows that $`T`$ is the result of applying jeu de taquin slides to $`P`$, and thus $`T`$ is the rectification of $`P`$. We show how Tesler’s algorithm may be regarded as a special case of hopscotch and give a combinatorial (as opposed to geometric) proof that his algorithm computes the rectification of $`P`$.
In these games of Tesler, a vertical strip (full of $``$’s) is moved through a tableau $`P`$, and some entries of $`P`$ are removed (forming the $``$’s). We describe the conjugate (replacing columns by rows) of Tesler’s rightward shift game in terms of a local rule. The leftward shift game is the reverse of this procedure.
###### Definition 6.4.
An almost standard tableau $`P`$ is a tableau $`\beta ^0\beta ^1\mathrm{}\beta ^k`$, where each $``$ in the chain is either a cover $``$ or an equality. Given an almost standard tableau $`P:\beta ^0\beta ^1\mathrm{}\beta ^k`$, Tesler’s rightward shift game constructs another almost standard tableau $`𝒮(P):\alpha ^0\mathrm{}\alpha ^k`$ where each $`\beta ^i/\alpha ^i`$ is a horizontal strip. This begins with $`\alpha ^0:=\beta ^0`$. If we have constructed $`\alpha ^0,\mathrm{},\alpha ^{i1}`$, then we have the (partial) growth diagram
$$\begin{array}{cc}\alpha ^{i1}& \\ \beta ^{i1}& \beta ^i\hfill \end{array}.$$
We construct $`\alpha ^i`$ as follows.
1. If $`\beta ^i=\beta ^{i1}`$, then we set $`\alpha ^i:=\alpha ^{i1}`$.
2. If $`\beta ^i/\beta ^{i1}`$ is a single box in the $`r`$th row and if the horizontal strip $`\beta ^{i1}/\alpha ^{i1}`$ has no boxes in rows less than $`r`$, then we set $`\alpha ^i:=\alpha ^{i1}`$.
3. If $`\beta ^i/\beta ^{i1}`$ is a single box in the $`r`$th row and if the horizontal strip $`\beta ^{i1}/\alpha ^{i1}`$ has boxes in rows less than $`r`$, then we choose the largest row $`r^{}`$ of $`\beta ^{i1}/\alpha ^{i1}`$ less than $`r`$ and let $`\alpha ^i/\alpha ^{i1}`$ be a box in that row.
By Remark 6.2, the similarity between these rules, particularly (3), and those for hopscotch is evident.
###### Example 6.5.
We give a completed example of $`𝒮`$ applied to a standard tableau (the second row below):
$$\begin{array}{cccccccccc}421\hfill & \mathit{421}\hfill & \mathit{431}\hfill & \mathit{431}\hfill & \mathit{432}\hfill & \mathit{432}\hfill & \mathit{442}\hfill & \mathit{542}\hfill & \mathit{552}\hfill & \mathit{552}\hfill \\ 421\hfill & 431\hfill & 432\hfill & 442\hfill & 4421\hfill & 5421\hfill & 5422\hfill & 5522\hfill & 5532\hfill & 5542\hfill \end{array}$$
We display this in terms of tableaux.
$`\stackrel{𝒮}{–-⟶}`$
Informally, we move the entries in the tableaux in order northeast to the nearest available space. The 1 vacates the tableau, the 2 moves into the empty space where the 1 was, the 3 vacates, the 4 moves where the 2 was, the 5 vacates, the 6 moves where the 3 was, the 7 moves where the 5 was, the 8 moves where the 7 was, and the 9 vacates.
Let $`r(P)`$ be the integers in $`P`$ but not in $`𝒮(P)`$, that is, the set of indices $`i`$ where $`\beta ^{i1}\beta ^i`$ but $`\alpha ^{i1}=\alpha ^i`$. In the example above, $`r(P)=\{1,3,5,9\}`$. Then Tesler’s rectification algorithm runs as follows: Given a standard skew tableau $`P`$, form tableaux $`P_0,\mathrm{},P_k`$ and sets $`r_1,\mathrm{},r_k`$ recursively, initially setting $`P_0:=P`$ and then
$$P_i:=𝒮(P_{i1})\text{and}r_i:=r(P_{i1}).$$
###### Proposition 6.6 (, Theorem 8.13).
The tableau $`P_k`$ is empty and $`r_1,\mathrm{},r_k`$ form the rows of the rectification $`S`$ of $`P`$.
We relate this to hopscotch, first showing that $`𝒮`$ is equivalent to a particular $``$-move applied to $`P`$, and then we show how to compute the tableau $`S`$ of Proposition 6.6 using hopscotch. This will give a new, combinatorial proof of this result of Tesler.
Let $`P:\beta ^0\beta ^1\mathrm{}\beta ^k`$ be a standard tableau where each shape $`\beta ^i`$ has $`n`$ nonnegative parts (appending 0’s if needed). Let $`b`$ be the first (largest) part of $`\beta ^k`$. Prepend $`\mathrm{}`$ and append 0 to each shape $`\beta ^i`$ and consider it to be a stable shape. Set $`\alpha ^0=\beta ^0`$ and define
$$\gamma ^0:=bb_1b_2\mathrm{}b_n>\overline{\mathrm{}},$$
where $`\beta ^0=\mathrm{}>b_1\mathrm{}b_n0`$. Then the stable horizontal strip $`\beta ^0/\gamma ^0`$ consists of two semi-infinite rows of boxes, one beginning in column $`b+1`$ in the first row, and one ending in column 0 in row $`n+1`$. Set $`Q:=\beta ^0/\gamma ^0`$ and $`(T;U):=(P;Q)`$. Define shapes $`\gamma ^1,\mathrm{},\gamma ^k`$ and $`\alpha ^0,\mathrm{},\alpha ^k`$ so that the standard tableau $`T`$ above is $`\gamma ^0\gamma ^1\mathrm{}\gamma ^k`$ and $`𝒮(P)=\alpha ^0\mathrm{}\alpha ^k`$. For each $`i=0,1,\mathrm{},k`$, set $`a_i:=\mathrm{\#}\{j<i\alpha ^j=\alpha ^{j+1}\}`$.
###### Lemma 6.7.
For each $`i=0,\mathrm{},k`$ the stable shape $`\gamma ^i`$ is
(6.1)
$$b+a_i\alpha _1^i\mathrm{}\alpha _n^i>\overline{\mathrm{}}.$$
In particular, $`r(P)`$ is the first row of $`T`$ and $`𝒮(P)`$ consists of the remaining rows of $`T`$.
Proof. We prove this by induction on $`i`$, the case $`i=0`$ being the definition of $`\gamma ^0`$. Suppose that (6.1) holds for $`\gamma ^0,\mathrm{},\gamma ^{i1}`$. We compare the $`i`$th steps of $``$ and $`𝒮`$. Since $`P`$ is standard, case (1) of Definition 6.4 for $`𝒮`$ does not occur. If we are in case (2) of Definition 6.4, then $`\alpha ^i=\alpha ^{i1}`$ and $`a_i=1+a_{i1}`$. Since the stable horizontal strip $`\beta ^{i1}/\gamma ^{i1}`$ has no boxes in rows between its first row and the row $`r`$ of the single box $`\beta ^i/\beta ^{i1}`$, the stable shapes $`\gamma ^i`$ and $`\gamma ^{i1}`$ differ only in the first row, by Remark 6.2. Finally, case (3) of Definition 6.4 is equivalent to hopscotch.
Figure 8 shows an example of this hopscotch move where we fill the stable horizontal strip $`Q`$ with $``$’s. Note the similarity with Example 6.5.
To compute the rectification tableau $`S`$ of Proposition 6.6 using hopscotch, we iterate the hopscotch move of Lemma 6.7, letting the initial tableau $`P`$ extend sufficiently many stable horizontal strips of the form $`Q`$ above. Given $`P`$ with $`\beta ^i`$, $`b`$, and $`b_i`$ as in the paragraph before Lemma 6.7, set $`l:=\mathrm{min}\{k,n\}`$. Define the stable tableau $`Q:=\alpha ^0\alpha ^1\mathrm{}\alpha ^l`$ by
$$\alpha ^i:=\mathrm{}^i>b^{li}b_1\mathrm{}b_n0^i>\overline{\mathrm{}}^{li}.$$
If we prepend $`\mathrm{}^l`$ and append $`0^l`$ to each partition $`\beta ^i`$, then the standard tableau $`P`$ extends the stable tableau $`Q`$. Set $`(T;U):=(P;Q)`$. Figure 9 displays this action of hopscotch on the tableau $`P`$ of Figure 8, rectifying it to $`T`$.
The tableau $`T`$ has partition shape and is in the initial $`l`$ rows and columns greater than $`b_n`$. By Theorem 5.6, $`T`$ is Knuth-equivalent to $`P`$ and hence is the rectification of $`P`$. Successively applying Lemma 6.7 and analysing the effect of hopscotch on the columns greater than $`b_n`$, we obtain a combinatorial proof that Tesler’s algorithm rectifies tableaux.
## 7. Jeux de Tableaux
In Section 1 we asserted that the new tableaux operations of $``$, $`𝒞`$, and $``$ have ramifications for ordinary tableaux, even though we generalised ordinary tableaux in order to define these operations. We describe four different algorithms to rectify a column-strict tableau, supplementing the classical jeu de taquin of Schützenberger . The first is constructed from the internal (row) insertion of Sagan and Stanley , and the second is similarly derived from column insertion . The last two are related to hopscotch. One, which we call row extraction, is the column-strict tableaux version of Tesler’s rectification algorithm of Section 6.2, and the other, which we call column extraction, is adapted from Stroomer’s column sliding algorithm .
In this section, tableaux all have (skew) partition shape, in the traditional sense: non-negative integer parts, and the initial part of a partition $`\lambda `$ is $`\lambda _1`$.
### 7.1. Internal Insertion Games
Let $`P`$ be a column-strict tableau of shape $`\lambda /\mu `$ with $`\mu ,\lambda `$ partitions. We assume that $`P`$ has entries in both the first row and first column. A(n) (inside) cocorner is an entry of $`P`$ with no neighbours in $`P`$ to the left or above. An internal row insertion on $`P`$ as introduced in begins with a cocorner of $`P`$. That entry is removed from $`P`$ and inserted into the subsequent rows of $`P`$ using Schensted (row) insertion. The resulting tableau is Knuth-equivalent to $`P`$.
The row insertion game begins with a column-strict tableau $`P`$ of shape $`\lambda /\mu `$ and a row $`j`$ with $`\mu _j=0`$. The game proceeds by successive internal row insertions beginning with cocorners in rows $`<j`$, and it ends when there are no such cocorners. The result is a tableau $`T`$ of partition shape (occupying rows $`j`$) Knuth-equivalent to $`P`$, that is, the rectification of $`P`$. Thus the result of the row insertion game is independent of the particular sequence of cocorners chosen. We illustrate this process in Figure 10. In each tableau in that sequence, we shade the cocorner and insertion path that creates the subsequent tableau.
We do the same with column insertion. An internal column insertion on a tableau $`P`$ begins with a cocorner of $`P`$. That entry is removed from $`P`$ and inserted into subsequent columns of $`P`$ using column insertion , and the resulting tableau is Knuth-equivalent to $`P`$.
The column insertion game begins with a column-strict tableau $`P`$ of shape $`\lambda /\mu `$ and proceeds by successive internal column insertions beginning with cocorners in columns $`1,\mathrm{},\mu _1`$, and it ends when there are no such cocorners. The result is a tableau $`T`$ of partition shape (occupying columns $`>\mu _1`$) Knuth-equivalent to $`P`$, that is, the rectification of $`P`$. Thus the result of the column insertion game is independent of the particular sequence of cocorners chosen. We illustrate this process in Figure 11. In each tableau in that sequence, we shade the cocorner and insertion path that creates the subsequent tableau.
### 7.2. Row Extraction
Row extraction is the extension of the rectification algorithm of Tesler, $`𝒮`$, described in Section 6.2, to column-strict tableaux. After describing this extension, we show how it may be used to compute the rectification of a column-strict skew tableau $`P`$. Define the cross order on the cells of a shape or tableau by $`cd`$ if the cell $`c`$ lies in the same column as, or a subsequent column to, $`d`$ and is in a previous row.
Let $`P:\mu =\lambda ^0\lambda ^1\mathrm{}\lambda ^k=\lambda `$ be a tableau where $`\lambda ^i/\lambda ^{i1}`$ is a horizontal strip of $`i`$’s. Row extraction creates an initially vacant horizontal strip of $``$’s and moves it successively past each horizontal strip $`\lambda ^i/\lambda ^{i1}`$ of $`P`$, possibly getting larger as it goes by the creation of new $``$’s, which replace entries of $`P`$. The result is a tableau $`(P)`$ with inner border $`\mu `$ extended by a horizontal strip of $``$’s whose outer border is $`\lambda `$. The collection of all replaced entries of $`P`$ is a multiset $`r(P)`$.
We describe how to move a horizontal strip of $``$’s past a horizontal strip $`\lambda ^i/\lambda ^{i1}`$ extending it. This proceeds left-to-right through $`\lambda ^i/\lambda ^{i1}`$, moving each entry $`i`$ as follows. The current entry $`i`$ in $`\lambda ^i/\lambda ^{i1}`$ is interchanged with the maximal $``$ smaller than it in the cross order, and if there is no such $``$, then that $`i`$ is removed and replaced by a new $``$. The horizontal strip of $``$’s is initially empty, and remains so until encountering the first nonempty horizontal strip $`\lambda ^i/\lambda ^{i1}`$, at which point all $`i`$’s are logically removed and replaced by $``$’s. We give an example in Figure 12 of an intermediate stage of the algorithm, where two $`i`$’s are removed.
This process is essentially Tesler’s algorithm $`𝒮`$ as described in Section 6.2, up to standard renumbering. The column-strict extension of Tesler’s rectification algorithm proceeds as follows: Given $`P`$, form tableaux $`P_0`$, $`P_1`$, $`\mathrm{}`$, $`P_k`$ and multisets $`r_1,\mathrm{},r_k`$ by $`P_0:=P`$, and for each $`i=1,\mathrm{},k`$,
$$P_i:=(P_{i1})\text{and}r_i:=r(P_{i1}).$$
Then $`P_k=\mathrm{}`$, and an analysis as in Section 6.2 shows that the multisets $`r_1,\mathrm{},r_k`$ form the rows of the rectification of $`P`$. Figure 13 shows this procedure applied to the tableau $`P`$ of Figures 10 and 11. Each row in the figure is one application of $``$, and the multiset of entries removed in each step is displayed to the right.
### 7.3. Column Extraction
The last algorithm of column extraction is strangest of all. Its discovery in February 1992 and a desire to find a common framework with jeu de taquin was the genesis of this paper. Column extraction is constructed from the column sliding algorithm of Stroomer , which we first describe.
Let $`P`$ be a column-strict tableau with entries in $`[k]`$, and let $`b`$ be an inner corner of $`P`$. A column slide beginning at $`b`$ is given by moving an empty box $`b`$ through $`P`$ as follows. Initially $`b`$ switches with the first 1 greater in the cross order in $`P`$. Then the empty box switches with the first 2 in the cross order, and so on, concluding when $`b`$ switches with a $`k`$. It is only permissible to begin a column slide at $`b`$ when there exists a chain of entries in the cross order labelled $`1,2,\mathrm{},k`$. In the figure below, it is not permissible to begin a column slide at the inner corner marked with an $`X`$, but it is permissible to begin a column slide at the inner corner marked with an $``$. We shade the paths of both the impermissible column slide beginning with $`X`$ and the permissible column slide beginning with $``$, and then in the right tableau show the result of the permissible column slide.
###### Lemma 7.1.
The result $`T`$ of a permissible column slide on a tableau $`P`$ is Knuth-equivalent to $`P`$.
Proof. Consider the complement $`P^C`$ writing $`P^C`$ above the tableau $`P`$. Then an inner corner $`b`$ of $`P`$ is an outer cocorner of $`P^C`$ from which we could begin a reverse column insertion . Either an entry of $`P^C`$ will be bumped from $`P^C`$ or this reverse insertion will result in a new box on the inner border of $`P^C`$. Stroomer showed that in the first case, it is not permissible to begin a column slide at $`b`$ on $`P`$, and in the second case, the column slide is permissible, and $`T^C`$ is the tableau obtained from the (internal) column insertion. Thus $`P^C`$ and $`T^C`$ are Knuth-equivalent, and so by Theorem 2.2, $`P`$ and $`T`$ are Knuth-equivalent.
This proof shows that a permissible column slide is complementary to an internal column insertion, and thus may be computed using hopscotch moves. Indeed, the left and right tableaux in Figure 14 are respectively the second and first rows in Figure 7. We illustrate the complementary nature of a permissible column slide and an internal insertion. The (reverse) column insertion in the tableau below on the left beginning with the shaded box results in the tableau on the right below. These tableaux are the complements of the tableaux in Figure 14.
Observe that if the first column of the tableau $`P`$ is full, that is, consists of all the entries in $`[k]`$, then any inner corner of $`P`$ can initiate a permissible column slide. Also note that if $`T`$ is the result of a permissible column slide beginning at an inner corner $`b`$, then we may initiate a permissible column slide at any inner corner of $`T`$ smaller than $`b`$ in the cross order.
###### Definition 7.2.
Let $`P`$ be a tableau with entries in $`[k]`$. Form the tableau $`P^{}`$ by placing a full column of the entries $`1,\mathrm{},k`$ to the left of, and beginning in the same row as, the first column of $`P`$. The cell $`b`$ just above the rightmost column of $`P^{}`$ is an inner corner of $`P^{}`$, and so it is permissible to initiate a column slide in that cell. The cell just above $`b`$ is an inner corner of the resulting tableau, and we may begin another column slide with this cell. Repeating this procedure at most $`k`$ times, we obtain a tableau $`T^{}`$ whose rightmost column is full, and we delete this column to obtain the tableau $`𝒳(P)`$.
By construction, $`P`$ and $`𝒳(P)`$ have the same content. More is true.
###### Lemma 7.3.
The tableaux $`P`$ and $`𝒳(P)`$ are Knuth-equivalent.
We will prove this lemma at the end of this section. We illustrate this operation on the tableau of our running example.
Iterating $`𝒳`$ sufficiently many times rectifies a tableau.
###### Theorem 7.4.
Let $`P`$ be a column-strict tableau. Let $`j`$ be the number of columns of $`P`$ that begin in rows higher than the first column of $`P`$. Then the $`j`$th iterate $`𝒳^j(P)`$ of column extraction applied to $`P`$ is the rectification of $`P`$.
Proof. The number of columns of $`𝒳(P)`$ is at most the number of columns of $`P`$. Also, if the first $`i`$ columns of $`P`$ begin in the same row, then the first $`i+1`$ columns of $`𝒳(P)`$ will begin in that same row, unless $`𝒳(P)`$ has fewer than $`i+1`$ columns. Thus $`𝒳^j(P)`$ has all of its columns beginning in the same row, and so it has partition shape. Since $`P`$ is Knuth-equivalent to $`𝒳^j(P)`$ by Lemma 7.3, we see that $`𝒳^j(P)`$ is the rectification of $`P`$.
Figure 16 illustrates the algorithm of Theorem 7.4, applying $`𝒳`$ twice more to the example of Figure 15.
Our proof of Lemma 7.3 uses Knuth equivalence $``$ of words in the alphabet $`[k]`$ . Given two words $`w,v`$, let $`w.v`$ be their concatenation. Given a column-strict tableau $`P`$, let $`\omega (P)`$ be its (column) word, which is the entries of $`P`$ listed from the bottom to top in each column, starting in the leftmost column and moving right. We use the result of Schützenberger that tableaux $`P`$ and $`Q`$ are Knuth-equivalent if and only if $`\omega (P)\omega (Q)`$. Set $`C:=k\mathrm{}2.1`$, the word of a full column. We first prove a lemma concerning Knuth equivalence and $`C`$, or rather commutation and cancellation in the plactic monoid of Lascoux and Schützenberger .
###### Lemma 7.5.
Let $`w,v`$ be any words in the alphabet $`[k]`$. Then
1. $`C.ww.C`$.
2. $`C.wC.v`$ if and only if $`wv`$.
Proof. For (i), it suffices to consider the case when $`w`$ is a single number $`i`$. Then $`i.C`$ is the column word of a tableau of skew shape $`2^k/1^{k1}`$ whose rectification has partition shape $`(2,1^{k1})`$ and word $`C.i`$.
The reverse direction of (ii) is trivial. Suppose $`C.wC.v`$ and let $`P,Q`$ be tableaux of partition shape with $`\omega (P)w`$ and $`\omega (Q)v`$. Then $`C.\omega (P)C.\omega (Q)`$. But for any tableau $`T`$, $`C.\omega (T)`$ is the word of a tableau $`T^{}`$ of partition shape obtained from $`T`$ by placing a full column to its left. Since there is a unique tableau of partition shape in any Knuth equivalence class, we must then have $`C.\omega (P)=C.\omega (Q)`$, thus $`P^{}=Q^{}`$, and so $`P=Q`$. But this implies $`wv`$.
Proof of Lemma 7.3. Let $`P,P^{},T^{}`$ be as in Definition 7.2 and set $`T=𝒳(P)`$. Then
$$C.\omega (P)=\omega (P^{})\omega (T^{})=\omega (T).CC.\omega (T).$$
The first Knuth-equivalence follows from Lemma 7.1, as $`T^{}`$ is obtained from $`P^{}`$ by column slides. The second Knuth-equivalence is Statement (i) of Lemma 7.5. Since $`C.\omega (P)C.\omega (T)`$, we deduce that $`\omega (P)\omega (T)`$, by (ii) of Lemma 7.5, and so $`P`$ is Knuth equivalent to $`𝒳(P)`$. |
warning/0002/quant-ph0002030.html | ar5iv | text | # 1 Introduction
## 1 Introduction
The aim of this paper is to make an introduction to Fine’s interpretation of quantum mechanics and to show how it can solve the EPR–Bell problem. In the real spin correlation experiments the measured quantum probabilities are identified with relative frequencies taken on a selected sub-ensemble of the emitted particle pairs: only those particles are taken into account which are coincidentally detected in the two wings. The detection/emission inefficiency is usually ascribed to random detection errors occurring independently in the two wings, and treated by the “enhancement hypothesis”, which assumes that the relative frequencies measured on the randomly selected sub-ensemble are equal to the ones taken on the whole statistical ensemble of emitted particle pairs.
Fine’s *prism model*<sup>(1)</sup> is a local hidden variable theory, in which the detection inefficiency is an effect not only of the random errors in the analyzer + detector equipment, but is also the manifestation of a predetermined hidden property of the particles.<sup>2</sup><sup>2</sup>2 This conception of hidden variable goes back to Einstein (Ref. 4, Chapter 4). I present one of Fine’s prism models for the EPR experiment and compare it with the recent experimental results.<sup>(2)</sup> As we shall see, it works well in case of the $`2\times 2`$ spin-correlation experiments. There appeared, however, a theoretical demand to embed the $`2\times 2`$ prism models into a large $`n\times n`$ prism model reproducing all potential $`2\times 2`$ sub-experiments. This demand was motivated by the idea that the real physical process does not know which directions are chosen in an experiment. On the other hand, it seemed that in the known prism models of the $`n\times n`$ spin-correlation experiment the efficiencies tended to zero, if $`n\mathrm{}`$, which contradicts what we expect of actual experiments.<sup>(3)</sup> This problem is investigated in the last part of the paper.
## 2 Real EPR experiments
Figure 1 shows an Aspect-type spin-correlation experiment. The analyzers can be set in orientation $`a`$ or $`a^{}`$ on the left hand side, and $`b`$ or $`b^{}`$ on the right. Denote $`A`$, $`A^{}`$, $`B`$ and $`B^{}`$ the corresponding “spin-up” detection-events. $`p(A|a)`$, for instance, denotes the conditional probability of $`A`$, given that the measurement set-up in the left wing is $`a`$.
Now, from the assumption that there exists a hidden variable $`\lambda `$ satisfying the screening off condition,
$$p\left(AB|ab\lambda \right)=p\left(A|a\lambda \right)p\left(B|b\lambda \right)$$
one can derive<sup>(4)</sup> the well known Clauser-Horne inequality:
$$1\underset{CH}{\underset{}{\left\{\begin{array}{c}p\left(AB|ab\right)+p\left(AB^{}|ab^{}\right)p\left(A^{}B|a^{}b\right)\\ +p\left(A^{}B^{}|a^{}b^{}\right)p\left(A^{}|a^{}\right)p\left(B|b\right)\end{array}\right\}}}0$$
(1)
According to the standard views, inequality (1) is violated in the real spin-correlation experiments, hence, the argument goes, any local hidden variable theory for the EPR experiment is excluded. For example, in case of spin-$`\frac{1}{2}`$ particles, if the directions $`a`$, $`a^{}`$, $`b`$, $`b^{}`$ are coplanar with $`\mathrm{}(a,b)=\mathrm{}(a,b^{})=\mathrm{}(a^{},b^{})=120^{}`$ and $`\mathrm{}(a^{},b)=0`$, then the Clauser-Horne expression in (1) is
$$CH=\frac{3}{8}+\frac{3}{8}0+\frac{3}{8}\frac{1}{2}\frac{1}{2}=\frac{3}{8}$$
(2)
There is, however a serious loophole in the real experiments. Compare the original apparatus configuration used for Bell’s 1971 proof with the one used in the real Aspect experiment (Fig. 1 and 2). The original configuration contains two ‘event-ready’ detectors, which signal both arms that a pair of particles has been emitted. So, the statistics are taken on the ensemble of particle pairs emitted by the source. In the real experiments, however, instead of the event-ready detectors, a four-coincidence circuit detects the ‘emitted particle-pairs’. This method yields to a *selected* statistical ensemble: only those pairs are taken into account, which coincidentally fire one of the left and one of the right detectors. Denote $`\left[A\right]`$ the event that there is any detection in the left wing with analyzer set-up $`a`$, that is, either the “up” detector or the “down” detector fires. Similarly, $`\left[A\right]\left[B\right]`$ denotes the corresponding double detection. So, what we actually observe is the violation of the following inequality:
$$1\left\{\begin{array}{c}p\left(AB|ab\left[A\right]\left[B\right]\right)+p\left(AB^{}|ab^{}\left[A\right]\left[B^{}\right]\right)\\ p\left(A^{}B|a^{}b\left[A^{}\right]\left[B\right]\right)+p\left(A^{}B^{}|a^{}b^{}\left[A^{}\right]\left[B^{}\right]\right)\\ p\left(A^{}|a^{}\left[A^{}\right]\right)p\left(B|b\left[B\right]\right)\end{array}\right\}0$$
(3)
If the selection procedure were *completely random* then the observed relative frequencies on the selected ensemble would be equal to the ones taken on the original ensemble, that is,
$`p\left(AB|ab\left[A\right]\left[B\right]\right)`$ $`=`$ $`p\left(AB|ab\right)`$
$`p\left(AB^{}|ab^{}\left[A\right]\left[B^{}\right]\right)`$ $`=`$ $`p\left(AB^{}|ab^{}\right)`$
$`\text{etc}.`$
(*enhancement* *hypothesis*) and the violation of inequality (3) would imply the violation of (1), in accordance with Bell’s point of view:
> … it is hard for me to believe that quantum mechanics works so nicely for inefficient practical set-ups and is yet going to fail badly when sufficient refinements are made. (Ref. 5, p. 154)
This is indeed the case if non-detections are caused by independent random errors in the detector+analyser equipment.
## 3 Fine’s interpretation
Arthur Fine approaches the detection inefficiency problem in a different way:
> … the efficiency problem ought not to be dismissed as merely one of biased statistics and conspiracies, for the issue it raises is fundamental. Can a hidden variable theory of the very type being tested explain the statistical distributions, inefficiencies and all, actually found in the experiments? If so then we would have a model (or theory) of the experiment that explains why the samples counted yield the particular statistics that they do. (Ref. 6, p. 465)
This conception of hidden variable was first realized in Fine’s *prism models*<sup>(1)</sup> for the $`2\times 2`$ spin-correlation experiment. Prism model is a local, deterministic hidden variable theory, in which the hidden variables predetermine not only the outcomes of the corresponding measurements, but also predetermine whether or not an emitted particle arrives to the detector and becomes detected. In other words, the measured observables can take on a new “value” corresponding to an inherent “no show” or defectiveness.
As an example, consider a prism model reproducing the quantum mechanical probabilities in (2). Figure 3 shows a parameter space $`\mathrm{\Lambda }\lambda `$ consisting of disjoint blocks of measure $`\frac{3}{32}`$ and $`\frac{1}{32}`$ respectively. A point of $`\mathrm{\Lambda }`$ (a value of the parameter $`\lambda `$) predetermines all events in question. Therefore, each EPR event can be represented as a subset of $`\mathrm{\Lambda }`$. For instance, assume that $`\lambda =\lambda _{\mathrm{example}}`$. Then, an $`a`$-measurement on the left particle produces neither event “up” nor event “down”, while if an $`a^{}`$-measurement is performed then the outcome is “down”. In the right wing, if we perform a $`b`$-measurement then the outcome is “up”, and if the $`b^{}`$-measurement is performed, the outcome is “down”. Consequently, in case, for example, we perform an $`a`$-measurement on the left particle and a $`b`$-measurement on the right one, then there is no coincidence registered, and the particle pair in question does not appear in the statistics of the measurement. On the contrary, if we perform an $`a^{}`$-measurement on the left particle and a $`b`$-measurement on the right one, then there is a coincidence registered and the counter of the total number of events as well as the $`B`$-counter count. Thus, the hidden parameter governs the whole process in such a way that the observed relative frequencies reproduce the probabilities measured in the experiment:
$`{\displaystyle \frac{\mu \left(A\right)}{\mu \left([A]\right)}}`$ $`=`$ $`{\displaystyle \frac{\mu \left(A^{}\right)}{\mu \left([A^{}]\right)}}={\displaystyle \frac{\mu \left(B\right)}{\mu \left([B]\right)}}={\displaystyle \frac{\mu \left(B^{}\right)}{\mu \left([B^{}]\right)}}={\displaystyle \frac{\frac{12}{32}}{\frac{24}{32}}}={\displaystyle \frac{1}{2}}`$
$`{\displaystyle \frac{\mu \left(AB\right)}{\mu \left([A][B]\right)}}`$ $`=`$ $`{\displaystyle \frac{\mu \left(AB^{}\right)}{\mu \left([A][B^{}]\right)}}={\displaystyle \frac{\mu \left(A^{}B^{}\right)}{\mu \left([A^{}][B^{}]\right)}}={\displaystyle \frac{\frac{6}{32}}{\frac{16}{32}}}={\displaystyle \frac{3}{8}}`$
$`{\displaystyle \frac{\mu \left(A^{}B\right)}{\mu \left([A^{}][B]\right)}}`$ $`=`$ $`{\displaystyle \frac{0}{\frac{16}{32}}}=0`$
## 4 Compatibility with the actual EPR experiments
As we have seen, the basic idea of the Einstein–Fine interpretation is that some elements of the statistical ensemble of identically prepared quantum systems (characterized by a quantum state $`W`$) do not produce outcome at all when one perform the measurement of a quantum observable $`A`$. Such systems are called $`A`$*-defective* in Fine’s terminology. In connection with this basic feature of the model, one can investigate some important characteristics of the above Einstein–Fine model of the EPR experiment, and compare them with the similar characteristics of the actual EPR experiments:
$`R^A`$ $`=`$ $`{\displaystyle \frac{\text{number of non}A\text{defective systems}}{\text{total number of systems}}}`$
$`R^A^{}`$ $`=`$ $`{\displaystyle \frac{\text{number of non}A^{}\text{defective systems}}{\text{total number of systems}}}`$
$`R^B`$ $`=`$ $`{\displaystyle \frac{\text{number of non}B\text{defective systems}}{\text{total number of systems}}}`$
$`R^B^{}`$ $`=`$ $`{\displaystyle \frac{\text{number of non}B^{}\text{defective systems}}{\text{total number of systems}}}`$
$`R^{AB}`$ $`=`$ $`{\displaystyle \frac{\text{number of non}A\text{defective }\&\text{ non}B\text{defective systems}}{\text{total number of systems}}}`$
$`R^{AB^{}}`$ $`=`$ $`{\displaystyle \frac{\text{number of non}A\text{defective }\&\text{ non}B^{}\text{defective systems}}{\text{total number of systems}}}`$
$`R^{A^{}B}`$ $`=`$ $`{\displaystyle \frac{\text{number of non}A^{}\text{defective }\&\text{ non}B\text{defective systems}}{\text{total number of systems}}}`$
$`R^{A^{}B^{}}`$ $`=`$ $`{\displaystyle \frac{\text{number of non}A^{}\text{defective }\&\text{ non}B^{}\text{defective systems}}{\text{total number of systems}}}`$
In case of the above example:
$$\begin{array}{ccc}\hfill R^A=R^A^{}=R^B=R^B^{}& =& 75\%\hfill \\ \hfill R^{AB}=R^{AB^{}}=R^{A^{}B}=R^{A^{}B^{}}& =& 50\%\hfill \end{array}$$
(4)
If any of the similar rates in a real experiment were higher than the corresponding one in (4), Fine’s interpretation would be *experimentally* refuted.
There are principal obstacles to an event ready detection, therefore we cannot have a precise information about the total number of systems. In one of the best experiments of the last years<sup>(2)</sup> the *estimated* rates are the following<sup>3</sup><sup>3</sup>3 I would like to thank G. Weihs and A. Zeilinger for the private communications about many interesting details of the experiment.:
$$\begin{array}{ccc}\hfill R^A=R^A^{}=R^B=R^B^{}& =& 5\%\hfill \\ \hfill R^{AB}=R^{AB^{}}=R^{A^{}B}=R^{A^{}B^{}}& =& 0.25\%\hfill \end{array}$$
(5)
So the prism model is in accordance with this experiment.
It can be (and probably is) the case that this very low detection/emission rate is mostly caused by the external random detection errors, different from the prism mechanism. In order to separate these two sources of inefficiency, consider a new characteristic of the model:
$$r_A^{AB}=\frac{\text{number of non}A\text{defective }\&\text{ non}B\text{defective systems}}{\text{number of non}A\text{defective systems}}$$
$$r_A^{AB^{}}=\frac{\text{number of non}A\text{defective }\&\text{ non}B^{}\text{defective systems}}{\text{number of non}A\text{defective systems}}$$
$$\mathrm{}$$
(6)
$$r_B^{}^{A^{}B^{}}=\frac{\text{number of non}A^{}\text{defective }\&\text{ non}B^{}\text{defective systems}}{\text{number of non}B^{}\text{defective systems}}$$
In our prism model:
$$r_A^{AB}=r_A^{AB^{}}=r_A^{}^{A^{}B}=r_A^{}^{A^{}B^{}}=r_B^{AB}=r_B^{}^{AB^{}}=r_B^{A^{}B}=r_B^{}^{A^{}B^{}}=66,66\%$$
The experiment by Weihs *et al*.<sup>(2)</sup> had a particular new feature: In the two wings independent data registration was performed by each observer having his own atomic clock, synchronized only once before each experiment cycle. A time tag was stored for each detected photon in two separate computers at the observer stations and the stored data were analyzed for coincidences long after measurements were finished. Due to this method of data registration, it was possible to count the rates in (6). Again, if any of these rates were higher than 66,66%, Fine’s interpretation wouldn’t be tenable. However the experimental values were only around 5%.
## 5 The prism model of an $`n\times n`$ spin correlation experiment
Let us turn now to a serious objection to Fine’s approach. In the EPR experiment we consider only $`2\times 2`$ different possible directions ($`\stackrel{}{a},\stackrel{}{a}^{},\stackrel{}{b},\stackrel{}{b}^{}`$). If nature works according to Fine’s prism model then there must exist, in principle, a larger $`n\times n`$ prism model reproducing all potential $`2\times 2`$ sub-experiments. It is because nature does not know about how the experiment is designed. So, in the final analysis, there is no such a thing as $`2\times 2`$ prism model of the $`2\times 2`$ spin-correlation experiment. If we want to describe a $`2\times 2`$ spin-correlation experiment with a prism model, then there must exist a large $`n\times n`$ (if not $`\mathrm{}\times \mathrm{})`$ prism model behind it.
The general schema of the prism model of a spin-correlation experiment is the following. In both wings one considers $`n`$ different possible events:
$$\stackrel{\mathrm{left}}{\stackrel{}{\underset{\left[A_1\right]=\left[A_2\right]}{\underset{}{A_1,A_2}},\underset{\left[A_3\right]=\left[A_4\right]}{\underset{}{A_3,A_4}},\mathrm{}\underset{\left[A_{n1}\right]=\left[A_n\right]}{\underset{}{A_{n1},A_n}}}}\stackrel{\mathrm{right}}{\stackrel{}{\underset{\left[B_1\right]=\left[B_2\right]}{\underset{}{B_1,B_2}},\underset{\left[B_3\right]=\left[B_4\right]}{\underset{}{B_3,B_4}},\mathrm{}\underset{\left[B_{n1}\right]=\left[B_n\right]}{\underset{}{B_{n1},B_n}}}}$$
(7)
$`A_1`$ denotes the event that the left particle has spin “up” along direction $`\stackrel{}{a}_1`$. $`A_2`$ denotes the event that the left particle has spin “down” along direction $`\stackrel{}{a}_1`$. Similarly, $`A_3`$ denotes the event that the left particle has spin “up” along direction $`\stackrel{}{a}_3`$ and $`A_4`$ denotes the event that the left particle has spin “down” along direction $`\stackrel{}{a}_3`$, etc. (We will also use the following notation: $`\stackrel{}{a}_2=\stackrel{}{a}_1,\stackrel{}{a}_4=\stackrel{}{a}_3,\mathrm{}\stackrel{}{a}_{2k}=\stackrel{}{a}_{2k1},\mathrm{}\stackrel{}{b}_{2k}=\stackrel{}{b}_{2k1}`$.) There are $`\frac{n}{2}`$ different directions on both sides. We also assume the following logical relationships:
$$\begin{array}{ccccccc}\hfill \left[A_1\right]=\left[A_2\right]& =& A_1A_2\hfill & & \hfill \left[B_1\right]=\left[B_2\right]& =& B_1B_2\hfill \\ \hfill \left[A_3\right]=\left[A_4\right]& =& A_3A_4\hfill & & \hfill \left[B_3\right]=\left[B_4\right]& =& B_3B_4\hfill \\ & \mathrm{}& & & & \mathrm{}& \\ \hfill \left[A_{n1}\right]=\left[A_n\right]& =& A_{n1}A_n\hfill & & \hfill \left[B_{n1}\right]=\left[B_n\right]& =& B_{n1}B_n\hfill \\ \hfill A_1A_2& =& 0\hfill & & \hfill B_1B_2& =& 0\hfill \\ \hfill A_3A_4& =& 0\hfill & & \hfill B_3B_4& =& 0\hfill \\ & \mathrm{}& & & & \mathrm{}& \\ \hfill A_{n1}A_n& =& 0\hfill & & \hfill B_{n1}B_n& =& 0\hfill \end{array}$$
(8)
that is, $`\left[A_1\right]`$ (which is equal to $`\left[A_2\right]`$) denotes the event that the left particle is predetermined to produce any outcome if the $`\stackrel{}{a}_1`$ direction is measured. The quantum probabilities are reproduced in the following way:
$`tr\left(WP_{A_i}\right)=q_i`$ $`=`$ $`{\displaystyle \frac{p\left(A_i\right)}{p\left(\left[A_i\right]\right)}}`$
$`tr\left(WP_{B_i}\right)=q_i^{}`$ $`=`$ $`{\displaystyle \frac{p\left(B_i\right)}{p\left(\left[B_i\right]\right)}}`$ (9)
$`tr\left(WP_{A_i}P_{B_j}\right)=q_{ij}`$ $`=`$ $`{\displaystyle \frac{p\left(A_iB_j\right)}{p\left(\left[A_i\right]\left[B_j\right]\right)}}`$
The quantum probabilities $`q_1,q_2,\mathrm{}q_1^{},q_2^{},\mathrm{}q_{ij},\mathrm{}`$ are the only fix numbers in the model.
The experimental setup shows the following simple and natural symmetries:
* None of the left and right wings is privileged.
* There is no privileged direction among the possible polariser positions.
Consequently, all physically plausible prism model have to satisfy these symmetry conditions, which imply the following two constraints:
$$p\left(\left[A_i\right]\right)=\omega \text{for all }1in$$
(10)
where $`\omega `$ is some uniform efficiency for all directions on both sides, and
$$p\left(\left[A_i\right]\left[B_j\right]\right)=\sigma \left(\mathrm{}(\stackrel{}{a}_i,\stackrel{}{b}_j)\right)\text{for all }1i,jn$$
(11)
where $`0\sigma \left(\mathrm{}(\stackrel{}{a}_i,\stackrel{}{b}_j)\right)\omega `$ is an arbitrary function of the angel $`\mathrm{}(\stackrel{}{a}_i,\stackrel{}{b}_j)`$.
Thus, the prism-model resolution of the EPR-Bell problem requires the existence of $`n\times n`$ prism models of the above type. On the other hand, this requirement appears to be a serious objection to Fine’s program. The reason is that in all the known $`n\times n`$ prism models the efficiency $`\omega `$ tends to zero if $`n\mathrm{}`$, which contradicts the recent experimental results.<sup>(7-8)</sup> Moreover, Fine has shown<sup>(3)</sup> that this is true for the class of $`n\times n`$ prism models satisfying certain symmetry conditions called *Exchangeability, Indifference* and *Strong Symmetry*. They are complex conditions, too complex to briefly recall the definitions. Although they do not express some natural and obvious symmetries of the experimental setup, they are instanced in all the known prism models.
If all physically plausible prism models had to satisfy the Exchangeability, Indifference and Strong Symmetry conditions, the problem of zero efficiency would mean a serious objection to Fine’s interpretation. Fortunately, this is not the case. In my http://arXiv.org/abs/quant-ph/0012042 I shown the existence of $`\mathrm{}\times \mathrm{}`$ prism models satisfying the symmetry conditions (S1) and (S2), whereas *the efficiencies are reasonably high.*
## Acknowledgments
The author wish to thank Professor A. Fine for his stimulating comments and suggestions. The research was supported by the OTKA Foundation, No. T015606 and T032771.
## References
1. A. Fine, “Some local models for correlation experiments”, *Synthese*, 50, 279 (1982)
2. G. Weihs, T. Jennewin, C. Simon, H. Weinfurter, and A. Zeilinger, “Violation of Bell’s Inequality under Strict Einstein Locality Conditions”, *Phys. Rev. Lett.*, 81, 5039 (1998)
3. A. Fine, “Inequalities for Nonideal Correlation Experiments”, *Foundations of Physics*, 21, 365 (1991)
4. B. C. van Frassen, “The Charybdis of Realism: Epistemological Implications of Bell’s Inequality”, in *Philosophical Consequences of Quantum Theory – Reflections on Bell’s Theorem*, J. T. Cushing and E. McMullin (eds.), University of Notre Dame Press, Notre Dame, (1989)
5. J. S. Bell, *Speakable and unspeakable in quantum mechanics*, Cambridge University Press, Cambridge, (1987)
6. A. Fine, “Correlations and Efficiency: Testing Bell Inequalities”, *Foundations of Physics*, 19, 453 (1989)
7. W. D. Sharp and N. Shank, , “Fine’s prism models for quantum correlation statistics”, *Philosophy of Science*, 52, 538 (1985)
8. T. Maudlin, *Quantum Non-Locality and Relativity**Metaphysical Intimations of Modern Physics*, Aristotelian Society Series, Vol. 13, Blackwell, Oxford, (1994) |
warning/0002/cond-mat0002442.html | ar5iv | text | # Relationship between spiral and ferromagnetic states in the Hubbard model in the thermodynamic limit
## I introduction
The itinerant ferromagnetism has been one of the most fundamental problems in condensed matter physics, and we are still some way from a full understanding of the problem. For the ferromagnetism in the Hubbard model, a simplest possible model for interacting itinerant electrons, there are various studies, which are pioneered by Gutzwiller, Hubbard, and Kanamori. For the Hubbard model on finite-size systems, it recently extended to rigorous proofs that spin-independent Coulomb interaction can indeed result in fully-polarized ferromagnetic ground states without degeneracy for appropriate conditions and/or appropriate models. Apart from these, there is a body of numerical studies for finite systems, where various authors have shown that finite Hubbard models with appropriate band fillings and boundary conditions have fully-polarized ferromagnetic ground states.
On the other hand, for infinite systems, the itinerant nature of electrons makes the problem interesting and subtle. The most important reason is the existence of the spiral-spin(SP) state, a spin-singlet state that has the spin-spin correlation length as large as the sample size, which is known to accompany the fully-polarized ferromagnetic(F) state in various finite-size itinerant models. By accompanying we mean SP and F states are close in energy in finite systems, where the ground state changes between them when the boundary condition is twisted.
The SP state does not seem to be an accident of the Hubbard model, since SP states are found not only in electron systems with strong charge-charge interactions, but also in those with spin-dependent interactions such as double-exchange or Kondo lattice models. So we believe that it is a general feature of strongly correlated systems for the F state to be accompanied by the SP state. It is rather surprising, since the latter has the total $`S=0`$ while the former has $`S=S_{\mathrm{max}}`$. Thus it is desirable to resolve this question to understand the ferromagnetism in itinerant systems in the thermodynamic limit. Thus the problem we have to clarify is the relation between the SP and F states in the thermodynamic limit (do those states degenerate or coalesce, etc), which we address in this paper.
First, we show that the SP state is eligible as the ground state (i.e., degenerate with the F state in energy in the thermodynamic limit) when the F state is the ground state. We check this for one-dimensional (1D) $`t`$-$`t^{}`$ Hubbard model with infinitely large coulomb interaction taken as a typical example and calculate the energies of the SP and F states as a function of the inverse system size with the density matrix renormalization group (DMRG) method.
Next, we discuss whether the symmetry of the SP state is broken in the thermodynamic limit, following the arguments by Koma and Tasaki on symmetry breaking and finite-size effects in quantum many-body systems on a lattice. We show that (i) the SP state possesses a kind of order parameter, (ii) although the state does not break the global SU(2) symmetry in the finite system, it does so in the thermodynamic limit by making a linear combination of other states that are degenerate in that limit.
We have studied excitation spectrum of the SP state by calculating one-particle spectral function, dynamical spin and charge susceptibilities for the SP state in various finite-size models, including 1D $`t`$-$`t^{}`$ Hubbard model, Tasaki’s flat band model, and Hubbard ladder. We have found that the excitation spectrum of these two states in the finite size system are almost identical.
The present results suggest that we should regard the SP state as identical with the F state although they have the total spins in the opposite extremes ($`S=0`$ for the SP state, $`S=S_{\mathrm{max}}`$ for F) in the finite size system. These properties can be exploited to determine the magnetic phase diagram from the calculation of the total spin of the finite-size system in an unambiguous fashion. Otherwise we have to worry about not only the total spin, but the spin-spin correlation to identify SP states in determining the phase diagram.
The paper is organized as follows. In Sec. II, we first confirm that the SP and F states have degenerate energy for $`L\mathrm{}(L`$: sample size) in that the energy per site is identical between them. DMRG is used to calculate the difference in the total energy, $`\mathrm{\Delta }EE_{\mathrm{SP}}E_\mathrm{F}`$, between the SP state and the F state as a function of system size $`L`$ for the $`t`$-$`t^{}`$ Hubbard model with $`U=\mathrm{}`$ as a typical example. The extrapolate $`\mathrm{\Delta }E(L\mathrm{})`$ is always finite (or zero depending on the boundary condition even in the thermodynamic limit), i.e., $`\mathrm{\Delta }E`$ per site does indeed vanish. It implies that when the F state is the ground state, the SP state is also eligible as the ground state if these states are translation invariant. Then we move on to the discussion on symmetry breaking of the SP state along the Koma-Tasaki’s argument. In Sec. III, we calculate the single-particle spectral function and dynamical spin and charge susceptibilities for various models. Our results suggest that the SP state and the F state are identical state in the thermodynamic limit.
In Sec. IV, we take an appropriate combination of the boundary condition and the number of electrons to make the SP state higher in energy than the SF state, and the phase diagram for the one-dimensional $`t`$-$`t^{}`$ Hubbard model, Tasaki model, and 2-leg Hubbard ladder model is obtained with exact diagonalization of small systems. We have found that the phase boundaries rapidly converge to those for larger systems obtained with DMRG. We may expect that it can be exploited to the determination of the phase diagram of 2D or higher dimensional systems. A summary of the present study is given in Sec. VI.
## II Discussion on the thermodynamic limit
### A The energy difference between the SP state and the F state in the thermodynamic limit
1D $`t`$-$`t^{}`$ Hubbard model is one of the simplest models that is thought to exhibit ferromagnetism for sufficiently strong electron-electron interaction. Müller-Hartmann suggested that in the low-density limit the ground state of $`t`$-$`t^{}`$ Hubbard model is ferromagnetic. Daul and Noack have carried out highly accurate DMRG calculation to conclude that there is an extensive ferromagnetic phase in the phase diagram. Since the $`t`$-$`t^{}`$ Hubbard model is numerically tractable, we study the energy difference of the SP and F states by taking this model with $`U=\mathrm{}`$ as an example.
The $`t`$-$`t^{}`$ Hubbard Hamiltonian is given by
$``$ $`=`$ $`t{\displaystyle \underset{i=1}{\overset{L}{}}}{\displaystyle \underset{\sigma }{}}(c_{i\sigma }^{}c_{i+1,\sigma }+\mathrm{H}.\mathrm{c}.)`$
$`+t^{}{\displaystyle \underset{i=1}{\overset{L}{}}}{\displaystyle \underset{\sigma }{}}(c_{i\sigma }^{}c_{i+2,\sigma }+\mathrm{H}.\mathrm{c}.)+U{\displaystyle \underset{i=1}{\overset{L}{}}}n_in_i,`$
where $`c_{i\sigma }^{}`$ creates an electron at site $`i`$ with spin $`\sigma (=,)`$, $`t`$ is the nearest-neighbor hopping, and $`t^{}`$ the next nearest-neighbor hopping, $`U`$ the Hubbard repulsion, and $`n_{i\sigma }c_{i\sigma }^{}c_{i\sigma }`$. In this section we take $`t=1`$, $`t^{}=0.2`$ and $`U=\mathrm{}`$.
As we can check by a simple exact diagonalization, the ground state of the system has always $`S_{\mathrm{tot}}=0`$ for any value of $`U`$ at least for system sizes up to 12 sites when a periodic boundary condition(PBC) is adopted. This is to be contrasted with the phase diagram obtained by Daul and Noack with an open boundary condition(OBC), in which a wide ferromagnetic region is found for $`t^{}=0.2`$ and large $`U`$.
Let us first confirm that the SP and F states are degenerate in the sense that their difference in the total energy behaves as
$`\mathrm{\Delta }EE_{\mathrm{singlet}}E_\mathrm{F}\mathrm{finite},`$
i.e., the energy per site vanishes,
$`\mathrm{\Delta }E/L0,`$
where $`E_{\mathrm{singlet}}`$ is the lowest energy within the $`S_{\mathrm{tot}}=0`$ sector, namely the energy of the SP state.
We take OBC, because the DMRG becomes most accurate for this condition. As Daul and Noack showed, the ground state of $`t`$-$`t^{}`$ Hubbard model is ferromagnetic for $`t^{}=0.2`$ and $`U\mathrm{}`$ in OBC at least up to 50 sites. To estimate $`\mathrm{\Delta }E`$, we must calculate the energy of the ground state of the $`S_{\mathrm{tot}}=0`$ sector which has higher energy than the ferromagnetic ground state. Therefore, we add the term $`\lambda 𝐒_{\mathrm{tot}}^2`$ to the original Hamiltonian (1) to selectively shift the states with higher total spin to higher energies by turning on $`\lambda >0`$, while conserving the SU(2) symmetry. Here we set $`\lambda =1`$. In Fig. 1, we show the spin correlation function, $`S_i^zS_j^z`$ ($`i=L/2,j=L/2+1,\mathrm{},L3`$) of the spin singlet state for $`L=32`$ sites, $`n=0.75`$ as an example. We can see that the spin correlation wave length is as large as the system size. Namely, the spin singlet state is the SP state.
Since DMRG is a variational procedure, the energy of the SP state is an upper bound, while the energy of the F state is calculated exactly because the F state does not feel the on-site coulomb interaction. Hence we overestimate $`\mathrm{\Delta }E`$. To minimize this overestimation, we must calculate the energy of the SP state as accurately as possible. This is the reason why we have set $`U=\mathrm{}`$ and $`t^{}=0.2`$. Namely, the exclusion of double occupancies reduces the Hilbert space drastically, while a small value of $`|t^{}|`$ reduces the truncation error.
Using the finite-size algorithm in DMRG, we calculated up to $`L=48`$ for the density of electrons $`n=0.5`$, sweeping the system about 20 times to improve the wave function. We store the density matrix at each step to construct good initial vector for each super-block diagonalization. We have kept up to $`m=500`$ states per block at each step, where the convergence is checked by comparing the results for $`m=300,400,500`$. The truncation error is smaller than $`10^5`$, which is small enough to enable us to extrapolate $`\mathrm{\Delta }E(L\mathrm{})`$.
In Fig. 2(solid line), we show the results for $`\mathrm{\Delta }EE_{\mathrm{SP}}E_\mathrm{F}`$ as a function of inverse system size $`1/L`$ for electron density $`n=0.5`$. We can see that all the points for $`L>24`$ fall upon a linear dependence on $`1/L`$, from which we can extract $`\mathrm{\Delta }E(L\mathrm{})`$. For all the densities studied, $`\mathrm{\Delta }E(L\mathrm{})`$ indeed remains finite ($`10^3`$).
In order to check whether the result is not an accident for $`U=\mathrm{}`$, we can introduce the effect of large but finite $`U`$ as an (antiferromagnetic) exchange interaction $`J>0`$,
$$J\underset{i=1}{}P_G^1\left(𝐒_𝐢𝐒_{𝐢+\mathrm{𝟏}}\frac{1}{4}n_in_{i+1}\right)P_G,$$
(1)
added to the $`U=\mathrm{}`$ $`t`$-$`t^{}`$ Hubbard model. Here $`S_i\frac{1}{2}_{\beta \gamma }c_{i\beta }^{}\stackrel{}{\sigma }_{\beta \gamma }c_{i\gamma }`$ is the spin operator, where $`\stackrel{}{\sigma }`$ is the Pauli matrices and $`P_G`$ denotes the Gutzwiller projection operator. In Fig. 2 the results for $`J=\pm 0.2`$ are superposed. We can see that $`\mathrm{\Delta }E`$ does not change drastically even if we introduce the effect of finite $`U`$.
This kind of $`\mathrm{\Delta }E`$ may depend on the boundary condition, which can be indeed the case with finite systems. For instance, Kusakabe and Aoki have shown for a finite two-dimensional Hubbard model that Nagaoka’s ferromagnetic state (one hole in the half-filled band with $`U=\mathrm{}`$) alternate with an SP state as the boundary condition is changed from periodic to anti-periodic. In other words a level crossing takes place between the two states as an Aharonov-Bohm magnetic flux is introduced to twist the boundary condition. Such a boundary effect can affect the ‘energy gap’ even in the thermodynamic limit in general.
Thus we compare in Fig. 3 the ordinary OBC (the same as Fig. 2, solid line) and the boundary condition in which we further turn off the nearest-neighbor transfers ($`t`$) at either end of the system in the $`t`$-$`t^{}`$ model for $`\mathrm{\Delta }E(L\mathrm{})`$ at $`n=0.5`$. We can see that $`\mathrm{\Delta }E(L\mathrm{})`$ becomes zero (negligibly small) in the latter case. Therefore, we can not exclude a possibility for which the SP state has a lower energy than the F state in some appropriate boundary condition. In other words, the problem which state has a lower energy is a very subtle problem.
### B The definition of the ground state
As we have seen, it is meaningless in the thermodynamic limit to identify the ground state by studying small ‘finite’ differences in the total energy between the candidates, so that a totally different point of view is required. The mathematical definition of the ground state in the thermodynamic limit is reviewed in Koma and Tasaki’s article and here we follow them. For simplicity, we focus on the case of PBC.
We first recapitulate the definition. Let $`A`$ be an arbitrary local operator that acts on a finite number of sites. We define $`\rho `$ as
$`\rho (A)\underset{\mathrm{\Lambda }Z^d}{lim}\rho _\mathrm{\Lambda }(A)`$
for each $`A`$, and we call $`\rho `$ as ‘state’ in the thermodynamic limit. Here $`\mathrm{\Lambda }`$ is a finite ($`L\times L\times \mathrm{}L\times L`$) $`d`$-dimensional hypercubic lattice while $`Z^d`$ is an infinite $`d`$-dimensional lattice, and we have defined
$`\rho _\mathrm{\Lambda }(\mathrm{})\mathrm{Tr}_{\mathrm{H}_\mathrm{\Lambda }}[(\mathrm{})\stackrel{~}{\rho }_\mathrm{\Lambda }],`$
where the trace is taken over the Hilbert space $`\mathrm{H}_\mathrm{\Lambda }`$ spanned over $`\mathrm{\Lambda }`$, and $`\stackrel{~}{\rho }_\mathrm{\Lambda }`$ is an arbitrary density matrix on $`\mathrm{H}_\mathrm{\Lambda }`$. We also assume a Hamiltonian $`_\mathrm{\Lambda }`$ that is local as defined by
$`_\mathrm{\Lambda }={\displaystyle \underset{x\mathrm{\Lambda }}{}}h_x.`$
where $`h_x`$ is a local component such as $`h_x=_\sigma (c_{x\sigma }^{}c_{x+1,\sigma }+\mathrm{H}.\mathrm{c}.)+Un_xn_x`$ for the Hubbard model. We then define the ground-state energy density $`\epsilon _0`$ as
$`\epsilon _0\underset{\mathrm{\Lambda }Z^d}{lim}\underset{\mathrm{\Phi }_\mathrm{\Lambda }\mathrm{H}_\mathrm{\Lambda },\mathrm{\Phi }=1}{inf}{\displaystyle \frac{1}{|\mathrm{\Lambda }|}}\mathrm{\Phi }_\mathrm{\Lambda }|_\mathrm{\Lambda }|\mathrm{\Phi }_\mathrm{\Lambda },`$
where $`\mathrm{\Phi }`$ is a state in $`_\mathrm{\Lambda }`$ and $`|\mathrm{\Lambda }|`$ is the norm of $`\mathrm{\Lambda }`$. A state $`\omega `$ in the infinite-volume limit is said to be a ground state if it satisfies
$`\omega (h_x)=\epsilon _0`$
for arbitrary $`x`$.
In the present context, we have seen in the previous section that the numerical results for the Hubbard model imply that the energy per site of the SP state and the F state is identical in the thermodynamic limit for any boundary condition, especially PBC. In this boundary condition the system has a translational invariance, so that an SP state in a finite system should have
$`T(x)|\mathrm{SP}_\mathrm{\Lambda }=\mathrm{exp}(i\gamma )|\mathrm{SP}_\mathrm{\Lambda }`$
under the translation $`T(x)`$ which translates the state by $`x`$. We can expect that the invariance holds in the thermodynamic limit as well, so that
$`\mathrm{SP}|h_{x_0+x}|\mathrm{SP}`$ $`=`$ $`\mathrm{SP}|T(x)h_{x_0}T^1(x)|\mathrm{SP}`$
$`=`$ $`\mathrm{SP}|h_{x_0}|\mathrm{SP}`$
$`=`$ $`\epsilon _0,`$
and we can conclude that the SP state in PBC is also the ground state in the thermodynamic limit.
### C Symmetry breaking in the spiral spin state
One important question is: does the SP state, despite its being spin singlet, accompany a symmetry breaking in the thermodynamic limit? For simplicity, we consider the 1D case. From the numerical results for the spin-spin correlation in finite systems, a natural quantity that is expected to become nonzero (in PBC) is
$$\frac{1}{L^2}\underset{i,j\mathrm{\Lambda }}{}\mathrm{cos}(\theta _i\theta _j)𝐒_i𝐒_j_{\mathrm{SP}}>0,$$
(2)
where $`\theta _j=2\pi j/L`$. We can then introduce an order parameter,
$`O_\mathrm{\Lambda }^\pm `$ $``$ $`{\displaystyle \frac{1}{L}}{\displaystyle \underset{i\mathrm{\Lambda }}{}}\mathrm{exp}(i\theta _i)S_i^\pm {\displaystyle \frac{1}{L}}{\displaystyle \underset{i\mathrm{\Lambda }}{}}o_i^\pm ,`$ (3)
$`O_\mathrm{\Lambda }^{(1)}`$ $`=`$ $`{\displaystyle \frac{1}{2}}(O_\mathrm{\Lambda }^++O_\mathrm{\Lambda }^{})`$ (4)
$``$ $`{\displaystyle \frac{1}{L}}{\displaystyle \underset{i\mathrm{\Lambda }}{}}\mathrm{cos}\theta _iS_i^x+\mathrm{sin}\theta _iS_i^y,`$ (5)
$`O_\mathrm{\Lambda }^{(2)}`$ $`=`$ $`{\displaystyle \frac{1}{2i}}(O_\mathrm{\Lambda }^+O_\mathrm{\Lambda }^{})`$ (6)
$``$ $`{\displaystyle \frac{1}{L}}{\displaystyle \underset{i\mathrm{\Lambda }}{}}\mathrm{sin}\theta _iS_i^x+\mathrm{cos}\theta _iS_i^y,`$ (7)
$`O_\mathrm{\Lambda }^{(z)}`$ $``$ $`{\displaystyle \frac{1}{L}}{\displaystyle \underset{i\mathrm{\Lambda }}{}}S_i^z`$ (8)
where $`S^\pm =S^x\pm iS^y`$. Three operators $`O_\mathrm{\Lambda }^{(1)},O_\mathrm{\Lambda }^{(2)},O_\mathrm{\Lambda }^{(z)}`$, which obey a Lie algebra, satisfy the assumption required for Theorem 2.4 in Ref. . The expectation values of the squared $`O_\mathrm{\Lambda }^{(1)}`$ is
$`{}_{\mathrm{\Lambda }}{}^{}\mathrm{SP}|O_\mathrm{\Lambda }^{(1)2}|\mathrm{SP}_{\mathrm{\Lambda }}^{}`$ $`=`$ $`{\displaystyle \frac{2}{3L^2}}{\displaystyle \underset{i,j\mathrm{\Lambda }}{}}\mathrm{cos}(\theta _i\theta _j)𝐒_i𝐒_j_{\mathrm{SP}}`$
$`+{\displaystyle \frac{1}{2iL^2}}{\displaystyle \underset{i,j\mathrm{\Lambda }}{}}\mathrm{sin}(\theta _i\theta _j)S_i^+S_j^{}_{\mathrm{SP}}.`$
We can always take $`|\mathrm{SP}_\mathrm{\Lambda }`$ to be real in PBC since $`_\mathrm{\Lambda }`$ is real. $`S_i^\pm `$ can also be represented as a real matrix, so that $`S_i^+S_j^{}_{\mathrm{SP}}`$ is a real number, and we have $`S_i^+S_j^{}_{\mathrm{SP}}=S_j^+S_i^{}_{\mathrm{SP}}`$ so that the last line above vanishes. From the inequality eq.(2) we have
$`{}_{\mathrm{\Lambda }}{}^{}\mathrm{SP}|O_\mathrm{\Lambda }^{(1)2}|\mathrm{SP}_{\mathrm{\Lambda }}^{}`$ $`=`$ $`{}_{\mathrm{\Lambda }}{}^{}\mathrm{SP}|O_\mathrm{\Lambda }^{(2)2}|\mathrm{SP}_{\mathrm{\Lambda }}^{}`$
$`=`$ $`{\displaystyle \frac{2}{3L^2}}{\displaystyle \underset{i,j\mathrm{\Lambda }}{}}\mathrm{cos}(\theta _i\theta _j)𝐒_i𝐒_j_{\mathrm{SP}}>0.`$
On the other hand the order parameters, when not squared, satisfy
$`{}_{\mathrm{\Lambda }}{}^{}\mathrm{SP}|O_\mathrm{\Lambda }^{(1)}|\mathrm{SP}_{\mathrm{\Lambda }}^{}=_\mathrm{\Lambda }\mathrm{SP}|O_\mathrm{\Lambda }^{(2)}|\mathrm{SP}_\mathrm{\Lambda }=0.`$
Thus we have
$`\mathrm{\Psi }_\mathrm{\Lambda }^M{\displaystyle \frac{(O_\mathrm{\Lambda }^+)^M|\mathrm{SP}_\mathrm{\Lambda }}{(O_\mathrm{\Lambda }^+)^M|\mathrm{SP}_\mathrm{\Lambda }}}`$
as the ‘low-lying state’ in the Theorem 2.4 in Ref. , which asserts that a symmetry breaking can occur when some low-lying excitations whose energies approach that of the state in question like $`1/L`$ mix with it. Here $`\mathrm{\Psi }_\mathrm{\Lambda }^M`$ does have an energy that approaches to $`E_{\mathrm{SP}}`$ as
$`|(\mathrm{\Psi }_\mathrm{\Lambda }^M,H_\mathrm{\Lambda }\mathrm{\Psi }_\mathrm{\Lambda }^M)E_{\mathrm{SP}}|\mathrm{const}.\times ({\displaystyle \frac{M^2}{L}}).`$
We also have, according to Theorem 2.5,
$`\underset{k\mathrm{}}{lim}\underset{\mathrm{\Lambda }Z}{lim}(\mathrm{\Xi }_\mathrm{\Lambda }^k,O_\mathrm{\Lambda }^{(1)}\mathrm{\Xi }_\mathrm{\Lambda }^k)=\mathrm{const}.>0`$
where,
$`\mathrm{\Xi }_\mathrm{\Lambda }^k{\displaystyle \frac{1}{\sqrt{2k+1}}}\left[|\mathrm{SP}_\mathrm{\Lambda }+{\displaystyle \underset{M=1}{\overset{k}{}}}\left(\mathrm{\Psi }_\mathrm{\Lambda }^M+\mathrm{\Psi }_\mathrm{\Lambda }^M\right)\right].`$
Here, let us introduce the translation operator $`T(x)`$, which satisfies $`T(x)c_{x_0}^{}T^1(x)=c_{x_0+x}^{}`$. Since
$`T(x)O_\mathrm{\Lambda }^+T^1(x)`$ $`=`$ $`{\displaystyle \frac{1}{L}}{\displaystyle \underset{x_0\mathrm{\Lambda }}{}}\mathrm{exp}(i\theta _{x_0})T(x)S_{x_0}^+T^1(x)`$
$`=`$ $`{\displaystyle \frac{1}{L}}{\displaystyle \underset{x_0\mathrm{\Lambda }}{}}\mathrm{exp}(i\theta _{x_0})S_{x+x_0}^+`$
$`=`$ $`\mathrm{exp}(i\theta _x)O_\mathrm{\Lambda }^+,`$
$`T(x)\mathrm{\Psi }_\mathrm{\Lambda }^M=\mathrm{exp}(iM\theta _x)\mathrm{exp}(i\gamma )\mathrm{\Psi }_\mathrm{\Lambda }^M`$, where $`\gamma `$ is a constant. Thus, for $`o_i^+\mathrm{e}^{i\theta _i}S_i^+`$ appearing in eq. (3),
$`{}_{\mathrm{\Lambda }}{}^{}\mathrm{\Psi }_\mathrm{\Lambda }^M|o_{x+x_0}^+|\mathrm{\Psi }_\mathrm{\Lambda }^{M1}_{\mathrm{\Lambda }}^{}`$ $`=`$ $`{}_{\mathrm{\Lambda }}{}^{}\mathrm{\Psi }_\mathrm{\Lambda }^M|\mathrm{e}^{i\theta _{x+x_0}}S_{x+x_0}^+|\mathrm{\Psi }_\mathrm{\Lambda }^{M1}_{\mathrm{\Lambda }}^{}`$
$`=`$ $`{}_{\mathrm{\Lambda }}{}^{}\mathrm{\Psi }_\mathrm{\Lambda }^M|\mathrm{e}^{i\theta _x}T(x)o_{x_0}^+T^1(x)|\mathrm{\Psi }_\mathrm{\Lambda }^{M1}_{\mathrm{\Lambda }}^{}`$
$`=`$ $`\mathrm{e}^{(i\theta _x+iM\theta _xi(M1)\theta _x+i\gamma i\gamma )}`$
$`\times _\mathrm{\Lambda }\mathrm{\Psi }_\mathrm{\Lambda }^M|o_{x_0}^+|\mathrm{\Psi }_\mathrm{\Lambda }^{M1}_\mathrm{\Lambda }`$
$`=`$ $`{}_{\mathrm{\Lambda }}{}^{}\mathrm{\Psi }_\mathrm{\Lambda }^M|o_{x_0}^+|\mathrm{\Psi }_\mathrm{\Lambda }^{M1}_{\mathrm{\Lambda }}^{}.`$
Therefore, we have
$`\underset{k\mathrm{}}{lim}\underset{\mathrm{\Lambda }Z}{lim}{\displaystyle \frac{1}{|\mathrm{\Omega }|}}{\displaystyle \underset{x\mathrm{\Omega }}{}}\mathrm{\Xi }_\mathrm{\Lambda }^k|o_x^+|\mathrm{\Xi }_\mathrm{\Lambda }^k`$ $`=`$ $`\underset{k\mathrm{}}{lim}\underset{\mathrm{\Lambda }Z}{lim}{\displaystyle \frac{1}{|\mathrm{\Omega }|}}{\displaystyle \underset{x\mathrm{\Omega }}{}}\mathrm{\Xi }_\mathrm{\Lambda }^k|o_{x+x_0}^+|\mathrm{\Xi }_\mathrm{\Lambda }^k`$
$`=`$ $`\mathrm{const}.>0,`$
where $`\mathrm{\Omega }`$ is a finite subspace of $`\mathrm{\Lambda }`$ whose size is $`|\mathrm{\Omega }|`$. Hence,
$`\underset{k\mathrm{}}{lim}\underset{\mathrm{\Lambda }Z}{lim}{\displaystyle \frac{1}{|\mathrm{\Omega }|}}{\displaystyle \underset{x\mathrm{\Omega }}{}}\mathrm{\Xi }_\mathrm{\Lambda }^k|S_x^+|\mathrm{\Xi }_\mathrm{\Lambda }^k=\mathrm{const}.>0`$
Namely, although the SP state itself does not break the symmetry, it does so in in the thermodynamic limit by letting $`\{\mathrm{\Psi }_\mathrm{\Lambda }^M\}`$ mix with it in that limit. In other words, an arbitrarily large but finite region in an SP state has a broken SU(2) with a finite magnetization.
If we use identities $`S^{}=_{x_i\mathrm{\Lambda }}S_i^{}`$, $`[O^{},S^{}]=0`$, $`\mathrm{SP}|S^{}=0`$,
$`\mathrm{\Psi }_\mathrm{\Lambda }^M|\mathrm{F};S^z=M`$
$`=_\mathrm{\Lambda }\mathrm{SP}|(O_\mathrm{\Lambda }^{})^M(S^{})^{S_{\mathrm{max}}M}|\mathrm{F};S^z=S_{\mathrm{max}}_\mathrm{\Lambda }`$
$`=0.`$
Thus both of the SP state and $`\{\mathrm{\Psi }_\mathrm{\Lambda }^M\}`$ that mixes with it in the thermodynamic limit turn out to be orthogonal to the F state for finite systems.
Finally, let us discuss whether $`\mathrm{\Xi }`$ is the ground state in the thermodynamic limit by checking whether $`\mathrm{\Xi }|h_x|\mathrm{\Xi }=\epsilon _0`$ holds for arbitrary $`x`$. Since
$`{}_{\mathrm{\Lambda }}{}^{}\mathrm{SP}|(O_\mathrm{\Lambda }^{})^Mh_{x_0+x}(O_\mathrm{\Lambda }^+)^M|\mathrm{SP}_{\mathrm{\Lambda }}^{}`$
$`=_\mathrm{\Lambda }\mathrm{SP}|(O_\mathrm{\Lambda }^{})^MT(x)h_{x_0}T^1(x)(O_\mathrm{\Lambda }^+)^M|\mathrm{SP}_\mathrm{\Lambda }`$
$`=_\mathrm{\Lambda }\mathrm{SP}|(O_\mathrm{\Lambda }^{})^Mh_{x_0}(O_\mathrm{\Lambda }^+)^M|\mathrm{SP}_\mathrm{\Lambda }`$
$`=\epsilon _0,`$
$`(\mathrm{\Psi }_\mathrm{\Lambda }^M,h_x\mathrm{\Psi }_\mathrm{\Lambda }^M)=\epsilon _0`$ does indeed hold for arbitrary $`x`$ in the limit of $`L\mathrm{}`$. This means $`\mathrm{\Xi }`$ is the ground state.
## III The excitation spectrum of the spiral spin state
Having shed light in the context of Koma-Tasaki’s argument, we may now conjecture that the SP state and the F state are identical in the thermodynamic limit. Thus let us have a closer look at the SP state as compared with the F state in three one-dimensional models, i.e., the $`t`$-$`t^{}`$ Hubbard model, Tasaki’s flat band model, and the two-leg Hubbard ladder. To determine whether the SP and F states are identical in the thermodynamic limit, we must show that the relation $`\omega _{\mathrm{SP}}(A)=\omega _{\mathrm{SF}}(A)`$ holds for arbitrary local operator $`A`$. Although it is impossible to confirm this for all the possible $`A`$, if the F state and the SP state are identical in that limit, the excitation spectrum of these state in finite size system are expected to be almost identical. To confirm this, we look specifically at the single-particle spectral function and the dynamical spin and charge correlation functions of the SP state. Indeed we will show that since the nature of the SP state in finite systems is almost identical to that of the F state, we cannot distinguish the SP state from the F state by looking at these quantities for sufficiently large system.
### A $`t`$-$`t^{}`$ Hubbard model
In the $`t`$-$`t^{}`$ Hubbard model where the electron transfer extends to next nearest neighbors ($`t^{}`$), the ground state is $`S_{\mathrm{tot}}=0`$ for any value of $`U`$ in PBC at least for sizes up to 12 sites as mentioned in the previous section. Let us first look at the spin-spin correlation function, $`\mathrm{\Phi }_G|S_i^zS_j^z|\mathrm{\Phi }_G`$, to confirm that the $`S_{\mathrm{tot}}=0`$ ground state for large enough $`UU_C`$ is the SP state. We show in Fig. 4 the result for 10 electrons in a 12-site ring (two holes) with $`U=20`$ or 40. For $`U=20`$, a short-range antiferromagnetic spin-spin correlation is observed for $`S_iS_{i+1}<0`$. By contrast, for $`U=40`$ (or larger, which falls upon the ferromagnetic region in the phase diagram; see Fig. 14 below), the spin-spin correlation has a wave length as large as the system size, which is the same behavior of the lowest $`S_{\mathrm{tot}}=0`$ state in OBC (Fig. 1).
A more detailed nature the states is encapsulated in the single-particle spectral function given by
$`A^\pm (k,\omega )={\displaystyle \frac{1}{\pi }}\mathrm{Im}\mathrm{\Phi }_G|\gamma _{k\sigma }^{}{\displaystyle \frac{1}{\omega \pm (E_0)i0}}\gamma _{k\sigma }^\pm |\mathrm{\Phi }_G,`$ (9)
where $`A^+(A^{})`$ denote the electron addition (removal) spectrum with $`\gamma _{k\sigma }^+c_{k\sigma }^{}`$, $`\gamma _{k\sigma }^{}(\gamma _{k\sigma }^+)^{}`$, $`|\mathrm{\Phi }_G`$ and $`E_0`$ the ground state and its energy, respectively. We have numerically calculated this quantity. In Fig. 5, we show the result for the parameter values adopted in Fig.4 with $`U=40`$. We can see that the Luttinger relation, $`k_F=\pi n/2`$ (with $`n=10/12=0.83`$ here), for the spin unpolarized electrons is violated in favor of $`k_F=\pi n`$, which would be expected for fully spin-polarized fermions. This behavior in the single particle spectral function is consistent with the nearly ferromagnetic nature of the SP state. Such a behavior was first found for the $`t`$-$`t^{}`$-$`J`$ model by Eder and Ohta, and then given an interpretation subsequently.
The $`k`$-dependence of the spectral function may be analyzed by evoking the argument by Doucot and Wen(DW), who studied the infinite-$`U`$ Hubbard model on a two-dimensional square lattice with two holes to find a trial state that gives an energy lower by $`1/L^2`$ than Nagaoka’s ferromagnetic state. The key idea of DW is based on the following intuition. Holes behave as free fermions for on-site interactions when the background spin state is ferromagnetic (or nearly so). Assume that fully polarized electrons take an open-shell configuration (in which there is a degeneracy in the free-electron configurations) in $`k`$-space in PBC. If the background spin state is changed from F to the SP spin texture, a fictitious gauge field corresponding to half flux quantum is generated. This shifts the $`k`$-points by half the $`k`$-point spacing, which lowers the kinetic energy because the polarized electrons now take a closed-shell configuration.
It is a nontrivial question whether the DW’s trial state may be applied to reproduce the spinless-fermion-like behavior in $`A(k,\omega )`$ found here for the 1D system. Let us assume, for simplicity, that the holes hop in a rigid spin background, although in the original paper DW considered the holes dressed by spin waves. The energy of the hole in the $`t`$-$`t^{}`$ Hubbard model is then given as
$`\epsilon `$ $`=`$ $`2t\mathrm{cos}\left(k\pm {\displaystyle \frac{\pi }{L}}\right)+2t^{}\mathrm{cos}2\left(k\pm {\displaystyle \frac{\pi }{L}}\right),`$ (10)
where $`k(=2\pi N/L,0NL,N`$ is an integer) is wave number and $`\pm `$ corresponds to the sign of the fictitious flux. We can see that the fitting in Fig. 5 is remarkably good. To be precise, the transfer integrals in the original DW are taken to be effective ones ($`tt\mathrm{cos}(\pi /L),t^{}t^{}\mathrm{cos}(2\pi /L)`$) to take care of the reduction in the transfer between slightly twisted spins, but the present result is better fitted with these reduction factors omitted (dashed curves).
In the addition spectrum an almost dispersionless band of low intensity peaks is seen (at around $`E=2`$ in Fig. 5). If the ground state were the fully polarized F state, then such a band is expected, since we can add an opposite spin at any $`k`$-point. A remnant of such a band again suggests the nearly ferromagnetic nature of the SP state.
A next question is whether the SP state persists for more than two holes. In Fig. 6, we show the spin-spin correlation for 6 electrons on a 12-site lattice (six holes, or $`n=0.5`$; quarter filled) in PBC for $`U=4,6`$. We can see that for $`U=6`$, the spin-spin correlation is again spiral. To investigate whether DW’s picture is valid for such a high hole doping, we have calculated the single-particle spectrum for $`U=6`$ in Fig. 7. The spectrum for the SP state is fitted well by the energy dispersion defined by the eq.(10), again $`t`$ and $`t^{}`$ not reduced, even though the hole concentration is as large as $`n_h=0.5`$. This is surprising, since the assumption that holes are nearly free would be valid only for small enough doping.
We next question whether the SP state is connected adiabatically to the antiferromagnetic state as we decrease the Hubbard $`U`$. We show the ground-state energy as a function of $`1/U`$ for 6 electrons on a 12-site ring in Fig. 8. We can clearly identify a level crossing appearing as a cusp around $`U5`$, which indicates that a transition (rather than an adiabatic connection) occurs within the $`S=0`$ space from the antiferromagnetic phase to the SP phase.
We now turn to the dynamical spin and charge correlation functions
$`N(k,\omega )`$ $`=`$ $`{\displaystyle \frac{1}{\pi }}\mathrm{Im}\mathrm{\Phi }_G|N_k{\displaystyle \frac{1}{\omega +(E_0)i0}}N_k|\mathrm{\Phi }_G,`$
$`S(k,\omega )`$ $`=`$ $`{\displaystyle \underset{\alpha }{}}{\displaystyle \frac{1}{\pi }}\mathrm{Im}\mathrm{\Phi }_G|S_k^\alpha {\displaystyle \frac{1}{\omega +(E_0)i0}}S_k^\alpha |\mathrm{\Phi }_G,`$
where $`N_k`$ and $`S_k^\alpha `$ are Fourier transform of $`(n_in)`$ and $`S_i^\alpha (\alpha =x,y,z)`$, respectively. In Fig. 9, we show the results for the F state in APBC and the SP state in PBC for 10 electrons in 12 sites with $`U=40`$. We can see that the behavior of these functions for the SP state is surprisingly similar to those of the F state even for $`L=12`$. We can expect that they will become identical in the thermodynamic limit.
### B One-dimensional Tasaki model
Tasaki proposed a Hubbard model on a special class of lattice structures for which the lowest energy band is dispersionless (flat). He proved rigorously that the ground state is ferromagnetic for arbitrary interaction strength ($`0<U<\mathrm{}`$) when the flat band is half-filled.
To clarify whether metallic ferromagnetism can be realized in models having a non-half-filled, nearly-flat band, Sakamoto and Kubo investigated with DMRG the Hubbard model on a chain of triangles, which may be thought of as a realization of Tasaki’s model in 1D. They obtained results which suggest that the system exhibits metallic, fully-polarized ferromagnetism for sufficiently strong interaction. On the other hand, Watanabe and Miyashita found that an SP state is the ground state at least for system sizes up to $`L=12`$ for appropriate conditions, i.e., sufficiently large $`U`$, $`n>0.5`$, even number of electrons and PBC.
The Hamiltonian is given by
$``$ $`=`$ $`t{\displaystyle \underset{i=1}{\overset{L/2}{}}}{\displaystyle \underset{\sigma }{}}(c_{2i,\sigma }^{}c_{2(i+1),\sigma }+\mathrm{H}.\mathrm{c}.)`$
$`+\alpha t{\displaystyle \underset{i=1}{\overset{L/2}{}}}{\displaystyle \underset{\sigma }{}}[(c_{2i1,\sigma }^{}+c_{2i+1,\sigma }^{})c_{2i,\sigma }+\mathrm{H}.\mathrm{c}.]`$
$`+\beta t{\displaystyle \underset{i=1}{\overset{L/2}{}}}{\displaystyle \underset{\sigma }{}}n_{2i1,\sigma }+U{\displaystyle \underset{i=1}{\overset{L}{}}}n_in_i,`$
where $`t`$ is the transfer for the bottom bond in a triangle, $`\alpha t`$ for the bonds connecting bottom and top sites and $`\beta t`$ the site energy of a top site. Since a unit cell contains two sites, the single-electron spectrum consists of two bands with the dispersion relation, up to a constant,
$`\epsilon ^\pm `$ $`=`$ $`{\displaystyle \frac{t}{2}}\left[2\mathrm{cos}k\pm \sqrt{4\mathrm{cos}^2k+(8\alpha ^24\beta )\mathrm{cos}k+\beta ^2+8\alpha ^2}\right].`$
The lower band is then flat when $`\beta =\alpha ^22`$ is satisfied. Hereafter in this subsection, we take $`t=1`$, $`\alpha =2`$, $`\beta =0`$ (for which the lower band is dispersive) and $`L`$ even.
While Watanabe and Miyashita have shown that SP states appear for more than quarter-filled bands ($`0.5<n<1.0`$), we find here that the SP state is the ground state of $`S_{\mathrm{tot}}=0`$ sector for sufficiently large $`U`$ in the region $`n0.5`$ as well at least for $`L12`$. In Fig. 10, we show the spin correlation for 6 electrons on a 12-site lattice in PBC for $`U=4`$ and $`U=6`$. We can see that for the case of $`U=6`$ the spin correlation length is indeed as large as the system size. We have also calculated the single-particle spectral function, eq.(9), where we take $`\gamma _{i,\sigma }^+=c_{2i,\sigma }^{}`$ or $`\gamma _{i,\sigma }^+=c_{2i1,\sigma }^{}`$. If we combine them, they should contain the contributions from the two bands. Here $`\gamma _{i\sigma }^+`$ is a Fourier transform of $`\gamma _{k\sigma }^+`$. In Fig. 11, we show the result for 6 electrons on 12 site lattice, where the spectrum is fitted by a two-band extension of eq.(10),
$`\epsilon ^\pm =t\mathrm{cos}(k\pm {\displaystyle \frac{\pi }{L}})\pm t\sqrt{\mathrm{cos}^2(k\pm {\displaystyle \frac{\pi }{L}})+8\mathrm{cos}(k\pm {\displaystyle \frac{\pi }{L}})+8}.`$ (11)
Fig. 11 shows that the DW picture (almost free electrons hopping in a twisted spin background) is surprisingly good even for the two-band case.
### C Two-leg Hubbard ladder
Ferromagnetic ground state in 1D has also been obtained for ladder systems, where two chains are connected with an inter-chain transfer, $`t_{}`$. Liang and Pang suggested that for $`t_{}/t=1`$ and $`U=\mathrm{}`$, the F state is one of the ground states for $`n_h<0.22`$, while the ground state is spin singlet for $`n_h>0.4`$. Kohno presented rigorous results for $`U=\mathrm{}`$ 2-leg Hubbard ladder in the limit of large inter-chain hopping ($`t_{}/t\mathrm{}`$). He also studied the case for finite values of $`t_{}/t`$ with DMRG to obtain the phase diagram. His results are consistent with that of Liang and Pang. These studies assumed OBC.
On the other hand, it is intriguing to study whether an SP state appears in the ladders as well for PBC, and if so, whether $`A(k,\omega )`$ exhibits a DW-like behavior. The Hamiltonian is given by
$``$ $`=`$ $`t{\displaystyle \underset{i=1}{\overset{L}{}}}{\displaystyle \underset{\alpha ,\sigma }{}}(c_{i,\alpha \sigma }^{}c_{i+1,\alpha \sigma }+\mathrm{H}.\mathrm{c}.)`$
$`t_{}{\displaystyle \underset{i=1}{\overset{L}{}}}{\displaystyle \underset{\sigma }{}}(c_{i,1\sigma }^{}c_{i,2\sigma }+\mathrm{H}.\mathrm{c}.)+U{\displaystyle \underset{n=1}{\overset{L}{}}}{\displaystyle \underset{\alpha }{}}n_{i\alpha }n_{i\alpha },`$
where $`\alpha (=1,2)`$ labels the two legs of the ladder. We set $`t=1`$, $`U=\mathrm{}`$.
We have calculated the intra-chain spin-spin correlation $`S_{i\alpha }^zS_{j\alpha }^z`$ (Fig.12), and found that the ground state is indeed SP. We have also calculated the spectral function, eq.(9), where we now take $`\gamma _{k\sigma }^+=c_{k1\sigma }^{}+c_{k2\sigma }^{}`$ (creating an electron in the bonding band) or $`\gamma _{k\sigma }^+=c_{k1\sigma }^{}c_{k2\sigma }^{}`$ (antibonding). In Fig.13, we show the result for two holes (14 electrons on a $`8\times 2`$ ladder). We fit the spectrum by the dispersion,
$`\epsilon =2t\mathrm{cos}\left(k\pm {\displaystyle \frac{\pi }{L}}\right)\pm t.`$ (12)
Again the picture of two holes hopping in a twisted spin background gives an accurate description.
## IV Magnetic phase diagram from finite-size studies
Our results in the previous sections suggest that we should regard an SP state as being ferromagnetic rather than non-magnetic even though $`S_{\mathrm{tot}}=0`$. So we finally come to the problem of how to determine the magnetic phase diagram from finite-size studies in the light of the SP state. As we have stressed, the ground state of finite systems in a certain boundary condition always has $`S_{\mathrm{tot}}=0`$ no matter how $`U`$ is strong in some one-dimensional models as exemplified by the $`t`$-$`t^{}`$ Hubbard model and the Hubbard ladder. Even when a ferromagnetic state appears for finite systems, the magnetic region in the phase diagram can shrink for some system size as shown by Sakamoto and Kubo for the 1D Tasaki’s model. This is due to the appearance of the SP state, which are encountered if we assume PBC for even number of electrons, or more generally, if we assume a boundary condition for which the fully-polarized electrons take an open-shell configuration in the ground state.
Since the SP state may be regarded as ferromagnetic as elaborated in previous sections, we have then to distinguish the SP state from nonmagnetic states by calculating the spin correlation, etc. Here we propose that this difficulty can be readily overcome by taking an appropriate combination of the boundary condition and the electron number for which the SP state is excluded from the ground state — this enables us to obtain a reliable magnetic phase diagram within finite-size studies by simply looking at $`S_{\mathrm{tot}}`$ without worrying about the existence of $`S_{\mathrm{tot}}=0`$ ferromagnetic-like states. We stress here that this procedure is allowed because the SP state may be regarded as ferromagnetic: if the SP state were a distinct state, then the exclusion of SP states would lead to a missed phase transition.
Here we illustrate this with an appropriate boundary condition that makes the F state take a closed-shell configuration. We present the phase diagram thus obtained for the models employed in section III, i.e., 1D $`t`$-$`t^{}`$ Hubbard model, Tasaki’s model, and 2-leg Hubbard ladder. In all cases, the results accurately coincide with those obtained by DMRG calculation for much larger systems with OBC. In other words, the result converges to the thermodynamic limit rapidly when we concentrate on the F states. We may expect that this method should be applicable to the determination of the phase diagram of 2D or higher dimensional systems.
### A $`t`$-$`t^{}`$ Hubbard model
We first observe that, in the one-electron energy band of the $`t`$-$`t^{}`$ Hubbard model,
$`\epsilon (k)=2t\mathrm{cos}k+2t^{}\mathrm{cos}(2k),`$
the band minimum at $`k=0`$ splits into double minima for $`t^{}>0.25`$ as $`t^{}`$ is increased. Hereafter in this section, we set $`t=1`$.
We first look at the case of $`t^{}=0.2<0.25`$. The band has a single minimum, so that an even number of fully-polarized electrons take a closed-shell configuration for all densities $`n<1`$ if we assume APBC to put the $`k`$-points symmetrically about $`k=0`$. In Fig. 14, we show the exact diagonalization results for 4, 6, 8, or 10 electrons in a 12-site system ($`0.33n0.83`$), and 4, 6, or 8 electrons in a 10-site system. We have also plotted the DMRG result for the phase boundary obtained by Daul and Noack for a 50-site system with OBC. We can immediately see that the results obtained with APBC for fairly small systems is close to that obtained for larger ones.
If we move on to the case of a larger $`t^{}=0.8`$, an even number of fully-polarized electrons only take a closed-shell configuration for APBC when the Fermi energy is higher than $`2t+2t^{}`$. When the Fermi energy is lower, the situation depends on how the $`k`$-points are located around the double-minimum dispersion in the given boundary condition, so we concentrate on the former case. In Fig. 15, we show the results for 8 or 10 electrons in a 12-site system and 8 electrons in a 10-site system. we can see that the result for small systems with APBC is very close to the DMRG result for a 50-site system.
### B One-dimensional Tasaki’s model
The phase diagram against $`U`$ and $`\beta `$ (that controls the dispersion of the band) of the 1D Tasaki’s model (a chain of triangles) was obtained by Sakamoto and Kubo for various band fillings ($`n=1/2,1/4,3/8`$) by means of the DMRG method with PBC. As they mention, the phase boundary between $`S_{\mathrm{tot}}=S_{\mathrm{max}}`$ and $`S_{\mathrm{tot}}=0`$ varies strongly with the system size for small systems, so that a reliable phase diagram can be obtained only by going to sufficiently large ($`L=32`$) systems.
In Fig. 16, we show the phase boundary determined by exact diagonalization for $`n=1/2`$ and system size as small as 8 or 12 sites with APBC. Here we fix $`\alpha =2`$, so that the lower band becomes flat for $`\beta =2`$. Again, the result excellently agrees with that of 32 sites with PBC.
### C Two-leg Hubbard ladder model
As mentioned in Sec. IIIC, Kohno obtained the phase diagram of $`U=\mathrm{}`$ Hubbard ladder by means of the DMRG method. According to his results, there is a wide region of partially ferromagnetic phase in the neighborhood of paramagnetic phase.
Let us consider whether we can reproduce the phase boundary between non-magnetic phase and partially polarized phase by assuming APBC for small systems. We consider the case of 6 and 8 electrons on a $`6\times 2`$ lattice and 4 and 6 electrons on a $`4\times 2`$ lattice. We performed calculation for $`0.125t_{}/(t_{}+t)0.909`$, and found the ferromagnetic-non-magnetic transition only for the case of 8 electrons on a $`6\times 2`$ lattice, around $`t_{}/(t_{}+t)0.45`$. This is consistent with Kohno’s the result.
All these results indicate that the phase boundary between the ferromagnetic and the antiferromagnetic states is reasonably close to the thermodynamic limit for small systems when a care is taken (i.e., letting fully-polarized electrons take closed shell configurations).
## V Summary
To summarize, we have studied the relation between the fully-polarized ferromagnetic state and the spiral spins state, a spin-singlet spin state that has a spin correlation length as large as the system size, which accompanies the fully-polarized ferromagnetic state in a number of electron correlation models. As a typical example, we calculate the energy of the SP state and the F state for one-dimensional $`U=\mathrm{}`$ $`t`$-$`t^{}`$ Hubbard model, It suggests that the SP state is also the ground state in the thermodynamic limit when the F state is the ground state. Following the argument by Koma and Tasaki, we have also indicated how a symmetry becomes broken in the SP state in the thermodynamic limit, where some ‘low-lying states’ become hybridized.
We have then characterized the SP state, an itinerant magnetic state, by calculating the one-particle spectral function. The result is interpreted in a picture in which holes move almost freely in a twisted spin configuration. We have also calculated the dynamical spin and charge correlation functions to find that their behaviors in the SP state are similar to those of the F state.
From these we have conjectured that we should regard the SP state and the F state are equivalent in the thermodynamic limit, even though the SP state (spin-singlet) and the F state (fully-polarized) are opposite extremes in terms of $`𝐒_{\mathrm{tot}}`$. We have then shown that the magnetic phase diagram can be determined accurately from finite-size studies by taking appropriate boundary condition (that depends on the number of electrons which are accommodated in the shells in the non-interacting case) that pushes up the energy of the SP state. This enables us to concentrate on the F states, i.e., to simply look at the change in $`S_{\mathrm{tot}}`$. This is permissible since the SP and F states are regarded to be equivalent, so that we can concentrate on either of them. We have obtained the phase diagram in this way for 1D $`t`$-$`t^{}`$ Hubbard model, Tasaki’s model and 2-leg Hubbard ladder. We have found that the phase boundary between the ferromagnetic state and the non-magnetic state determined in such a way accurately coincides with those obtained by the DMRG calculation for much larger systems. Since the method does not depend on the dimensionality, we may expect that it should be applicable to the determination of the phase diagram in two or higher dimensional systems.
## VI Acknowledgments
We would like to thank Kazuhiko Kuroki, Koichi Kusakabe, Tohru Koma for illuminating discussions. R.A. is supported by Japan Society for Promotion of Science. Numerical calculation were done on FACOM VPP 500/40 at the Supercomputer center, Institute for Solid State Physics, University of Tokyo. |
warning/0002/cond-mat0002063.html | ar5iv | text | # Characterization of One-Dimensional Luttinger Liquids in Terms of Fractional Exclusion Statistics
## I Introduction
It is well-known that the Landau theory of Fermi liquids fails to describe most of one-dimensional (1-d) interacting many-body systems. To provide a substitute, Haldane proposed, years ago, the concept of the Luttinger liquid , defined by a set of low-lying excitations and critical exponents of the asymptotic correlation functions. Like Fermi liquids, there is a (pseudo-)Fermi surface for the quasiparticle-like excitations in Luttinger liquids, so that the classification of low-lying excitations is similar to that in Fermi liquids. However, the exponents of the asymptotic correlation functions (at low temperature) are distinct from those for Fermi liquid theory. For one-component systems (without internal degrees of freedom), the low-energy or low-temperature behavior of a Luttinger liquid is controlled by a single parameter, the Haldane controlling parameter. It controls not only all exponents, but also the velocity ratios between different types of elementary excitations. The Fermi liquid theory is a special case of the Luttinger liquids with the Haldane parameter $`\lambda =1`$.
In recent years, the failure of Landau’s theory of Fermi liquids to describe several newly discovered strongly correlated electron systems have revived the interests in the theory of Luttinger liquids. Among other questions, compared with Fermi liquids, one would like very much to know the answer to the following questions:
* What is the physical meaning of Haldane’s controlling parameter? Or more precisely, how to use physical properties of low-lying excitations to characterize the concept of Luttinger liquids?
* Does Haldane’s theory of Luttinger liquids possess universality in one dimension, just like Landau’s theory of Fermi liquids in three dimensions? Or equivalently, in terms of modern language of renormalization group, does the Luttinger liquids describe the infrared (or low-energy) fixed points in 1-d systems?
* In what directions could one expect to go for generalizing the concept of Luttinger liquids to higher than one dimensions?
In short, a characterization of 1-d Luttinger liquids, other than using a bunch of excitations and exponents, is in demand for gaining more insights and looking for possible generalization.
To achieve this, let us recall what motivated Landau’s concept of Fermi liquids, which is known to describe an infrared fixed point (or a universality class) of interacting electron systems. The basic idea behind it is based on the following organizing principle for interacting many-body systems: At low temperature, the low-lying excited states of an interacting many-body system above a stable ground state can be viewed as consisting of weakly coupled elementary excitations. Here ”weakly coupled” only means that the total energy can be written as a sum of single-particle (dressed) energies, while the dispersion of the dressed energy may well depend on the total particle number, a signal of remnant interactions between the quasiparticles. According to Landau, the ground state and the low-lying excited states of a Fermi liquid are approximately, to a good accuracy at sufficiently low temperature, described by those of an ideal Fermi gas with dressed energy for the quasiparticles.
We note the significant role played by the ideal Fermi gas distribution (with dressed energy) in this description of Fermi liquids. Actually it is the ideal Fermi gas that gives a characterization to the Fermi liquid fixed point, and a meaning to the universality of the concept of Fermi liquids. This inspires us to try to give a characterization of the 1-d Luttinger liquids along a similar line of thoughts, namely using a properly generalized concept of exclusion statistics, of which a special case is the usual Fermi statistics. Because the concept of quantum statistics in statistical mechanics is independent of the dimensionality of a system, a characterization of 1-d infrared fixed points using statistics, if successful, would shed light on how to generalize to higher dimensional non-Fermi liquids.
Fortunately, a generalization as such has been available recently, under the name of fractional exclusion statistics (FES). It is based on a new combinatoric rule for the many-body state counting , which is essentially an abstraction and generalization of Yang-Yang’s state counting in 1-d soluble many-body models. FES has been shown to be applicable to elementary excitations in a number of exactly solvable models for strongly correlated systems , anyons in the lowest Landau level , and quasiparticle excitations in the fractional quantum Hall effect . The thermodynamics of the so-called generalized ideal gas (GIG) associated with FES have been studied in a general framework.
Inspired by these results, the thoughts along the lines indicated in the above paragraphs have led two of present authors to propose that at least for some strongly correlated systems or non-Fermi liquids, their low-energy or low-temperature fixed point may be described by a GIG associated with FES, similar to the way that of the Fermi liquid fixed point by the ideal Fermi gas . As a testimony to this proposition, a sketchy proof was given in that short letter that the low-$`T`$ critical properties of the 1-d Luttinger liquids are exactly reproduced by those of 1D ideal excluson gas(IEG), if one identifies the Haldane parameter of the former with the statistics parameter $`\lambda `$ of the latter. (We call the particles obeying the FES without mutual statistics exclusons). Threfore IEG can be used to describe the fixed points of the Luttinger liquids. In this paper, we will present our results obtained in in details, much of which was not published before.
A main tool we use in this study is bosonization of the 1-d excluson systems at low $`T`$, à la Tomonaga and Mattis and Lieb . To bosonize an IEG system is a little bit tricky, because at low temperature the linearized dispersion of dressed energy versus pseudo-momentum has different slope outside and inside the pseudo-Fermi sea: There is ‘refraction’ at both pseudo-Fermi points. In spite of this, we still manage to construct well-defined density fluctuation operators that obey the $`U(1)`$ current algebras and physically describe free phonons. Then, the Tomonaga-Mattis-Lieb bosonization applies, resulting a bosonized effective field theory, in agreement with Haldane’s harmonic fluid description of the Luttinger liquid . Then the asymptotic correlation functions and their exponents can be systematically calculated. In this way, the critical properties of IEG reproduce those of the Luttinger liquids.
An important consequence of our bosonization is that the low energy behavior of IEG is controlled by an orbifold conformal field theory (CFT) with central charge $`c=1`$ and compactified radius $`R=\sqrt{1/\lambda }`$. This variant of $`c=1`$ CFT is not the ordinary $`c=1`$ CFT compactified on a circle $`S^1`$, rather it is compactified on an orbifold $`S^1/Z_2`$, which is topologically an interval . The differences arise due to different selection rules for vertex operators, that constrain quantum numbers of possible quasiparticle excitations in the system. In the usual literature this difference quite often is overlooked. Only within the orbifold CFT the IEG with statistics parameter $`\lambda =1`$ recovers ideal Fermi gas, as it should be. Also the two classes of $`c=1`$ CFT’s have different duality relation; only the one in orbifold CFT reproduces the known particle-hole duality in IEG, $`\lambda 1/\lambda `$, as given in . (For the details and more elaboration, see below. We note that a similar situation happens for the Calogero-Sutherland (C-S) model: The low-energy effective field theory for the bosonic and fermionic C-S models belongs to, respectively, the above-mentioned two classes of $`c=1`$ CFT.)
The fact that the low-$`T`$ behavior of IEG is controlled by a conformally invariant theory is significant, implying that indeed IEG provides a characterization of infrared fixed points, having the conformal invariance as required by renormalization group. We have also studied the effects of mutual statistics between different pseudomomenta and of the Luttinger-type (density-density) interactions among exclusons. In either case, the low-$`T`$ behavior is controlled by an effective statistics $`\lambda _{eff}`$ for excitations near the Fermi points, the same way as $`\lambda `$ in the case of IEG. In one dimension both the momentum-independent part of interactions and change in chemical potential $`\mu `$ are relevant perturbations , leading to a continuous shift in the fixed-point line parameterized by $`\lambda `$. All these will be explained in details in the present paper.
To make this paper self-contained, we devote the next two sections, Sec. II and Sec. III, to reviewing the Luttinger liquid theory and the GIG associated with FES, respectively. In Sec. IV, we discuss the low-energy behavior of the IEG system and achieve its bosonization. In Sec. V, the generalization to the GIG with mutual statistics as well as the non-ideal gas with FES are provided. The last section is dedicated to conclusions and discussions.
## II Luttinger liquid
The Luttinger liquid, which describes a very large class of one-dimensional interacting many-body systems, is introduced because of the infrared divergence of certain vertices in the Fermi liquid description of the 1-d systems. Some pioneering works have been done in the Luttinger model before the Luttinger liquid concept . The model has been exactly solved by using the bosonization technique . Haldane re-solved the model with the following important observations:
(i) Besides a linearized spectrum of non-zero mode excitations, i.e., the density fluctuations (sound waves), there are two kinds of zero mode excitations, single-particle excitations by adding extra particles to the system and persistent currents by making Galileo boosts.
(ii) There is a fundamental relation among the velocities of these three types of excitations
$$v_s=\sqrt{v_Nv_J},$$
(1)
where $`v_s`$ is the sound velocity, $`v_J`$ the current velocity and $`v_N`$ a velocity related to the change in particle number. The velocity ratios define a controlling parameter, $`e^{2\phi }`$, by
$$v_N=v_se^{2\phi },v_J=v_se^{2\phi }.$$
(2)
(iii) The above defined controlling parameter measures the essential renormalized coupling constant, and is the unique parameter that determines the exponents of power-law decay in the zero-temperature correlation functions.
Based on these observations, Haldane defined the Luttinger liquids as 1-d systems that have similar behavior (i)-(iii) at low temperature just like the Luttinger model. In this way the Luttinger liquids are characterized through their excitations and the exponents of the asymptotic correlation functions.
To be more precise, recall that the Luttinger model describes a one-dimensional interacting fermion system with the Hamiltonian
$$H=𝑑x|\psi |^2+\frac{1}{2}𝑑x𝑑yV(xy)\rho (x)\rho (y).$$
(3)
In the low energy limit, the Hamiltonian (3) can be bosonized as
$$H=v_s\underset{q}{}|q|b_q^{}b_q+\frac{1}{2}(\pi /L)(v_NM^2+v_JJ^2),$$
(4)
where $`b_q`$ are the standard boson annihilation operators, and $`M`$ and $`J`$ the operators corresponding to adding extra particle and boosting persistent currents, whose eigenvalues obey the following selection rule,
$$(1)^J=(1)^M.$$
(5)
The total momentum of the system also has a bosonized form
$$P=[k_F+(\pi /)M]J+\underset{q}{}qb_q^{}b_q,$$
(6)
with $`k_F`$ being the Fermi momentum.
Eqs.(1,2,4-6) turned to be universally valid for the description of the low-energy properties of gapless interacting one-dimensional spinless fermion systems even for those not exactly soluble with a conserved current $`J`$. This universality class is named as the Luttinger liquid by Haldane . The Luttinger liquid has a model-independent representation, namely the harmonic fluid description , which is convenient for calculating the correlation functions. The results of the harmonic fluid representation are listed in the Appendix, for later use to be compared with our bosonization theory of the IEG.
The Haldane theory of Luttinger liquids is based on the significant observation that the low-$`T`$ behavior of the Luttinger model is universal. Naturally arises the question: Why is it so? In this paper we intend to answer this question by pointing out a profound coincidence of the low-$`T`$ behavior of the Luttinger model and that of ideal excluson gas (IEG), i.e., ideal gas of particles obeying fractional exclusion statistics: The universality of the former is due to that of the latter.
## III Generalized ideal gas
In quantum mechanics, there are two ways to define the statistics of particles. One is in terms of the symmetry of the many-body wave function under particle exchange. The other is based on the state counting. Here we are interested in the latter definition. As is well-known, bosons and fermions have different countings for many-body states, or different statistical weights $`W`$: The number of quantum states of $`N`$ particles occupying a group of $`G`$ states is, for bosons and fermions respectively, given by
$$W_b=\frac{(G+N1)!}{N!(G1)!},\mathrm{or}W_f=\frac{G!}{N!(GN)!}.$$
(7)
A simple interpolation between bosons and fermions is given by
$$W=\frac{[G+(N1)(1\lambda )]!}{N![G\lambda N(1\lambda )]!},$$
(8)
with $`\lambda =0`$ corresponding to bosons and $`\lambda =1`$ to fermions. The physical meaning of this equation is the following: By assumption, the statistical weight remains to be a single combinatoric number, so one can count the states by thinking of the particles effectively either as bosons or as fermions, with the effective number of available single-particle states being linearly dependent on the particle number:
$$G_{eff}^{(b)}=G\lambda (N1),\mathrm{or}G_{eff}^{(f)}=G(1\lambda )(N1).$$
(9)
Obviously, for genuine bosons (or fermions), $`G_{eff}^{(b)}`$ (or $`G_{eff}^{(f)}`$) is independent of the particle number. In all other cases, either of the two $`G_{eff}`$ is linearly dependent on the particle number. This is the defining feature of the FES. The statistics parameter $`\lambda `$ tells us, on the average, how many single-particle states that a particle can exclude others to occupy. A proper understanding of this has been discussed in . Thus, the expression (8) for the statistical weight, $`W`$, formulates a generalized Pauli exclusion principle, as first recognized by Haldane .
It is easy to generalize this state counting to more than one species, labeled by the index $`i`$:
$$W=\underset{i}{}\frac{[G_i+N_i1\underset{j}{}\lambda _{ij}(N_j\delta _{ij})]!}{(N_i)![G_i1_j\lambda _{ij}(N_j\delta _{ij})]!}.$$
(10)
Here $`G_i`$ is the number of states when the system consists of only a single particle of species $`i`$. By definition, the diagonal $`\lambda _{ii}`$ is the “self-exclusion” statistics of species $`i`$, while the non-diagonal $`\lambda _{ij}`$ (for $`ij`$) is the mutual-exclusion statistics. Note that $`\lambda _{ij}`$, which Haldane called statistical interactions, may be asymmetric in $`i`$ and $`j`$. The interpretation is similar to that of the one-species case: The number of available single-particle states for species $`i`$, in the presence of other particles, is again linearly dependent on particle numbers of all species:
$`G_{eff,i}^{(b)}=G_i{\displaystyle \underset{j}{}}\lambda _{ij}(N_j\delta _{ji}),`$ (11)
$`\mathrm{or}`$ (12)
$`G_{eff,i}^{(f)}=G_{eff,i}^{(b)}+N_i1.`$ (13)
The definition (8) or (10) starts with a postulated form for the statistical weight, and thus is more direct and convenient for the purpose of formulating quantum statistical mechanics. One of us has first formulated the quantum statistical mechanics by proposing the notion of generalized ideal gas(GIG): A GIG satisfies the following two conditions: (i) The total energy (eigenvalue) is always of the form of a simple sum, in which the $`i`$-th term is linear in the particle number $`N_i`$:
$$E=\underset{i}{}N_i\epsilon _i^0,$$
(14)
with $`\epsilon _i^0`$ identified as the energy of a particle of species $`i`$; (ii) The state-counting (10) for statistical weight $`W`$ is applicable. When there are no statistical interactions (i.e., $`\lambda _{ij}=0`$ for $`ij`$), we have the usual ideal gas, which we call as IEG.
With the assumptions (14) and (10), the thermodynamics of a GIG can be worked out by the usual techniques in statistical mechanics. Consider a grand canonical ensemble at temperature $`T`$ and with chemical potential $`\mu _i`$ for species $`i`$, whose partition function is given by
$$Z=\underset{\{N_i\}}{}W(\{N_i\})\mathrm{exp}\{\underset{i}{}N_i(\mu _i\epsilon _i^0)/T\}.$$
(15)
As usual, we expect that for very large $`N_i`$, the summation has a very sharp peak around the set of most-probable (or mean) particle numbers $`\{\overline{N}_i\}`$. Using the Stirling formula, introducing the average “occupation number per state” defined by $`n_i\overline{N}_i/G_i`$, and maximizing
$$\frac{}{n_i}\left[\mathrm{ln}W+\underset{i}{}G_in_i(\mu _i\epsilon _i^0)/T\right]=0,$$
(16)
one obtains the equations that determine the most-probable distribution of $`n_i`$
$$\underset{j}{}(\delta _{ij}w_j+g_{ij})n_j=1,$$
(17)
with $`g_{ij}\lambda _{ij}G_j/G_i`$, and $`w_i`$ being determined by the functional equations
$$(1+w_i)\underset{j}{}\left(\frac{w_j}{1+w_j}\right)^{\lambda _{ji}}=e^{(\epsilon _i^0\mu _i)/T}.$$
(18)
The thermodynamic potential $`\mathrm{\Omega }=T\mathrm{ln}Z`$ and the entropy $`S`$ are then given by
$`\mathrm{\Omega }`$ $``$ $`PV=T{\displaystyle \underset{i}{}}G_i\mathrm{log}{\displaystyle \frac{1+n_i\underset{j}{}g_{ij}n_j}{1_jg_{ij}n_j}}`$ (19)
$`=`$ $`T{\displaystyle \underset{i}{}}G_i\mathrm{ln}(1+w_i^1);`$ (20)
$`S`$ $`=`$ $`{\displaystyle \underset{i}{}}G_i\left\{n_i{\displaystyle \frac{\epsilon _i^0\mu _i}{T}}+\mathrm{ln}{\displaystyle \frac{1+n_i\underset{j}{}g_{ij}n_j}{1_jg_{ij}n_j}}\right\}`$ (21)
$`=`$ $`{\displaystyle \underset{i}{}}G_i\left\{n_i{\displaystyle \frac{\epsilon _i^0\mu _i}{T}}+\mathrm{ln}(1+w_i^1)\right\}.`$ (22)
Other thermodynamic functions follow straightforwardly. As usual, one can easily verify that the fluctuations, $`(\overline{N_{i}^{}{}_{}{}^{2}}\overline{N_i}^2)/\overline{N_i}^2`$, of the occupation numbers are negligible, which justifies the validity of the above approach.
## IV Bosonization of 1-d ideal excluson gas
Let us first consider the simplest case, the 1-d IEG without internal degrees of freedom. We expect to obtain a continuous interpolation between the usual ideal Bose and ideal Fermi gas. Moreover, we want to show that the low-energy behavior of the IEG reproduces that of the Luttinger liquid and, therefore, provides a better characterization of the infrared fixed points associated with the Luttinger liquid.
### A Ideal Excluson Gas
Consider a GIG of $`N`$ particles on a ring with size $`L`$. Single-particle states are labeled by pseudo-momenta $`k_i`$. The total energy and momentum are given by
$$E=k_i^2,P=k_i.$$
(23)
According to (13), in the thermodynamic limit the hole density, $`\rho _a(k,T)`$, (or the density of available single-particle states) is linearly dependent on the particle density, $`\rho (k,T)`$. By definition, the statistics interaction matrix is given by
$$\lambda (k_i,k_j)=\frac{\mathrm{\Delta }\rho _a(k_i)}{\mathrm{\Delta }\rho (k_j)}.$$
(24)
Or in the thermodynamic limit, one has
$$\lambda (k,k^{})=\delta \rho _a(k)/\delta \rho (k^{}).$$
(25)
The system is called an IEG of statistics $`\lambda `$ (with no mutual statistics between different momenta), if
$$\lambda (k_i,k_j)=\lambda \delta (k_ik_j),$$
(26)
or (13) reads
$$\rho (k_j)=\frac{1}{2\pi }+\frac{1}{L}(1\lambda )\underset{ij}{}\delta (k_jk_i)\rho (k_i)\mathrm{\Delta }k,$$
(27)
which, in the thermodynamic limit, can be simply written as
$$\rho _a(k,T)+\lambda \rho (k,T)=\rho _0(k,T),$$
(28)
where $`\rho _0(k)1/2\pi `$ is the bare density of single-particle states. Thus, $`\lambda =1`$ corresponds to fermions, and $`\lambda =0`$ to bosons. The thermodynamic potential, now reads, in terms of (20)
$$\mathrm{\Omega }=\frac{T}{2\pi }_{\mathrm{}}^{\mathrm{}}𝑑k\mathrm{ln}(1+w(k,T)^1),$$
(29)
with the function $`w(k,T)\rho _a(k)/\rho (k)`$ satisfying an algebraic equation,
$$w(k,T)^\lambda [1+w(k,T)]^{1\lambda }=e^{(k^2\mu )/T}.$$
(30)
Firstly, we consider the ground state, in which the particles are distributed in a finite and origin-symmetric interval in the pseudo-momentum space. The (pseudo-)Fermi momentum is defined by
$$k_F^2=\mu $$
(31)
and its value is fixed by the average particle density $`\overline{d}_0=N_0/L`$ in the ground state,
$$_{k_F}^{k_F}𝑑k\rho (k)=\overline{d}_0.$$
(32)
Because holes are absent in the ground state, the particle density in the ground state is easily obtained from (28),
$$\rho (k)=\{\begin{array}{cc}\frac{1}{2\pi \lambda },\hfill & \mathrm{for}|k|<k_F;\hfill \\ 0,\hfill & \mathrm{for}|k|>k_F.\hfill \end{array}$$
(33)
Hence, one has
$$k_F=\pi \lambda \overline{d}_0,\mu =(\pi \lambda \overline{d}_0)^2.$$
(34)
Then the ground state energy and momentum are given by
$`{\displaystyle \frac{E_0}{L}}={\displaystyle _{k_F}^{k_F}}𝑑k\rho (k)k^2={\displaystyle \frac{1}{3}}\pi ^2\lambda ^2\overline{d}_0^3,`$ (35)
$`P_0={\displaystyle _{k_F}^{k_F}}𝑑k\rho (k)k=0.`$ (36)
Now let us examine possible excitations in an IEG. First there are density fluctuations due to particle-hole excitations, i.e., sound waves with velocity (see the next subsection)
$$v_s=v_F2k_F.$$
(37)
Besides, by adding extra $`M`$ particles to the ground state, one can create particle excitations, and by Galileo boosts a persistent current. We observe that the velocities of these three classes of elementary excitations in IEG also satisfy the fundamental relation (1). Indeed, shifting $`N_0`$ to $`N=N_0+M`$, the change in the ground state energy is
$`\delta _ME_0`$ $`=`$ $`{\displaystyle \frac{1}{3}}\pi ^2\lambda ^2(N/L)^3{\displaystyle \frac{1}{3}}\pi ^2\lambda ^2(N_0/L)^3`$ (38)
$`=`$ $`\pi ^2\overline{d}_0^2M+\pi (\lambda k_F)M^2+O(M^3/L^3),`$ (39)
while a persistent current, created by the boost of the Fermi sea $`kk+\pi J/L`$, leads to the energy shift
$`\delta _JE_0`$ $`=`$ $`{\displaystyle _{k_F+\pi J/L}^{k_F+\pi J/L}}𝑑k\rho (k)k^2{\displaystyle _{k_F}^{k_F}}𝑑k\rho (k)k^2`$ (40)
$`=`$ $`\pi (k_F/\lambda )J^2.`$ (41)
Therefore the total change in energy, due to charge and current excitations, is
$$\delta E_0\mu M=\frac{\pi }{2L}v_F(\lambda M^2+\lambda ^1J^2).$$
(42)
The total momentum change due to the current excitations is
$$\delta P_0=\underset{k}{}\frac{\pi J}{L}=\pi (\overline{d}_0+\frac{M}{L})J.$$
(43)
If we denote the variation in free energy as $`\delta F_0=\delta E_0\mu M`$, and identify $`\lambda `$ as the controlling parameter $`e^{2\phi }`$ in the Luttinger liquid theory, (42) just recuperates the zero-mode contributions in (4). Comparing (42) with (4) we identify the velocities $`v_N`$ and $`v_J`$ to be
$$v_N=v_F\lambda ,v_J=v_F/\lambda ,$$
(44)
Then we see the velocity relation (1), i.e., $`v_s=\sqrt{v_Nv_J}`$, that Haldane used to characterize the Luttinger liquids, is satisfied in IEG. The selection rule (5) also holds for the IEG, since the system should correspond to the ideal Fermi gas if $`\lambda =1`$.
Encouraged by this relationship between the IEG and Luttinger liquids, we want to calculate the critical exponents of IEG to see whether they reproduce those of the Luttinger liquids. This motivates to develop a bosonization for the density fluctuations in IEG.
### B Low Energy Limit and Bosonization
Following Yang and Yang, we introduce the dressed energy $`ϵ(k,T)`$ by writing
$$w(k,T)=e^{ϵ(k,T)/T}.$$
(45)
The point is that the grand partition function $`Z_G`$, corresponding to the thermodynamic potential (29), is of the form of that for an ideal system of fermions with a complicated, $`T`$-dependent energy dispersion given by the dressed energy:
$$Z_G=\underset{k}{}(1+e^{ϵ(k,T)/T}).$$
(46)
However, this fermion representation is not very useful, because of the implicit $`T`$-dependence of the dressed energy. To simplify, we consider the low-$`T`$ limit. By using the dressed energy, (30) reads
$$ϵ(k,T)=k^2\mu T(1\lambda )\mathrm{ln}(1+e^{ϵ(k,T)/T}).$$
(47)
Because there is no singularity in $`ϵ(k,T)`$ at $`T=0`$, the zero temperature dressed energy is given by
$$ϵ(k)=\{\begin{array}{cc}(k^2k_F^2)/\lambda ,\hfill & |k|<k_F,\hfill \\ k^2k_F^2,\hfill & |k|>k_F.\hfill \end{array}$$
(48)
Denote
$$ϵ(k,T)=ϵ(k)+\stackrel{~}{ϵ}(k,T),$$
(49)
where
$$\stackrel{~}{ϵ}(k,0)=0.$$
(50)
In the low-$`T`$ limit, one has
$$ϵ(k,T)=\{\begin{array}{cc}\frac{k^2\mu }{\lambda }(\lambda ^11)T\mathrm{ln}(1+e^{|ϵ(k)|/T}),\hfill & |k|<k_F,\hfill \\ (k^2\mu )(1\lambda )T\mathrm{ln}(1+e^{|ϵ(k)|/T}),\hfill & |k|>k_F,\hfill \end{array}.$$
(51)
Hence,
$$\stackrel{~}{ϵ}(k)=\{\begin{array}{cc}(1\lambda ^1)T\mathrm{ln}(1+e^{|ϵ(k)|/T}),\hfill & |k|<k_F,\hfill \\ (\lambda 1)T\mathrm{ln}(1+e^{|ϵ(k)|/T}),\hfill & |k|>k_F.\hfill \end{array}$$
(52)
For low energies, one can consider only the excitations around the Fermi surface,
$`{\displaystyle \frac{\mathrm{\Omega }(T)}{L}}`$ $``$ $`{\displaystyle \frac{T}{2\pi }}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑k\mathrm{ln}(1+e^{ϵ(k)/T})`$ (53)
$`+`$ $`{\displaystyle \frac{1}{2\pi }}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑k{\displaystyle \frac{\stackrel{~}{ϵ}(k,T)}{1+e^{ϵ(k)/T}}}`$ (54)
$``$ $`{\displaystyle \frac{1}{2\pi }}{\displaystyle _{k_F}^{k_F}}𝑑kϵ(k){\displaystyle \frac{T}{\pi \lambda }}{\displaystyle _{k_F\delta }^{k_F}}𝑑k\mathrm{ln}(1+e^{|ϵ(k)|/T})`$ (56)
$`{\displaystyle \frac{T}{\pi }}{\displaystyle _{k_F}^{k_F+\delta }}𝑑k\mathrm{ln}(1+e^{|ϵ(k)|/T}),`$
where the first term on the right hand side of the last equality is recognized as $`\mathrm{\Omega }(0)/L`$. The cut-off $`\delta `$ is of order $`O(T/v_s)`$ (actually, a few times of $`T/v_s`$). Mathematically, we take the limit of $`T0`$ followed by $`\delta 0`$. Using the integral formula
$$_0^{\mathrm{}}𝑑x\mathrm{ln}(1+e^x)=\frac{\pi ^2}{12},$$
we have the low-$`T`$ thermodynamic potential
$$\frac{\mathrm{\Omega }(T)}{L}\frac{\mathrm{\Omega }(0)}{L}=\frac{\pi T^2}{6v_s},$$
(57)
which implies that the theory is cut-off independent at low temperature. Notice that $`F=\mathrm{\Omega }\mu N`$. Because we only consider the particle-hole excitations near the Fermi surface contribute to thermal excitations, $`N(T)N(0)=0`$, which can be checked by an explicit calculation in terms of the definition of $`\rho (k,T)`$. Thus, we have
$$\frac{F(T)}{L}\frac{F(0)}{L}=\frac{\mathrm{\Omega }(T)}{L}\frac{\mathrm{\Omega }(0)}{L}=\frac{\pi T^2}{6v_s}.$$
(58)
This means that the low energy behavior of the IEG is controlled by a $`c=1`$ CFT. This result can be verified by a finite-size scaling in the spatial direction,
$$\frac{F_L(0)}{L}\frac{F(0)}{L}=\frac{\pi v_s}{6L^2},$$
(59)
where $`F_L(0)`$ is the zero temperature free energy for a system with size $`L`$. (For details, see .)
The above relation agrees with the finite-size scaling of a conformally invariant system with central charge $`c=1`$. So we want to see whether the low-energy effective theory of the IEG is really a CFT. Let us start with the grand partition function (46). At low temperature, the solution (51) leads to $`\stackrel{~}{ϵ}(k,T)=O(Te^{|ϵ|/T})`$, so one can simply replace $`ϵ(k,T)`$ with $`ϵ(k)`$ in the grand partition function:
$$Z_G\underset{k}{}(1+e^{\beta ϵ(k)}).$$
(60)
Note that the dressed energy with $`k`$ outside the Fermi points $`\pm k_F`$ has a slope different from that with $`k`$ inside $`\pm k_F`$. The former is $`\frac{2\pi }{L}`$ and the latter $`\frac{2\pi \lambda }{L}`$. It is necessary to keep this in mind for writing down the correct ground state wave functions of the excluson system. Now that the dispersion $`ϵ(k)`$ is $`T`$-independent, the grand partition function in the low-$`T`$ limit can be expressed in a fermionic representation as
$$Z_G=\mathrm{Tr}e^{\beta H_{\mathrm{eff}}},$$
(61)
where the effective Hamiltonian is given by
$$H_{\mathrm{eff}}=\underset{k}{}ϵ(k)c_k^{}c_k,$$
(62)
where $`c_k^{}`$ are fermionic creation operators. We also see that $`ϵ(k_F)=0`$, which can be used to define the Fermi momentum.
Physically, it is the phonon excitations that dominate the low-energy behavior of the system. In the low-$`T`$ limit, it is enough to consider the density fluctuations only near the Fermi points, $`k\pm k_F`$, where the left- and right-moving sectors are separable and decoupled:
$$H_{\mathrm{eff}}=H_++H_{}.$$
(63)
Besides this, another important simplification for excitations near Fermi points in the low-$`T`$ limit is that their energy, $`H_\pm `$, has a linearized dispersion:
$$ϵ_\pm (k)=\{\begin{array}{cc}\pm v_F(kk_F),\hfill & |k|>k_F,\hfill \\ \pm v_F(kk_F)/\lambda ,\hfill & |k|<k_F.\hfill \end{array}$$
(64)
We note the ‘refractions’ at $`k=\pm k_F`$, which implies to create a particle with pseudo-momentum $`k`$ and to create a hole with $`k^{}`$ cost different energies, even if $`|kk_F|=|k^{}k_F|`$. The reason for this is that $`k`$ is not the actual momentum carried by $`c_k^{}`$, as we will see soon.
The key thing for bosonization is to construct a density fluctuation operator. Taking into account the different slopes for dressed energy inside and outside the Fermi points, the density fluctuation operator at $`kk_F`$ is constructed as follows:
$`\rho _q^{(+)}={\displaystyle \underset{k>k_F}{}}:c_{k+q}^{}c_k:+{\displaystyle \underset{k<k_F\lambda q}{}}:c_{k+\lambda q}^{}c_k:`$ (65)
$`+{\displaystyle \underset{k_F\lambda q<k<k_F}{}}:c_{\frac{kk_F}{\lambda }+k_F+q}^{}c_k:`$ (66)
for $`q>0`$. A similar density operator $`\rho _q^{()}`$ can also be defined at $`kk_F`$,
$`\rho _q^{()}={\displaystyle \underset{k<k_F}{}}:c_{kq}^{}c_k:+{\displaystyle \underset{k>k_F+\lambda q}{}}:c_{k\lambda q}^{}c_k:`$ (67)
$`+{\displaystyle \underset{k_F+\lambda q>k>k_F}{}}:c_{\frac{k+k_F}{\lambda }k_Fq}^{}c_k:`$ (68)
To define the normal ordering we write, e.g.,
$$c_k=\{\begin{array}{cc}c_k,\hfill & k>k_F,\hfill \\ d_k^{},\hfill & k<k_F,\hfill \end{array}$$
(69)
where $`d_k^{}`$ is understood as a creation operator of a hole. Then normal ordering is done as usual: putting the annihilation operators to the right of the creation ones. Hence we have, e.g.,
$`\rho _q^{(+)}={\displaystyle \underset{k>k_F}{}}:c_{k+q}^{}c_k:+{\displaystyle \underset{k<k_F\lambda q}{}}:d_{k+\lambda q}d_k^{}:`$ (70)
$`+{\displaystyle \underset{k_F\lambda q<k<k_F}{}}:c_{\frac{kk_F}{\lambda }+k_F+q}^{}d_k^{}:`$ (71)
Within the Tomonaga approximation , in which commutators are taken to be their ground-state expectation value, we obtain
$`[\rho _q^{(\pm )},\rho _q^{}^{(\pm )}]0|[\rho _q^{(\pm )},\rho _q^{(\pm )}]|0`$ (72)
$`={\displaystyle \underset{k_F\lambda q<k<k_F}{}}0|c_{k+\lambda q}c_{k+\lambda q^{}}^{}|0`$ (73)
$`=\delta _{q,q^{}}{\displaystyle \underset{k_F\lambda q<k<k_F}{}}1={\displaystyle \frac{L}{2\pi }}q\delta _{q,q^{}}`$ (74)
Also, the commutators between $`H_{\mathrm{eff}}`$ and $`\rho _q^{(\pm )}`$ are
$$[H_\pm ,\rho _q^{(\pm )}]0|[H_\pm ,\rho _q^{(\pm )}]|0=\pm v_Fq\rho _q^{(\pm )}.$$
(75)
(74) and (75) describe 1-d free phonons with the sound velocity $`v_s=v_F`$ (so we have proved (37)). Introducing normalized boson annihilation operators
$$b_q=\sqrt{2\pi /qL}\rho _q^{(+)},\stackrel{~}{b}_q=\sqrt{2\pi /qL}\rho _q^{()}$$
(76)
and adding back the zero mode contributions, the bosonized Hamiltonian satisfying (74) is given by
$$H_B=v_s\{\underset{q>0}{}q(b_q^{}b_q+\stackrel{~}{b}_q^{}\stackrel{~}{b}_q)+\frac{1}{2}\frac{\pi }{L}[\lambda M^2+\frac{1}{\lambda }J^2]\},$$
(77)
which agrees with the bosonized Hamiltonian (4) in the Luttinger liquid theory.
In passing, we make a comment on linearization of the dressed energy dispersion. When we did this, we changed the ground state energy, because we assumed that for all $`k`$ the spectrum is linear in $`k`$. However, we changed neither the ground state wave function, nor the low-$`T`$ physics. On the other hand, the linearized spectrum was valid only for phonon excitations, it has nothing to do with the zero-mode excitations. So, after the linearized phonon part of the Hamiltonian is bosonized, we had to add back the zero-mode excitations.
The construction of the bosonized momentum operator is a bit more tricky, because $`c_k^{}`$ does not carry a momentum $`k`$. Each term in (66) should carry the same momentum $`q`$, therefore the fermion created by $`c_k^{}`$ carries a dressed momentum $`p`$, which is related to $`k`$ by
$`p(k)=\{\begin{array}{ccc}kk_F+(k_F/\lambda ),\hfill & k>k_F,\hfill & \\ k/\lambda ,\hfill & |k|<k_F,.\hfill & \\ k+k_F(k_F/\lambda ),\hfill & k<k_F.\hfill & \end{array}`$ (81)
In terms of this variable, the linearized dressed energy $`ϵ(p)`$ is of a simple form: $`ϵ_\pm (p)=\pm v_s(pp_F)`$, with $`p_F=k_F/\lambda `$. The bosonized total momentum operator, corresponding to the fermionized $`P=_kp(k)c_k^{}c_k`$, is
$$P=\underset{q>0}{}q(b_q^{}b_q\stackrel{~}{b}_q^{}\stackrel{~}{b}_q)+\pi (\overline{d}_0+M/L)J.$$
(82)
We see that the fundamental velocity relation, the bosonized Hamiltonian and momentum, and the selection rule of the quantum numbers in the Luttinger liquid theory can all be reproduced in IEG if we identify
$$\lambda e^{2\phi }.$$
(83)
To say that IEG can be used to characterize the renormalization group fixed points of Luttinger liquids, we still need to check the conformal invariance of the bosonized theory of IEG, and to verify the critical properties of IEG reproduce those of the Luttinger liquids.
### C Effective Field Theory and Conformal Invariance
To check conformal invariance, we need to rewrite the above bosonized effective Halmitonian (77) into a form of field theory in coordinate space. Employing the Fourier transformation, the density operator can be written as
$`\rho (x)`$ $`=`$ $`\rho _R(x)+\rho _L(x),`$ (84)
$`\rho _R(x)`$ $`=`$ $`{\displaystyle \frac{M_R}{L}}+{\displaystyle \underset{q>0}{}}\sqrt{{\displaystyle \frac{q}{2\pi L\lambda }}}(e^{iqx}b_q+e^{iqx}b_q^{}),`$ (85)
$`\rho _L(x)`$ $`=`$ $`{\displaystyle \frac{M_L}{L}}+{\displaystyle \underset{q>0}{}}\sqrt{{\displaystyle \frac{q}{2\pi L\lambda }}}(e^{iqx}\stackrel{~}{b}_q+e^{iqx}\stackrel{~}{b}_q^{}),`$ (86)
where $`M_{R,L}`$ are given by $`M=M_R+M_L`$ and $`\stackrel{~}{b}_q=b_q`$ for $`q>0`$.
The boson field $`\varphi (x)`$, which is conjugated to $`\rho (x)`$ and satisfies
$$[\varphi (x),\rho (x^{})]=i\delta (xx^{}),$$
(87)
is given by
$`\varphi (x)`$ $`=`$ $`\varphi _R(x)+\varphi _L(x),`$ (88)
$`\varphi _R(x)`$ $`=`$ $`{\displaystyle \frac{\varphi _0}{2}}+{\displaystyle \frac{\pi J_Rx}{L}}+i{\displaystyle \underset{q>0}{}}\sqrt{{\displaystyle \frac{\pi \lambda }{2qL}}}(e^{iqx}b_qe^{iqx}b_q^{}),`$ (89)
$`\varphi _L(x)`$ $`=`$ $`{\displaystyle \frac{\varphi _0}{2}}+{\displaystyle \frac{\pi J_Lx}{L}}+i{\displaystyle \underset{q>0}{}}\sqrt{{\displaystyle \frac{\pi \lambda }{2qL}}}(e^{iqx}\stackrel{~}{b}_qe^{iqx}\stackrel{~}{b}_q^{}),`$ (90)
with $`J=J_R+J_L`$. We have to assign the quantum numbers such that there are only two independent each other in $`M_{R,L}`$ and $`J_{R,L}`$. A consistent choice is
$$M_R=J_R,M_L=J_L.$$
(91)
Then,
$$J=J_R+J_L,M=J_RJ_L.$$
(92)
Here $`\varphi _0`$ is an angular variable conjugated to $`M`$: $`[\varphi _0,M]=i`$. The Hamiltonian (4) becomes
$$H=\frac{1}{2}_0^L𝑑x[\pi v_N\rho (x)^2+v_J/\pi (_x\varphi (x))^2],$$
(93)
or by a field rescaling,
$$H=\frac{v_s}{2\pi }_0^L𝑑x[\mathrm{\Pi }(x)^2+(_xX(x))^2],$$
(94)
where
$$\mathrm{\Pi }(x)=\pi \lambda ^{1/2}\rho (x),X(x)=\lambda ^{1/2}\varphi (x).$$
(95)
With $`X(x,t)=e^{iHt}X(x)e^{iHt}`$, the Lagrangian density reads
$$=\frac{v_s}{2\pi }_\alpha X(x,t)^\alpha X(x,t).$$
(96)
This is the Lagrangian density of a free scalar field theory in $`1+1`$-dimensions. Writing the corresponding operators as the functionals of the scalar field $`X(x,t)`$, all correlation functions can be obtained by using the propagators of $`X_R(x,t)`$ and $`X_L(x,t)`$,
$`X_R(x,t)X_R(0,0)={\displaystyle \frac{1}{4}}\mathrm{ln}(xv_st),`$ (97)
$`X_L(x,t)X_L(0,0)={\displaystyle \frac{1}{4}}\mathrm{ln}(x+v_st).`$ (98)
The statistics of an operator in the theory can also be inferred by the commutators of the scalar fields,
$$[X_{R,L}(x),X_{R,L}(x^{})]=\pm \frac{i\pi }{4}\theta (xx^{}).$$
(99)
We recognize that $``$ (96) is the Lagrangian of a $`c=1`$ CFT, consistent with the finite-size scaling (58). Alternatively, it is easy to check that the theory is invariant under the conformal transformations generated by a set of the Virasoro generators
$$L_m=\frac{1}{2}\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}\alpha _{nm}\alpha _n,\stackrel{~}{L}_m=\frac{1}{2}\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}\stackrel{~}{\alpha }_{nm}\stackrel{~}{\alpha }_n,$$
(100)
where the oscillators $`\alpha _m=m^{1/2}b_q`$ and $`\stackrel{~}{\alpha }_m=m^{1/2}\stackrel{~}{b}_q^{}`$ for $`m=qL/2\pi >0`$ being integers. $`\alpha _0=(\pi /2L)^{1/2}[J\lambda ^{1/2}M\lambda ^{1/2}]`$ and $`\stackrel{~}{\alpha }_0=(\pi /2L)^{1/2}[J\lambda ^{1/2}+M\lambda ^{1/2}]`$. The generators obey the Virasoro algebra with the central charge $`c=1`$,
$$[L_m^{\mathrm{tot}},L_n^{\mathrm{tot}}]=(mn)L_{m+n}^{\mathrm{tot}}+\frac{1}{12}(m^3m)\delta _{m+n,0},$$
(101)
with $`L_m^{\mathrm{tot}}=L_m+\stackrel{~}{L}_m`$.
Since $`\varphi _0`$ is an angular variable, there is a hidden invariance in the theory under $`\varphi \varphi +2\pi `$. The field $`X`$ is thus said to be “compactified” on a circle, with a radius that is determined by the exclusion statistics :
$$XX+2\pi R,R^2=1/\lambda .$$
(102)
Noting the selection rule (5), the Hamiltonian has a duality
$$\lambda 1/\lambda ,MJ,$$
(103)
which has referred to the particle-hole duality . Using the CFT terminology, this duality is represented as the duality of the compactified radii,
$$R1/R.$$
(104)
We note that this is different from the duality relation $`R2/R`$, in the usual $`c=1`$ CFT compactifed on a circle. Actually, according to the standard terminology in CFT , our selection rule (5) and duality relation (104) make what we obtained above a $`c=1`$ CFT compactified on an orbifold $`S^1/Z_2`$, i.e., a circle folded by a reflection about a diameter, which topologically is a semi-circle or an interval. This difference can also be seen from the grand partition function: Using the identification between $`H_{\mathrm{eff}}`$ and $`L_0^{\mathrm{tot}}`$, i.e., $`H_{\mathrm{eff}}=v_sL_0^{\mathrm{tot}}`$, the grand partition function of IEG (in the low-$`T`$ limit) can be rewritten as
$$Z_G=Tr_{}[q^{L_0}\overline{q}^{\stackrel{~}{L}_0}],$$
(105)
where $`q=e^{iv_s\tau }`$ with $`\tau =i\beta =i/T`$. Thus, the selection rule (5) severely constrain the allowed values for the eigenvalues of $`L_0`$ and $`\stackrel{~}{L}_0`$. It makes the CFT we obtained have an unusual spectrum and duality relation, corresponding to the $`c=1`$ orbifold CFT . In next subsection we will see that because of the difference in the selection rules, the statistics of the allowed charge-1 operators in the two classes of CFT’s are not the same.
We note that a similar situation happens for the CFT that describes the low-$`T`$ behavior of the Calogero-Sutherland (C-S) model . This model has two different versions, with the long-range interactions being among bosons or among fermions, respectively. At low temperature, the two versions have different selection rules for the zero-mode quantum numbers, thus leading to different CFT’s: The low-$`T`$ CFT for the bosonic C-S model is the usual $`c=1`$ CFT compactified on a circle, which has been studied extensively in the literatures ; while for the fermionic C-S model the low-$`T`$ limit gives rise to the $`c=1`$ orbifold CFT. This is because the selection rule for zero modes severely constrains the spectrum of the system, i.e., possible quantum numbers of the allowed excitations. (For details, see ref. .) Thus, only the fermionic (not the bosonic) C-S model respects a duality relation $`\lambda 1/\lambda `$ that coincides with the particle-hole duality in IEG .
### D Correlation Functions
The CFT description of the IEG offers a better understanding for the space of quantum states in the theory. States $`V[X]|0`$ or operators $`V[X]`$ are allowed only if they respect the invariance (102),
$$V[X+2\pi R]V[X],$$
(106)
with a given boundary condition restriction. Here, a Fermi or a Bose operator obeys the periodic boundary condition (PBC). So quantum numbers of quasiparticles are strongly constrained, in particular by the selection rule for zero-mode quantum numbers. For example, the primary fields obeying the PBC in the CFT are given by
$`\varphi _{M,J}(x)`$ $``$ $`f(J,X^0):e^{i(M\lambda ^{1/2}+J/\lambda ^{1/2})X_R(x)}`$ (107)
$`\times `$ $`e^{i(M\lambda ^{1/2}J/\lambda ^{1/2})X_L(x)}:,`$ (108)
$`f(J,X^0)`$ $`=`$ $`e^{iJ(\lambda ^{1/2}\lambda ^{1/2})X_R^0}e^{iJ(\lambda ^{1/2}\lambda ^{1/2})X_L^0}`$ (109)
where the prefactor $`f(J,X^0)`$ makes the fields satisfy the PBC, $`M`$ and $`J`$ eigenvalues of the number and current operators, and $`X^0=\pi Mx/L`$. The field carries the charge $`M`$ and current $`J`$. The conformal dimensions of the fields are
$`h`$ $`=`$ $`{\displaystyle \frac{1}{2}}[(M\lambda ^{1/2}+J/\lambda ^{1/2})^2+(M\lambda ^{1/2}J/\lambda ^{1/2})^2]`$ (110)
$`=`$ $`M^2\lambda +J^2\lambda ^1.`$ (111)
The statistics of the field can be calculated by using (99) and the statistics factors are
$`\mathrm{exp}\{i{\displaystyle \frac{\pi }{4}}[(M\lambda ^{1/2}+J/\lambda ^{1/2})^2(M\lambda ^{1/2}J/\lambda ^{1/2})^2]\}`$ (112)
$`=`$ $`(1)^{MJ}.`$ (113)
Consider the charge-1 primary fields, with $`M=1`$. Therefore, they can only be fermions since $`J=`$ odd due to the selection rule. The general charge-1 fermion operator is a linear combination of the charge-1 primary fields. A careful construction of the allowed fermion field with unit charge leads to
$`\mathrm{\Psi }_F^{}(x,t)=\rho (x)^{1/2}`$ $`{\displaystyle \underset{m=\mathrm{}}{\overset{\mathrm{}}{}}}e^{iO_m}:e^{i(\lambda ^{1/2}+(2m+1)/\lambda ^{1/2})X_R(x_{})}:`$ (115)
$`:e^{i(\lambda ^{1/2}(2m+1)/\lambda ^{1/2})X_L(x_+)}:,`$
where the prefactor $`f`$ has been suppressed and the hermitian, constant-valued operators $`O_m`$ satisfy
$$[O_m,O_m^{}]=i\pi (mm^{}).$$
(116)
The multi-sector density operator is the linear combination of those primary fields with $`M=0`$ and $`J=`$even,
$`\widehat{\rho }(x)`$ $`=`$ $`\mathrm{\Psi }_F^{}(x)\mathrm{\Psi }_F(x)`$ (117)
$`=`$ $`\rho (x){\displaystyle }{}_{m}{}^{}:\mathrm{exp}\{i2m[X_R(x)X_L(x)]/\lambda ^{1/2}\}:.`$ (118)
All the secondary fields in the CFT follow by considering the sound wave contribution to the conformal weight of the fields.
The correlation functions can easily be calculated by using the CFT techniques. For examples, the density-density and single particle correlation functions are as follows,
$`\widehat{\rho }(x,t)\widehat{\rho }(0,0)`$ $``$ $`\overline{d}_0^2[1+{\displaystyle \frac{1}{(2\pi \overline{d}_0)^2\lambda }}({\displaystyle \frac{1}{x_R^2}}+{\displaystyle \frac{1}{x_L^2}})`$ (119)
$`+`$ $`{\displaystyle \underset{m=1}{\overset{\mathrm{}}{}}}A_m{\displaystyle \frac{1}{[x_Rx_L]^{m^2/\lambda }}}\mathrm{cos}(2\pi \overline{d}_0mx)],`$ (120)
and
$`G(x,t)\mathrm{\Psi }_F^{}(x,t)\mathrm{\Psi }_F(0,0)`$ (121)
$`\overline{d}_0{\displaystyle \underset{m=\mathrm{}}{\overset{\mathrm{}}{}}}`$ $`B_m{\displaystyle \frac{1}{x_R^{(\lambda ^{1/2}+(2m+1)\lambda ^{1/2})^2/4}}}`$ (124)
$`{\displaystyle \frac{1}{x_L^{(\lambda ^{1/2}(2m+1)\lambda ^{1/2})^2/4}}}`$
$`e^{i(2\pi (m+\lambda /2)\overline{d}_0x+\mu t)},`$
where $`x_{R,L}=xv_st`$ and $`A_m`$ and $`B_m`$ regularization-dependent constants.
Usually a physical quantity, e.g., a boson field, satisfies the periodic boundary conditions (PBC). Hence, a charge-1 bosonic excitations are not allowed in the theory, because it is anti-periodic. However, as we know, an anyon field needn’t to obey the PBC. So in the theory, there may be allowed anyonic excitations. A charge-1 anyonic (or exclusonic) operator is a primary field that does not obey the PBC,
$$\mathrm{\Psi }_\lambda ^{}(x)=:\mathrm{\Psi }_F^{}(x)e^{i(\lambda ^{1/2}\lambda ^{1/2})(X_R(x)X_L(x))}:.$$
(125)
The anyon commutation relation is easy to check:
$$\mathrm{\Psi }_\lambda ^{}(x)\mathrm{\Psi }_\lambda ^{}(x^{})e^{i\pi \lambda \mathrm{sgn}(xx^{})}\mathrm{\Psi }_\lambda ^{}(x^{})\mathrm{\Psi }_\lambda ^{}(x)=0,\mathrm{for}xx^{}.$$
(126)
In other words, the anyon field carries a fractional current. Or by the $`MJ`$-duality, the anyon with integer $`J`$ carries a fractional charge. The correlation function of the single-anyon reads
$`G(x,t;\lambda )`$ $``$ $`\mathrm{\Psi }_\lambda ^{}(x,t)\mathrm{\Psi }_\lambda (0,0)`$ (127)
$`\overline{d}_0{\displaystyle \underset{m=\mathrm{}}{\overset{\mathrm{}}{}}}`$ $`B_m^a`$ $`{\displaystyle \frac{1}{x_R^{(m+\lambda )^2/\lambda }}}{\displaystyle \frac{1}{x_L^{m^2/\lambda }}}e^{i(2\pi (m+\lambda /2)x+\mu t)},`$ (128)
This correlation function coincides with the asymptotic one in the Calogero-Sutherland model. We see that
(i) if $`\lambda =1`$, (128) consists with (124);
(ii) there are no boson excitations ($`\lambda =0`$) because $`G(x,t;0)=0`$;
(iii) and moreover, $`\lambda >0`$ is implied since (128) will diverge at the long distance if $`\lambda <0`$.
(iv) Look at $`m=0`$. The critical exponents can be reads out,
$$\eta _f=\lambda +\lambda ^1,\eta _\lambda =2\lambda .$$
Thus,
$$\begin{array}{cc}\eta _f>\eta _\lambda ,\hfill & \mathrm{if}\lambda <1;\hfill \\ \eta _f<\eta _\lambda ,\hfill & \mathrm{if}\lambda >1.\hfill \end{array}$$
(v) The multi-sector density operator for exclusons is the same as that of the fermion.
The single-hole state, i.e. $`\mathrm{\Psi }_{1/\lambda }^{}|0\mathrm{\Psi }_\lambda (\lambda \lambda ^1)|0`$, with charge $`1/\lambda `$ alone is not allowed. The minimum allowed multi-hole state is given by
$$\mathrm{\Psi }_{1/\lambda }^{}(x_1)\mathrm{}\mathrm{\Psi }_{1/\lambda }^{}(x_p)|0$$
if $`\lambda =p/q`$ is rational. One may obtain, e.g. ,
$$[\mathrm{\Psi }_{1/\lambda }^{}(x,t)]^p[\mathrm{\Psi }_{1/\lambda }(0,0)]^p[G(x,t;1/\lambda )]^p.$$
A more interesting allowed operator is what creates $`q`$ particle excitations accompanied by $`p`$ hole excitations:
$$\widehat{n}(x,t)=[\mathrm{\Psi }_\lambda ^{}(x,t)]^q[\mathrm{\Psi }_{1/\lambda }^{}(x,t)]^p.$$
We note the similarity of this operator to Read’s order parameter for fractional quantum Hall fluids (in bulk). Its correlation function can be calculated by using Wick’s theorem:
$$\widehat{n}(x,t)\widehat{n}(0,0)[G(x,t;\lambda )]^q[G(x,t;1/\lambda )]^p.$$
(129)
If the contribution from the $`m=0`$ sector dominates, then one gets
$$\widehat{n}(x,t)\widehat{n}(0,0)(xv_st)^{(p+q)}$$
.
## V Two Extensions
Now we proceed to go beyond IEG. Two extensions will be discussed in this section: The one-component GIG with the mutual statistics, and the non-ideal gas with the Luttinger-type interactions. In either case, we will show that the low-temperature behavior is that of an IEG, controlled by a single ”effective statistics” parameter $`\lambda _{eff}`$, whose value depends on the mutual statistics and the coupling constants in the interactions.
### A Generalized ideal gas with mutual statistics
We turn to discussing the effects of mutual statistics. Consider a GIG with the statistics matrix (25) in momentum space given by
$$g(kk^{})=\delta (kk^{})+\mathrm{\Phi }(kk^{}).$$
(130)
Here $`\mathrm{\Phi }(k)=\mathrm{\Phi }(k)`$ is a smooth function. $`\mathrm{\Phi }(kk^{})`$ stands for mutual statistics between particles with different momenta; for IEG $`\mathrm{\Phi }(k)=(\lambda 1)\delta (k)`$. The thermodynamic properties of GIG is given by eq. (29), but now $`w(k,T)`$ satisfies integral equation which, in terms of the dressed energy (45), is of the form
$$ϵ(k,T)=ϵ_0(k)+T_{\mathrm{}}^{\mathrm{}}\frac{dk^{}}{2\pi }\mathrm{\Phi }(kk^{})\mathrm{ln}(1+e^{ϵ(k^{},T)/T}),$$
(131)
where $`ϵ_0(k)k^2\mu `$. In the low-$`T`$ limit, it can be proven by the iteration that $`ϵ(k,T)=ϵ(k)+O(T^2/v_s)`$, where $`ϵ(k)`$ is the zero-temperature dressed energy given below. At $`T=0`$, the Fermi momentum $`k_F`$ is determined by
$$ϵ(\pm k_F)=0.$$
(132)
Introduce
$`(\alpha \beta )[k_F,k_F]`$ $``$ $`{\displaystyle _{k_F}^{k_F}}{\displaystyle \frac{dk}{2\pi }}\alpha (k)\beta (k),`$ (133)
$`(\mathrm{\Phi }\alpha )(k;k_F,k_F]`$ $``$ $`{\displaystyle _{k_F}^{k_F}}{\displaystyle \frac{dk^{}}{2\pi }}\mathrm{\Phi }(kk^{})\alpha (k^{}).`$ (134)
Then both $`\rho (k)`$ and $`ϵ(k)`$ in the ground state satisfy an integral equations like
$$\alpha (k)=\alpha _0(k)(\mathrm{\Phi }\alpha )(k;k_F,k_F].$$
(135)
The dressed momentum $`p(k)`$ is related to $`\rho (k)`$ by
$$dp(k)=2\pi \rho (k)dk,p(k)=p(k).$$
(136)
The ground state energy is given by
$$E_0/L=(ϵ_0\rho )[k_F,k_F].$$
(137)
Using the equation satisfied by $`\rho (k)`$, it can be expressed by the dressed energy
$$E_0/L=(ϵ\rho _0)[k_F,k_F].$$
(138)
The above equations are of the same form as those in the thermodynamic Bethe ansatz , hence the Luttinger-liquid relation , $`v_s=\sqrt{v_Nv_J}`$, remains true. A simple proof is sketched as follows. The sound velocity is well-known:
$$v_s=ϵ(p_F)/p_F.$$
(139)
The charge velocity is given by
$$v_N=v_sz(k_F)^2,$$
(140)
where the dressed charge $`z(k)`$ is given by the solution to the integral equation
$$z(k)=1(\mathrm{\Phi }z)(k;k_F,k_F].$$
(141)
This relation can be easily derived from the definitions
$`v_N`$ $`=`$ $`L\mu /N_0,`$ (142)
$`z(k)`$ $`=`$ $`\delta ϵ(k)/\delta \mu .`$ (143)
To create a persistent current, let us boost the Fermi sea by
$$\pm k_F\pm k_F+\mathrm{\Delta },$$
(144)
where $`\mathrm{\Delta }=z(k_F)/L\rho (k_F)`$. Then the total energy of the state with the persistent current is
$`E_\mathrm{\Delta }/L`$ $`=`$ $`(ϵ_0\rho _\mathrm{\Delta })[k_F+\mathrm{\Delta },k_F+\mathrm{\Delta }]`$ (145)
$`=`$ $`(ϵ_\mathrm{\Delta }\rho _0)[k_F+\mathrm{\Delta },k_F+\mathrm{\Delta }],`$ (146)
where
$$\rho _\mathrm{\Delta }(k)=\rho _0(k)(\mathrm{\Phi }\rho _\mathrm{\Delta })(k;k_F+\mathrm{\Delta },k_F+\mathrm{\Delta }]$$
(147)
and
$$ϵ_\mathrm{\Delta }(k)=ϵ_0(k)(\mathrm{\Phi }ϵ_\mathrm{\Delta })(k;k_F+\mathrm{\Delta },k_F+\mathrm{\Delta }].$$
(148)
Now, using the last expression for $`E_\mathrm{\Delta }`$ and substituting $`ϵ_\mathrm{\Delta }`$ in (V A), we have
$`E_\mathrm{\Delta }/L=(ϵ\rho _0)[k_F,k_F]+{\displaystyle \frac{\mathrm{\Delta }^2}{2}}ϵ^{}(k_F)`$ (149)
$`\{\rho _0(k_F)+(\rho _02\pi F)(k_F;k_F,k_F]\}`$ (150)
$`{\displaystyle \frac{\mathrm{\Delta }^2}{2}}ϵ^{}(k_F)\{\rho _0(k_F)+(\rho _02\pi F)(k_F;k_F,k_F]\}.`$ (151)
Here $`F(k,k^{})`$ is determined by
$$F(k,k^{})=\frac{1}{2\pi }\mathrm{\Phi }(k,k^{})\frac{1}{2\pi }_{k_F}^{k_F}𝑑k^{\prime \prime }\mathrm{\Phi }(k,k^{\prime \prime })F(k^{\prime \prime },k^{}).$$
(152)
On the other hand, we note that the equation for $`\rho _0(k)`$ can be rewritten as
$$\rho (k)=\rho _0(k)(\rho _02\pi F)(k;k_F,k_F].$$
(153)
Thus, we have
$$E_\mathrm{\Delta }E_0=L\mathrm{\Delta }^2ϵ^{}(k_F)\rho (k_F)=(2\pi /L)v_sz(k_F)^2.$$
(154)
This verifies $`v_J=v_sz(k_F)^2`$. In view of eq. (44), at low energies, the GIG looks like an IEG with
$$\lambda _{eff}=z(k_F)^2.$$
(155)
It can be shown that it is the effective statistics (155) that controls the low-$`T`$ critical properties of GIG, as $`\lambda `$ does for IEG. Linearization near the Fermi points and bosonization of the low-energy effective Hamiltonian go the same way as before for IEG. The only difference now is that the slope of the linearized dispersion for the dressed energy $`ϵ_\pm (k)=\pm ϵ^{}(k_F)(kk_F)+\mu =\pm v_s(p(k)p_F)+\mu `$, is smooth at $`k\pm k_F`$. So bosonization is standard and the bosonized Hamiltonian is the same as eq. (77) for IEG, only with $`\lambda `$ replaced by $`\lambda _{eff}`$. However, before going to the bosonization we need an effective Hamiltonian of the fermions with the dressed energy. Unlike the IEG, in the GIG case, $`ϵ(k,T)=ϵ(k)+O(T^2/v_s)`$. Now, we work out the $`T`$-expansion of $`ϵ(k,T)`$ explicitly in the low-$`T`$ limit:
$$ϵ(k,T)=ϵ(k)+\stackrel{~}{ϵ}(k,T)+O(T^3/v_s^2).$$
(156)
One finds that
$$\stackrel{~}{ϵ}(k,T)=\frac{\pi T^2}{6ϵ^{}(k_F)}f(k),$$
(157)
with the function $`f`$ determined by
$`f(k)`$ $`=`$ $`\mathrm{\Phi }(k_Fk)(\mathrm{\Phi }f)(k;k_F,k_F]`$ (158)
$`=`$ $`\mathrm{\Phi }(k_Fk)(\mathrm{\Phi }\mathrm{\Phi })(k;k_F,k_F]`$ (160)
$`+(\mathrm{\Phi }(\mathrm{\Phi }\mathrm{\Phi })(k;k_F,k_F]+\mathrm{}`$
Note that the equation that $`\rho (k)`$ obeys can be rewritten as
$`{\displaystyle \frac{\rho (k)}{\rho _0}}=1{\displaystyle _{k_F}^{k_F}}dk^{}\{\mathrm{\Phi }(kk^{})`$ (161)
$`+(\mathrm{\Phi }\mathrm{\Phi })(k^{};k_F,k_F](\mathrm{\Phi }(\mathrm{\Phi }\mathrm{\Phi })(k^{};k_F,k_F]+\mathrm{}`$ (162)
Integrating (160) over $`k`$ and comparing with (162), one has
$$_{k_F}^{k_F}\frac{dk}{2\pi }f(k)=1\frac{\rho (k_F)}{\rho _0},$$
(163)
and then
$$_{k_F}^{k_F}\frac{dk}{2\pi }\stackrel{~}{ϵ}(k,T)=\frac{\pi T^2}{6ϵ^{}(k_F)}(12\pi \rho (k_F)).$$
(164)
Substituting (156) into the thermodynamic potential (29), we have
$$\frac{\mathrm{\Omega }(T)}{L}=\frac{T}{2\pi }_{\mathrm{}}^{\mathrm{}}𝑑k\mathrm{ln}(1+e^{ϵ(k)/T})+\frac{1}{2\pi }_{\mathrm{}}^{\mathrm{}}\frac{dk}{1+e^{ϵ(k)/T}}\stackrel{~}{ϵ}(k,T).$$
(165)
In the low-$`T`$ limit, the first term in the last equation gives
$$\frac{\mathrm{\Omega }(0)}{L}\frac{\pi T^2}{6ϵ^{}(k_F)},$$
with
$$\frac{\mathrm{\Omega }(0)}{L}=\frac{1}{2\pi }_{k_F}^{k_F}𝑑kϵ(k).$$
and the second term is approximately given by (164). Thus,
$$\frac{\mathrm{\Omega }(T)}{L}\frac{\mathrm{\Omega }(0)}{L}=\frac{\pi T^2}{6v_s},$$
(166)
which proves the central charge $`c=1`$ CFT behavior of the theory at the low energy.
We may also confirm this from the finite size scaling in the spatial direction. To see this, we consider the discrete version of the equation in which the density $`\rho _L(k_i)`$ obeys
$$\rho _L(k_i)=\frac{1}{2\pi }\underset{ji}{}\mathrm{\Phi }(k_ik_j).$$
(167)
Using the relation between discrete sum and integration
$`{\displaystyle \frac{1}{L}}{\displaystyle \underset{n=N_1}{\overset{N_2}{}}}f({\displaystyle \frac{I_n}{L}})={\displaystyle _{(N_1+1/2)/L}^{(N_21/2)/L}}𝑑xf(x)`$ (168)
$`+{\displaystyle \frac{1}{24L^2}}[f^{}((N_11/2)/L)f^{}((N_2+1/2)/L)]`$ (169)
$`+O(1/L^3),`$ (170)
one has
$`\rho _L(k){\displaystyle \frac{1}{2\pi }}(\mathrm{\Phi }\rho _L)(k;k_F,k_F]`$ (171)
$`{\displaystyle \frac{1}{24L^2}}{\displaystyle \frac{1}{\rho (k_F)}}\left[{\displaystyle \frac{\mathrm{\Phi }(kk^{})}{dk^{}}}\right]_{k_F}`$ (172)
$`+{\displaystyle \frac{1}{24L^2}}{\displaystyle \frac{1}{\rho (k_F)}}[{\displaystyle \frac{\mathrm{\Phi }(kk^{})}{dk^{}}}]_{k_F}.`$ (173)
Denote
$$\rho _L=\rho +\rho _1,$$
(174)
where $`\rho (k)`$ is of the order $`O(1/L^0)`$ and $`\rho _1(k)`$ the order $`O(1/L^2)`$. Then, $`\rho (k)`$ is as defined and $`\rho _1(k)`$ is determined by
$`\rho _1(k)={\displaystyle \frac{1}{24L^2}}{\displaystyle \frac{1}{\rho (k_F)}}\left[{\displaystyle \frac{\mathrm{\Phi }(kk^{})}{dk^{}}}\right]_{k_F}`$ (175)
$`+{\displaystyle \frac{1}{24L^2}}{\displaystyle \frac{1}{\rho (k_F)}}\left[{\displaystyle \frac{\mathrm{\Phi }(kk^{})}{dk^{}}}\right]_{k_F}(\mathrm{\Phi }\rho _1)(k;k_F,k_F],`$ (176)
The corresponding thermodynamic potential reads
$`{\displaystyle \frac{\mathrm{\Omega }_L(0)}{L}}={\displaystyle \frac{1}{L}}{\displaystyle \underset{i}{}}ϵ_0(k({\displaystyle \frac{I_i}{L}}))`$ (177)
$`={\displaystyle _{k_F}^{k_F}}𝑑k\rho _L(k)ϵ_0(k)+{\displaystyle \frac{1}{24L^2\rho (k_F)}}[ϵ_0^{}(k)|_{k_F}ϵ_0^{}(k)|_{k_F}]`$ (178)
$`={\displaystyle _{k_F}^{k_F}}𝑑k\rho (k)ϵ_0(k)+{\displaystyle _{k_F}^{k_F}}𝑑k\rho _1(k)ϵ_0(k)`$ (179)
$`+{\displaystyle \frac{1}{24L^2\rho (k_F)}}[ϵ_0^{}(k)|_{k_F}ϵ_0^{}(k)|_{k_F}].`$ (180)
The first term of the last equation is $`\mathrm{\Omega }(0)/L`$ and the rest, using (176), can be written as
$`{\displaystyle \frac{1}{24L^2\rho (k_F)}}{\displaystyle \frac{}{k}}(ϵ_0(k)+(1)(\mathrm{\Phi }ϵ_0)`$ (181)
$`+(1)^2((\mathrm{\Phi }\mathrm{\Phi })ϵ_0)+\mathrm{})(k;k_F,k_F])_{k=k_F}`$ (182)
$`+{\displaystyle \frac{1}{24L^2\rho (k_F)}}{\displaystyle \frac{}{k}}(ϵ_0(k)+(1)(\mathrm{\Phi }ϵ_0)`$ (183)
$`+(1)^2((\mathrm{\Phi }\mathrm{\Phi })ϵ_0)+\mathrm{})(k;k_F,k_F])_{k=k_F}.`$ (184)
Recall the equation that $`ϵ(k)`$ obeys, one has immediately,
$$\frac{\mathrm{\Omega }_L(0)}{L}\frac{\mathrm{\Omega }(0)}{L}=\frac{\pi }{12L^2}\frac{ϵ^{}(k)_{k_F}ϵ^{}(k)_{k_F}}{2\pi \rho (k_F)}=\frac{\pi v_s}{6L^2}.$$
(185)
as desired.
Similar to the case of IEG, we also could have a fermion representation of the grand partition function with the temperature-dependent spectrum. To derive the low-energy effective theory, however, one rewrites the thermodynamic potential (165) in the low-$`T`$ limit as
$$\frac{\mathrm{\Omega }(T)}{L}\frac{\mathrm{\Omega }(0)}{L}2T\rho (k_F)I(k_F,T),$$
(186)
where
$`I(k_F,T)`$ $`=`$ $`{\displaystyle _{k_F\delta }^{k_F+\delta }}𝑑k\mathrm{ln}(1+e^{|ϵ(k)|/T})`$ (187)
$`=`$ $`{\displaystyle _{p_F\delta }^{p_F+\delta }}{\displaystyle \frac{dp}{2\pi }}\rho (p)\mathrm{ln}(1+e^{|ϵ(k((p))|/T}).`$ (188)
That is,
$`{\displaystyle \frac{\mathrm{\Omega }(T)}{L}}`$ $`=`$ $`{\displaystyle _{k_F}^{k_F}}{\displaystyle \frac{dk}{2\pi }}ϵ(k)`$ (189)
$``$ $`{\displaystyle \frac{T}{2\pi }}{\displaystyle _{p_F+\delta }^{p_F\delta }}𝑑p\mathrm{ln}(1+e^{|ϵ(k(p))|/T}),`$ (190)
where $`p`$ is the physical (dressed) momentum. The grand partition function reads
$$Z_G\underset{k^{}}{}(1+e^{\beta ϵ(k(k^{}))}),$$
(191)
where $`k^{}=k`$ for $`|k|<k_F\delta `$ and $`k^{}=p`$ for $`|k|>k_F\delta `$. Now, we can have an effective Hamiltonian because $`ϵ(k)`$ is $`T`$-independent,
$$H_{\mathrm{eff}}=\underset{k^{}}{}ϵ(k(k^{}))c_k^{}^{}c_k^{}.$$
(192)
Similar to the IEG case, the low-$`T`$ excitations can be considered by taking the linear approximation near the Fermi points, and after bosonization the zero-temperature excitations should be added back. The way to bosonize the linear Hamiltonian is also similar to the case of IEG. Because the dressed energy is smooth at the Fermi points now, the bosonization is even simpler. The density fluctuation operators are simply given by
$`\rho _q^{(+)}`$ $`=`$ $`{\displaystyle \underset{kk_F}{}}:c_{p+q}^{}c_p:`$ (193)
$`\rho _q^{()}`$ $`=`$ $`{\displaystyle \underset{kk_F}{}}:c_{pq}^{}c_p:.`$ (194)
The commutators among $`\rho _q^{(\pm )}`$ and $`H_\pm `$ are
$`[\rho _q^{(\pm )},\rho _q^{}^{(\pm )}]`$ $``$ $`0|[\rho _q^{(\pm )},\rho _q^{(\pm )}]|0={\displaystyle \underset{p_Fq<p<p_F}{}}0|c_{p+q}c_{p+q^{}}^{}|0`$ (195)
$`=`$ $`\delta _{q,q^{}}{\displaystyle \underset{p_Fq<p<p_F}{}}1={\displaystyle \frac{L}{2\pi }}q\delta _{q,q^{}}`$ (196)
and
$$[H_\pm ,\rho _q^{(\pm )}]0|[H_\pm ,\rho _q^{(\pm )}]|0=\pm v_Fq\rho _q^{(\pm )}.$$
(197)
Introducing the normalized bosonic annihilation operators
$$b_q=\sqrt{2\pi /qL}\rho _q^{(+)},\stackrel{~}{b}_q=\sqrt{2\pi /qL}\rho _q^{()}$$
(198)
and adding back the zero-mode contributions, the bosonized Hamiltonian satisfying (197) is given by
$$H_B=v_s\{\underset{q>0}{}q(b_q^{}b_q+\stackrel{~}{b}_q^{}\stackrel{~}{b}_q)+\frac{1}{2}\frac{\pi }{L}[\lambda _{\mathrm{eff}}M^2+\frac{1}{\lambda _{\mathrm{eff}}}J^2]\},$$
(199)
which agrees with the bosonized Hamiltonian (4) in the Luttinger liquid theory. We see that with $`\lambda `$ replaced by $`\lambda _{\mathrm{eff}}`$, the bosonized Hamiltonian for the GIG is the same as that for the IEG. So, all consequences we have obtained from the bosonized Hamiltonian in the IEG case can be applied to the GIG case. Especially, there is an (allowed) $`\mathrm{\Psi }_{\lambda _{eff}}^{}`$ describing the particle excitation near the Fermi surface with both anyon and exclusion statistics being $`\lambda _{eff}`$. In this sense, one may say that the effect of mutual statistics is to renormalize the statistics matrix.
Here we remark that in IEG, $`\mathrm{\Phi }(k,k^{})=(\lambda 1)\delta (kk^{})`$ is not smooth, so the dressed charge has a jump at $`k_F`$: $`z(k_F^+)=1`$ and $`z(k_F^{})=\lambda ^1`$ for $`k_F^\pm =k_F\pm 0^+`$. The general Luttinger-liquid relation is of the form
$$v_N=v_s[z(k_F^+)z(k_F^{})]^1,v_J=v_sz(k_F^+)z(k_F^{}).$$
(200)
### B Non-ideal Gas
Finally, we examine non-ideal gases, e.g., with general Luttinger-type density-density interactions,
$`H`$ $`=`$ $`H_{\mathrm{eff}}+H_I,`$ (201)
$`H_I`$ $`=`$ $`{\displaystyle \frac{\pi }{L}}{\displaystyle \underset{q0}{}}[U_q(\rho _q\rho _q^{}+\stackrel{~}{\rho }_q\stackrel{~}{\rho }_q^{})+V_q(\rho _q\stackrel{~}{\rho }_q^{}+\stackrel{~}{\rho }_q\rho _q^{})],`$ (202)
where $`H_{\mathrm{eff}}`$ is given by (192) describing a GIG, and $`\rho _q`$ and $`\stackrel{~}{\rho }_q`$ are the excluson density fluctuations near $`\pm k_F`$ respectively. After bosonization, the total Hamiltonian remains bilinear in densities:
$`H`$ $`=`$ $`H_B+H_I`$ (203)
$`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{q>0}{}}q[(v_s+U_q)(b_q^{}b_q+\stackrel{~}{b}_q^{}\stackrel{~}{b}_q+b_qb_q^{}+\stackrel{~}{b}_q\stackrel{~}{b}_q^{})`$ (206)
$`+V_q(b_q^{}\stackrel{~}{b}_q^{}+b_q\stackrel{~}{b}_q+\stackrel{~}{b}_q^{}b_q^{}+\stackrel{~}{b}_qb_q)]+{\displaystyle \frac{1}{2}}{\displaystyle \frac{\pi }{L}}[v_NM^2+v_JJ^2]`$
$`+{\displaystyle \frac{\pi }{L}}[U_0(M_R^2+M_L^2)+2V_0M_RM_L]{\displaystyle \underset{q>0}{}}v_sq`$
Using the Bogoliubov transformation, the Hamiltonian can be easily diagonalized
$`H`$ $`=`$ $`{\displaystyle \underset{q>0}{}}\omega _q(a_q^{}a_q+\stackrel{~}{a}_q^{}\stackrel{~}{a}_q)+{\displaystyle \frac{1}{2}}(\pi /L)[\stackrel{~}{v}_NM^2+\stackrel{~}{v}_JJ^2]+_0,`$ (207)
$`_0`$ $`=`$ $`{\displaystyle \underset{q>0}{}}(\omega _qv_sq),`$ (208)
where
$`a_q^{}`$ $`=`$ $`\mathrm{cosh}\stackrel{~}{\phi }_0b_q^{}\mathrm{sinh}\stackrel{~}{\phi }_0\stackrel{~}{b}_q^{},`$ (209)
$`\stackrel{~}{a}_q^{}`$ $`=`$ $`\mathrm{cosh}\stackrel{~}{\phi }_0\stackrel{~}{b}_q^{}\mathrm{sinh}\stackrel{~}{\phi }_0b_q^{},`$ (210)
and the renormalized velocities are
$`v_s`$ $``$ $`\stackrel{~}{v}_s=|(v_s+U_0)^2V_0^2|^{1/2},`$ (211)
$`v_N`$ $``$ $`\stackrel{~}{v}_N=\stackrel{~}{v}_se^{2\stackrel{~}{\phi }_0},`$ (212)
$`v_J`$ $``$ $`\stackrel{~}{v}_J=\stackrel{~}{v}_se^{2\stackrel{~}{\phi }_0}.`$ (213)
with the controlling parameter $`\stackrel{~}{\phi }_0`$ determined by
$$\mathrm{tanh}(2\stackrel{~}{\phi }_0)=\frac{v_Jv_N2V_0}{v_J+v_N+2U_0}.$$
(214)
Thus, the Luttinger-liquid relation ((44) survives with $`\lambda _{eff}`$ of GIG renormalized to
$$\stackrel{~}{\lambda }_{eff}=e^{2\stackrel{~}{\phi _0}}.$$
(215)
Note that the new fixed point depends both on the position of the Fermi points and on the interaction parameters $`U_0`$ and $`V_0`$, leading to “non-universal”exponents.
## VI Discussions and conclusions
In conclusion, we have shown that 1-d IEG (without mutual statistics) exactly reproduces the low-energy and low-$`T`$ properties of (one-component) Luttinger liquids. This gives rise to the following physical picture: At low temperature, the Luttinger liquids can be approximately thought of as an IEG consisting of quasiparticle excitations. Introducing mutual statistics or/ and Luttinger-type interactions among these excitations only shifts the value of $`\lambda _{eff}`$. Thus the essence of Luttinger liquids is to have an IEG obeying FES as their fixed point. This is our characterization of Luttinger liquids in terms of FES.
In this way, we have explicitly answered the three questions raised in the introduction about Luttinger liquids:
* The physical meaning of the Haldane’s controlling parameter is the quasiparticle’s effective statistics, $`\lambda _{eff}`$.
* The Luttinger liquids, more precisely, the IEG, indeed describe the infrared (or low-energy) fixed points in 1-d systems, since their effective field theory at low energy is conformally invariant. However, these fixed points are not isolated; they form a fixed-point line. Both the chemical potential and coupling constants are relevant perturbations that can drive the fixed point to move along the line, corresponding to the ”renormalization” of the effective statistics $`\lambda _{eff}`$ and leading to ”non-universal” exponents.
* It is conceivable that some strongly correlated systems, exhibiting non-Fermi liquid behavior, in two or higher dimensions may also be characterized as having a GIG with appropriate statistics matrix as their low-energy or low-temperature fixed point. This is because the concept of exclusion statistics is independent of spatial dimeniosnality of the system.
Moreover, we also showed that the effective field theory of 1-d IEG is a CFT with central charge $`c=1`$ and compactified radius $`R=\sqrt{1/\lambda }`$. The particle-hole duality of the exclusons implies the CFT has an unusual duality $`R1/R`$, meaning that the CFT belongs to a new variant of the $`c=1`$ CFT’s, i.e the ones that are compactified on an ”orbifold” $`S^1/Z_2`$ rather than on a circle. Physically, the differences are due to different constraints on the zero-mode quantum numbers. The CFT explanation makes a better understanding of the single-particle operators, especially, the anyonic (or exclusonic ) ones. Also, the CFT techniques provide a systematic way to calculate the correlation functions.
Finally we observe several additional implications of this work: 1) Our bosonization and operator derivation of CFT at low energies or in low-$`T`$ limit can be applied to Bethe ansatz solvable models, including the long-range (e.g., Calogero-Sutherland) one . 2) Here we have only consider one-species cases, i.e., with excitations having no internal quantum numbers such as spin. Our bosonization and characterization of Luttinger liquids are generalizable to GIG with multi-species, with the effective statistics matrix related to the dressed charge matrix . 3) The chiral current algebra in eqs. (74) and (75)with $`\lambda =1/m`$ coincides with that derived by Wen for edge states in $`\nu =1/m`$ fractional quantum Hall fluids. So these edge states and their chiral Luttinger-liquid fixed points can be described in terms of chiral IEG.
This work was supported in part by the U.S. NSF grant PHY-9309458, PHY-9970701 and NSF of China.
## A The harmonic fluid description
In coordinate space, there is a harmonic fluid description of the Luttinger liquid. Instead of the $`\theta `$-$`\varphi `$ representation that Haldane originally used, we prefer the right-left-moving representation. The density operator can be written as the Fourier transformations
$`\rho (x)`$ $`=`$ $`\rho _R(x)+\rho _L(x),`$ (A1)
$`\rho _R(x)`$ $`=`$ $`{\displaystyle \frac{M_R}{L}}+{\displaystyle \underset{q>0}{}}\sqrt{{\displaystyle \frac{q}{2\pi Le^{2\phi }}}}(e^{iqx}b_q+e^{iqx}b_q^{}),`$ (A2)
$`\rho _L(x)`$ $`=`$ $`{\displaystyle \frac{M_L}{L}}+{\displaystyle \underset{q>0}{}}\sqrt{{\displaystyle \frac{q}{2\pi Le^{2\phi }}}}(e^{iqx}\stackrel{~}{b}_q+e^{iqx}\stackrel{~}{b}_q^{}),`$ (A3)
where $`M_{R,L}`$ are given by $`M=M_R+M_L`$ and $`\stackrel{~}{b}_q=b_q`$ for $`q>0`$.
The boson field $`\varphi (x)`$, which is conjugated to $`\rho (x)`$ and satisfies
$$[\varphi (x),\rho (x^{})]=i\delta (xx^{}),$$
(A4)
is given by
$`\varphi (x)=\varphi _R(x)+\varphi _L(x),`$ (A5)
$`\varphi _R(x)={\displaystyle \frac{\varphi _0}{2}}+{\displaystyle \frac{\pi J_Rx}{L}}+i{\displaystyle \underset{q>0}{}}\sqrt{{\displaystyle \frac{\pi e^{2\phi }}{2qL}}}(e^{iqx}b_qe^{iqx}b_q^{}),`$ (A6)
$`\varphi _L(x)={\displaystyle \frac{\varphi _0}{2}}+{\displaystyle \frac{\pi J_Lx}{L}}+i{\displaystyle \underset{q>0}{}}\sqrt{{\displaystyle \frac{\pi e^{2\phi }}{2qL}}}(e^{iqx}\stackrel{~}{b}_qe^{iqx}\stackrel{~}{b}_q^{}),`$ (A7)
with $`J=J_R+J_L`$. We have to assign the quantum numbers such that there are only two independent variables in $`M_{R,L}`$ and $`J_{R,L}`$. A consistent choice is
$$M_R=J_R,M_L=J_L.$$
(A9)
Then,
$$J=J_R+J_L,M=J_RJ_L.$$
(A10)
Here $`\varphi _0`$ is an angular variable conjugated to $`M`$: $`[\varphi _0,M]=i`$. The Hamiltonian (4) becomes
$$H=\frac{1}{2}_0^L𝑑x[\pi v_N\rho (x)^2+v_J/\pi (_x\varphi (x))^2],$$
(A11)
or by a field rescaling,
$$H=\frac{v_s}{2\pi }_0^L𝑑x[\mathrm{\Pi }(x)^2+(_xX(x))^2],$$
(A12)
where
$$\mathrm{\Pi }(x)=\pi e^\phi \rho (x),X(x)=e^\phi \varphi (x).$$
(A13)
With $`X(x,t)=e^{iHt}X(x)e^{iHt}`$, the Lagrangian density reads
$$=\frac{v_s}{2\pi }_\alpha X(x,t)^\alpha X(x,t),$$
(A14)
which describes a free scalar field theory in $`1+1`$-dimensions. |
warning/0002/hep-th0002017.html | ar5iv | text | # Effects of Magnetic Fields on String Pair Creation
## 1 Introduction and summary
One way to unravel some nonperturbative features of quantum systems is to study their dynamics in the presence of background fields. A classic example is Schwinger’s calculation of the rate at which electron-positron pairs are created by a constant electric field.
More recently, the dynamics of linearly extended objects placed in various backgrounds became of interest in the context of string theory. In this direction, open strings coupled through their end-points to external electromagnetic fields have been studied by various authors . In particular, the Born-Infeld effective action for the external field was obtained in , whereas discussed the possibility of string pair creation through the Schwinger mechanism. Ref. considered bosonic strings in weak electric fields, whereas Ref. studied both bosonic and fermionic strings, for electric field strengths up to the value (in natural units) of the string tension. The above quantum effects were studied by calculating a one-loop effective action (whose imaginary part is relevant for the pair creation process) in a background electric field; additional magnetic fields were left out of this type of analysis, although they affect the pair creation rate.
The present note is an extension and a generalization of the seminal work . We study the creation of open strings in presence of a constant, but otherwise arbitrary, electromagnetic field strength $`F_{\mu \nu }`$, and compute the pair production rate (PPR) in this general situation. We display the various corrections with respect to the pure electric field background result, which is recovered as a particular case.
Moreover, our results will be obtained via path integral methods. Their use offers, even in the case of pure electric fields, several advantages over the Hamiltonian methods used in : the derivation is self-contained; it does not rely on previous work concerning the spectrum of the open string in external fields, but rather encodes it in the final result. In this way one also bypasses a somehow puzzling aspect of canonical methods, namely the use of a complex valued substraction for the vacuum, and hence for the spectrum, of the string in electric fields. For arbitrary external electromagnetic fields, which are our concern, path integrals allow a direct calculation - without previous diagonalization of $`F_{\mu \nu }`$ \- of the PPR. This is difficult to envisage through canonical methods. Still, to compare easily with the work of , we will choose to first block-diagonalize $`F_{\mu \nu }`$. This will ask for a semiquantitative analysis of the dependence of the eigenvalues of $`F_{\mu \nu }`$ on its non-diagonal components, which is of general interest for physics in more than four dimensions. It will turn out that one cannot find the eigenvalues analytically for a spacetime dimension $`D10`$.
We will also provide a transparent derivation of the Born-Infeld (BI) effective action, complementary to the original work of . In particular, the connection between the BI action and the finite spatial extension of the string is clearly seen in this formalism, as already noticed in (Ref. anticipated our present discussion in the case of pure electric fields). Beyond being somehow at the center of the calculation, the BI term produces an enhancement of the PPR, a qualitative stringy effect somehow unexpected in a bosonic theory (where the magnetic field usually decreases the pair production rate). This effect becomes particularly important in the dissipative limit of string theory described in . We treat in detail only the case of bosonic strings, which displays all the interesting aspects of the situation; the results for the case of superstrings are only briefly mentioned. Our methods and results might be of some use in the light of D-branes , for instance in the context of non-commutative low energy actions for strings stretched between branes with electromagnetic fluxes on them . The results obtained in this note were checked by using the boundary state formalism .
Last but not least, the open string spectrum in presence of a magnetic field contains tachyonic excitations , which destabilize the theory. These instabilities, like the similar tachyonic mode appearing in Yang-Mills theories , can be traced back to the nonminimal coupling of the electromagnetic field to excitations with spin greater than one . Our expedient remedy will be to ’Higgs’ the theory <sup>1</sup><sup>1</sup>1I am grateful to C.Bachas for useful discussions on this point by stretching the string between some different, parallel, D-branes. The Dirichlet boundary conditions along the coordinates orthogonal to the branes then produce an additional mass term, which can overcome the destabilizing contribution of the magnetic field.
## 2 Generalities
i) Vacuum decay rate. We briefly review one way to calculate the rate at which pairs are tunneled out of the vacuum, in presence of an external field.
The vacuum free energy in an external field $`F_{\mu \nu }`$ is the logarithm of the vacuum-to-vacuum transition amplitude $`e^{iW(F_{\mu \nu })_{(vac)}}=<0|e^{i\widehat{H}\times (time)}|0>`$. For a static field, $`W_{vac}(F_{\mu \nu })=_{vac}(F_{\mu \nu })\times (time)`$, where $`(time)`$ is the total time interval. Reexpressing the vacuum free energy à la Schwinger, and taking the Hamiltonian to be the one for open strings, one obtains
$$W_{vac}=_0^{\mathrm{}}\frac{dt}{t}Tre^{t\widehat{H}_{string}}.$$
(1)
The trace is evaluated by means of a (suitably normalised) path integral
$$Tre^{t\widehat{H}_{string}}=DX_0D\stackrel{}{X}e^{S(X_0,\stackrel{}{X},F_{\mu \nu })}.$$
(2)
The action $`S(X_0,\stackrel{}{X},F_{\mu \nu })`$ (discussed below) includes the external electromagnetic field, which couples to the end-points of the string. Due to that coupling, the vacuum energy $`_{vac}(F_{\mu \nu })`$ gets an imaginary part, $`\frac{\mathrm{\Gamma }}{2}`$, which induces the vacuum decay. The decay rate per unit volume $`\gamma =\frac{\mathrm{\Gamma }}{V}`$ reads
$$\gamma =2Im_0^{\mathrm{}}\frac{dt}{t}^{}DX_0D\stackrel{}{X}e^{S(X_0,\stackrel{}{X},F_{\mu \nu })}.$$
(3)
The prime means that we have factored out the zero mode part of the action (upon integration, it gives precisely the space-time volume).
For a bosonic open string living on a Euclidean world-sheet, coupled to a $`U(1)`$ gauge field, the action to be used in eqs.(2,3) reads:
$`S`$ $`=`$ $`{\displaystyle \frac{T}{2}}{\displaystyle _0^t}𝑑\tau {\displaystyle _0^l}𝑑\sigma [({\displaystyle \frac{X^\mu }{\tau }})^2+({\displaystyle \frac{X^\mu }{\sigma }})^2]`$
$`iq_1F_{\mu \nu }{\displaystyle _0^t}𝑑\tau \left[X_\mu {\displaystyle \frac{X_\nu }{\tau }}\right]_{\sigma =0}iq_2F_{\mu \nu }{\displaystyle _0^t}𝑑\tau \left[X_\mu {\displaystyle \frac{X_\nu }{\tau }}\right]_{\sigma =l},`$
where $`T`$ denotes the string tension, whereas $`q_1`$ and $`q_2`$ are the magnitudes of the charges situated at the end-points of the string. This action being invariant under a rescaling of both $`\sigma `$ and $`\tau `$ by a factor of $`l`$, we set $`l=1`$.
ii) Free path integral. We first review the zero electromagnetic field case . If $`F_{\mu \nu }0`$, eq.(2) factorizes into products of free path integrals along each space-time direction. For a generic uncoupled coordinate $`X`$ we have to evaluate $`DXe^S`$, with $`S=\frac{T}{2}_0^t𝑑\tau _0^1𝑑\sigma [(\frac{X}{\tau })^2+(\frac{X}{\sigma })^2]`$. Taking the boundary conditions to be periodic along $`\tau `$, $`X(t+\tau ,\sigma )=X(\tau ,\sigma )`$, and Neumann along $`\sigma `$ , $`\frac{X}{\sigma }|_{\sigma =0,1}=0,`$ which means expanding $`X(\tau ,\sigma )`$ as follows
$$X(\tau ,\sigma )=\underset{nZ}{}\underset{kN}{}X_{nk}cos(k\pi \sigma )exp(2\pi in\frac{\tau }{t}),$$
(5)
one gets the result
$$Z_{free}DXe^S=\sqrt{\frac{T}{4\pi }}e^{\frac{\pi }{6}\frac{1}{t}}\underset{n1}{}\frac{1}{1e^{4\pi n\frac{1}{t}}}.$$
(6)
Through a modular transformation one can reexpress (6) as follows:
$$Z_{free}=\sqrt{\frac{T}{4\pi }\frac{2}{t}}e^{\frac{\pi }{24}t}\underset{k1}{}\frac{1}{1e^{\pi kt}}.$$
(7)
iii) ’Higgsing’ the theory. Alternatively, one could use Dirichlet boundary conditions for the string’s ends, by obliging them to stay on two parallel D-branes situated at a relative distance $`d`$, i.e. ask $`X(\sigma =0)=0`$ and $`X(\sigma =1)=d`$. The corresponding mode expansion is
$$X(\tau ,\sigma )=d\sigma +\underset{nZ}{}\underset{k>0}{}X_{nk}\mathrm{sin}(k\pi \sigma )exp(2\pi in\frac{\tau }{t}),$$
(8)
and it modifies the final result of the path integration in only one way: an additional exponential factor $`e^{\frac{t}{2}Td^2}`$ appears in (6). This is equivalent to a mass term $`\frac{1}{2}Td^2`$, whose role is to stabilize the theory in presence of a magnetic field. This term can be inserted at any step of the calculation; we will come back to it later.
## 3 Eigenvalues
A non-zero background field $`F_{\mu \nu }`$ can be block-diagonalized in any dimension $`D`$: $`F_{01}=F_{10}=`$, $`F_{23}=F_{32}=_1`$, $`F_{45}=F_{54}=_2`$, etc., and $`F_{ij}=0`$ for $`ij\underset{¯}{+}1`$. $``$ represents the electric-like eigenvalue of $`F_{\mu \nu }`$, whereas the $``$’s are the magnetic-like ones. The PPR rate being a relativistic invariant, one can calculate it either by first diagonalizing $`F_{\mu \nu }`$ and path integrating subsequently, or by path integrating directly. The second approach will be mentioned later (eq. 21). We will use a block-diagonalized $`F_{\mu \nu }`$ and we will obtain the PPR rate as a function of its eigenvalues. It is thus of interest to study their dependence on the initial, in general non-diagonal, field strength.
The most important effect of non-diagonal $`F_{ij}`$’s is to change $``$ \- which enters the exponential factor of the PPR. Inserting an imaginary factor in front of the electric components $`F_{0j}`$, in order to use the Euclidean metric in the eigenvalue equation $`det(F_{\mu \nu }\eta _{\mu \nu }\lambda )=0`$, $`F_{\mu \nu }`$ reads, in $`D=4`$,
$$F=\left(\begin{array}{cccc}0& iE& 0& 0\\ iE& 0& b& 0\\ 0& b& 0& B\\ 0& 0& B& 0\end{array}\right).$$
(9)
$`B`$ and $`b`$ are the components of the magnetic field parallel, respectively orthogonal, to the electric field $`E`$; $``$, the real eigenvalue of $`F`$, is
$$2(^2)_{1,2}=\sqrt{(B^2+b^2E^2)^2+4E^2B^2}(B^2+b^2E^2).$$
(10)
For $`b=E`$, $``$ decreases with respect to the case $`b=0`$; it vanishes if $`B=0`$. In five dimensions, for a field strength tensor with non-zero components
$$F_{01}=iE,F_{12}=b,F_{34}=B,$$
(11)
and if $`b=E`$, the electric-like eigenvalue is zero for any $`B`$. This happens because, inside the matrix given by (11), the block containing $`E`$ and $`b`$ does not have lines or columns in common with the block containing the $`B`$’s.
In general, for an electric field parallel to the $`x`$-axis and a purely magnetic block partially diagonalized (with the $`F_{1j}`$’s left in their original form)
$$F=\left(\begin{array}{ccccccc}0& iE& 0& 0& 0& 0& \\ iE& 0& b_1& b_2& b_3& b_4& \\ 0& b_1& 0& B_1& 0& 0& \\ 0& b_2& B_1& 0& 0& 0& \\ 0& b_3& 0& 0& 0& B_2& \\ 0& b_4& 0& 0& B_2& 0& \\ .& .& .& .& .& .& .\end{array}\right),$$
(12)
the magnetic components fall into two classes: 1) The $`F_{1j}`$’s, called here $`b`$’s; they decrease the electric eigenvalue $``$ and might even cancel it (as in eq.11), unless they sit above a $`2\times 2`$ block containing a non-zero $`B`$ (as in eq.9). 2) The other components, the $`B`$’s; they do not influence $``$ in absence of the $`b`$’s; for non-zero $`b`$’s, they temper the reduction of $``$ those produce. At fixed $`F_{1j}`$’s, $``$ grows when the $`B`$’s are increased. This is easily seen up to $`D=5`$ space-time dimensions and, in principle, also up to $`D=9`$, by finding analytic expressions for the eigenvalues. In ten dimensions the characteristic equation of the matrix $`F`$ becomes of degree five and is not any more solvable by radicals. We have tested numerically various cases for $`D`$ from $`6`$ to $`10`$, and the conclusions above held. The $`F_{1j}`$’s reduce $``$ and the production rate, whereas the other magnetic components temper their decreasing effect if their corresponding planes intersect. This is probably true in any space-time dimension. Moreover, in higher dimensions the effect of a given variation of one single $`F_{ij}`$ is less important than a similar change in lower dimensions. The other numerous components provide a kind of inertia.
One can also fill the empty off-diagonal magnetic part of $`F`$ (e.g. the $`4\times 4`$ matrix containing $`B_1`$ and $`B_2`$ in (12)). Increasing those components might decrease or increase $``$, but their influence is small, being supressed by at least one order of magnitude with respect to their initial variation.
## 4 Path integral evaluation
Once $`F_{\mu \nu }`$ is diagonalized, we are left with path integrals along pairs of coupled coordinates. At this point, we could just use the results obtained in , for pure electric fields. Nevertheless, we prefer to make a more direct, path integral, calculation. Beyond being of intrinsic interest, the evaluation of the path integral which follows is a different way to make one of the few known nonperturbative, nonsupersymmetric, but still interesting, calculations in string theory. Choosing the $`01`$ plane, we evaluate
$$Z_{01}=DX_0DX_1e^{S(X_0,X_1,F_{01})},$$
(13)
the action $`S(X_0,X_1,F_{01})`$ being eq. (2) now restricted to $`\mu =0,1`$:
$`S`$ $`=`$ $`{\displaystyle \frac{T}{2}}{\displaystyle _0^t}𝑑\tau {\displaystyle _0^1}𝑑\sigma [({\displaystyle \frac{X^\mu }{\tau }})^2+({\displaystyle \frac{X^\mu }{\sigma }})^2]`$
$`iq_1E_1{\displaystyle _0^t}𝑑\tau \left[X_0{\displaystyle \frac{X_1}{\tau }}\right]_{\sigma =0}iq_2E_2{\displaystyle _0^t}𝑑\tau \left[X_0{\displaystyle \frac{X_1}{\tau }}\right]_{\sigma =1}.`$
We remark that we can treat in this way the more general case in which different field strengths $`E_{1,2}`$ are applied to the two string end-points. We will subsequently include the two charges $`q_1`$ and $`q_2`$ into the field strength value through the more compact notation $`q_{1,2}E_{1,2}E_{1,2}`$.
Developing in the same interaction-independent Fourier basis as in the free case (5), the action (4) becomes $`S=S(0)+_{n=1}^{\mathrm{}}S(n)`$, with $`S(0)=\frac{t}{2}_{k>0}[(X_{0k}^0)^2\frac{T}{2}\pi ^2k^2(X_{0k}^1)^2\frac{T}{2}\pi ^2k^2]`$ and $`S(n>0)=𝐗^{}A𝐗`$. The term $`𝐗^{}A𝐗`$ encodes the modes which couple due to the electric field:
$$𝐗^{}=(X_{n,0}^0,X_{n,1}^0,\mathrm{}X_{n,k}^0,\mathrm{},X_{n,0}^1,X_{n,1}^1,\mathrm{}X_{n,k}^1,\mathrm{}),$$
$$A=\left(\begin{array}{ccccccccccccc}a_0& 0& 0& & & & C_1& C_2& C_1& & & & \\ 0& a_1& 0& & & & C_2& C_1& C_2& & & & \\ & & & & & & & & & & & & \\ D_1& D_2& D_1& & & & b_0& 0& 0& & & & \\ D_2& D_1& D_2& & & & 0& b_1& 0& & & & \\ & & & & & & & & & & & & \end{array}\right),$$
(15)
where $`C_1=D_1=2\pi n(E_1+E_2)`$, $`C_2=D_2=2\pi n(E_1E_2)`$, whereas $`a_0=Tt\frac{4\pi ^2n^2}{t^2}`$, $`a_{k>0}=T\frac{t}{2}(\frac{4\pi ^2}{t^2}n^2+\pi ^2k^2)`$, and $`b_k=a_k`$, $`k0`$. The appearance of nonzero $`a_k`$ terms for $`k1`$ is due to the finite spatial extension of the string. One can prove that
$`Det(A)`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{\mathrm{}}{}}}a_ib_i\times [1C_1D_1({\displaystyle \underset{(ij)even}{}}{\displaystyle \frac{1}{a_ib_j}})C_2D_2({\displaystyle \underset{(ij)odd}{}}{\displaystyle \frac{1}{a_ib_j}})`$ (16)
$`+(C_1^2C_2^2)(D_1^2D_2^2)({\displaystyle \underset{(ij)odd,(kl)odd}{}}{\displaystyle \frac{1}{a_ia_jb_kb_l}})].`$
$`_{i=1}^{\mathrm{}}a_ib_i`$ corresponds to the uncoupled coordinates, for which we can use (6). Using now the identities (valid for some complex $`x`$)
$`{\displaystyle \frac{1}{2x^2}}`$ $`+`$ $`{\displaystyle \underset{k=2,4,6\mathrm{}}{}}{\displaystyle \frac{1}{x^2+k^2}}={\displaystyle \frac{\pi }{4x}}\mathrm{coth}({\displaystyle \frac{\pi x}{2}})`$ (18)
$`{\displaystyle \underset{k=1,3,5\mathrm{}}{}}{\displaystyle \frac{1}{x^2+k^2}}={\displaystyle \frac{\pi }{4x}}\mathrm{tanh}({\displaystyle \frac{\pi x}{2}}),`$
we obtain the following partition function:
$$Z_{01}=\frac{T}{4\pi }e^{\frac{\pi }{3}\frac{1}{t}}\underset{n1}{}\frac{[(1E_1^2)(1E_2^2)]^1}{[1\frac{(1+E_1)(1+E_2)}{(1E_1)(1E_2)}e^{4\pi n\frac{1}{t}}][1\frac{(1E_1)(1E_2)}{(1+E_1)(1+E_2)}e^{4\pi n\frac{1}{t}}]}.$$
(19)
The string tension $`T`$ was absorbed into $`E`$: $`E/TE`$.
Using the convenient notation $`ϵ=ϵ_1+ϵ_2`$, with $`ϵ_j=arcthE_j,\text{ }j=1,2`$, as well as $`\zeta `$-function regularizing the divergent product $`_{n1}[(1E_1^2)(1E_2^2)]`$ via $`_{k1}1=lim_{s0}\zeta (s)=\frac{1}{2}`$, we finally obtain
$$Z_{01}=\frac{T}{4\pi }e^{\frac{\pi }{3}\frac{1}{t}}\sqrt{(1E_1^2)(1E_2^2)}\underset{n=1}{\overset{\mathrm{}}{}}\frac{1}{[1e^{2ϵ4\pi n\frac{1}{t}}][1e^{2ϵ4\pi n\frac{1}{t}}]}.$$
(20)
Under the $`\zeta `$-function regularization, the divergent infinite product, to which all the string oscillation modes along $`\sigma `$ do contribute, has metamorphosed into the Born-Infeld term $`\sqrt{(1E_1^2)(1E_2^2)}`$. This happens not only in the case of globally neutral strings but rather, as usual from T-duality/D-branes arguments, each string end-point has associated with it a separate BI action, no matter what its charge is.
We quote now the result of a longer calculation, performed in $`D`$ dimensions without previous diagonalization of $`F_{\mu \nu }`$. It says that the partition function is just the one without electromagnetic field present, times a product of BI-like factors:
$$Z_{(F_{\mu \nu })}=Z_{(F=0)}\times \underset{n=1}{\overset{\mathrm{}}{}}[Det(\eta _{\mu \nu }cth(2\pi n/t)F_{\mu \nu })]^1.$$
(21)
This shows the importance of the BI action for the whole problem treated here: (21) encodes all the effects of the external field, in particular the way it distorts the open string spectrum. Eq. (21) arises from the determinant of a matrix analogous to (15) (now containing $`D\times D`$ infinite blocks) upon summing over the $`\sigma `$-oscillators of the string. Remarkably enough, after extensive use of (18,18), each infinite off-diagonal block gets replaced by a simple term $`F_{\mu \nu }cth(2\pi n/t)`$, whereas on the principal diagonal one gets the metric tensor $`\eta _{\mu \nu }`$.
It is sometimes useful to recast (20) in a different form (switching from the closed string channel to the open string one). Using the transformation properties of the Dedekind eta function $`\eta (x)=e^{i\pi \frac{x}{12}}_1^{\mathrm{}}(1e^{2\pi inx})=\frac{1}{\sqrt{ix}}\eta (\frac{1}{x})`$, and of the first $`\mathrm{\Theta }`$-function $`\mathrm{\Theta }_1(v|\tau )=2q^{\frac{1}{8}}\mathrm{sin}\pi v_{n=1}^{\mathrm{}}(1q^n)(1e^{2\pi iv}q^n)(1e^{2\pi iv}q^n)=e^{i\frac{v^2}{\tau }}\frac{1}{\sqrt{i\tau }}\mathrm{\Theta }_1(\frac{v}{\tau }|\frac{1}{\tau })`$, (20) becomes
$$Z_{01}=(E_1+E_2)\frac{T}{4\pi }e^{\frac{\pi }{12}t}\frac{e^{\frac{t}{2\pi }ϵ^2}}{sin(ϵ\frac{t}{2})}\underset{n>1}{}\frac{1}{[1e^{(iϵ\pi n)t}][1e^{(iϵ+\pi n)t}]}.$$
(22)
A linear $`(E_1+E_2)`$ factor appears in front instead of the BI term. Both (20) and (22) display poles, which signal the string pair production. They are
$$t_p(k)=\frac{2k\pi }{ϵ};\text{ }k=0,1,2,\mathrm{}$$
(23)
The path integrals along the other coupled directions are obtained from (20) or (22). The $`Z_{23}`$ path integral, for the $`2`$ and $`3`$ directions coupled by a magnetic eigenvalue $`B`$, for instance, is obtained by replacing $`ϵ_{1,2}if_{1,2}`$ ($`i=\sqrt{1}`$ ) in (20) . Now $`f_{1,2}=arctgB_{1,2}`$ and $`f=f_1+f_2`$. $`Z_{23}`$ does not exhibit poles, as expected:
$`Z_{23}`$ $`=`$ $`{\displaystyle \frac{T}{4\pi }}e^{\frac{\pi }{3}\frac{1}{t}}\sqrt{(1+B_1^2)(1+B_2^2)}{\displaystyle \underset{n>1}{}}{\displaystyle \frac{1}{[1e^{2if4\pi n\frac{1}{t}}][1e^{2if4\pi n\frac{1}{t}}]}}`$ (24)
$`=`$ $`{\displaystyle \frac{T}{4\pi }}e^{\frac{\pi }{12}t}(B_1+B_2){\displaystyle \frac{e^{\frac{t}{2\pi }f^2}}{sh(f\frac{t}{2})}}{\displaystyle \underset{n>1}{}}{\displaystyle \frac{1}{[1e^{(f\pi n)t}][1e^{(f+\pi n)t}]}}.`$
One has to take into account also the uncoupled directions and the ghosts. The contribution of one free coordinate, denoted $`Z_{free}`$, has been displayed in (6, 7); the ghosts cancel the stringy part of two free coordinates and give
$$z_g^2=(Z_{free})^2\times \frac{T}{2\pi t}.$$
(25)
Using eqs. (6) or (7), (25), (20) or (22), and (24), one can write down the whole partition function, $`Z=z_g^2\times (Z_{free})^d\times Z_{01}\times Z_{23}\times Z_{45}\times \mathrm{}`$, in $`D`$ dimensions out of which $`d`$ are left uncoupled by the electromagnetic field.
## 5 Pair production rate
From (6,22,25), and evaluating $`Im\frac{dt}{t}Z_{01}(t)`$, one obtains the pair production rate in presence of an electric field in D dimensions (compare to )
$$\gamma (E)\underset{k=1}{\overset{\mathrm{}}{}}\gamma _k=\pi \underset{k}{}()^{k+1}(E_1+E_2)e^{kϵ}\frac{T}{4\pi }\frac{ϵ}{k\pi }[Z_{free}]^{D2},$$
(26)
with $`Z_{free}`$ given by eq. (7), in which $`t`$ is replaced by $`\frac{2k\pi }{ϵ}`$, cf. (23). The term $`\gamma _1`$ is the dominant one. This result gets corrected in two ways in presence of a magnetic field. First, the electric-like eigenvalue $``$ may change in presence of other components of $`F_{\mu \nu }`$, as already discussed (thus we assume $`F_{\mu \nu }`$ to be in block-diagonal form, with $`E`$). Second, the presence of non-zero magnetic-like eigenvalues changes the form of the production rate, a point to which we now turn our attention. We consider only one non-zero magnetic eigenvalue (the analysis proceeds identically and independently for several $`B`$’s). Then each term $`\gamma _k`$ in the production rate (26) gets corrected, $`\gamma _k(E,B)=\gamma _k(E)\times \delta _k`$, with the following correction factor
$$\delta _k(f)=\sqrt{1+\frac{B^2}{T^2}}\underset{n=1}{\overset{\mathrm{}}{}}\frac{(1x^n)^2}{(1e^{2if}x^n)(1e^{2if}x^n)}=\frac{1}{\mathrm{cos}f}\underset{n=1}{\overset{\mathrm{}}{}}\frac{(1x^n)^2}{(12x^n\mathrm{cos}2f+x^{2n})}.$$
(27)
$`x`$ is evaluated at the pole of order $`k`$: $`x=e^{4\pi /t_p(k)}`$. We now discuss the behaviour of $`\delta _1\delta `$ as a function of the magnetic field.
In the absence of the BI term, $`\delta 1`$; the magnetic field would always decrease the PPR, as is the case for bosonic (Klein-Gordon) point-particles. Nevertheless, this $`\sqrt{1+(\frac{B}{T})^2}`$, stringy, contribution triggers a qualitative change; the pair production can be enhanced by the magnetic field, although by a small factor: $`\sqrt{2}`$ at most per string end-point if we do not allow the fields to be greater than the string tension (if we drop this restriction the PPR increases indefinitely, as $`B\mathrm{}`$). Although this enhancement is quite tiny, we think it is somehow unexpected in a bosonic theory, and further understanding of it would be required.
In the figure above, the coordinate axes are $`x=e^{4\pi /t_p}`$, $`f=arctg(B)`$, with $`t_p=2\pi /ϵ`$. The part of the figure at the left of the curve connecting $`(0,0)`$ and $`(x0.1335,f=\pi /4)`$ corresponds to an increase of the PPR, with a maximum at $`(x=0,f=\pi /4)`$, where an enhancement by a factor of $`\sqrt{2}`$ is obtained. On the above mentioned curve (and for $`f=0`$) the PPR equals the pure electric field one, whereas on its right it is smaller.
A dissipative, non-relativistic, limit of string theory can be obtained , by taking the velocity $`v`$ of propagation of excitations along space-like coordinates to be much smaller than the velocity of light $`c=1`$ along the time coordinate. In this case, the enhancement due to the BI term may become dramatic: it is of the form $`\sqrt{1+(\frac{B}{vT})^2}`$, with $`v1`$.
In the case of supersymmetric strings the tachyon disappears from the free spectrum, hence also the exponential increase it would produce (see eqs. 7, 26). For zero magnetic field the fermionic contribution was evaluated in and contains the sum over the even spin structures of powers of $`\mathrm{\Theta }`$\- and $`\eta `$-functions; in particular, it cancels the $`e^{kϵ}`$ factor in (26). An additional magnetic field further corrects the contribution of each of the spin structures $`s`$ by a factor
$$\delta _k^{(s)}(f)=\underset{n>0}{}\frac{(1+()^s2x^n\mathrm{cos}2f+x^{2n})}{(1()^sx^n)^2},s=2,3,4,$$
(28)
with $`n`$ integer for $`s=2`$ and half-integer for $`s=3,4`$; $`x`$ is evaluated at the $`k`$-th pole, as in (27). These corrections can increase the PPR, independently of the BI term. This effect is just the fermionic string counterpart of what happens in the case of pointlike Dirac fermions. The BI factor remains unchanged and no further poles appear in the path integral. The PPR reads
$$\gamma _{superstring}(ϵ,f)=\underset{k=1}{\overset{\mathrm{}}{}}\gamma _k\delta _k(f)\times \underset{s=2,3,4}{}[Z_{free}^{(s)}]^{D2}\times \delta _k^{(s)}(iϵ)\times \delta _k^{(s)}(f),$$
(29)
where $`\gamma _k`$, $`\delta _k`$, and $`\delta _k^{(s)}`$ are the ones from eqs.(26), (27), and (28), respectively. $`Z_{free}^{(s)}=\sqrt{\mathrm{\Theta }_s(t)/\eta (t)}`$ is the contribution per uncoupled fermionic direction.
We have now to remember that if we stretch the string along one direction \- in order to stabilize the theory - an additional factor $`e^{\frac{t_p}{2}Td^2}=e^{\frac{\pi }{ϵ}Td^2}`$ appears. It corresponds to a mass term $`\frac{1}{2}Td^2`$, whose role is to compensate the possibly negative contribution of (in string tension units):
$$\alpha ^{}m^2=\frac{f^2}{2\pi ^2}\frac{f}{\pi }(n1/2)+(n1)1,$$
(30)
$`n=a_1^{}a_1`$ being the number operator for the string modes lying on the first Regge trajectory. Thus, it is enough to take $`\frac{1}{2}Td^2=1`$; this provides a damping factor for the PPR of the form $`e^{\frac{\pi }{ϵ}}`$. Nevertheless, for big electric fields ($`E1`$, or $`ϵ\mathrm{}`$) this does not influence much the PPR, and the previous conclusions hold.
One final remark is in order. Using eqs. (24) and (7) one sees that in $`D=26`$ and in the limit $`t\mathrm{}`$ the leading term in the partition function is of the form $`e^{\pi t(\frac{f^2}{2\pi ^2}\frac{f}{2\pi }+1)}`$, modulo a prefactor. The coefficient of $`(t)`$ in the exponent gives precisely the lower tachyonic mass ($`n=0`$) in eq. (30), as it should. This check confirms that the path integral encodes \- at least in principle - all the information about the spectrum.
Acknowledgments
I am grateful to C. Bachas for stimulating discussions, in particular for suggesting the T-dual of the ’Higgs mechanism’ used here. A discussion with E. Gava is acknowledged. I thank R. Iengo for introducing me to the subject, for very useful discussions and encouragement, and for a careful reading of the manuscript. |
warning/0002/nlin0002044.html | ar5iv | text | # 1 Geometry of the rippled channel.
Quantum-classical correspondence for local density of states and eigenfunctions of a chaotic periodic billiard
G. A. Luna-Acosta\*, J. A. Méndez-Bermúdez, and F. M. Izrailev
Instituto de Física, Universidad Autónoma de Puebla, Puebla, Pue. 72570, Apartado Postal J - 48, México.
## Abstract
Abstract
Classical-quantum correspondence for conservative chaotic Hamiltonians is investigated in terms of the structure of the eigenfunctions and the local density of states, using as a model a 2D rippled billiard in the regime of global chaos. The influence of the observed localized and sparsed states in the quantum-classical correspondence is discussed.
PACS: 5.45.+b; 03.65.-w; 03.20.
Keywords: Quantum-classical correspondence, chaotic billiards.
* Corresponding author. E-mail: gluna@sirio.ifuap.buap.mx
The subject of this Letter concerns the quantum-classical correspondence of autonomous Hamiltonian systems with classical chaotic dynamics. As a paradigm for this class of systems we consider the motion of particles in the 2D rippled billiard depicted in Fig. 1. It is known \[1-3\] that the dynamics of a classical particle in this billiard undergoes the generic transition to chaos (regular-mixed-global) as the amplitude $`a`$ is increased ($`d`$ is fixed). Thus the results that can be obtained by studying this particular model are applicable to a wide class of systems. A modern physical realization of this rippled billiard of finite length is a mesoscopic electron wave guide. In its classical transport properties, such as resistivity, were related to its dynamical properties, yielding a transport signature of chaos. The analysis of the quantum motion in the infinitely long rippled channel (the periodic rippled channel) is useful for the understanding of universal features of electronic band structures of real crystals , propagation in periodic waveguides , quantum wires \[6-8\] and films . Signatures of chaotic diffusion in the band spectra of a similar periodic billiard have been investigated in Ref. . In the energy band structure of the quantum version of the periodic rippled channel was calculated and certain aspects of the quantum-classical correspondence were investigated in terms of the Husimi distributions.
The standard signature of chaos in quantum mechanics based on the random matrix conjecture showed a qualitative agreement for the energy level spacing statistics of the periodic rippled channel . However, a more detailed study revealed important deviations. Motivated by this discrepancy, which may be found in other chaotic systems, here we use a novel approach that has been suggested in Ref. to analyze the quantum-classical correspondence of complex systems. Such an approach relies on the study of the structure of the local density of states (LDOS), and the shape of the eigenfunctions (SEF) in the basis of non-interacting particles (quasi-particles) . It has been confirmed that the LDOS and SEF have well defined classical counterparts in the chaotic regime in dynamical systems such as that of two interacting spin particles and the three orbital schematic shell model , see also very recent works in . Below we extend these studies to chaotic billiards by treating the degrees of freedom as independent particles and incorporating the effects of the boundaries into the Hamiltonian operator. We show how to construct the classical counterparts of LDOS and SEF for chaotic billiards. We shall demonstrate that these quantities determine the global shape of the quantum LDOS and SEF, and that the phenomenon of quantum localization within the energy shell manifests itself as strong fluctuations about the global shape. We remark here that this localization occurs not in configuration space but in energy (or momentum) space (for chaotic billiards this kind of localization was first investigated in Ref. ).
We limit our studies here to the case of global classical chaos which occurs in wide ($`d>2\pi `$) channels . The control parameter that determines the degree of chaoticity, obtained by the standard linearization of its mapping, is given by $`K=\frac{2da}{\pi }`$. In agreement with Chirikov’s overlapping criteria our numerical experiments show global chaos for $`K\stackrel{>}{}1`$. We have compared the phase space dynamics for various combinations of $`a`$ and $`d`$ yielding the same value of $`K>1`$, and found their Poincare maps to be indistinguishable from each other. In what follows we take ($`a/2\pi ,d/2\pi )=(.06,1)`$, i.e., $`K1.5`$.
The quantum mechanical description involves the solution of the Schrödinger equation $`\widehat{H}=\frac{\mathrm{}^2}{2}(\widehat{p}_x^2+\widehat{p}_y^2)`$ subject to the boundary conditions $`\psi (x,y)=0`$ at $`y=0`$ and $`y=d+a\mathrm{cos}x`$. It is both convenient and illuminating to go to the new coordinates $`u=x`$, and $`v=\frac{yd}{d+acos(x)}`$, for in these coordinates the boundary conditions become simpler: $`\mathrm{\Psi }(u,v)=0`$ at $`v=0,d`$. On the other hand, the Hamiltonian acquires a much more complicated form,
$$\widehat{H}=\frac{\mathrm{}^2}{2}g^{1/4}\widehat{P}_\alpha g^{\alpha \beta }g^{1/2}\widehat{P}_\beta g^{1/4},\alpha ,\beta =u,v,$$
(1)
which is simply the kinetic energy expressed in covariant form . The momentum is now given by $`\widehat{P}_\alpha =i\mathrm{}[_\alpha +\frac{1}{4}_\alpha ln(g)]=i\mathrm{}g^{1/4}_\alpha g^{1/4}`$, where $`g^{\alpha \beta }`$ is the metric tensor, and $`g=Det(g_{\alpha \beta })=[1+ϵ\mathrm{cos}u]^2`$, $`ϵa/d`$, is the metric (see for details). Thus, we have formally transformed the model of one-particle motion in the rippled billiard to that of two interacting particles (identified with the two degrees of freedom). In this way, the complexity of the boundary in the original model is incorporated into the interaction potential.
Since the Hamiltonian is periodic in $`u`$, the energy eigenstates satisfy Bloch’s theorem: $`\mathrm{\Psi }_E(u,v)=\mathrm{exp}(iku)\mathrm{\Phi }_k(u,v)`$ with $`\mathrm{\Phi }_k(u+2\pi ,v)=\mathrm{\Phi }_k(u,v)`$ and Bloch wave vector $`k=k(E)`$. For an infinite channel, $`k`$ takes a continuous range of values and we choose the first Brillouin zone to lie in the interval $`[\frac{1}{2}k\frac{1}{2}]`$. Further, we can expand $`\mathrm{\Phi }_k(u,v)`$ in a double Fourier series and write the $`\alpha ^{th}`$ eigenstate of energy $`E^\alpha (k)`$ as
$$\mathrm{\Psi }^\alpha (u,v;k)=\underset{m=1}{\overset{\mathrm{}}{}}\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}C_{mn}^\alpha (k)\varphi _{mn}^k(u,v)$$
(2)
where
$$\varphi _{mn}^k(u,v)=<u,vm,n>_k=\pi ^{1/2}g^{1/4}sin(m\pi v)e^{i(k+n)u}.$$
(3)
The factor $`\pi ^{1/2}g^{1/4}`$ comes from the orthonormality condition in the curvilinear coordinates $`(u,v)`$. For future reference we shall call $`m`$ the mode (or channel) number and $`n`$ the Brillouin Zone number. Note that $`\varphi _{mn}^k`$ are the eigenstates of the unperturbed momenta squared (and, therefore, of the unperturbed Hamiltonian): $`\widehat{P}_v^2\varphi _{mn}^k=\mathrm{}^2(m\pi )^2\varphi _{mn}^k`$, $`\widehat{P}_u^2\varphi _{mn}^k=\mathrm{}^2(k+n)^2\varphi _{mn}^k`$. The problem requires the diagonalization of the matrix $`H_{l,l^{}}(k)=<l\widehat{H}l^{}>_k`$, where $`l>_k=m,n>_k`$. Note that for each pair $`(m,n)`$ we associate a number $`l`$. Even though the energy spectra is independent of the choice of the function $`l(m,n)`$, the structure of the matrix clearly depends on it. For our purposes, it is essential (as it will become clear below) that the unperturbed Hamiltonian basis be ordered in increasing energy: $`E_{l+1}^0(k)>E_l^0(k)`$, where $`H^0\varphi _l^k=E_l^0(k)\varphi _l^k`$, $`E_l^0(k)=\frac{\mathrm{}^2}{2}((n+k)^2+(m\pi )^2)`$. This defines the choice of assignment $`(m,n)l`$. Once the matrix $`H_{l,l^{}}(k)`$ has been diagonalized, its (exact) eigenstates $`\mathrm{\Psi }^\alpha (k)=C_l^\alpha (k)\varphi _l^k`$ are also re-ordered in energy $`(E^{\alpha +1}E^\alpha )`$. We adopt the convention that the Greek superindex (Latin subindex) denotes the exact (unperturbed) state. The amplitudes $`C_l^\alpha (k)`$ form the state vector matrix; its elements along the row $`\alpha `$ are the components of the $`\alpha ^{th}`$ state in the representation of the unperturbed re-ordered basis, and the elements along the column $`l`$ give the unperturbed state $`l`$ expanded in the re-ordered perturbed basis. We shall study the structure of both, rows and columns, the latter being essential in the construction of the local density of states (LDOS) of the system.
In the previous work we studied the level spacing statistics of the rippled channel in connection with the Random Matrix Theory conjecture. Using the Bloch momentum $`k`$ as an external parameter we found that the level spacing statistics is the same for all values of $`k`$ (except $`k0`$) within the Brillouin zone. Thus, here we arbitrarily take the value $`k=0.1`$ to explore the structure of the exact eigenstates and from now on we drop, for economy of notation, the label $`k`$.
Fig. 2 shows three consecutive high-energy eigenstates ($`C_j^\alpha `$, $`\alpha `$=2078, 2079, and 2080) in the re-ordered unperturbed basis. This sequence illustrates the kind of eigenstates that can occur typically; in a given sequence for this range of values of $`\alpha `$ roughly 15 percent are states such as the one shown in Fig. 2b, and about 60 percent are like those of Figs. 2a and 2c. The rest are intermediate between these two types: “sparse” states characterized by a few dominant peaks in a background of smaller components. Eigenstates can be characterized by various localization measures, such as the “entropy localization length” $`l_{}`$ and the “inverse participation length” $`l_{ipr}`$ (see, e.g., Ref. ) defined by
$$l_{}=2.08\mathrm{exp}(),l_{ipr}=3/𝒫,$$
(4)
where
$$=\underset{\alpha =1}{\overset{N}{}}\psi ^\alpha (E)^2ln\psi ^\alpha (E)^2,and𝒫=\underset{\alpha =1}{\overset{N}{}}\psi ^\alpha (E)^4.$$
(5)
These two measures are found to be two orders of magnitude smaller for the eigenstate 2079 than for the eigenstates 2078 and 2080 shown in Fig. 2. Extensive numerical results show that these lengths fluctuate wildly as a function of the number of the eigenstate $`\alpha `$, or, equivalently, as a function of its corresponding energy $`E^\alpha `$. As for extremely localized states such as that of Fig. 2b, they are expected to play an important role in transport since they are practically the same as the unpertubed (regular) states of the model. In the energy spectrum, they are located at the edges of the cells defined by different values of the Brillouin Zone number $`n`$. Moreover, these states have the largest group velocity $`E/k`$; see for details.
Our main interest here is to relate the structure of the matrix $`C_j^\alpha `$ to properties of the corresponding classical system. Since $`C_j^\alpha =<\mathrm{\Psi }^\alpha \varphi _j>`$ as a function of $`j`$ (or equivalently, $`E_j^0`$) is the projection of the exact state onto the states of the unperturbed system, a classical counterpart of $`C_j^\alpha ^2`$ as a function of $`j`$ can be obtained by projecting the (chaotic) trajectories generated by the Hamiltonian $`H`$ onto the unperturbed Hamiltonian $`H^0`$, where $`H=H^0+V`$ and $`V`$ is the perturbation. Specifically, this can be done by substituting the orbit $`\mathrm{\Phi }(t)(x(t),y(t),p_x(t),p_y(t))`$ of a typical initial condition generated by $`H`$ with energy $`E`$ onto $`H^0`$. Then the energy of the unperturbed Hamiltonian varies in time, $`E^0(t)`$, fluctuating within a range $`\mathrm{\Delta }E`$ which defines the width of the so-called “energy shell.” After a sufficiently long time a distribution of energies $`W_{cl}(E^0/E)`$ can then be obtained from $`E^0(t)`$ for a single (chaotic) orbit. Alternatively, $`W_{cl}(E^0/E)`$ can be calculated for shorter times but taking several orbits. The equivalence between these two ways is expected if motion generated by the full Hamiltonian is globally chaotic.
The above discussion outlines the general procedure that claims that $`W_{cl}(E^0/E)`$ is the classical counterpart of $`C_j^\alpha ^2`$ as a function of $`E_j^0`$. It has been applied \[15-19\] in a straight forward manner in systems where $`H^0`$ and $`V`$ can readily be identified. For chaotic billiards (with Neumann or Dirichlet boundary conditions) the Hamiltonian of both unperturbed and perturbed systems is simply the kinetic energy and the non-integrability comes from the boundary conditions. In order to incorporate the perturbation into the Hamiltonian operator, it is necessary to perform an appropriate coordinate transformation in the same way as it was done above for the quantum model. To accomplish this, we take the quantum Hamiltonian, Eq. (1), and commute all the operators, transforming them to c-numbers,
$$H=\frac{1}{2}g^{\alpha \beta }P_\alpha P_\beta ,$$
(6)
subject to the “flat” boundary conditions in $`(u,v)`$ coordinates. Expanding this expression, we obtain
$$H=\frac{1}{2}\left[P_u^2\frac{2ϵv\xi _u}{1+ϵ\xi }P_uP_v+\frac{1+(ϵv\xi _u)^2}{(1+ϵ\xi )^2}P_v^2\right],$$
(7)
where $`\xi cosu,\xi _ud\xi /du`$, and $`ϵa/d`$. The canonical transformation from the old to new variables is given by $`u=x`$, $`v=yd/(d+a\xi )`$, $`P_u=P_x+ϵy\xi _xP_y/(1+ϵ\xi )`$, and $`P_v=(1+ϵ\xi )P_y`$. We now separate the Hamiltonian $`H=H^0+V`$ analogously to the quantum model: $`H^0\frac{1}{2}\left[P_u^2+P_v^2\right];VHH^0`$.
To compare the classical distribution $`W_{cl}(E^0/E)`$ with the quantum one $`C_j^\alpha ^2`$ (as a function of $`E_j^0`$) we take the energy of the classical Hamiltonian to be equal to the energy of the state $`\mathrm{\Psi }^\alpha >`$ under consideration. Since the classical phase space dynamics is independent of the energy of the particle (the kinetic energy can be rescaled) we arbitrarily take it to be equal to one and introduce an effective Planck’s constant $`\mathrm{}_{eff}`$.
In Fig. 3 we show the classical and quantum distributions (shapes) of four eigenstates. The areas under the classical and quantum curves have been normalized to 1. These states were chosen from four different energy regions (different values of $`\mathrm{}_{eff}`$) and they represent typical states in each region of about one hundred eigenstates. These are typical in the sense that most states in each interval are characterized by typical localization lengths in that range. For example in the region of $`700<\alpha <800`$, the length $`l_{}`$ oscillates around $`l_{}=200`$ with only 8 eigenstates with $`l_{}<50`$ and the maximum of $`l_{}`$ is about 340. That is, the states of Fig. 3 are neither extremely delocalized nor extremely localized.
A simple inspection of Fig. 3 indicates that for low energies few unperturbed components participate in the construction of the exact eigenstate and as the energy increases, the participation of more and more components increases. Based on the Weyl formula and the condition that the de Broglie wavelength be equal to the amplitude of the ripple, one obtains an estimate for the number of levels defining the perturbative border $`N_p`$. For the geometrical parameters of the channel ($`a/2\pi =.06,d/2\pi =1`$) we have $`N_p139`$. Thus, for low energies, as exemplified by Fig. 3a, there is hardly any resemblance between the classical and quantum distributions, since at such energies the system is still in the perturbative regime and far from the semiclassical domain. It is also clear that as the energy increases the global shape of the quantum distribution is closer to the classical distribution. It is remarkable that the range of energy components that participate in the construction of the perturbed eigenstate can be deduced from the classical distribution. This range of energies defines the energy shell of the eigenstate under consideration. Note that the slight assymmetry observed in the range of the quantum distribution with respect to $`E^0=1`$ is the same as in the classical distribution and can be explained analytically, see details in Ref. .
For states with largest values of $`l_{ipr}`$ or $`l_{}`$, our numerical experiments show a much better quantum-classical correspondence. This occurs because as more components are needed to form the perturbed state, then the number of components with amplitudes above the classical shape decreases. Even so, strong fluctuations of the quantum distribution with respect to the classical shape do not dissappear completely even for the case of extremely delocalized (within the energy shell) states, see Fig. 3d.
We now try to smooth these fluctuations by calculating the averaged distribution, also known as the average shape of the eigenfunction (SEF)
$$W(E_j^0/E^\alpha )=\underset{\alpha ^{}}{}C_j^\alpha ^{}^2\delta (EE^\alpha ^{}),$$
(8)
where the sum is taken over a number of eigenstates in the neighborhood of the exact state $`\mathrm{\Psi }^\alpha >`$.
Another quantum quantity that shall be compared with its classical counterpart is the so called “local density of states” (LDOS) or “strength function”, widely used in solid state and nuclear physics. It is defined as
$$\omega (E^\alpha /E_j^0)=\underset{j^{}}{}C_j^{}^\alpha ^2\delta (E^0E_j^{}),$$
(9)
where the summation is done over a number of perturbed states in the neighborhood of the unperturbed state $`j>`$. This quantity gives information about the evolution of the system. In particular, the spreading width of this function determines the energy range associated with its “decay” into other states due to the interaction. The classical counterpart of the LDOS is constructed in much the same way as for the “classical eigenfunction” distribution $`W_{cl}(E^0/E)`$ except now we project the trajectories of the unperturbed Hamiltonian $`H^0`$ onto the exact Hamiltonian $`H`$ .
Figures 4a and 4b show both distributions, the average SEF and the LDOS, together with their classical counterparts. Even though the agreement between classical and averaged quantum eigenfunction shapes has improved compared to the individual shapes of Fig. 3, the quantum fluctuations are still relatively large. Moreover, the averaging procedure reveals a three-peak structure for both, classical and quantum models. The origin of these peaks can be undestood by a detailed analysis of the classical trajectories. In particular for SEF, it was found that the chaotic trajectories dweel for a long time in the neighborhood of stable and unstable period-one orbits. The unstable (stable) period-one orbit is defined by $`P_u=0`$ and $`x=0(\pi )`$. The right peak corresponds to motion perpendicular to the $`x`$ axis at $`x=\pi `$ (unstable orbit). Similarily, the left one results from the stable orbit at $`x=0`$, see details in Ref. . In contrast, the central peak is produced when the trajectory is parallel to the $`x`$-axis. A similar analysis explains the origin of the structure of the classical LDOS . One can detect a difference in the structure of the LDOS in comparison to that of the SEF, mainly, in the central region. This fact may be explained by the difference between time averages and phase space averages for finite times (see also discussion in Ref. ). It is remarkable that the quantum distributions are sensitive to the existence of these periodic orbits. In general there is an overall good correspondence between classical and quantum distributions (as exemplified by Fig. 4) but with large fluctuations due to the existence of localized and sparse states within the energy shell.
The kind of localization which we observed in this study occurs in a finite energy range determined by the classical energy shell. On the other hand, in time-periodic models, localization happens in unbounded momentum or energy space, and chaotic eigenstates are dense in the unperturbed basis. On the contrary, localization in conservative systems reveals itself in two different ways, one is very strong (perturbative) localization, see example in Fig. 2b, and the second one appears in the form of sparse eigenstates, see Fig. 3a. So far, it is not clear whether in conservative systems, the former type of localization can also appear.
In summary, we have analyzed quantum-classical correspondence for the eigenstates and the local spectral density of states for non-integrable billiards, using as a paradigm a 2D channel with periodic sinusoideal boundaries. We have shown that the global shape of the quantum quantities can be determined from their classical ones. In particular, the range of the energy shell and its global shape can be deduced from the classical distributions. It was also shown that the existence of localized and sparse states within the energy shell give rise to strong quantum fluctuations. These states are also responsible for the observed deviations from random matrix theory predictions for level spacing statistics and are expected to play an important role in the transport properties of the system. We remark that the comparison was done here for quantum and purely classical quantities. Thus, it is very interesting to develop a kind of semiclassical approach for the shape of eigenstates and the LDOS, based on the trace formula (for time-periodic pertubation, this problem has been addressed in Ref. ).
Acknowledgements
This work was partially supported by CONACYT (Mexico) grant, No. 26163-E and NSF-CONACYT, grant No. E120.3341. |
warning/0002/nucl-th0002034.html | ar5iv | text | # 1 Introduction
## 1 Introduction
In a high-energy heavy-ion collision, the fate of a $`J/\psi `$ or $`\psi ^{}`$ produced in the collision depends sensitively on its dissociation cross sections with ordinary light hadrons. If we are to use $`J/\psi `$ production as an indicator of the formation of a quark-gluon plasma in a high-energy heavy-ion collision, as suggested by Matsui and Satz , we should concurrently evaluate the competing process of $`J/\psi `$ absorption by hadrons that are produced during the collision \[2-10\].
The dissociation of a $`J/\psi `$ by light mesons has been examined previously, but the cross sections from various analyses unfortunately differ considerably, due largely to different assumptions regarding the dominant scattering mechanism \[11-18\]. A description of the present calculations is found in Ref. .
Charmonium dissociation processes can presumably be described in terms of the fundamental quark and gluon interactions but are of greatest phenomenological interest at energy scales in the resonance region. For this reason, we advocate the use of the known quark-gluon forces to specify the underlying scattering amplitude, which must then be convolved with the explicit nonrelativistic quark model hadron wavefunctions of the initial and final mesons to obtain the dissociation cross section.
## 2 The model
The model of Barnes and Swanson, which we shall use to describe these processes, is a quark interchange model, with a quark-quark interaction taken from quark potential models . In this model the meson-meson scattering amplitude is given by the sum of four quark line diagrams, which are shown in Fig. 1. These are evaluated as overlap integrals of quark model wavefunctions, using the “Feynman rules” given in App. C of Ref. . This method has previously been applied successfully to the closely related no-annihilation scattering channels $`I=2`$ $`\pi \pi `$ , $`I=3/2`$ $`K\pi `$ , $`I=0,1`$ S-wave $`KN`$ scattering , and the short-range repulsive NN interaction .
Fig. 1. Quark-exchange diagrams included in the calculation.
The diagrams of Fig. 1 are the ‘prior’ forms; there are also ‘post’ forms, in which the interaction takes place after the interchange of the quarks. In the non-relativistic formulation one can show that these two forms are equivalent , although we find small differences between them when we use relativistic kinematics for the initial and final mesons.
We take the interaction $`H_{ij}(r_{ij})`$ between quarks $`i`$ and $`j`$, represented by the curly line in Fig. 1, to be the conventional Coulomb plus linear plus spin-spin hyperfine interaction,
$`H_{ij}(r_{ij})={\displaystyle \frac{\lambda (i)}{2}}{\displaystyle \frac{\lambda (j)}{2}}\left\{{\displaystyle \frac{\alpha _s}{r_{ij}}}{\displaystyle \frac{3b}{4}}r_{ij}{\displaystyle \frac{8\pi \alpha _s}{3m_im_j}}S_iS_j{\displaystyle \frac{\sigma ^3}{\pi ^{3/2}}}e^{\sigma ^2r_{ij}^2}+V_{con}\right\},`$ (1)
where $`\alpha _s`$ is the strong coupling constant, $`b`$ is the string tension, $`m_i`$ and $`m_j`$ are the masses of the interacting constituents, and $`\sigma `$ is the range parameter in the Gaussian-regularized contact hyperfine term. (For antiquarks, the generator $`\lambda /2`$ is, as usual, replaced by $`\lambda ^T/2`$.) We used the parameter set
$`\alpha _s`$ $`=`$ $`0.58,b=0.18\mathrm{GeV}^2,\sigma =0.897\mathrm{GeV},`$
$`m_u`$ $`=`$ $`m_d=0.345\mathrm{GeV},m_c=1.931\mathrm{GeV},V_{con}=0.612\mathrm{GeV},`$ (2)
which gives a reasonably accurate description of the $`q\overline{q}`$ meson spectrum and, when applied to scattering, also gives an $`I=2`$ $`\pi \pi `$ phase shift, which is in good agreement with experiment. This Hamiltonian was used in the nonrelativistic Schrödinger equation to generate $`q\overline{q}`$ wavefunctions, which were in turn used with the same Hamiltonian in the diagrams of Fig. 1 to evaluate the scattering amplitudes and cross sections. We also used a second set of parameters, $`\alpha _s=0.594`$, $`b=0.162`$ GeV<sup>2</sup>, $`\sigma =0.897`$ GeV, $`m_u=m_d=0.335`$ GeV, $`m_c=1.6`$ GeV and a flavor-dependent $`V_{con}`$, found by fitting a large set of experimental masses, to test the sensitivity of our results to parameter variations.
## 3 Cross section for the dissociation of $`J/\psi `$ and $`\psi ^{}`$ by $`\pi `$
The dissociation of $`J/\psi `$ and $`\psi ^{}`$ by pions of sufficient energy can lead to many different open-charm channels. In $`\pi +J/\psi `$ and $`\pi +\psi ^{}`$ collisions the channels with the lowest thresholds are $`D\overline{D^{}}`$, $`D^{}\overline{D}`$, and $`D^{}\overline{D^{}}`$.
Fig. 2. Total and partial $`\pi +J/\psi `$ (Fig. 2$`a`$) and $`\pi +\psi ^{}`$ (Fig. 2$`b`$) dissociation cross sections.
In Fig. 2 we show the cross sections for the dissociation of $`J/\psi `$ and $`\psi ^{}`$ by $`\pi `$ as a function of the initial kinetic energy in the center-of mass system, $`E_{KE}=\sqrt{s}m_Am_B`$, where $`A`$ and $`B`$ are the colliding particles. The total cross sections are shown as thick solid curves and are the means of the ‘prior’ and ‘post’ results. The estimated errors due to the post-prior discrepancy and parameter variations are indicated by bands in the figures. The total cross section for the dissociation of $`J/\psi `$ by $`\pi `$ has an initial kinetic energy threshold of 0.65 GeV, and is about 1 mb just above threshold (Fig. 2$`a`$). This $`\pi +J/\psi `$ cross section is rather smaller than the 7 mb obtained by Martins $`etal.`$ , who used the same formalism but made different assumptions about the confining interaction.
The threshold for $`\pi +\psi ^{}\{D\overline{D^{}},D^{}\overline{D}\}`$ is only 0.05 GeV. The total cross section for the dissociation of $`\psi ^{}`$ by $`\pi `$ has local maxima at $`E_{KE}0.1`$ GeV and $`0.22`$ GeV, where the cross section is predicted to be, respectively, 6.2(0.8) mb and 4.6(1.8) mb (Fig. 2$`b`$). Thus, we find that the low energy $`\pi +J/\psi `$ dissociation cross section should be rather small, but in contrast we predict the $`\pi +\psi ^{}`$ cross section to be quite large. This is in qualitative agreement with earlier analyses of $`J/\psi `$ and $`\psi ^{}`$ suppression in pA, O-A, and S-U collisions .
## 4 Cross section for the dissociation of $`J/\psi `$ and $`\psi ^{}`$ by $`\rho `$
We next calculate the dissociation cross sections for $`J/\psi `$ and $`\psi ^{}`$ in their collisions with $`\rho `$ mesons. A $`\rho `$ meson is a spin-1 particle, so the total spin of $`\rho +\{J/\psi \mathrm{or}\psi ^{}\}`$ is $`S_{tot}=0,1,`$ and $`2`$. The lowest mass final reaction products can be $`(D,\overline{D})`$ with $`(S_3,S_4)=(0,0)`$ and $`S_{tot}=0`$, $`(D,\overline{D^{}})`$ with $`(S_3,S_4)=(0,1)`$ or $`(D^{},\overline{D})`$ with $`(S_3,S_4)=(1,0)`$ for $`S_{tot}=1`$, and $`(D^{},\overline{D}^{})`$ with $`(S_3,S_4)=(1,1)`$, for which $`S_{tot}`$ can be 0, 1 or 2.
Note that since the reaction $`\rho +J/\psi D\overline{D}`$ is exothermic, this cross section diverges as $`1/|\stackrel{}{v}_{\rho (J/\psi )}|`$ near threshold, independent of the detailed assumptions regarding the scattering mechanism. For the other channels the thresholds lie above $`m_\rho +m_{J/\psi }`$, so those subprocesses are endothermic and have zero cross section at threshold. The total dissociation cross section we find is shown as a solid line in Fig. 3. Numerically it is quite large, about 11(3) mb for an initial-state kinetic energy of 0.1 GeV, decreasing to 6(2) mb at a kinetic energy of 0.2 GeV.
Fig. 3. Cross sections for the dissociation of $`J/\psi `$ by $`\rho `$ into the six lowest open-charm channels.
## 5 Dependence of the dissociation cross sections on the $`\rho `$ mass
The collision of two pions leads to isospin $`I=0,`$ 1 and 2 states. The $`I=1`$ cross section at low energies is dominated by the $`\rho (770)`$, which has a width of 0.150 GeV. Because of this large width, $`\rho `$ mesons can be formed significantly far off the energy shell in $`\pi \pi `$ collisions.
A $`J/\psi `$ interacting with off-shell $`\rho `$ mesons will have a range of different effective thresholds, so it is useful to calculate these dissociation cross sections as a function of the $`\rho `$ meson mass. For this calculation we assume that the off-shell $`\rho `$ has the same spatial wavefunction as on-shell, and that only the kinematics of the reaction change. The results are shown in Fig. 4 for the $`J/\psi `$ dissociation cross section as a function of the center-of-mass kinetic energy above the lowest threshold, $`E_{KE}^{}`$. Here, $`E_{KE}^{}=\sqrt{s}m_\rho m_{J/\psi }E_{th}(0)`$, where $`E_{th}(0)`$ is the lowest threshold, and is given by $`(m_D+m_{\overline{D}}m_{J/\psi }m_\rho )\mathrm{\Theta }(m_D+m_{\overline{D}}m_{J/\psi }m_\rho )`$.
Figure 4 shows that the cross section increases from about 1 mb to 3 mb as $`m_\rho `$ increases from 0.45 GeV to 0.55 GeV. At $`m_\rho =0.65`$ GeV, the reaction $`\rho `$+$`J/\psi `$ becomes exothermic, and the cross section completely changes character; it diverges as $`1/\sqrt{E_{KE}^{}}`$ near $`E_{KE}^{}=0`$.
The results in Fig. 4 show that the dissociation cross section for $`\rho +J/\psi `$ is quite sensitive to the $`\rho `$ mass; higher mass $`\rho `$ mesons have much larger cross sections near threshold.
Fig. 4. The total $`\rho +J/\psi `$ dissociation cross section for different values of $`m_\rho `$.
We have carried out the corresponding calculations for the $`\rho +\psi ^{}`$ dissociation cross sections as a function of the $`\rho `$ meson mass; the results are shown in Fig. 5 as a function of $`E_{KE}^{}=\sqrt{s}m_\rho m_\psi ^{}E_{th}(0)`$, where the lowest threshold is $`E_{th}(0)=(m_D+m_{\overline{D}}m_\psi ^{}m_\rho )\mathrm{\Theta }(m_D+m_{\overline{D}}m_\psi ^{}m_\rho )`$. The cross section diverges as $`1/\sqrt{E_{KE}^{}}`$ near $`E_{KE}^{}=0`$ for the $`\rho `$ masses considered in Fig. 5. There are additional contributions from the $`D^{}\overline{D}`$ and $`D^{}\overline{D}^{}`$ channels. For a given $`E_{KE}^{}`$ and $`m_\rho `$, the $`\psi ^{}`$ dissociation cross section is larger than the $`J/\psi `$ cross section, with the exception of $`m_\rho =770`$ MeV and $`E_{KE}^{}>0.2`$ GeV, at which the two cross sections are nearly equal.
The results of Fig. 5 show that the dissociation cross section for $`\psi ^{}`$ by an off energy shell $`\rho `$ is quite sensitive to $`m_\rho `$.
Fig. 5. The total dissociation cross sections for $`\psi ^{}`$ by $`\rho `$ for different $`m_\rho `$.
## 6 Conclusions and future applications
The quark-interchange scattering model of Barnes and Swanson has been supported by extensive comparisons with experimental data in channels such as $`I=2`$ $`\pi \pi `$ scattering, $`I=3/2`$ $`K\pi `$ scattering, $`I=0,1`$ S-wave $`KN`$ scattering, the small $`\pi +J/\psi `$ and large $`\pi `$+$`\psi ^{}`$ dissociation cross sections, and indirectly through its incorporation of a quark model Hamiltonian that gives a good description of the experimental meson spectrum. The results for $`\rho +J/\psi `$ and $`\rho +\psi ^{}`$ dissociation cross sections we have found using this model are presumably qualitatively correct, since they depend mainly on the endothermic or exothermic nature of these processes. The more detailed numerical predictions certainly need to confront experimental data through a detailed study of $`J/\psi `$ and $`\psi ^{}`$ production and evolution in high-energy heavy-ion collisions.
It may also be possible to infer the $`\rho `$ dissociation cross sections from studies of heavy-ion collisions. $`\rho `$ mesons are produced in $`\pi \pi `$ collisions and they also decay inversely into two pions. The production and the decay of $`\rho `$ mesons can readily reach an equilibrium and the the law of mass action applies. As a consequence, the density of $`\rho `$ mesons should vary as the square of the pion density. In a heavy-ion collision the density of pions produced varies approximately as the path length $`L`$ of the colliding nuclei . Thus, the density of $`\rho `$ mesons depends approximately quadratically on the path length $`L`$. Charmonium absorption by these $`\rho `$ mesons thus leads to a term $`c\times \sigma (\rho +J/\psi )\times L^2`$ in the exponential absorption index. This quadratic dependence will be more evident if the cross section $`\sigma (\rho `$+$`J/\psi )`$ for $`\rho +J/\psi `$ dissociation is large, as we have found here. This term will also be enhanced as we increase the radii of the colliding nuclei, for example in going from S+U to Pb+Pb collisions. The degree to which the observed anomalous suppression in Pb+Pb collisions is due to $`\rho +J/\psi `$ dissociation will require a careful quantitative study, taking into account the variation of the dissociation cross section on the off energy shell $`\rho `$ mass. In any case, the possible importance of such a nonlinear suppression term should be considered in the attempt to identify the creation of a quark-gluon plasma through suppression of $`J/\psi `$ production.
## Acknowledgments
This research was supported by the Division of Nuclear Physics, DOE, under Contract No. DE-AC05-96OR21400 managed by Lockheed Martin Energy Research Corp. ES acknowledges support from the DOE under grant DE-FG02-96ER40944 and DOE contract DE-AC05-84ER40150 under which the Southeastern Universities Research Association operates the Jefferson National Accelerator Facility. TB acknowledges additional support from the Deutsche Forschungsgemeinschaft DFG under contract Bo 56/153-1. The authors would also like to thank C. M. Ko and S. Sorensen for useful discussions. |
warning/0002/quant-ph0002089.html | ar5iv | text | # Operational criterion and constructive checks for the separability of low rank density matrices.
## I Introduction
Entanglement is one of the quantum properties with no classical counterpart. It is closely connected to fundamental questions of quantum mechanics, and to physical phenomena which are important for quantum information processing . The relevance of entanglement effects was first demonstrated for pure states. However, in realistic physical situations one deals usually with mixed states, in which pure state entanglement has been significantly weakened by noise. In order to overcome the problems caused by noise (i.e. in order to reduce it) the idea of distillation of entanglement in spatially separated laboratories was introduced . It was proved then that for bipartite systems of low dimensional Hilbert space $`𝒞^M\times 𝒞^N`$ or simply $`M\times N`$ (namely systems with $`M=2`$ and $`N=2`$ or $`3`$) mixed state entanglement can always be distilled into its pure form. However, it turned out that in higher dimensional systems ($`MN>6`$) bound entanglement – which cannot be distilled, as opposed to free entanglement – exists.
Unlike in the case of pure states, it is in general very difficult to know whether a given mixed state is entangled (inseparable) or non-entangled (separable). According to the definition a state supported on a Hilbert space $`_{AB}=_A_B`$ is separable if and only if it can be written in (or approximated by) the form
$$\varrho =\underset{i=1}{\overset{k}{}}p_i|e_i,f_ie_i,f_i|,\underset{i}{}p_i=1,$$
(1)
where $`|e_i,f_i`$ stands here for a normalized vector $`|e_i|f_i_A_B`$. In finite dimensional cases, the ones we will be concerned here, the approximation part is not necessary, as for any separable state one can always find a set $`\{|e_i,f_i\}`$ of product vectors for which $`k\mathrm{dim}(_{AB})^2`$ in the above formula .
Several necessary conditions for separability are known: Werner derived a condition based on the analysis of local hidden variables models and mean value of the, so called, flipping operator ; the Horodeckis proposed a necessary criterion based on $`\alpha `$-entropy inequalities . Peres demonstrated that the partial transpose $`\varrho ^{T_A}`$ of the matrix $`\varrho `$, defined as $`m,\mu |\varrho ^{T_A}|n,\nu =n,\mu |\varrho |m,\nu `$ for any fixed product basis $`|n,\nu |e_n_A|e_\nu _B`$, must be still a legitimate density matrix if $`\varrho `$ is separable . This operationally friendly, necessary condition, called positive partial transpose (PPT) condition, turned out to be very strong.
Soon after Peres result, a general connection between positive map theory and separability was established in , where necessary and sufficient separability conditions were derived in terms of positive maps. In particular it implied that for systems of low dimensions ($`MN6`$) the PPT condition is also sufficient for separability. It was also shown that this is not the case for systems of higher dimensions ($`MN>6`$). Later on explicit counterexamples of entangled states with PPT property were provided by means of another separability criterion, based on the analysis of the range of the density matrix (cf. ). It was then shown that they represent bound entanglement . Let us note that on mathematical grounds there were examples, provided earlier , of elements of positive matrices cones which can be treated as prototypes of PPT entangled states.
Sufficient conditions for separability are also known. We remark that the results of readily imply that any state close enough to the completely random state $`\pi `$ is separable. Thus, as quantified in , any mixture $`\stackrel{~}{\varrho }=(1p)\varrho +p\pi `$ in a $`M\times N`$ system is separable if $`p(1+2/MN)^1`$ or, in other words, as we wish to make explicit here, a full rank mixed state is separable provided its smallest eigenvalue be greater than or equal to $`(2+MN)^1`$.
On the other hand the analysis of the range of the density matrices, first applied in the separability criterion , led to an algorithm for the optimal decomposition of mixed states into a separable and an inseparable part , and to a systematic method of constructing examples of PPT entangled states and peculiar positive maps . Also, the technique of diminishing the rank of a PPT density matrix by subtraction of selected product vectors, which was worked out in , turned out to be very useful. This and other techniques have allowed quite recently to study operational necessary and sufficient separability conditions for states of a $`2\times N`$ system . In particular it has been shown that:
($`i`$) all PPT states of rank smaller than $`N`$ are separable;
($`ii`$) the separability of generic states such that $`r(\varrho )+r(\varrho ^{T_A})3N`$ reduces to analyzing the roots of some complex polynomials (a constructive separability criterion was derived, thus providing also the decomposition of such separable states into pure product states);
($`iii`$) states invariant under partial transpose, and those that are not “very different” from their partial transpose are necessarily separable.
This paper can be considered an extension and generalization of Ref. . The results ($`i`$) and ($`ii`$) obtained there for $`2\times N`$ systems are here generalized non-trivially to the case of $`M\times N`$ systems ($`MN`$). We show, namely, that
* any state $`\varrho `$ supported on $`M\times N`$ ($`MN`$) and with rank $`r(\varrho )N`$ is separable iff its partial transpose is positive;
* separability of generic PPT density matrices with $`r(\varrho )+r(\varrho ^{T_A})2MNMN+2`$ reduces to solving a system of coupled polynomial equations.
In both cases a pure product state decomposition for separable states is obtained.
Throughout this paper we make use of the following definition: we say that a state $`\rho `$ acting on $`M\times N`$ is supported on $`M\times N`$ if this is the smallest product Hilbert space on which $`\rho `$ can act. Let us introduce the local ranks $`r(\varrho _A)`$ and $`r(\varrho _B)`$, where $`\varrho _{A,B}`$ Tr$`{}_{A,B}{}^{}\varrho `$ are the reduced density operators. It immediately follows from the first of the above results that there is no PPT bound entanglement of rank 3. Indeed, a rank $`3`$ state $`\rho `$ either has at least one of the local ranks $`r(\rho _A)`$ and $`r(\rho _B)`$ greater than $`3`$, and in this case is distillable (i.e., $`\rho ^{T_A}`$ is not positive), or else can be supported on a $`MN6`$ or on a $`3\times 3`$ system, and thus is separable. This implies in particular that the bound entangled states constructed with the UPB method and those based on the chess-board structure of eigenvectors are optimal with respect to their ranks.
For our second main result, concerned with those PPT density matrices for which the sum of ranks satisfies $`r(\varrho )+r(\varrho ^{T_A})2MNMN+2`$, we identify the eligible product vectors (that is, those that can appear in decomposition (1) if $`\rho `$ is separable) with the solutions of a system of coupled polynomial equations. We analyze these equations, which are arguably expected to have only a finite number of solutions. For this case we present a constructive (i.e. leading to a product state decomposition) method to check separability. Also for the same case we discuss an alternative, constructive method to check separability numerically. These checks represent a necessary and sufficient condition for separability.
We wish to remark the importance of having separability conditions for low rank density matrices, especially in relation to unsolved problems concerning the nature of bound entanglement (BE). Notice that such conditions are of great value when trying to construct states with BE. Among the open questions we encounter the existence of BE having a non-positive partial transpose NPT (see ). Also, whether a finite or a vanishing amount of free entanglement is required to asymptotically create bound entangled states. There are in addition several conjectures concerning bound entanglement (see ) among them the ones connected to capacities of quantum channels and bound entanglement assisted distillation. Finally, we have been recently able to establish a general connection between low rank bound entangled states and positive maps. This connection allows for a systematic construction of independent linear maps in arbitrary dimensions, including $`2\times N`$, where the procedures based on unextendible product bases do not work . The discussion of this connection will be presented elsewhere.
This paper is organized as follows: we start by generalizing some needed results of related to diminishing the rank of $`\rho `$ by subtracting projectors on product vectors; in Section III we present our theorem about the separability of states with rank $`N`$; in Section IV the necessary and sufficient separability conditions for generic matrices with $`r(\varrho )+r(\varrho ^{T_A})2MNMN+2`$ are formulated, and discussed in the context of $`3\times 3`$ systems; finally, Section V contains our conclusions and acknowledgments.
## II Diminishing of the rank - generalizations
Before we turn to the main results of this paper we need to generalize some of those presented in .
Consider a state $`\varrho `$ of a $`M\times N`$ system satisfying $`\varrho ^{T_A}0`$. Throughout this paper $`K(X),R(X),k(X)`$, and $`r(X)`$ denote the kernel, the range, the dimension of the kernel, and the rank of the operator $`X`$, respectively. By $`\{|a_i\}_{i=1}^M`$ and $`\{|b_i\}_{i=1}^N`$ we will denote orthonormal basis in $`_A`$ and $`_B`$, and by $`|e^{}`$ we will denote the complex conjugated vector of $`|e`$ in the orthonormal basis $`|1_A,\mathrm{},|M_A`$ in which we perform the partial transposition; that is, if $`|e=_{i=1}^M\alpha _i|i`$ then $`|e^{}=_{i=1}^M\alpha _i^{}|i`$.
For the time being we do not require $`MN`$. The following Lemma is a generalization of Lemma 6 of Ref. proved there for $`M=2`$:
Lemma 1.- If $``$ $`|f𝒞^N`$ such that $`|a_i,fK(\varrho )`$ for $`i=1,\mathrm{},M1`$, then either (i) $`|a_M,fK(\varrho )`$ or (ii)
$`\varrho |a_M,f=|a_M,g`$ (2)
$`\varrho ^{T_A}|a_{M}^{}{}_{}{}^{},f=|a_{M}^{}{}_{}{}^{},g`$ (3)
for some $`|g𝒞^N`$.
Proof.- From the assumptions we have immediately $`\varrho ^{T_A}|a_{i}^{}{}_{}{}^{},f=0`$ ($`i=1,\mathrm{},M1`$). In particular $`|h𝒞^N`$ we have $`a_{M}^{}{}_{}{}^{},h|\varrho ^{T_A}|a_{i}^{}{}_{}{}^{},f=0`$ or, equivalently, $`a_i,h|\varrho |a_M,f=0`$. Since $`|h`$ is arbitrary we have either statement (i) or $`\varrho |a_M,f=|a_M,g`$ for some $`|g0`$. The second case needs further analysis. In a similar way we can prove that either $`\varrho ^{T_A}|a_{M}^{}{}_{}{}^{},f=0`$ (which is still equivalent to the statement (i)) or $`\varrho ^{T_A}|a_{M}^{}{}_{}{}^{},f=|a_{M}^{}{}_{}{}^{},g^{}`$ for some $`|g^{}0`$. It remains to prove that $`|g^{}=|g`$. Indeed $`|g^{}=a_{M}^{}{}_{}{}^{}|\varrho ^{T_A}|a_{M}^{}{}_{}{}^{},f=a_M|\varrho |a_M,f=|g`$.
The second lemma below is also a generalization of the results from Ref. :
Lemma 2.- If $`\varrho `$ satisfies the assumptions of Lemma 1, and the possibility (ii) of Lemma 1 holds, then
$$\varrho _1=\varrho \lambda |a_M,ga_M,g|,$$
(4)
where $`\lambda a_M,g|\varrho ^1|a_M,g`$ and
(i) $`\varrho _1`$ is a PPT state with $`r(\varrho _1)=r(\varrho )1`$ and $`r(\varrho _1^{T_A})=r(\varrho ^{T_A})1`$.
(ii) $`\varrho _1`$ is supported either on a $`(M1)\times (N1)`$ or on a $`M\times (N1)`$.
(iii) $`\varrho _1`$ is separable iff $`\varrho `$ is separable.
Proof .- Following Corollary 1 and Lemma 2 from Ref. , we observe that
$$\varrho _1=\varrho \frac{|a_M,ga_M,g|}{a_M,g|\varrho ^1|a_M,g}$$
(5)
is positive, and that $`r(\varrho _1)=r(\varrho )1`$. Then (i) follows from taking into account that since
$`\lambda ^1`$ $`=`$ $`a_M,g|\varrho ^1|a_M,g=`$ (6)
$`g|f`$ $`=`$ $`a_{M}^{}{}_{}{}^{},g|\varrho _{}^{T_A}{}_{}{}^{1}|a_{M}^{}{}_{}{}^{},g,`$ (7)
we have that
$$\varrho _{1}^{}{}_{}{}^{T_A}=\varrho ^{T_A}\frac{|a_{M}^{}{}_{}{}^{},ga_{M}^{}{}_{}{}^{},g|}{a_{M}^{}{}_{}{}^{},g|\varrho _{}^{T_A}{}_{}{}^{1}|a_{M}^{}{}_{}{}^{},g},$$
(8)
so that also $`\varrho _1^{T_A}`$ is positive and $`r(\varrho _1^{T_A})=r(\varrho ^{T_A})1`$. From assumptions on the kernel of $`\varrho `$ it follows that all vectors $`\{|a_i,f\}_{i=1}^M`$ belong to the kernel of $`\varrho _1`$, hence $`\varrho _1`$ can be embedded into a $`M\times (N1)`$ space; on the other hand since $`\varrho _{1,A}=\varrho _A\lambda g|g|a_Ma_M|`$, $`r(\varrho _{1,A})`$ must be either $`M=r(\varrho _A)`$ or $`M1`$, which finishes the proof of (ii). In order to prove (iii) let us assume that $`\varrho `$ is separable and let us show that also $`\varrho _1`$ is so (if $`\varrho _1`$ is separable then obviously $`\varrho `$ is also separable). Since $`\varrho |a_i,f=0`$ ($`i=1,\mathrm{},M1`$), we can always write
$$\varrho =|e_i,f_ie_i,f_i|+|a_Ma_M|\eta ,$$
(9)
where $`f|f_i=0`$ and $`\eta `$ is a positive operator acting on $`𝒞^N`$. If we impose $`|a_M,g=\varrho |a_M,f`$ we obtain $`|g=\eta |f`$, and therefore $`|gR(\eta )`$. We can write
$$\varrho _1=|e_i,f_ie_i,f_i|+|a_Ma_M|(\eta \lambda |gg|),$$
(10)
so that if we show that the operator $`(\eta \lambda |gg|)0`$ then we have that $`\rho _1`$ is separable. Using (7) we have that such an operator is
$$\eta \frac{1}{g|f}|gg|=\eta \frac{1}{g|\eta ^1|g}|gg|,$$
(11)
and that therefore it is positive (cf. Lemma 1 in ). $`\mathrm{}`$
## III All rank $`N`$ PPT states supported on a $`M\times N`$ system ($`MN`$) are separable
In this section we generalize the following theorem proved in Ref. :
Theorem.-\[Theorem 1 of Ref. \]
Let $`\varrho `$ be a PPT state of rank $`N`$ supported on a $`2\times N`$ space. Then $`\varrho `$ is separable and can be written as
$$\varrho =\underset{i=1}{\overset{N}{}}|e_i,f_ie_i,f_i|$$
(12)
with all $`\{|f_i\}`$ linearly independent.
We will express a density matrix in terms of its reduced operators $`i_A|\varrho |j_A`$ acting on $`_B`$. For instance, we will write (12) as
$`\varrho =\left[\begin{array}{cc}\stackrel{~}{A}& \stackrel{~}{B}^{}\\ \stackrel{~}{B}& \stackrel{~}{C}\end{array}\right],`$ (15)
where $`\stackrel{~}{A}1|\rho |10`$, $`\stackrel{~}{C}2|\rho |20`$ and $`\stackrel{~}{B}2|\rho |1`$. More generally, a $`M\times N`$ density matrix will be expressed as
$`\varrho =\left[\begin{array}{cccc}E_{11}& E_{12}& \mathrm{}& E_{1M}\\ E_{12}^{}& E_{22}& \mathrm{}& \mathrm{}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ E_{1M}^{}& \mathrm{}& \mathrm{}& E_{MM}\end{array}\right],`$ (20)
where now $`E_{ij}i_A|\varrho |j_A`$. We start by using the previous Theorem to prove:
Lemma 3.- Let $`\varrho `$ be a rank $`N`$ PPT state supported on a $`2\times N`$ space. Then after a reversible local filtering operation <sup>*</sup><sup>*</sup>*A reversible transformation is a transformation that can be reversed with nonzero probability. A local filtering in Bob’s side $`IV`$ is then reversible iff the operator $`V`$ can be inverted, i.e. iff $`V^1`$ exists. the state is proportional to the matrix (called hereafter $`2\times N`$ canonical form):
$`\mathrm{\Sigma }\left[\begin{array}{cc}B^{}B& B^{}\\ B& I\end{array}\right]=[BI]^{}[BI],`$ (23)
with $`B`$ normal, i.e. $`[B,B^{}]=0`$.
Proof .- We write the density matrix (12) in the form (15). Because there is only a finite number of $`|e_i`$ in (12) we can always find a vector $`|a`$ such that $`a|e_i0`$ for all $`i`$. Let this $`|a`$ be the second element of the orthonormal basis in Alice’s space, i.e. $`|2=|a`$. The matrix
$$\stackrel{~}{C}=a|\varrho |a=\underset{i=1}{\overset{N}{}}|a|e_i|^2|f_if_i|$$
has then maximal rank $`N`$, since the $`|f_i`$ are linearly independent. Taking the local filter $`V=(\sqrt{\stackrel{~}{C}})^1`$ on Bob’s side (this corresponds to sandwitching the state between $`IV`$ and $`IV^{}`$) we obtain
$`\stackrel{~}{\varrho }=\left[\begin{array}{cc}AB^{}& \\ BI& \end{array}\right],`$ (26)
which is still positive and PPT (because any local operation preserves the PPT property ). We can write
$$\stackrel{~}{\varrho }=\mathrm{\Sigma }+\mathrm{diag}[\mathrm{\Delta },0],$$
(27)
where the positive matrix $`\mathrm{\Sigma }`$ from the expression (23) has rank $`r(\mathrm{\Sigma })=N`$ as its kernel $`K(\mathrm{\Sigma })`$ has at least dimension $`N`$ containing all vectors of the type
$$|\varphi _f=|1|f+|2|Bf,$$
(28)
while its range has at least dimension $`N`$ due to the identity entry on the diagonal. Notice that diag$`[\mathrm{\Delta },0]`$ is also positive, because positivity of $`\stackrel{~}{\varrho }`$ implies that $`\mathrm{\Delta }=AB^{}B0`$ . Now, since in addition $`r(\stackrel{~}{\rho })=r(\mathrm{\Sigma })`$, we also have that $`R(\stackrel{~}{\rho })=R(\mathrm{\Sigma })R(`$diag$`[\mathrm{\Delta },0])`$, that is $`K(`$diag$`[\mathrm{\Delta },0])K(\mathrm{\Sigma })`$. But $`K(\mathrm{\Sigma })`$ is spanned by the states (28), for which then $`\varphi _f|`$diag$`[\mathrm{\Delta },0]|\varphi _f=0`$, which finally implies $`\mathrm{\Delta }|f=0|f`$. This ends the proof of the fact that $`\mathrm{\Delta }=0`$, or in another words that $`A=B^{}B`$. This proves therefore the canonical form (23), but not yet the normality of $`B`$. The latter property can be simply proven from the positivity of $`\stackrel{~}{\varrho }^{T_A}`$, which implies that $`BB^{}B^{}B0`$ . The latter (positive) operator has at the same time null trace and therefore it must vanish. Thus $`B`$ is normal as stated.$`\mathrm{}`$
Let us prove now the generalization of Lemma 3 to the case of $`3\times N`$ systems ($`N3`$), and then to the $`M\times N`$ case, where $`MN`$ from now on.
Lemma 4.- Let $`\varrho `$ be a PPT state of rank $`N`$ in a $`3\times N`$ space. Let the reduced state $`\varrho _B`$ and the entry $`E_{33}`$ in some local basis have also the same rank $`N`$. Then $`\varrho `$ can be transformed using some reversible local transformation to the canonical form:
$$\varrho [C,B,I]^{}[C,B,I]$$
(29)
where $`C,B`$ are normal and $`[B,C^{}]=[B,C]=0`$.
Note that in Lemma 4, in contrast to Lemma 3, we assume that in some basis $`r(E_{33})=N`$. Later on, in Theorem 1, we will prove that this assumption is always satisfied.
Proof .- In order to obtain the identity matrix $`I`$ at the diagonal we use an analogous reversible local filter to the one used in the proof of Lemma 3. After that we readily obtain the form
$`\varrho \stackrel{~}{\varrho }=\left[\begin{array}{ccc}C^{}C& D^{}& C^{}\\ D& B^{}B& B^{}\\ C& B& I\end{array}\right].`$ (33)
with both $`B`$ and $`C`$ normal and some unknown $`D`$. Indeed, expression (33) as well as the normality of $`B`$ and $`C`$ follow from the fact that after a local projection by projectors $`P_kI(|kk|+|33|)I`$, $`k=1,2`$ we get a $`2\times N`$ state satisfying the assumptions of Lemma 3.
Now, notice that $`\mathrm{\Psi }_f|\stackrel{~}{\varrho }|\mathrm{\Psi }_f=0`$ for $`|\mathrm{\Psi }_f|2|f|3|Bf`$. Since $`\stackrel{~}{\varrho }0`$ we have that
$$0=\stackrel{~}{\varrho }|\mathrm{\Psi }_f=|1|D^{}fC^{}Bf,$$
(34)
which, as $`f`$ is arbitrary, leads to $`D^{}=C^{}B`$. Thus formula (29) holds. Finally e shall use the latter as well as normality of $`B`$ and $`C`$ to prove that $`[B,C^{}]=[B,C]=0`$. We have
$`\varrho ^{T_A}\stackrel{~}{\varrho }^{T_A}=\left[\begin{array}{ccc}C^{}C& B^{}C& C\\ C^{}B& B^{}B& B\\ C^{}& B^{}& I\end{array}\right],`$ (38)
and we can check that for any $`|f𝒞^N`$ and for $`|\mathrm{\Phi }_f|2|f|3|B^{}f`$, $`\mathrm{\Phi }_f|\stackrel{~}{\varrho }^{T_A}|\mathrm{\Phi }_f=0`$. As $`\varrho `$ is PPT this implies that
$$\varrho ^{T_A}|\mathrm{\Phi }_f=|1|[B^{},C]f$$
(39)
must vanish. Since the above equation holds for arbitrary $`|f`$ we have immediately that $`[B,C^{}]=[C,B^{}]^{}=0`$ Normality of $`B`$ and of $`C^{}`$ implies that these operators can be decomposed as a complex linear combination of projectors into eigenvectors. That they commute means that they actually have the same eigenvectors, and thus so do $`B`$ and $`C`$, i.e. $`[B,C]=0`$.$`\mathrm{}`$
Lemma 5.- Any PPT state supported on a $`M\times N`$ space ($`MN`$) satisfying that (i) $`r(\varrho )=N`$, (ii) in some product basis $`r(E_{ii})=N`$ for some $`i`$, can be transformed after a reversible local transformation to the canonical form:
$$\varrho Z^{}Z=[C_1,\mathrm{},C_{M1},I]^{}[C_1,\mathrm{},C_{M1},I]$$
(40)
with $`[C_i,C_j^{}]=[C_i,C_j]=0`$, $`i,j=1,\mathrm{},M1`$.
Proof .- It follows easily from the application of Lemmas 3 and 4. In particular one has to use the local projections $`PI=(|kk|+|MM|)I`$, $`1k<M`$, $`P^{}I=(|mm|+|m^{}m^{}|+|MM|)I`$, $`1m<m^{}<M`$.$`\mathrm{}`$
As an immediate consequence we have
Lemma 6 .- Any PPT state supported on a $`M\times N`$ space ($`MN`$) satisfying that (i) $`r(\varrho )=N`$, (ii) in some product basis $`r(E_{ii})=N`$ for some $`i`$, is separable and can be expressed as
$$\stackrel{~}{\varrho }=\underset{i=1}{\overset{N}{}}|e_i,f_ie_i,f_i|,$$
(41)
where the $`\{|e_i\}`$ are possibly unnormalized and the $`\{|b_i\}`$ are linearly independent.
Proof .- We make use of Lemma 5. It is easy to see that the matrix $`Z^{}Z`$ has nonzero eigenvectors of the form $`|a_i,b_i`$. Here $`|b_i`$ is the $`i`$th common eigenvector of all operators $`C_j`$, $`C_k^{}`$ while $`|e_i^{}=[c_i^{(1)},\mathrm{},c_i^{(M1)},1]`$ is a row of all $`i`$th eigenvalues of matrices $`C_1,\mathrm{},C_{M1},I`$. Thus, after some reversible local transformation the state $`\rho `$ becomes:
$$\stackrel{~}{\varrho }=\underset{i=1}{\overset{N}{}}|e_i,b_ie_i,b_i|,$$
(42)
where the $`\{|b_i\}`$ are orthonormal. Reversing the previous local filtering we obtain (41) $`\mathrm{}`$.
Remark .- The above procedure gives a constructive algorithm to decompose any state which satisfies the assumptions of the Lemma.
The main disadvantage of the above results is that all of them contain assumptions about $`r(E_{ii})=N`$ for some $`i`$ and for some product basis, which as we have mentioned, are not necessary. Our main theorem is free of that assumption (i.e., it shows that such $`E_{ii}`$ always exists). To prove it we have to use induction with respect to $`M+N=K`$ and use the previous Lemmas. We consider only $`r(\varrho )=N`$, as a PPT state supported on $`M\times N`$ cannot have smaller rank. Indeed, since $`r(\varrho _B)=N`$, if $`r(\varrho )<N`$ then $`\varrho `$ is distillable, which implies that $`\varrho ^{T_A}`$ is not positive .
Theorem 1.- All rank $`N`$ PPT states $`\varrho `$ supported on $`M\times N`$ are separable.
Proof.- We will prove that in some product basis we have $`r(E_{ii})=N`$ for some $`i`$. Separability of $`\varrho `$ will follow from the previous Lemmas. Let us observe that the Theorem and the latter fact are true for $`M=2`$ and arbitrary $`N2`$. In particular they are true for $`M+N=K=4`$ and 5. Let us assume that they hold for $`M+NK`$. We shall now demonstrate that they also hold for $`M+N=K+1`$.
To this aim let us consider the case of $`\varrho `$ supported on a $`(M+1)\times N`$ space, with $`M+1N`$, $`r(\varrho _A)=M+1`$, $`r(\varrho _B)=N`$. and $`M+N=K`$. In an orthonormal, product basis representation the state $`\varrho `$ has the form of a $`(M+1)\times (M+1)`$ matrix with $`N\times N`$ entries
$`\varrho =\left[\begin{array}{cccc}E_{11}& E_{12}& \mathrm{}& E_{1,M+1}\\ E_{12}^{}& E_{22}& \mathrm{}& \mathrm{}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ E_{1,M+1}^{}& \mathrm{}& \mathrm{}& E_{M+1,M+1}\end{array}\right].`$ (47)
Let us consider the following $`M\times M`$ submatrix of $`\varrho `$
$`W(\varrho )W=\left[\begin{array}{cccc}E_{22}& E_{23}& \mathrm{}& E_{2,M+1}\\ E_{23}^{}& E_{22}& \mathrm{}& \mathrm{}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ E_{2,M+1}^{}& \mathrm{}& \mathrm{}& E_{M+1,M+1}\end{array}\right],`$ (52)
resulting after removing the first row and the first column from the representation (47). As the latter action can be achieved by a local projection on Alice’s side, $`W`$ is an unnormalized PPT state acting in $`M\times N`$. For Bob’s reduced matrix $`W_B=Tr_B(W)`$, we shall consider two alternative possibilities: (i) $`r(W_B)=N`$, (ii) $`r(W_B)<N`$.
In case (i) we must have $`r(W)=N`$ as otherwise we would have that the global rank is less than one of the local ranks, resulting in distillability of $`\varrho `$, ergo in violation of the PPT condition . But it means that $`W`$ is a PPT state supported on $`M\times N`$ with global and local rank equal to $`N`$. According to the induction assumption, it is thus separable and for some product basis has an entry $`E_{ii}`$ for some $`i=2,\mathrm{},M+1`$ with the rank $`N`$. But then $`\varrho `$ has an entry $`E_{ii}`$ with rank $`N`$ in the same product basis, and from Lemmas 5 and 6 it follows immediately that it is separable.
Consider now case (ii). If $`W`$ has $`r(W_B)<N`$, then obviously there exist a sequence of product vectors $`|a_i,fK(W)`$, $`i=2,\mathrm{},M+1`$. It is immediate to check that they must belong to kernel of $`\varrho `$. That means that the assumptions of the Lemma 1 are fulfilled. The possibility (i) of this Lemma cannot hold because otherwise one could embed $`\varrho `$ into $`(M+1)\times (N1)`$ space, and $`r(\varrho _B)`$ would be $`N1`$ instead of $`N`$. The possibility (ii) of Lemma 1 means that $`\varrho `$ can be written in the form (cf. Lemma 2)
$$\varrho =\varrho ^{}+\lambda |1,g1,g|$$
(53)
where $`\varrho ^{}`$ is a rank $`N1`$ PPT state supported either on a $`(M+1)\times (N1)`$ subspace or on a $`M\times (N1)`$ subspace, $`\lambda ^11,g|\rho ^1|1,g`$ and $`|1,g1,g|`$ is an unnormalized projector onto a product state such that $`\varrho ^1|1,g`$ is orthogonal to $`R(\varrho ^{})`$.
At the same time it must hold that $`r(\varrho _B^{})=N1`$, since i) Bob’s space has now only $`N1`$ dimensions, ii) $`r(\varrho _B^{})`$ cannot be smaller than $`N1`$, since $`N=r(\varrho _B=\varrho _B^{}+|gg|)`$ and $`|gg|`$ can increase at most in one unit the rank of $`\varrho _B^{}`$. All that means that the matrix $`\varrho ^{}`$ fulfills the induction assumption as $`(M+1)+(N1)=K`$ (or $`M+(N1)=K1`$ ) and $`r(\varrho ^{})=r(\varrho _B^{})`$, ergo it is separable and has in some product basis $`|a_i,b_j`$ the entry $`E_{ii}^{}=a_i|\varrho ^{}|a_i`$ with rank $`N1`$. Lemma 6 implies then that $`\varrho `$ ($`=\varrho ^{}+\lambda |1,g1,g|`$) can be decomposed into
$$\underset{i=1}{\overset{N1}{}}|e_i,f_ie_i,f_i|+\lambda |1,g1,g|,$$
(54)
where $`|g`$ is linearly independent from the set of (also linearly independent) vectors $`|b_i`$. Since there is only a finite number of projectors in the decomposition above, we can always find a vector $`|a`$ in Alice’s space such that $`e_i|a01|e`$. Including such a vector in a product basis to express $`\varrho `$ we will obtain the wished rank $`N`$ element $`a|\varrho |a`$. This completes the proof of the induction step, and by induction completes thus the proof of the theorem. $`\mathrm{}`$
## IV Separability criteria for rank$`(\varrho )+`$rank$`(\varrho ^{T_A})2MNMN+2`$
In this section we generalize the results obtained for $`2\times N`$ systems in Ref. . The idea is that a PPT density operator $`\varrho `$ with $`r(\varrho )+r(\varrho ^{T_A})2MNMN+2`$ may have a finite number of product vectors $`|e_i,f_i`$ in its range, such that $`|e_i^{},f_iR(\varrho ^{T_A})`$. These product vectors are the only possible candidates to appear in decomposition (1) . Finding them requires solving a system of polynomial equations. First we show how to solve these equations in a generic case (namely when the coefficients of such equations do not happen to satisfy a large series of conditions, which amounts to having only a finite number of solutions), and once all the product states $`\{|e_i,f_i\}_{i=0}^{L<\mathrm{}}`$ have been obtained, we present an algorithmic method to check whether $`\rho `$ is separable. This is done in a finite number of computational steps, and thus solves operationally the problem of separability for states with $`r(\varrho )+r(\varrho ^{T_A})2MNMN+2`$ and finite $`L`$.
### A Eligible product vectors.
Let the linearly independent vectors $`|K_i`$, $`|\stackrel{~}{K}_i`$ form a basis in the kernel of $`\varrho `$ and in the kernel of $`\varrho ^{T_A}`$, respectively:
$`K(\varrho )`$ $`=`$ $`\text{span}\{|K_i,i=1,\mathrm{},k(\varrho )\}`$ (55)
$`K(\varrho ^{T_A})`$ $`=`$ $`\text{span}\{|\stackrel{~}{K}_i,i=1,\mathrm{},k(\varrho ^{T_A})\}`$ (56)
We consider here the case when $`k(\varrho )+k(\varrho ^{T_A})M+N2`$. We can always expand $`|K_i`$ and $`|\stackrel{~}{K}_i`$ in an orthonormal basis in Alice’s space
$`|K_i`$ $`=`$ $`{\displaystyle \underset{m=0}{\overset{M}{}}}|m,k_i^m,`$ (57)
$`|\stackrel{~}{K}_i`$ $`=`$ $`{\displaystyle \underset{m=0}{\overset{M}{}}}|m,\stackrel{~}{k}_i^m.`$ (58)
A product vector $`|e,f`$ belonging to the range $`R(\varrho )`$ must be orthogonal to all $`|K_i`$; simultaneously, if its partial complex conjugation belongs to $`R(\varrho ^{T_A})`$, $`|e^{},f`$ must be orthogonal to all $`|\stackrel{~}{K}_i`$. Thus the eligible product vectors are the solutions of $`k(\varrho )+k(\varrho ^{T_A})`$ equations, namely
$`K_i|e,f=0,i=1,\mathrm{},k(\varrho ),`$ (59)
$`\stackrel{~}{K}_i|e^{},f=0,i=1,\mathrm{},k(\varrho ^{T_A}).`$ (60)
Let us now expand $`|e`$ in the above formula as:
$$|e=\left[\begin{array}{c}\alpha _1\\ \mathrm{}\\ \alpha _M\end{array}\right].$$
(61)
We restrict ourselves to $`\alpha _1=1`$. The reason is that we expect to find only a finite number $`L`$ of inequivalent vectors $`|e_i,f_i`$ that fulfill the requirements. A generic choice of an orthonormal basis $`\{|a_i\}`$ in Alice’s space will imply that $`1|a_i0`$ for all $`i=1,\mathrm{},L`$. In this basis $`\alpha _1`$ can be set equal to 1.
Equations (60) can be rewritten as follows:
$$A(\alpha _1,\mathrm{},\alpha _M;\alpha _1^{},\mathrm{},\alpha _M^{})|f=0,$$
(62)
where the $`[k(\varrho )+k(\varrho ^{T_B})]\times N`$ matrix $`A`$ is defined as follows:
$`A(\alpha _1,\mathrm{},\alpha _M;\alpha _1^{},\mathrm{},\alpha _M^{})\left[\begin{array}{c}_{m=1}^M\alpha _mk_1^m|\\ \mathrm{}\\ _{i=1}^M\alpha _mk_{k(\varrho )}^m|\\ _{i=1}^M\alpha _m^{}\stackrel{~}{k}_1^m|\\ \mathrm{}\\ _{i=1}^M\alpha _m^{}\stackrel{~}{k}_{k(\varrho ^{T_A})}^m|\end{array}\right]`$ (69)
$`\left[\begin{array}{c}D_{k(\varrho )\times N}(\alpha )\\ \stackrel{~}{D}_{k(\varrho ^{T_B})\times N}(\alpha ^{})\end{array}\right].`$ (72)
If (62) holds for some $`|f0`$ and $`|e0`$, this means that for the corresponding set of $`\alpha `$’s the rank of $`A`$ is smaller than $`N`$. Therefore, in order to identify eligible product vectors we have to require that at most $`N1`$ rows of $`A`$ be linearly independent vectors. In what follows we restrict ourselves to the limiting case $`k(\varrho )+k(\varrho ^{T_A})=M+N2`$, the others containing more restrictions and consequently less solutions than this.
Let us then take $`N1`$ rows of $`A`$, say the first ones, and let us require that each of the remaining $`M1`$ rows be linearly dependent of them. Recall that we can use the $`M1`$ variables $`\alpha _2,\mathrm{},\alpha _M`$ in order to achieve this. Then, parameter counting strongly suggests that we need to fix all the $`M1`$ $`\alpha `$’s in order to make $`A`$ have rank smaller than $`N`$, this corresponding to a zero measure set of points in the $`\alpha `$-space $`[\alpha _1=1,\alpha _2,\mathrm{},\alpha _M]`$. We will in addition relate the number of solutions to the roots of complex polynomials, which under generic conditions have only a finite number of roots. Numerical experience acquired for the $`2\times N`$ case further supports the expectation that the number of solutions be typically finite.
### B Generic Polynomials.
Let us discuss a bit further sufficient conditions for the existence of a finite set of solutions, while presenting a systematic method to find them once the conditions are fulfilled. This method also works for $`k(\varrho )+k(\varrho ^{T_A})<M+N2`$ by just adding more equations.
Matrix $`A`$ will have at most rank $`N1`$ after requiring that all its rank $`N`$ minors vanish. At risk of finally finding more solutions than just those of equations (62), we can impose that only $`M1`$ of these minors vanish. The reason for doing so is that this will already allow us to prove that only a finite number of product vectors fulfill (60) under some generic circumstances. Thus we consider the determinant of $`N\times N`$ submatrices of $`A`$ formed by taking its first $`N1`$ rows and then also one of the $`M1`$ remaining ones. We shall denote these minors by $`F_i(z_1,\mathrm{},z_{2M})`$, $`i=1,\mathrm{},M1`$, where $`z_j\alpha _j`$ and $`z_{j+M}\alpha _j^{}`$ ($`i=1,\mathrm{},M`$) will be taken as $`2(M1)`$ independent variables ($`z_1=z_{M+1}1`$). Again, this will only imply that when we now set
$`F_i(z_1,\mathrm{},z_{2M})`$ $`=`$ $`0,`$ (73)
for $`i=1,\mathrm{},M1`$, some of the solutions we find do not correspond to product vectors, although all the $`|e_i,f_i`$ we look for are among the solutions of (73).
We have $`2(M1)`$ variables $`z_i`$ and the same number $`2(M1)`$ of polynomial equations for them, $`M1`$ coming from the minors $`F_i(z_1,z_{2M})=0`$ and the remaining $`M1`$ from its complex conjugation, which are inequivalent to the first ones as the variables are mapped according to $`z_iz_{i+M}`$, $`i=1,\mathrm{},M`$, under complex conjugation.
No theorem exists for complex polynomials $`P(\stackrel{}{\alpha },\stackrel{}{\alpha }^{})`$ which allows us to know the number of roots they have. However, in a generic case, namely when $`P(\stackrel{}{\alpha },\stackrel{}{\alpha }^{})`$ is not proportional to its complex conjugate, we can prove that only a finite number of solutions exist. In a method to find such roots was developed for polynomials depending on one $`\alpha `$ and its complex conjugate. Accordingly, from $`P(\alpha ,\alpha ^{})`$ another polynomial $`Q(\alpha )`$ containing all the roots of $`P`$ was obtained. Such a method admits a generalization to the present case, which we shall discuss later on by means of an example when analyzing states of a $`3\times 3`$ system. As already mentioned, we were not able to determine when a density matrix $`\varrho `$ will lead to a set of non-generic polynomials. However, we expect this to be rarely the case. In what follows we will assume that the polynomials derived from $`\varrho `$ are generic, and that therefore there is only a finite number of product vectors that can appear in (1).
### C Separability criterion.
When the number of solutions of equation (60) is finite, we can formulate a necessary and sufficient separability condition which follows from the following general theorem:
Theorem 2 (see also ) .- A state $`\varrho `$ of rank $`r(\varrho )`$ is separable iff it can be written as a convex combination of at most $`\mathrm{min}\{r(\varrho )^2,r(\varrho ^{T_A})^2\}`$ linearly independent projectors $`|e_i,f_ie_i,f_i|`$ onto product vectors.
Proof.- The inverse implication is obvious. For the direct implication we will assume, without loss of generality, that $`r(\varrho )r(\varrho ^{T_A})`$. Caratheodoris’ theorem tells us then that $`\varrho `$ can be expressed as a convex combination of $`r(\varrho )^2`$ product projectors,
$$\varrho =\underset{i=1}{\overset{r(\varrho )^2}{}}p_i|e_i,f_ie_i,f_i|.$$
(74)
Suppose these projectors are not linearly independent. This means we can find $`_ic_i|e_i,f_ie_i,f_i|=0`$ with at least some non-vanishing $`c_i`$. Set $`\lambda \mathrm{min}\{p_i/c_i\}`$. Then the decomposition
$$\varrho =\underset{i=1}{\overset{r(\varrho )^2}{}}(p_i\lambda c_i)|e_i,f_ie_i,f_i|,$$
(75)
corresponds also to a convex combination of the previous projectors $`|e_i,f_ie_i,f_i|`$, but with at least one of the terms having vanishing weight. Now, if the remaining projectors do not form yet a linearly independent set, we can repeat the same procedure and get rid of another product projector. This can be iterated until expressing $`\varrho `$ as a convex combination of linearly independent product projectors. $`\mathrm{}`$
Consequently, once we obtain all product vectors $`|e_i,f_iR(\varrho )`$ such that $`|e_i^{},f_iR(\varrho ^{T_A})`$, $`i=1,\mathrm{},L<\mathrm{}`$, we can find out whether $`\varrho `$ is separable by proceeding as follows:
* We build all possible maximal subsets of linearly independent projectors $`|e_i,f_ie_i,f_i|`$ (with at least $`L_0\mathrm{max}\{r(\varrho ),r(\varrho ^{T_A})\}`$ elements). Notice that there is only a finite number of subsets.
* For each of these subsets we express $`\varrho `$ as a linear combination of projectors in the subset.
* If this is possible, then we have to see whether the coefficients of the linear combination are all positive.
We immediately have:
Separability criterion.- $`\varrho `$ is separable iff all coefficients are non-negative in (at least) one of the linear combinations described above.
### D Numerical methods.
We notice that for a $`\varrho `$ with just a finite, but large number $`L`$ of eligible product vectors it may be impractical to construct all possible subsets of linearly independent product projectors, as described above. In this case the linear programming theory has developed various methods to try to find out a solution to whether $`\varrho `$ can be expressed as a linear combination, with positive weights, of the over complete but finite set of projectors $`|e_i,f_ie_i,f_i|`$. We propose, however, to use for this aim the best separable approximation (BSA) method, developed by us in Ref. . It has nice physical analogies also for non-separable states, providing the expansion
$$\varrho =\varrho _s+(1\lambda )\delta \varrho ,$$
where $`\varrho _s=_i\mathrm{\Lambda }_iP_i`$ is a separable state, $`\lambda =_i\mathrm{\Lambda }_i`$ is maximal, and finally $`\delta \varrho `$ is a state that does not have any product vector in its range. The paper describes an efficient algorithm for finding such expansion, by optimizing each of the $`\mathrm{\Lambda }_i`$ individually, and each of the pairs $`\mathrm{\Lambda }_i`$, $`\mathrm{\Lambda }_j`$ with respect to $`\mathrm{\Lambda }_i+\mathrm{\Lambda }_j`$. For the purpose of checking if a given matrix is separable, the BSA method of Ref. is sufficient; in the context of the present paper it is interesting to introduce here a generalization of the results of to the PPT states :
Lemma 7.- Let $`\varrho `$ be a PPT state. For a given set of $`P_i=|e_i,f_ie_i,f_i|`$, such that the product vectors $`|e_i,f_iR(\delta \varrho )`$, such that $`|e_i^{},f_iR(\delta \varrho ^{T_A})`$, there exists the best separable approximation of $`\varrho `$, in the form
$$\varrho =\varrho _s+(1\lambda )\delta \varrho ,$$
where $`\varrho _s=_i\mathrm{\Lambda }_iP_i`$ is a separable state, $`\lambda =_i\mathrm{\Lambda }_i`$ is maximal, and finally both $`\delta \varrho 0`$, and $`\delta \varrho ^{T_A}0`$. Moreover, there does not exist a product vector $`|e,fR(\delta \varrho )`$, such that $`|e^{},fR(\delta \varrho ^{T_A})`$.
The proof of the above lemma is the same as the proof in Ref. . Similarly, one can find an efficient algorithm for finding the BSA, by requiring that:
* All $`\mathrm{\Lambda }_i`$ should be maximal, i.e.
$`\mathrm{\Lambda }_i=\mathrm{min}(e_i,f_i|(\rho {\displaystyle \underset{ji}{}}\mathrm{\Lambda }_jP_j)^1|e_i,f_i^1,`$ (76)
$`e_i^{},f_i|(\rho ^{T_A}{\displaystyle \underset{ji}{}}\mathrm{\Lambda }_jP_j^{T_A})^1|e_i^{},f_i^1).`$ (77)
* All pairs of $`\mathrm{\Lambda }_i,\mathrm{\Lambda }_j`$ should be maximized with respect to $`\mathrm{\Lambda }_i+\mathrm{\Lambda }_j`$. This requirement can also be expressed in an analytical form for $`\mathrm{\Lambda }`$’s, which will be presented elsewhere .
### E Example: $`3\times 3`$ system.
We end this section by describing with an example in a $`3\times 3`$ system how to estimate the number $`L`$ of eligible product vectors. This example illustrates how to generalize to several independent $`\alpha `$’s the method developed in .
Suppose $`r(\varrho )4`$ and $`r(\varrho ^{T_A})9`$. For $`r(\varrho )=4`$ and $`r(\varrho ^{T_A})=9`$ (least favorable case) we have that the matrix $`A`$ reads
$`A=\left[\begin{array}{c}k_1^1|+\alpha _2k_1^2|+\alpha _3k_1^3|\\ k_2^1|+\alpha _2k_2^2|+\alpha _3k_2^3|\\ k_3^1|+\alpha _2k_3^2|+\alpha _3k_3^3|\\ k_4^1|+\alpha _2k_4^2|+\alpha _3k_4^3|\\ k_5^1|+\alpha _2k_5^2|+\alpha _3k_5^3|\end{array}\right],`$ (83)
so that after constructing the $`3\times 3`$ submatrices $`A_{1,2,3}`$ by taking the first two rows of $`A`$ and one of the remaining rows at a time, we obtain three 3-rd order equations for $`\alpha _1`$ and $`\alpha _2`$:
$`F_1`$ $`=`$ $`detM_1{\displaystyle \underset{k=0}{\overset{3}{}}}\alpha _2^kP_3^k(\alpha _3)=0,`$ (85)
$`F_2`$ $`=`$ $`detM_2{\displaystyle \underset{k=0}{\overset{3}{}}}\alpha _2^kQ_3^k(\alpha _3)=0,`$ (86)
$`F_3`$ $`=`$ $`detM_3=0,`$ (87)
where $`P_s(x)`$ denotes a $`s`$-th order polynomial in $`x`$. By only using equations (85) and (86) we can obtain two quadratic equations in $`\alpha _2`$ as follows: on the one hand we multiply (85) by $`Q_3^3(\alpha _3)`$, (86) by $`P_3^3(\alpha _3)`$, and then subtract them, leading to
$$\underset{k=0}{\overset{2}{}}\alpha _2^kR_6^k(\alpha _3)=0;$$
(88)
on the other hand we multiply (85) by $`Q_3^0(\alpha _3)`$, (86) by $`P_3^0(\alpha _3)`$, again subtract them, and after dividing by $`\alpha _2`$ we obtain
$$\underset{k=0}{\overset{2}{}}\alpha _2^kS_6^k(\alpha _3)=0.$$
(89)
Finally, applying the same trick but now to equations (88) and (89), we obtain two linear equations for $`\alpha _2`$, from which a unique $`18`$-th order equation for $`\alpha _3`$ follows. Therefore there are at most $`18`$ different values of $`\alpha _3`$ which in principle could lead to an eligible product vector. For each such values one should now still solve the $`3`$ $`3`$rd order equations (85-87) for $`\alpha _2`$, and see which solutions survive, if any Notice that for a given $`\alpha _3`$ in principle we could find $`0,1,2`$ or $`3`$ valid values of $`\alpha _2`$. For simplicity we assume in the final estimation of the number $`L`$ of eligible product vectors that to each solution $`\alpha _3`$ there corresponds at most one valid $`\alpha _2`$. . Finally, for those triads $`[1,\alpha _2,\alpha _3]`$ which indeed fulfill (85-87) we can diagonalize $`A`$ and find the corresponding Bob’s local vector $`|f`$ in the kernel of $`A`$. We have obtained, thus, $`L18`$.
Before going into the conclusions we shall discuss briefly the question of the relative size of $`r(\varrho )`$ and $`r(\varrho ^{T_A})`$. It is natural to expect that this difference is not too big. However some naive intuitions must be abandoned (see ). Here we shall make the simple observation:
Observation .- Let $`\varrho `$ be a PPT state. If kernel of $`\varrho `$ contains the range of some PPT state $`\sigma `$, then the kernel of $`\varrho ^{T_A}`$ contains the range of $`\sigma ^{T_A}`$, so that $`r(\varrho ^{T_A})MNr(\sigma ^{T_A})`$.
The above observation about rank of $`\varrho `$ follows easily from the fact that Tr$`(AB)=`$Tr$`(A^{T_A}B^{T_A})`$. Note that $`\sigma `$ can be a separable state. In particular, if the kernel of $`\varrho `$ contains any system of $`n`$ orthogonal product vectors (in particular UPB set ) then $`r(\varrho ^{T_A})`$ can not exceed the value of $`MNn`$. The same holds if $`\sigma `$ from our observation is PPT bound entangled state defined as a UPB complement . The rank of the latter does not change under partial transpose, so again $`r(\varrho ^{T_A})`$ can not exceed the value of $`MNr(\sigma )`$. It can be also extended in other direction: taking $`\sigma `$ as a nontrivial PPT invariant state. Apart from all $`\sigma `$’s being complements of real UPB’s, there is an other nontrivial class (provided in ) of $`N\times N`$ states of that kind all having $`r(\sigma )=\frac{N(N1)}{2}+1`$. From the above discussion and the Theorem 1 we immediately know, for example, that for all the $`3\times 3`$ PPT entangled states with the kernel containing UPB complement both ranks: $`r(\varrho ^{T_2})`$ and $`r(\varrho ^{T_2})`$ must amount to either $`4`$ or $`5`$ so they cannot differ much from each other.
## V Conclusions
We have presented in this paper a relatively complete list of separability criteria for density matrices of low rank. There are several problems, however, which remain open and are worth further studies:
* In our analysis of the kernels of $`\varrho `$ and $`\varrho ^{T_A}`$ we have essentially used only those of their properties that are consequences of the dimensionality. On the other hand, it is expectable that both kernels are structurally related through the partial transpose operation. It would be important to investigate such relations, since it would probably automatically put much more stringent restrictions on the existence of separable matrices of low rank.
* All of the results of this paper can be generalized to the case of multipartite systems, and in particular $`3`$ partite systems. We have already obtained several results, but we leave a detailed and complete discussion of this problem to a separate publication. Let us just mention here that according to our studies: i) there are no rank $`N`$ PPT entangled states for $`N\times N\times N`$ systems; ii) In $`2\times 2\times 2`$ spaces PPT states of rank 4 are separable with respect to $`2\times 4`$ space of Alice and joint space of Bob and Charles, and posses a unique decomposition into a sum of 4 projectors onto product vectors in $`2\times 4`$ space; they are fully separable iff those 4 product vectors are at the same time product vectors in the sense of $`2\times 2\times 2`$; iii) In $`2\times 2\times 2`$ spaces generic PPT states with $`r(\varrho )+r(\varrho ^{T_A})+r(\varrho ^{T_B})+r(\varrho ^{T_C})4\times 82\times 2+1=29`$ have a finite number of product vectors in their range, such that the partial conjugates of those product vectors are in the corresponding ranges of partial transposes.
This work has been supported by Deutsche Forshungsgemeinschaft (SFB 407 and Schwerpunkt “Quanteninformationsverarbeitung”), the Österreichisher Fonds zur Fr̈derung der wissenschaftlichen Forschung (SFB P11), the European TMR network ERB-FMRX-CT96-0087, and the Institute for Quantum Information Gmbh. J. I. C. thanks the University of Hannover for hospitality. P. H. acknowledges the grant from Deutscher Akademisher Austauschdienst. We thank S. Karnas, A. Sanpera, J. Smolin, B. Terhal for fruitful discussions. We thank J. Werner for indicating to us relations to linear programming theory. |
warning/0002/nlin0002018.html | ar5iv | text | # Chaos or Noise — Difficulties of a Distinction
## I Introduction
It is a long debated question if and by what means we can distinguish whether an observed irregular signal is deterministically chaotic or stochastic . If the signal was obtained by iterating a certain model on a computer, we can give a definite answer, because we know the law which generated the signal.
In the case of time series recorded from experimental measurements, we are in a totally different situation. Indeed, in most of cases, there is no unique model of the “system” which produced the data. Moreover, we will see that knowing the character of the model might not be an adequate answer to the question of the character of the signal. For example, data of Brownian motion can be modeled by a deterministic regular process as well as by a deterministic chaotic or stochastic process, as we will show in section IV.
In principle, if we were able to determine the maximum Lyapunov exponent ($`\lambda `$) or the Kolmogorov-Sinai entropy per unit time ($`h_{KS}`$) of a data sequence, we would know, with no uncertainty, whether the sequence is generated by a deterministic law (in which case $`\lambda ,h_{KS}<\mathrm{}`$) or by a stochastic law (in which case $`\lambda ,h_{KS}\mathrm{}`$).
In spite of their conceptual relevance, there are evident practical problems with such quantities that are defined as infinite time averages taken in the limit of arbitrary fine resolution, since, typically, we have access only to a finite (and often very limited) range of scales. In order to cope with these limitations, in this paper we make use of the “finite size Lyapunov exponent” (FSLE) , a variant of the maximum Lyapunov exponent, and the $`(ϵ,\tau )`$-entropy per unit time , a generalization of the Kolmogorov-Sinai entropy per unit time. Basically, while for evaluating $`\lambda `$ and $`h_{KS}`$ one has to detect the properties of a system with infinite resolution, for determining the FSLE, $`\lambda (ϵ)`$, or the $`(ϵ,\tau )`$-entropy per unit time, $`h(ϵ,\tau )`$, the investigation on the system is performed at a finite scale $`ϵ`$, i.e. with a finite resolution. $`\lambda (ϵ)`$ gives us the average exponential rate of the divergence between close (on scale $`ϵ`$) trajectories of a system, and $`h(ϵ,\tau )`$ is the average rate of information needed for prediction. If properly defined, $`h(ϵ,\tau )\stackrel{ϵ,\tau 0}{}h_{KS}`$ and $`\lambda (ϵ)\stackrel{ϵ0}{}\lambda `$, if $`\lambda 0`$. Thus, if we have the possibility of determining the behavior of $`\lambda (ϵ)`$ or $`h(ϵ,\tau )`$ for arbitrarily small scales, as pointed out above, we could answer the original question about the character (deterministic or stochastic) of the law that generated the recorded signal.
However, the limits of infinite time and resolution, besides being unattainable when dealing with experimental data, may also result to be physically uninteresting. As a matter of fact, it is now clear that the maximum Lyapunov exponent and the Kolmogorov-Sinai entropy are not completely satisfactory for a proper characterization of the many faces of complexity and predictability of nontrivial systems, such as, for instance, intermittent systems or systems with many degrees of freedom . For example, in the case of the maximum Lyapunov exponent, one has to consider infinitesimal perturbations, i.e. infinitesimally close trajectories resp. infinite resolution. In systems with many degrees of freedom (e.g. turbulence) an infinitesimal perturbation means, from a physical point of view, that the differences $`\delta x_k=x_k^{}x_k`$ of the components, $`x_k^{}`$ and $`x_k`$, of the initially close state vectors $`𝐱^{}`$ and $`𝐱`$, have to be much smaller than the typical values $`\stackrel{~}{x}_k`$ of the variables $`x_k`$. If the $`\stackrel{~}{x}_k`$-s take very different values, then the concept of infinitesimal perturbation becomes physically unimportant, in the case one is interested only in the evolution of the components with the largest typical values (e.g., the large scales in a turbulent motion).
Taking into account all the limitations mentioned above, in particular the practical impossibility to reach arbitrarily fine resolution, we propose a different point of view on the distinction between chaos and noise. Neither it relies on a particular model for a given data set nor it ignores the fact that the character of a signal may depend on the resolution of the observation. Indeed $`h(ϵ,\tau )`$ (or equivalently $`\lambda (ϵ)`$) usually displays different behaviors as the range of scales is varied. According to these different behaviors, as it will become clear through the paper, one can define a notion of deterministic and stochastic behavior, respectively, on a certain range of scales.
In section II we recall the definitions of the $`(ϵ,\tau )`$-entropy and the finite size Lyapunov exponent. In section III we discuss how one can consistently classify the stochastic or chaotic character of a signal by using the information theoretic concepts as the $`(ϵ,\tau )`$-entropy and the redundancy and compare our approach with previous attempts. In section IV we discuss some examples showing that systems at the opposite in the realm of complexity can give similar results when analyzed from a time series point of view. Section V is devoted to a critical discussion of some recent, intriguing and (sometimes) controversial results for data analysis of “microscopic” chaos, in particular we comment the point of view to be adopted in interpreting the result of sec. IV. In section VI the reader finds some remarks on non trivial behaviors of high dimensional systems. Section VII summarizes and concludes the paper.
## II Two concepts for a resolution dependent time series analysis
### A $`(ϵ,\tau )`$-entropy and redundancy
In this section we recall the definition of the $`(ϵ,\tau )`$-entropy discussing its numerical computation, the possible technical problems as well as its properties.
We start with a continuous (in time) variable $`𝐱(t)\mathrm{I}\mathrm{R}^d`$, which represents the state of a $`d`$-dimensional system, we discretize the time by introducing a time interval $`\tau `$ and we consider the new variable
$$𝐗^{(m)}(t)=(𝐱(t),𝐱(t+\tau ),\mathrm{},𝐱(t+m\tau \tau )).$$
(1)
Of course $`𝐗^{(m)}(t)\mathrm{I}\mathrm{R}^{md}`$ and it corresponds to the discretized trajectory in a time interval $`T=m\tau `$.
Usually, in data analysis, the space where the state vectors of the system live is not known. Mostly, only a scalar variable $`u(t)`$ can be measured. In these cases one considers vectors $`𝐗^{(m)}(t)=(u(t),u(t+\tau ),\mathrm{},u(t+m\tau \tau ))`$, that live in $`\mathrm{I}\mathrm{R}^m`$ and allow a reconstruction of the original phase space, known as delay embedding in the literature . It can be viewed as a special case of (1).
We introduce now a partition of the phase space $`\mathrm{I}\mathrm{R}^d`$, using cells of length $`ϵ`$ in each of the $`d`$ directions. Since the region where a bounded motion evolves contains a finite number of cells, each $`𝐗^{(m)}(t)`$ can be coded into a word of length $`m`$, out of a finite alphabet:
$$𝐗^{(m)}(t)W^m(ϵ,t)=(i(ϵ,t),i(ϵ,t+\tau ),\mathrm{},i(ϵ,t+m\tau \tau )),$$
(2)
where $`i(ϵ,t+j\tau )`$ labels the cell in $`\mathrm{I}\mathrm{R}^d`$ containing $`𝐱(t+j\tau )`$. From the time evolution of $`𝐗^{(m)}(t)`$ one obtains, under the hypothesis of stationarity, the probabilities $`P(W^m(ϵ))`$ of the admissible words $`\{W^m(ϵ)\}`$. We can now introduce the $`(ϵ,\tau )`$-entropy per unit time, $`h(ϵ,\tau )`$ :
$`h_m(ϵ,\tau )`$ $`=`$ $`{\displaystyle \frac{1}{\tau }}[H_{m+1}(ϵ,\tau )H_m(ϵ,\tau )]`$ (3)
$`h(ϵ,\tau )`$ $`=`$ $`\underset{m\mathrm{}}{lim}h_m(ϵ,\tau )={\displaystyle \frac{1}{\tau }}\underset{m\mathrm{}}{lim}{\displaystyle \frac{1}{m}}H_m(ϵ,\tau ),`$ (4)
where $`H_m`$ is the block entropy of block length $`m`$:
$$H_m(ϵ,\tau )=\underset{\{W^m(ϵ)\}}{}P(W^m(ϵ))\mathrm{ln}P(W^m(ϵ)).$$
(5)
For sake of simplicity, we ignored the dependence on details of the partition. For a more rigorous definition one has to take into account all partitions with elements of size smaller than $`ϵ`$, and then define $`h(ϵ,\tau )`$ by the infimum over all these partitions (see e.g. ). In numerical calculations we circumvent this difficulty by using coverings instead of partitions (see below).
A concept which is complementary to the $`ϵ`$-entropy is the $`ϵ`$-redundancy (see e.g. ) which measures the amount of uncertainty on future observations which can be removed by the knowledge of the past, namely:
$`r_m(ϵ,\tau )={\displaystyle \frac{1}{\tau }}[H_1(ϵ,\tau )(H_{m+1}(ϵ,\tau )H_m(ϵ,\tau ,))],`$where $`H_1(ϵ)`$ estimates the uncertainty of the single outcome of the measurement, i.e. neglecting possible correlations in the signal. Alternatively, we can write the redundancy in the form
$$r_m(ϵ,\tau )=\frac{1}{\tau }H_1(ϵ,\tau )h_m(ϵ,\tau ),$$
(6)
which emphasizes the complementarity between redundancy and entropy. If the data are totally independent, i.e. IID, one has $`H_m(ϵ,\tau )=mH_1(ϵ)`$ and, therefore, $`r_m(ϵ,\tau )=0`$. On the opposite side, in the case of a periodic signal the redundancy is maximal $`r_m(ϵ,\tau )=H_1(ϵ,\tau )/\tau `$.
The Kolmogorov-Sinai entropy, $`h_{KS}`$, is obtained by taking the limit $`ϵ,\tau 0`$:
$$h_{KS}=\underset{\tau 0}{lim}\underset{ϵ0}{lim}h(ϵ,\tau ).$$
(7)
The KS-entropy is a dynamical invariant, i.e. it is independent of the used state representation (1), while this is not the case for the $`ϵ`$-entropy (4). To simplify the notation we drop the $`\tau `$ dependence in the following, apart from cases in which the $`\tau `$ dependency is explicitly considered as in section IV.
In a genuine deterministic chaotic system one has $`0<h_{KS}<\mathrm{}`$ ($`h_{KS}=0`$ for a regular motion), while for a random process $`h_{KS}=\mathrm{}`$. The entropies $`H_m(ϵ)`$ were above introduced using a partition and the usual Shannon entropy, however it is possible to arrive at the same notion starting from other entropy-like quantities, which are more suitable for numerical investigations. Following Cohen and Procaccia one can estimate $`H_m(ϵ)`$ as follows. Given a signal composed of $`N`$ successive records and the embedding variable $`𝐗^{(m)}`$, let us introduce the quantities:
$$n_j^{(m)}=\frac{1}{Nm1}\underset{ij}{}\mathrm{\Theta }(ϵ|𝐗^{(m)}(i\tau )𝐗^{(m)}(j\tau )|),$$
(8)
then the block entropy $`H_m(ϵ)`$ is given by
$$H_m^{(1)}(ϵ)=\frac{1}{(Nm)}\underset{j}{}\mathrm{ln}n_j^{(m)}(ϵ).$$
(9)
In practice $`n_j^{(m)}(ϵ)`$ is an approximation of $`P(W^m(ϵ))`$. From the numerical point of view the even more suited quantities are the correlation entropies
$$H_m^{(2)}(ϵ)=\mathrm{ln}\left(\frac{1}{Nm}\underset{j}{}n_j^{(m)}(ϵ)\right)H_m^{(1)}(ϵ),$$
(10)
where one approximates the Shannon entropy by the Renyi entropy of order $`q=2`$.
In the determination of $`h_{KS}`$ by data analysis, one has to consider some subtle points (see for a detailed discussion).
Let us just make some remarks about the general problems in the computation of the Kolmogorov-Sinai entropy from a time series of a deterministic system. The first point is the value of the embedding dimension $`m`$. Let us assume that the information dimension of the attractor of the deterministic system is $`D`$. In order to be able to observe a finite entropy, $`m`$ has to be larger than $`D`$, since the behavior of the entropies in the limit $`ϵ0`$ is
$$h_m(ϵ)=const+O(ϵ)h_{KS},$$
(11)
provided $`m>D`$ . The second relevant point is the fact that the saturation, i.e. the regime where the entropy $`h_m(ϵ)`$ does not depend on the length scale $`ϵ`$, can be observed only on length scales smaller than some $`ϵ_u`$. Thus it is possible to distinguish a deterministic signal from a random one only for $`ϵ<ϵ_u`$. Due to the finiteness of the data set there is a lower scale $`ϵ_l`$ below which no information can be extracted from the data. Taking into account the number of points of the series, $`N`$, it is possible to give the following relation between the embedding dimension, the $`KS`$-entropy, the information dimension and the saturation range $`ϵ_u/ϵ_l`$ :
$$\frac{ϵ_u}{ϵ_l}\left(Ne^{m\tau h_{KS}}\right)^{\frac{1}{D}},$$
(12)
where $`ϵ_u`$ and $`ϵ_l`$ are the upper and lower bounds of the interval of scales at which the deterministic character of a deterministic signal shows up. Note that this relation does not determine $`ϵ_u`$. For more details see .
If $`m`$ is not large enough and/or $`ϵ`$ is not small enough one can obtain misleading results, e.g. see sec. V.
The $`ϵ`$-entropy $`h(ϵ,\tau )`$ is well defined also for stochastic processes. Its dependence on $`ϵ`$ can give some insight in the underlying stochastic process . In the case of finite $`\tau `$ it is possible to define a saturation range: below some length scale $`ϵ_u(\tau )`$ we have
$$h_m(ϵ)=const\mathrm{ln}ϵ+O(ϵ)$$
(13)
But, the limit $`\tau 0`$ will lead to $`ϵ_u0`$, thus the saturation will disappear. As it was shown in , for some stochastic processes it is possible to give an explicit expression of $`h(ϵ,\tau )`$ in this limit. For instance, in the case of a stationary Gaussian process with spectrum $`S(\omega )\omega ^2`$ one has :
$$\underset{\tau 0}{lim}h(ϵ,\tau )\frac{1}{ϵ^2},$$
(14)
the same scaling behavior is also expected for Brownian motion . It can be recovered by looking at $`h(ϵ,\tau )`$ in a certain $`(ϵ,\tau )`$ region. See Ref. for a detailed derivation of (14). We have to stress that the behavior predicted by eq. (14) may be difficult to be experimentally observed due to problems related to the choice of $`\tau `$ (see sect. IV).
### B Finite Size Lyapunov Exponent
The Finite Size Lyapunov Exponent (FSLE) was originally introduced in the context of the predictability problem in fully developed turbulence . Such an indicator, as it will become clear below, is for some aspects the dynamical systems counterpart of the $`ϵ`$-entropy.
The basic idea of the FSLE is to define a growth rate for different sizes of the distance between a reference and a perturbed trajectory. In the following we discuss how the FSLE can be computed, by assuming to know the evolution equations. The generalization to data analysis is obtained following the usual ideas of “embedology” . First one has to define a norm to measure the distance $`ϵ(t)=\delta 𝐱(t)`$ between a reference and a perturbed trajectory. In finite dimensional systems the maximum Lyapunov exponent is independent of the used norm. But when one considers finite perturbations there could be a dependence on the norm (as for infinite dimensional systems). Having defined the norm, one has to introduce a series of thresholds starting from a very small one $`ϵ_0`$, e.g., $`ϵ_n=r^nϵ_0`$ ($`n=1,\mathrm{},P`$), and to measure the “doubling time” ($`T_r(ϵ_n)`$) at different thresholds. $`T_r(ϵ_n)`$ is the time a perturbation of size $`ϵ_n`$ takes to grow up to the next threshold, $`ϵ_{n+1}`$. The threshold rate $`r`$ should not be taken too large, in order to avoid the error to grow through different scales. On the other hand, $`r`$ cannot be too close to one, because otherwise the doubling time would be of the order of the time step in the integration (sampling time in data analysis) affecting the statistics. Typically, one uses $`r=2`$ or $`r=\sqrt{2}`$. For simplicity $`T_r`$ is called “doubling time” even if $`r2`$.
The doubling times $`T_r(ϵ_n)`$ are obtained by following the evolution of the distance $`\delta 𝐱(t)`$ from its initial value $`ϵ_{min}ϵ_0`$ up to the largest threshold $`ϵ_P`$. Knowing the evolution equations, this is obtained by integrating the two trajectories of the system starting at an initial distance $`ϵ_{min}`$. In general, one must choose $`ϵ_{min}ϵ_0`$, in order to allow the direction of the initial perturbation to align with the most unstable direction in the phase-space. Moreover, one must pay attention to keep $`ϵ_P<ϵ_{saturation}`$, so that all the thresholds can be attained ($`ϵ_{saturation}`$ is the typical distance of two uncorrelated trajectories).
The evolution of the error from the initial value $`ϵ_{min}`$ to the largest threshold $`ϵ_P`$ carries out a single error-doubling experiment. At this point the model trajectory is rescaled at a distance $`ϵ_{min}`$ with respect to the true one, and another experiment starts out. After $`𝒩`$ error-doubling experiments, we can estimate the expectation value of some quantity $`A`$ as:
$$A_e=\frac{1}{𝒩}\underset{i=1}{\overset{𝒩}{}}A_i.$$
(15)
This is not the same as taking the time average because each error doubling experiment may take a different time from the others. For the doubling time we have
$$\lambda (ϵ_n)=\frac{1}{T_r(ϵ_n)_e}\mathrm{ln}r;$$
(16)
for details, see . The method described above assumes that the distance between the two trajectories is continuous in time. This is not the case for maps or for discrete sampling in time, thus the method has to be slightly modified. In this case $`T_r(ϵ_n)`$ is defined as the minimum time at which $`ϵ(T_r)rϵ_n`$, and now we have
$$\lambda (ϵ_n)=\frac{1}{T_r(ϵ_n)_e}\mathrm{ln}\left(\frac{ϵ(T_r)}{ϵ_n}\right)_e.$$
(17)
It is worth to note that the computation of the FSLE is not more expensive than the computation of the Lyapunov exponent by the standard algorithm. One has simply to integrate two copies of the system and this can be done also for very complex simulations.
One can expect that in systems with only one positive Lyapunov exponent, one has $`\lambda (ϵ)h(ϵ)`$, see Ref. for details. Additionally it is shown in Ref. how it is possible to use the FSLE for characterizing the predictability (also from the data analysis point of view) of systems containing a slow component and a fast one. Let us comment on some advantages of the FSLE with respect to the $`(ϵ,\tau )`$-entropy. For the FSLE it is not necessary to introduce an $`ϵ`$-partition and, most important, at variance with the $`(ϵ,\tau )`$-entropy, the algorithmic procedure automatically finds the “proper time”; so that it is not necessary to decide on the right sampling time and to test the convergence at varying the words block size $`N`$. This point will be discussed in section IV C.
## III Classification by $`ϵ`$-dependence
In the previous section we discussed the $`ϵ`$-entropy and the FSLE as tools to characterize dynamical processes. Let us re-examine the question of distinguishing chaos and noise posed in the Introduction. Equations (11) and (13) allow us to make rigorous statements about the behavior of the entropy in the limit $`ϵ0`$. Then the behavior of the redundancy can be determined by using the relation to the entropy, given by (6) if we take into account that $`H_1(ϵ)\mathrm{ln}ϵ`$ for continuous valued non-periodic process. Both the behavior of the entropy and the redundancy are summarized in the following table .
deterministic ($`m>D`$) stochastic $`r_m(ϵ)\mathrm{}`$ $`h_m(ϵ)\mathrm{}`$ chaotic non-chaotic $`lim_m\mathrm{}h_m(ϵ)>0`$ $`lim_m\mathrm{}h_m(ϵ)=0`$ white noise colored noise $`r_m(ϵ)=0`$ $`r_m(ϵ)>0`$
The behavior of the FSLE in the limit $`ϵ0`$ is similar to that of the $`ϵ`$-entropy, $`h(ϵ)`$. It is worth noting that the FSLE defined through the doubling times (see sect. II B) is zero also if $`\lambda <0`$.
In all practical situations we have only a finite amount of data. Let us assume we have embedded the time series in a $`m`$-dimensional space, e.g. by time delay embedding. Then one can relate to this set of points an empirical measure $`\mu ^{}`$
$$\mu ^{}(𝐗^{(m)})=\frac{1}{N}\underset{i=1}{\overset{N}{}}\delta (𝐗^{(m)}𝐗^{(m)}(i\tau )).$$
(18)
This empirical measure $`\mu ^{}`$ approximates the true measure only on length scales larger than a finite length scale $`ϵ_l`$. This means that we cannot perform the limit $`ϵ0`$. Of course on a finite scale $`ϵ_l`$, both entropy and redundancy are always finite, therefore we are unable to decide which will reach infinity for $`ϵ0`$. But we can define stochastic and deterministic behavior of a time series at the length scale dependence of the entropy and the redundancy. Fig. 1 shows the typical behavior of the entropy $`h_m(ϵ)`$ and the redundancy $`r_m(ϵ)`$ in case of a deterministic model (2-dimensional chaotic map) and a stochastic model (autoregressive model of first order).
For a time series long enough, a “typical” system can show a saturation range for both the entropy and the redundancy. For decreasing length scales $`ϵ`$ with $`ϵ<ϵ_u`$ one observes the following behaviors
deterministic stochastic $`r_m(ϵ)\mathrm{ln}ϵ`$ $`h_m(ϵ)\mathrm{ln}ϵ`$ $`h_m(ϵ)const`$ $`r_m(ϵ)const`$
In addition, as far as stochastic behaviors are concerned, the $`ϵ`$-entropy can exhibit power laws on large scales, e.g. in the case of diffusion eq. (14) (see section IV and Ref. for further details).
If on some range of length scales either the entropy $`h_m(ϵ)`$ or the redundancy $`r_m(ϵ)`$ is a constant, we call the signal deterministic or stochastic on these length scales, respectively. Thus we have a practical tool to classify the character of a signal as deterministic or stochastic without referring to a model, and we are no longer obliged to answer the metaphysical question, whether the system which produced the data was a deterministic or a stochastic one.
Moreover, from this point of view, we are now able to give the notion of noisy chaos a clear meaning: chaotic scaling on large scales, stochastic scaling on small scales. We have also the notion of chaotic noise, namely stochastic scaling on large scales and deterministic scaling on small scales. These notions will become clear with the examples in the following sections.
The presented method for distinguishing between chaos and noise is a refinement and generalization of one of the first discussed methods to approach this problem: estimating the correlation dimension and taking a finite value as a sign for the deterministic nature of the signal . The main criticism to this approach is based on the work of Osborne and Provenzale , where they claimed that stochastic systems with a power-law spectrum will produce time series which exhibit a finite correlation dimension. A detailed discussion of this problem is beyond the scope of this paper, but a main step to clarify the problem was taken by Theiler . Firstly, he noted that the discussed signals were non-stationary and highly correlated with correlation times of the order of the length of the time series. From a conservative point of view one has to stop at this point with any attempt to calculate dimensions or entropies. If one proceeds nonetheless, Theiler showed that the result will depend on the number of data points and the length scale. If one has a sufficient number of data points one will encounter also for this kind of signals the embedding dimension which leads to the typical behavior as given in (13) for the entropy. Moreover, if one uses the usual time delay embedding like (1) in contrast to and , the result depends strongly on the chosen delay time.
We are aware that there are a lot of other attempts to distinguish chaos from noise discussed in the literature. They are based on the difference in the predictability using prediction algorithms rather than the estimating the entropy or they relate determinism to the smoothness of the signal . All these methods have in common that one has to choose a certain length scale $`ϵ`$ and a particular embedding dimension $`m`$. Thus they could, in principle, shed light at the interesting crossover scenarios we are going to describe in the next chapter.
## IV Difficulties in the distinction between chaos and noise: examples
In this section we analyze in a detailed way some examples which illustrate how subtle the transition from the large scales behavior to the small scales one can be; and thus the difficulties arising in distinguishing, only from data analysis, a genuine deterministic chaotic system from one with intrinsic randomness.
### A The diffusive regime
We first discuss the problems, due to the finite resolution, which one can have in analyzing experimental data. We consider the following map which generates a diffusive behavior on the large scales :
$$x_{t+1}=[x_t]+F\left(x_t[x_t]\right),$$
(19)
where $`[x_t]`$ indicates the integer part of $`x_t`$ and $`F(y)`$ is given by:
$$F(y)=\{\begin{array}{cc}(2+\mathrm{\Delta })y\hfill & \text{if}y[0,1/2[\hfill \\ (2+\mathrm{\Delta })y(1+\mathrm{\Delta })\hfill & \text{if}y]1/2,1].\hfill \end{array}$$
(20)
The maximum Lyapunov exponent $`\lambda `$ can be obtained immediately: $`\lambda =\mathrm{ln}|F^{}|`$, with $`F^{}=dF/dy=2+\mathrm{\Delta }`$. Therefore one expects the following scenario for $`h(ϵ)`$ (and for $`\lambda (ϵ))`$:
$$h(ϵ)\lambda \mathrm{for}ϵ<1,$$
(21)
$$h(ϵ)\frac{D}{ϵ^2}\mathrm{for}ϵ>1,$$
(22)
where $`D`$ is the diffusion coefficient, i.e.
$$\left(x_tx_0\right)^22Dt\mathrm{for}\mathrm{large}t.$$
(23)
Fig.s 3 and 4 show $`\lambda (ϵ)`$ and $`h(ϵ)`$ respectively. Let us briefly comment on a technical aspect. The numerical computation of $`\lambda (ϵ)`$ does not present any particular difficulties; on the other hand the results for $`h(ϵ)`$ depend on the used sampling time $`\tau `$. This can be appreciated by looking at Fig. 25b of the Gaspard and Wang review , where the power law behavior (22) in the diffusive region is obtained only if one considers the envelope of $`h_m(ϵ,\tau )`$ evaluated for different values of $`\tau `$; while looking at a single $`\tau `$ one has a rather not conclusive result. This is due to the fact that, at variance with the FSLE, when computing $`h(ϵ,\tau )`$ one has to consider very large $`m`$, in order to obtain a good convergence for $`H_m(ϵ)H_{m1}(ϵ)`$.
Because of the diffusive behavior, a simple dimensional argument shows that, by sampling the system every elementary time step, a good convergence holds for $`mϵ^2/D`$. Thus, for $`ϵ=10`$ and typical values of the diffusion coefficient $`D10^1`$, one has to consider enormous block size. A possible way out of this computational difficulty may be the following. We call $`l(=1)`$ the length of the interval $`[0,1]`$ where $`F(y)`$ is defined; if we adopt a coarse-grained description on a scale $`ϵ=l`$, i.e. we follow the evolution of the integer part of $`x_t`$, the dynamical system (19) is well described by means of a random walk , with given probability $`p_0`$ that in a time step $`s(=1)`$ the integer part of $`x`$ does not change: $`[x_{t+s}]=[x_t]`$, and probabilities $`p_\pm `$ that $`[x_{t+s}]=[x_t]\pm 1`$. A diffusive behavior means that the probability for a changing of $`[x_t]`$ by $`\pm k`$ in $`n`$ elementary time steps is given by
$$P(k,n)\frac{e^{\left(k^2/2\alpha n\right)\left(l^2/s\right)}}{\sqrt{2\pi \alpha n}}\sqrt{\frac{l^2}{s}},$$
(24)
with $`\alpha `$ a function of $`p_0`$. Now we increase the graining and identify all the $`x`$ values in a cell of size $`ϵ=L`$, with $`L`$ an integer multiple of $`l`$. If we observe the coarse-grained state of the system only every $`\tau (>s)`$ time steps and we define the variables $`\stackrel{~}{k}`$ and $`\stackrel{~}{n}`$, such that $`kl\stackrel{~}{k}L`$ and $`ns\stackrel{~}{n}\tau `$, we can write
$$P(\stackrel{~}{k},\stackrel{~}{n})\frac{e^{\left(\stackrel{~}{k}^2/2\alpha \stackrel{~}{n}\right)\left(L^2/\tau \right)}}{\sqrt{2\pi \alpha \stackrel{~}{n}}}\sqrt{\frac{L^2}{\tau }},$$
(25)
for the probability of finding the system $`\stackrel{~}{k}`$ $`L`$-cells apart after $`\stackrel{~}{n}`$ $`\tau `$-intervals. Thus, by choosing $`L^2/\tau =l^2/s`$, we expect that the sequences generated by checking the system either on a scale $`L`$ every $`\tau `$ steps or on a scale $`l`$ every elementary time step $`s`$, have the same statistics, in particular the same entropy, as a signature of a diffusive behavior. Note that if $`lim_m\mathrm{}H_m(ϵ,\tau =ϵ^2)/m`$ is constant at varying $`ϵ`$, as we found numerically, then the $`ϵ`$-entropy per unit time $`lim_m\mathrm{}H_m(ϵ,\tau )/m\tau `$ goes like $`1/ϵ^2`$.
Since the equality $`L^2/\tau =l^2/s`$ is assured by the choice $`ϵ=L=\gamma l`$ and $`\tau =\gamma ^2s`$, one can expect that, for a diffusion process, the following scaling relation hold: $`lim_m\mathrm{}H_m(ϵ,\tau )/m=lim_m\mathrm{}H_m(\gamma ϵ,\gamma ^2\tau )/m`$, with $`\gamma `$ an arbitrary scaling parameter. This scaling relation allows us to see why the power law behavior (14) is expected to be valid generally for the Brownian motion. Indeed, if we choose $`\gamma =1/ϵ`$ we have $`H_m(ϵ,\tau )=H_m(1,\tau /ϵ^2)`$ and, finally, taking the limit $`\tau 0`$ the $`ϵ`$-entropy is given by:
$$h(ϵ)=\underset{\tau 0}{lim}\frac{\left[H_{m+1}(1,\tau /ϵ^2)H_m(1,\tau /ϵ^2)\right]}{\tau }\underset{\tau 0}{lim}\frac{const\tau /ϵ^2+O(\tau ^2)}{\tau }\frac{1}{ϵ^2}$$
(26)
which is (14). Note that the first equality in (26) has been obtained by a Taylor expansion around $`\tau =0`$, and by noting that $`h(1,0)=0`$ otherwise the entropy for unit time will be infinite at finite $`ϵ`$ which is impossible.
### B Finite resolution effects
We consider now a stochastic system, namely a map with dynamical noise
$$x_{t+1}=[x_t]+G\left(x_t[x_t]\right)+\sigma \eta _t,$$
(27)
where $`G(y)`$ is shown in Fig. 2 and $`\eta _t`$ is a noise with uniform distribution in the interval $`[1,1]`$, and no correlation in time. As it can be seen from Fig. 2, the new map $`G(y)`$ is a piecewise linear map which approximates the map $`F(y)`$. When $`dG/dy<1`$, as is the case we consider, the map (27), in the absence of noise, gives a non-chaotic time evolution.
Now one can compare the chaotic case, i.e. eq. (19) with the approximated map (27) with noise. For example let us start with the computation of the finite size Lyapunov exponent for the two cases. Of course from a data analysis point of view we have to compute the FSLE by reconstructing the dynamics by embedding. However, for this example we are interested only in discussing the resolution effects. Therefore we compute the FSLE directly by integrating the evolution equations for two (initially) very close trajectories, in the case of noisy maps using two different realizations of the noise. In Fig. 3 we show $`\lambda (ϵ)`$ versus $`ϵ`$ both for the chaotic (19) and the noisy (27) map. As one can see the two curves are practically indistinguishable in the region $`ϵ>\sigma `$. The differences appear only at very small scales $`ϵ<\sigma `$ where one has a $`\lambda (ϵ)`$ which grow with $`ϵ`$ for the noisy case, remaining at the same value for the chaotic deterministic chase.
Both the FSLE and the $`(ϵ,\tau )`$-entropy analysis show that we can distinguish three different regimes observing the dynamics of (27) on different length scales. On the large length scales $`ϵ>1`$ we observe diffusive behavior in both models. On length scales $`\sigma <ϵ<1`$ both models show chaotic deterministic behavior, because the entropy and the FSLE is independent of $`ϵ`$ and larger than zero. Finally on the smallest length scales $`ϵ<\sigma `$ we see stochastic behavior for the system (27) while the system (19) still shows chaotic behavior.
### C Effects of finite block length
In the previous section we discussed the difficulties arising in classifying a signal as chaotic or stochastic because of the impossibility to reach an arbitrary fine resolution. Here we investigate the reasons, which make it difficult to distinguish a stochastic behavior from a deterministic non-chaotic one. In particular, we show that a non-chaotic deterministic system may produce a signal practically indistinguishable from a stochastic one, provided its phase space dimension is large enough.
The simplest way to generate a non-chaotic (regular) signal having statistical properties similar to a stochastic one is by considering the Fourier expansion of a random signal . One can consider the following signal:
$$x(t)=\underset{i=1}{\overset{M}{}}X_{0i}\mathrm{sin}\left(\mathrm{\Omega }_it+\varphi _i\right)$$
(28)
where the frequencies are equispaced discrete frequencies, i.e. $`\mathrm{\Omega }_i=\mathrm{\Omega }_0+i\mathrm{\Delta }\mathrm{\Omega }`$, the phases $`\varphi _i`$ are random variables uniformly distributed in $`[0,2\pi ]`$ and the coefficient $`X_{0i}`$ are chosen in such a way to have a definite power spectrum, e.g. a power law spectrum, which is a common characteristic of many natural signals. Of course (28) can be considered as the Fourier expansion of a stochastic signal only if one consider a set of $`2M`$ points such that $`M\mathrm{\Delta }\mathrm{\Omega }=\pi /\mathrm{\Delta }t`$, where $`\mathrm{\Delta }t`$ is the sampling time . Time series as (28) have been used to claim that suitable stochastic signals may display a finite correlation dimension , see the discussion in sect. III.
Here we adopt a slightly different point of view. The signal (28) can also be considered as the displacement of a harmonic oscillator linearly coupled to other harmonic oscillators. Indeed, it is well known since long times that a large ensemble of harmonic oscillators can originate stochastic-like behaviors . In particular, we refer to , where it was proved that an impurity of mass $`\mu `$ linearly coupled to a one-dimensional equal mass, $`\mu _0`$, chain of $`M`$ oscillators coupled by a nearest-neighbor harmonic interaction, in the limit of $`\mu \mu _0`$ and of infinite oscillators ($`M\mathrm{}`$), undergoes a Brownian motion. Practically our observable is given by the sum of harmonic oscillations as in eq. (28), where the frequencies $`\mathrm{\Omega }_i`$ have been derived in the limit $`\mu _0/\mu 1`$ by Cukier and Mazur . The phases $`\varphi _i`$ are chosen as uniformly distributed random variables in $`[0,2\pi ]`$ and the amplitudes $`X_{0i}`$ are chosen as follows:
$$X_{0i}=C\mathrm{\Omega }_i^1$$
(29)
where the $`C`$ is an arbitrary constant and the $`\mathrm{\Omega }`$ dependence is just to obtain a diffusive-like behavior. Note that for a signal of length $`2M`$ the random phases and the $`X_{0i}`$’s represent a initial condition of the $`M`$ oscillators, because their phase space is $`2M`$-dimensional.
In Fig 5a we show an output of the signal (28) and, for comparison, in Fig 5b we also show an artificial continuous time Brownian motion obtained integrating the following equation
$$\frac{\mathrm{d}x(t)}{\mathrm{d}t}=\xi (t)$$
(30)
where $`\xi (t)`$ is a Gaussian white noise, produced by a random number generator (the variance of the process is chosen as to mimic that obtained by (28)) . Because the random number generator uses a high entropic one-dimensional deterministic map, this is an example for a high entropic low dimensional system, which produces stochastic behavior. As it is possible to see the two signals appears to be very similar already at a first sight.
Now it is important to stress that if $`M<\mathrm{}`$ the signals obtained according to (28) cannot develop a true Brownian motion especially if one is interested in long time series. Indeed for a long enough record one should be able to recognize the regularities in the trajectory of $`x(t)`$. However, even if the time record is long enough, in order to give a definite answer about the value of the entropy one also needs very large embedding dimensions. The basic fact is that deterministic behavior can be observed only, if the embedding dimension $`m`$ is larger than the dimension of the manifold where the motion takes place, which is $`M`$ for $`M`$ harmonic oscillators. This means that although the entropy $`h_{KS}`$ is zero, the conditional entropies, $`h_m(ϵ,\tau )=(H_{m+1}(ϵ)H_m(ϵ))/\tau `$, for finite $`m`$ are nonzero and maybe even slowly decreasing for $`m>M`$. Moreover, one can encounter some quasi-convergences with respect to $`m`$ for $`m<M`$, if $`\tau `$ is large enough, i.e. the entropy can seem to be independent of $`m`$, e.g. see Fig. 6.
In Fig. 7a-b we show the $`ϵ`$-entropy, calculated by the Grassberger-Procaccia method . The deterministic signal (28) and the stochastic one (30) produce indeed very similar results. Note that we calculated $`h_m^{(2)}(ϵ,\tau )`$ instead of $`h_m^{(1)}(ϵ,\tau )`$ because it is used more often due to better statistics in most cases. However, in Fig. 8 one can see, that on the relevant length scales both entropies are almost equal.
As in Refs. we have considered different time delays, $`\tau `$, in computing the $`ϵ`$-entropy because of the problems discussed in the previous section. The power law behavior $`ϵ^2`$ for the $`ϵ`$-entropy is finally obtained only as the envelope of different computations with different delay times. The results for the FSLE calculated from the time series are shown in Fig. 9. Both the $`ϵ`$-entropy and the FSLE display the $`1/ϵ^2`$ behavior, which denotes that the signals can be classified as Brownian motion . It is worth noting that the FSLE computed from the time record is not too sensitive on the choice of the delay time $`\tau `$ and the embedding dimension $`m`$.
From this simple example is it easy to understand that the impossibility of reaching high enough embedding dimensions severely limits our ability to make definite statements about the ”true” character of the system which generated a given time series as well as the already analyzed problem of the lack of resolution.
## V Some remarks on a recent debate about “microscopic” chaos
The issue of the detection of “microscopic” chaos by data analysis has recently received some attention after a work of Gaspard et al. . Gaspard et al., from an entropic analysis of an ingenious experiment on the position of a Brownian particle in a liquid, claim to give an empirical evidence for microscopic chaos, i.e.: they claim to give evidence that the diffusive behavior observed for a Brownian particle is the consequence of chaos on the molecular scale. Their work can be briefly summarized as follows: from a long ($`1.510^5`$ data) record on the position of a Brownian particle they compute the $`ϵ`$-entropy with the Cohen-Procaccia method, described in section II, from which they obtain:
$$h(ϵ)\frac{D}{ϵ^2}$$
(31)
where $`D`$ is the diffusion coefficient. Then they assume that the system is deterministic and therefore, because of the inequality $`h(ϵ>0)h_{KS}`$, they conclude that the system is chaotic. However, from the results presented in the previous sections we can understand that their result does not give a direct evidence that the system is deterministic and chaotic. Indeed, the power law (31) can be produced in different ways:
1. a genuine chaotic system with diffusive behavior (as the map (20) of sect. IV A);
2. a non chaotic system with some noise, as the map (27);
3. a genuine Brownian system, as (30);
4. a deterministic linear non chaotic system with many degrees of freedom (as the example 28);
5. “complicated” non chaotic system as the Ehrenfest wind-tree model where a particle diffuses in a plane due to collisions with randomly placed, fixed oriented square scatters (as discussed by Cohen et al. in their comment to Gaspard et al. ).
It seems to us that the very weak points of Gaspard et al. are:
the explicit assumption that the system is deterministic
the neglecting of the limited number of data points and, therefore, of both limitations in the resolution and the block length.
The point (a) is crucial, without this assumption (even with an enormous data set) it is not possible to distinguish between 1) and 2). One has to say that in the cases 4) and 5) at least in principle it is possible to understand that the systems are “trivial” (i.e. not chaotic) but for this one has to use a very huge number of data. For example Cohen et al. estimated that in order to distinguish between 1) and 5) using realistic parameters of a typical liquid, the number of data points required has to be at least $`10^{34}`$. Let us remind that Gaspard et al. used $`1.510^5`$ points!
It seems to us that the distinction between chaos and noise makes sense only in very peculiar cases, e.g. very low dimensional systems. Nevertheless even in such a case an entropic analysis can be unable to recognize the “true” character of the system due to the lack of resolution. This is particularly evident in the comparison between the diffusive map (20) and the noisy map (27), if one only observes at scales $`ϵ>\sigma `$. According to the pragmatic proposal of classification discussed in sect. III one has that for $`\sigma ϵ1`$ both the system (20) and (27), in spite of their “true” character can be classified as chaotic, while for $`ϵ1`$ both can be considered as stochastic.
In this respect, the problem of the lack of resolution is even more severe for high entropic systems. One can roughly estimate the critical $`ϵ`$ ($`ϵ_u`$), below which the saturation can be observed, to be $`ϵ_u\mathrm{exp}(h_{KS})`$. Indeed $`\mathrm{exp}(h_{KS})`$ estimates the number of symbols, i.e. cells of the partition required to reconstruct the dynamics. Therefore, in the Brownian motion studied by Gaspard et al. , where the $`KS`$-entropy is expected to be proportional to the number of molecules present in the fluid, the possible small $`ϵ_u`$ is pushed on scales far from being reachable with the finest experimental resolution available.
## VI Discussions
Here we briefly review two examples, studied in details in , showing that high-dimensional systems can display non-trivial behaviors at varying the resolution scales. We believe this discussion can be useful in furtherly clarify the subtle aspects of the distinction between stochastic and deterministic behaviors in dynamical systems.
Olbrich et al. analyzed an open flow system described by unidirectionally coupled map lattice:
$$x_j(t+1)=(1\sigma )f(x_{j+1}(t))+\sigma x_j(t)$$
(32)
where $`j=1,\mathrm{},N`$ denotes the site of a lattice of size $`N`$, $`t`$ the discrete time and $`\sigma `$ the coupling strength. A detailed numerical study (also supported by analytical arguments) of the $`ϵ`$-entropy $`h(ϵ)`$ at different $`ϵ`$, in the limit of small coupling, gives the following scale-dependent scenario: for $`1ϵ\sigma `$ there is a plateau $`h(ϵ)=\lambda _s`$ where $`\lambda _s`$ is the Lyapunov exponent of the single map $`x(t+1)=f(x(t))`$. For $`\sigma ϵ\sigma ^2`$ another plateau appears at $`h(ϵ)2\lambda _s`$, and so on for $`\sigma ^{n1}ϵ\sigma ^n`$ one has $`h(ϵ)n\lambda _s`$. Similar results hold for the correlation dimension which increases step by step as the resolution increases, showing that the high-dimensionality of the system becomes evident only as $`ϵ0`$.
Therefore one has a strong evidence that the dynamics at different scales is basically ruled by a hierarchy of low-dimensional systems whose ”effective” dimension $`n_{eff}(ϵ)`$ increase as $`ϵ`$ decreases:
$$n_{eff}(ϵ)\left[\frac{\mathrm{ln}(1/ϵ)}{\mathrm{ln}(1/\sigma )}\right],$$
(33)
where $`[\mathrm{}]`$ indicates the integer part. In addition, for a resolution larger than $`ϵ`$, it is possible to find a suitable low-dimensional noisy system (depending on $`ϵ`$) which is able to mimic $`x_1(t)`$ given by eq. (32). It is interesting to note that, looking at $`h(ϵ)`$ on an extended range of values of $`ϵ`$, one observes for $`ϵ\sigma ^N`$
$$h(ϵ)\mathrm{ln}\frac{1}{ϵ}$$
(34)
i.e. the typical behavior of a stochastic process. Of course for $`ϵ\sigma ^N`$ one has to realize that the system is deterministic and $`h(ϵ)N\lambda _s`$. Even if this study mainly concerns the small unidirectional coupling, the diffusive and strong coupling case deserves further analysis, it represents a first step toward the understanding of the issue of data analysis of high-dimensional systems.
Let us now briefly discuss the issue of the macroscopic chaos, i.e. the hydro-dynamical-like irregular behavior of some global observable, with typical times much longer than the times related to the evolution of the single (microscopic) elements composing a certain system. This interesting kind of dynamical behavior has been studied in some recent works for globally coupled maps evolving according to the equation:
$$x_i(t+1)=(1\sigma )f_a(x_i(t))+\frac{\sigma }{N}\underset{j=1}{\overset{N}{}}f_a(x_j(t))$$
(35)
where $`N`$ is the total number of elements and $`f_a`$ is a nonlinear function depending on a parameter $`a`$. Cencini et al. (see also Shibata and Kaneko for a related work) studied the behavior of a global variable (i.e. the center of mass) using the FSLE analysis. Their results can be summarized as follows:
* at small $`ϵ(1/\sqrt{N})`$ one recovers the “microscopic” Lyapunov exponent, i.e. $`\lambda (ϵ)\lambda _{micro}`$
* at large $`ϵ(1/\sqrt{N})`$ one observe another plateau (corresponding to what we can call the “macroscopic” Lyapunov exponent) $`\lambda (ϵ)\lambda _{macro}`$ which can be much smaller than the microscopic one.
The above results suggest that at a coarse-grained level, i.e. $`ϵ1/\sqrt{N}`$, the system can be described by an “effective” hydro-dynamical equation (which in some cases can be low-dimensional), while the “true” high-dimensional character appears only at very high resolution, i.e.
$$ϵϵ_cO\left(\frac{1}{\sqrt{N}}\right).$$
The presence of two plateaux for the FSLE at different length scales is present in generic systems with slow and fast dynamics . The interesting fact is that in systems like (35) the two temporal scales are generated by the dynamics itself.
Let us stress that the behavior $`h(ϵ)=const`$ at small $`ϵ`$ and $`h(ϵ)`$ decreasing for larger $`ϵ`$ is not a peculiarity of the diffusive map (20). In typical high-dimensional chaotic systems one has $`h(ϵ)=h_{KS}O(N)`$ for $`ϵϵ_c`$ (where $`N`$ is the number of degrees of freedom and $`ϵ_c0`$ as $`N\mathrm{}`$) while for $`ϵϵ_c`$, $`h(ϵ)`$ decreases (often with a power law). From this point of view, the fact that in certain stochastic processes $`h(ϵ)ϵ^\alpha `$ can be indeed extremely useful for modeling such high-dimensional systems. As a relevant example we just mention the fully developed turbulence which is a very high-dimensional system whose non infinitesimal (the so-called inertial range) properties can be successfully mimicked in terms of a suitable stochastic process .
## VII Conclusions
We have shown how an entropic analysis at different resolution scales (in terms of $`ϵ`$-entropy and Finite Size Lyapunov Exponent) of a given data record allows us for a classification of the stochastic or chaotic character of a signal.
In practice, without any reference to a particular model, one can define the notion of deterministic or chaotic behavior of a system on a certain range of scales. In our examples we show that, according to the pragmatic classification proposed in Sect. III, one can consider (on a certain resolution) a system as random or deterministic independently from its “true” nature. At first glance this can appear disturbing, however, if one adopts a non “metaphysical” point of view there is an advantage in the freedom of modeling the behavior of the system at least if one is interested on a certain (non infinitesimal) coarse-graining property.
###### Acknowledgements.
We thank R. Hegger and T. Schreiber for useful discussions and suggestions. M.C., M.F. and A.V. have been partially supported by INFM (PRA-TURBO) and by the European Network Intermittency in Turbulent Systems (contract number FMRX-CT98-0175). |
warning/0002/hep-th0002204.html | ar5iv | text | # References
Diverse PBGS Patterns and Superbranes
E. Ivanov
Bogoliubov Laboratory of Theoretical Physics, JINR,
141 980, Dubna, Moscow Region, Russian Federation
This is a brief account of the approach to superbranes based upon the concept of Partial Breaking of Global Supersymmetry (PBGS).
1. Introduction. The view of superbranes as theories of partial spontaneous breaking of global SUSY (PBGS) received a considerable attention (see, e.g., \- ). The salient feature of this approach is that it deals with the Goldstone superfields living on the worldvolume superspace of unbroken SUSY and accommodating the superbrane physical degrees of freedom.
It is well known that the superbranes in the Green-Schwarz (GS) formulation break half of the target SUSY (see, e.g., ). Choosing the static gauge with respect to the worldvolume diffeomorphisms and killing half of the target $`\theta `$-coordinates by $`\kappa `$-symmetry, one ends up with the physical multiplet which comprises transverse target bosonic coordinates and the rest of $`\theta `$-coordinates (in the case of D-branes and their generalizations, the physical multiplets can include additional fields). The remaining $`\theta `$’s are shifted under half of the target SUSY, which suggests to interpret them as the relevant Goldstone fermions. Thus one half of the original SUSY, with respect to which the physical degrees of freedom form a supermultiplet, is unbroken, while the other is spontaneously broken, with the physical fermions as the Goldstone ones.
One can reverse the argument: adopt PBGS as the guiding principle and deduce superbranes just from it. In doing so, one can take advantage of the nonlinear realizations method - which provides the universal framework for treating spontaneously broken symmetries. In this method, the Goldstone (super)fields are identified with the parameters of the coset space of the full (super)symmetry group over its unbroken symmetry subgroup. The invariant actions of the Goldstone superfields constructed as integrals over the worldvolume superspace are reduced to the corresponding GS-type actions after passing to components and eliminating the auxiliary fields.
The nonlinear realizations method was worked out for constructing nonlinear sigma models of internal symmetries. It is less known that the bosonic $`p`$-branes in the physical (or “static”) gauge can also be treated in the nonlinear realizations language . To describe in this way some $`p`$-brane moving in $`D`$-dimensional Minkowski space, one should consider a nonlinear realization of the relevant Poincaré group $`𝒫_D=ISO(1,D1)`$, such that the residual unbroken (vacuum stability) symmetry group is the product of the Poincaré group of $`p+1`$-dimensional brane worldvolume and the rotation group of the transverse brane coordinates. The relevant coset manifold is $`𝒫_D/SO(1,p)SO(Dp1)`$. One splits the full $`D`$ space translation generator $`P_M`$, $`M=0,1,\mathrm{}D1`$, as
$$P_M(P_m,P_\mu ),m=0,1,\mathrm{}p;\mu =p+1,\mathrm{}D,$$
(1)
and associates with $`P_m,P_\mu `$ the worldvolume coordinate $`x^m`$ and the Goldstone field $`X^\mu (x)`$ as the coset parameters ($`X^\mu `$ becomes the transverse $`p`$-brane coordinate). Also, one introduces the Goldstone fields $`\mathrm{\Lambda }_m^\mu (x)`$ parametrizing the spontaneously broken part of the Lorentz group $`SO(1,D1)`$ (with the generators $`L_\mu ^m`$):
$$P_mx^m,P_\mu X^\mu (x),L_\mu ^m\mathrm{\Lambda }_m^\mu (x).$$
(2)
The application of the Cartan forms techniques augmented with some extra covariant constraints (the inverse Higgs ) gives rise to the following minimal invariant action
$$S_{br}d^{p+1}x\left(\sqrt{(1)^pg}1\right),g=\text{det}(\eta _{mn}_mX^\mu _nX^\mu ).$$
(3)
It is the static gauge form of the $`p`$-brane Nambu-Goto (NG) action.
As an illustration, let us consider the simplest example of the massive particle ($`0`$-brane) in $`D=2`$. The $`D=2`$ Poincaré group $`𝒫_{(2)}`$ involves two translation generators $`P_0,P_1`$ and the $`SO(1,1)`$ Lorentz generator $`L`$, with the only non-vanishing commutators
$$[L,P_0]=iP_1,[L,P_1]=iP_0.$$
(4)
In accord with the said above, we should construct a nonlinear realization of $`𝒫_{(2)}`$, with the one-dimensional “Poincaré group” generated by the generator $`P_0`$ as the stability subgroup. Thus we are led to place all generators into the coset (in this particular case it coincides with the full group):
$$G=\text{e}^{itP_0}\text{e}^{iX(t)P_1}\text{e}^{i\mathrm{\Lambda }(t)L}.$$
(5)
Here the wordline evolution parameter $`t`$ (time) is the coset coordinate associated with $`P_0`$, and the Goldstone fields are associated with the rest of generators. The group $`𝒫_{(2)}`$ acts as left shifts of $`G`$. The Cartan forms
$$G^1dG=i\omega _tP_0+i\omega _1P_1+i\omega _LL,$$
(6)
$`\omega _t=\sqrt{1+\mathrm{\Sigma }^2}dt+\mathrm{\Sigma }dX,\omega _1=\sqrt{1+\mathrm{\Sigma }^2}dX+\mathrm{\Sigma }dt,`$
$`\omega _L={\displaystyle \frac{1}{\sqrt{1+\mathrm{\Sigma }^2}}}d\mathrm{\Sigma },\mathrm{\Sigma }\text{sh}\mathrm{\Lambda }`$ (7)
by construction are invariant under this left group action. Next, we observe that the Lorentz Goldstone field $`\mathrm{\Sigma }(t)`$ can be traded for $`\dot{X}(t)`$ by the inverse Higgs constraint
$$\omega _1=0\mathrm{\Sigma }=\frac{\dot{X}}{\sqrt{1\dot{X}^2}}.$$
(8)
This constraint is covariant since $`\omega _1`$ is the group invariant (in the generic case, the coset Cartan forms undergo homogeneous rotation in their stability subgroup indices). Thus the obtained expression for $`\mathrm{\Sigma }`$ possesses correct transformation properties. Substituting it into the remaining Cartan forms we find
$$\omega _0=\sqrt{1\dot{X}^2}dt,\omega _L=\sqrt{1\dot{X}^2}\frac{d}{dt}\left(\frac{\dot{X}}{\sqrt{1\dot{X}^2}}\right)dt.$$
(9)
The simplest invariant action, the covariant length
$$S=\omega _0=𝑑t\sqrt{1\dot{X}^2},$$
(10)
is recognized, up to a renormalization factor of the dimension of mass, as the action of $`D=2`$ massive particle in the static gauge $`X^0(t)=t`$. The equation of motion for $`X(t)`$ can also be given the manifestly covariant form
$$\omega _L=0.$$
(11)
Actually, we could start from the action (10) with the original expression (7) for $`\omega _0`$ and reproduce (8) as the algebraic equation of motion for $`\mathrm{\Sigma }(t)`$. It is important to realize that the correct form of the action is recovered just because we have included the Lorentz Goldstone field into the coset.
In the generic case of $`p`$-brane in $`D`$-dimensional Minkowski space the construction is analogous. One equates to zero the Cartan forms $`\omega ^\mu `$ associated with the transverse translation generators, which is once again the manifestly covariant constraint, and in this way eliminates the Lorentz Goldstone fields $`\mathrm{\Lambda }^{\mu m}(x)`$:
$$\mathrm{\Lambda }_m^\mu _mX^\mu .$$
(12)
Substituting these expressions into the Cartan forms $`\omega ^m`$ which are covariant differentials of $`x^m`$, one finds that the invariant volume of $`x`$-space, i.e. the integral of the external product of these 1-forms, is just the $`p`$-brane static gauge NG action (3).
2. Superbranes from PBGS. The PBGS approach is the generalization of the above nonlinear realization view of branes to the superbranes. Though originally the phenomenon of partial breaking of global SUSY ($`N=2`$ down to $`N=1`$ in $`D=4`$) was studied without any reference to superbranes, the subsequent study - revealed the profound relationship between both concepts. It was shown in that $`N=1,D=4`$ superstring in the static gauge can be understood as the theory of partial breaking of $`N=1,D=4`$ SUSY to its $`N=(2,0),d=2`$ subgroup. The very existence of self-consistent GS type actions for superbranes, with the appropriate fermionic $`\kappa `$-symmetry, was inferred in the pioneering paper from the study of the partial breaking $`N=(1,0),D=6N=1,d=4`$. The same PBGS pattern was treated in from the $`D=4`$ perspective as the partial breaking of $`N=2`$ SUSY with two central charges down to $`N=1`$ one, with the systematic use of the nonlinear realizations and inverse Higgs effect techniques. It was shown that the self-consistent theory can be constructed in terms of the (covariantly) chiral and antichiral bosonic $`N=1`$ superfields which are the Goldstone ones corresponding to the central charges (that is, to the translation operators in fifth and sixth directions from the $`D=6`$ viewpoint). The fermionic Goldstone superfields associated with the spontaneously broken supertranslation generators are covariantly expressed in terms of these basic Goldstone superfields by the covariant constraints of the type (8). The relevant Goldstone superfields action, a nonlinear generalization of the standard free action of $`N=1`$ chiral superfields, was constructed in . In components, after eliminating the auxiliary fields, it yields just the static gauge form of the GS action for the $`N=(1,0),D=6`$ 3-brane in a flat background. An interesting new phenomenon was discovered in . It turned out that the same SUSY admits several different PBGS options depending on into which multiplet of unbroken SUSY one embeds the Goldstone fermion. In the $`N=2N=1,D=4`$ case, instead of placing this field into the chiral $`N=1`$ multiplet, one can place it into the abelian vector $`N=1`$ multiplet as a “photino”, by imposing the appropriate covariant constraints on the Goldstone fermionic superfield which generalize the Bianchi identities for the flat $`N=1`$ Maxwell superfield strength. The relevant invariant action, on the one hand, is $`N=2`$ extension of the Born-Infeld (BI) action with the hidden, nonlinearly realized half of supersymmetry, and, on the other, is (in components) a gauge-fixed form of the GS action of the “space-time filling” D3-brane. One more option is to embed the Goldstone fermion into the $`N=1`$ tensor (or linear) multiplet . The emerging brane is the super “L3-brane” in $`D=5`$ (in terminology of ), it is related to the $`N=(1,0),D=6`$ 3-brane via the familiar duality between $`N=1`$ tensor and chiral superfields.
The study of the PBGS patterns corresponding to partial breaking of SUSY with 16 supercharges ($`N=1,D=10`$; $`N=2,D=6`$; $`N=4,D=4`$; …) was initiated in our works . Let us dwell on this and related subjects in some detail.
3. Hypermultiplet as a Goldstone superfield. To describe the $`1/2`$ breaking of Poincaré SUSY with 16 supercharges, it is natural to start from the maximally symmetric situation where the broken and unbroken SUSY “live” as simple ones. This amounts to considering the PBGS option $`N=1,D=10N=(1,0),d=6`$ (we could equally choose $`N=(0,1)`$). From the $`d=6`$ perspective, $`N=1,D=10`$ is a central extension of $`N=(1,1)`$:
$$N=1,D=10\{Q_\alpha ^i,P_{\alpha \beta },S^{\beta a},Z^{ia}\},$$
(13)
where $`\alpha ,\beta =1,\mathrm{},4;i=1,2;a=1,2`$ are, respectively, the $`Spin(1,5)`$ indices and the $`SU(2)`$ doublet indices. The basic anticommutation relations read
$$\{Q_\alpha ^i,Q_\beta ^j\}=ϵ^{ij}P_{\alpha \beta },\{Q_\alpha ^i,S^{a\beta }\}=\delta _\alpha ^\beta Z^{ia},\{S^{a\alpha },S^{b\beta }\}=ϵ^{ab}P^{\alpha \beta }.$$
(14)
One also must add generators of the $`D=10`$ Lorentz group $`SO(1,9)\{M_{\alpha \beta \gamma \delta },T^{ij},T^{ab},K_{ia}^{\alpha \beta }\}`$, with $`M`$ and $`T`$ generating the $`d=6`$ Lorentz group $`SO(1,5)`$ and $`R`$-symmetry $`SU(2)\times SU(2)`$.
We wish $`N=(1,0),d=6`$ SUSY $`\{Q_\alpha ^i,P_{\alpha \beta },T^{ij},T^{ab},M_{\alpha \beta \gamma \delta }\}`$ to be unbroken. Like in the bosonic case of Sect.1, we are led to place the generators $`Q_\alpha ^i,P_{\alpha \beta },S^{\alpha a},Z^{ia},K_{\alpha \beta }^{ia}`$ into the coset and to treat the coset parameters associated with two first generators as the coordinates of $`N=(1,0),d=6`$ superspace, while the remaining ones as the Goldstone superfields:
$`Q_\alpha ^i\theta _i^\alpha ,P_{\alpha \beta }x^{\alpha \beta }.`$
$`S^{\alpha a}\mathrm{\Psi }_{\alpha a}(x,\theta ),Z^{ia}q_{ia}(x,\theta ),K_{\alpha \beta }^{ia}\mathrm{\Lambda }_{ia}^{\alpha \beta }(x,\theta ).`$ (15)
The coset element $`G`$ in the exponential parametrization is constructed like in eq. (5). The Cartan forms are defined by
$$G^1dG=\mathrm{\Omega }_Q+\mathrm{\Omega }_P+\mathrm{\Omega }_Z+\mathrm{\Omega }_S+\mathrm{\Omega }_K+\mathrm{},$$
(16)
where the subscripts denote the relevant generators. The forms $`\mathrm{\Omega }_{Q,P}`$ are the covariant differentials of the superspace coordinates, $`\mathrm{\Omega }_{Z,S,K}`$ are those of Goldstone superfields, they are homogeneously transformed under the supergroup left shifts. We shall be interested in the linearized structure of $`\mathrm{\Omega }_Z=\mathrm{\Omega }_Z^{ia}Z_{ia}`$ :
$$\mathrm{\Omega }_Z^{ia}=dq^{ia}+2\mathrm{\Lambda }_{\alpha \beta }^{ia}dx^{\alpha \beta }+\psi _\alpha ^ad\theta ^{i\alpha }+\mathrm{}.$$
(17)
Comparing it with the form $`\omega _1`$, eq. (7) of the toy example of Sect. 1, we observe that both forms have a similar structure. Hence, the superfields $`\mathrm{\Lambda },\mathrm{\Psi }`$ can be expressed through the basic Goldstone superfield $`q^{ia}`$ from the inverse Higgs constraint analogous to (8)
$`\mathrm{\Omega }_Z=0\mathrm{\Lambda }_{\alpha \beta }^{ia}`$ $`=`$ $`_{\alpha \beta }q^{ia}=_{\alpha \beta }q^{ia}+\mathrm{},`$ (18)
$`\mathrm{\Psi }_\alpha ^a`$ $`=`$ $`{\displaystyle \frac{1}{2}}_\alpha ^kq_k^a={\displaystyle \frac{1}{2}}D_\alpha ^kq_k^a+\mathrm{}.`$
Here $`D_\alpha ^i=/\theta _i^\alpha \frac{1}{2}\theta ^{i\beta }_{\beta \alpha }`$ and dots stand for nonlinear terms. Thus, $`q^{ia}`$ is the essential Goldstone superfield, analogue of $`X(t)`$ in the $`D=2`$ example. Its first bosonic component $`q^{ia}(x)`$ parametrizes the transverse directions in the $`D=10`$ Minkowski space, while the physical fermionic component is the Goldstino related to the spontaneously broken supertranslations (with the generator $`S^{a\alpha }`$). This field content is that of the scalar $`N=1,D=10`$ 5-brane.
The constraint (18) differs in a few important aspects from (8) and the similar inverse Higgs constraint considered in ref. . It not only eliminates the redundant Goldstone superfields, but also implies the differential constraint for $`q^{ia}(x,\theta )`$
$$_\alpha ^{(k}q^{i)a}=D_\alpha ^{(k}q^{i)a}+\mathrm{}=0.$$
(19)
It is just the nonlinear, “brane” generalization of the standard hypermultiplet constraint . The latter reduces the field content of $`q^{ia}`$ to the set of (8+8) components and puts them on shell:
$`q^{ia}(x,\theta )\varphi ^{ia}(x)+\theta ^{\alpha i}\psi _\alpha ^a(x)+x\text{-derivatives},`$ (20)
$`\mathrm{}\varphi ^{ia}(x)=0,^{\alpha \beta }\psi _\beta ^a=0.`$ (21)
Thus eq. (19) describes the on-shell dynamics of $`N=1,D=10`$ 5-brane as the natural generalization of the free hypermultiplet dynamics. Since the off-shell superfield action for the latter can be constructed only in harmonic superspace , it is reasonable to expect that there exists a brane generalization of the hypermultiplet harmonic superspace action. It would be $`N=(1,0),d=6`$ (or $`N=2,d=4`$ after dimensional reduction) analogue of the $`N=1`$ Goldstone superfield action of refs. . No systematic recipes are known so far for constructing PBGS actions. At present we are aware only of the appropriate fourth-order correction to the free hypermultiplet action. Even for the dimensionally-reduced cases, including the simplest case of $`N=2,D=5`$ superparticle, the full action is rather difficult to find. For the superparticle case it is known up to the sixth order in fields :
$$S_{br}^q=𝑑\zeta ^{(4)}q_a^+D^{++}q^{+a}+\alpha 𝑑Z(A^2+2\alpha AB^{++}B^{}+\mathrm{})$$
(22)
where
$$A=q_a^+D^{}q^{+a},B^{++}=q_a^+_tq^{+a},$$
$`q^{+a},D^{\pm \pm }`$ are the $`d=1`$ reduction of the analytic hypermultiplet superfield and harmonic derivatives , $`B^{}`$ is defined by
$$D^{++}B^{}D^{}B^{++}=0$$
and $`\alpha `$ is a coupling constant. It would be tempting to find out the geometric principle allowing to restore the whole action. Let us now apply to a simpler PBGS case where the complete off-shell action can be constructed .
4. N=1, D=4 supermembrane. This case corresponds to the partial breaking of $`N=1,D=4`$ SUSY to $`N=1,d=3`$ one. The $`N=1,D=4`$ superalgebra in the $`d=3`$ notation reads
$$\{Q_a,Q_b\}=P_{ab},\{Q_a,S_b\}=ϵ_{ab}Z,\{S_a,S_b\}=P_{ab},a,b=1,2.$$
(23)
The $`d=3`$ momentum operator $`P_{ab}`$ together with the central charge $`Z`$ form the $`D=4`$ translation operator. The vacuum stability subalgebra consists of the $`N=1,d=3`$ superalgebra (generators $`Q_a`$ and $`P_{ab}`$) and that of the $`d=3`$ Lorentz group $`SO(1,2)`$. The second SUSY, the central charge $`Z`$ and the generators $`K_{ab}`$ of the coset $`SO(1,3)/SO(1,2)`$ define spontaneously broken symmetries. The appropriate coset parameters are introduced as
$$Q_a\theta ^a,P_{ab}x^{ab},S_a\psi ^a(x,\theta ),Zq(x,\theta ),K_{ab}\lambda ^{ab}(x,\theta ).$$
(24)
After construction of the relevant Cartan forms and covariant elimination of the Goldstone superfields $`\psi ^a(x,\theta ),\lambda ^{ab}(x,\theta )`$ by the appropriate inverse Higgs constraints, one is left with $`q(x,\theta )`$ as the only unremovable Goldstone superfield. Its physical fields (one boson and two fermions) parametrize one transverse bosonic and two fermionic directions in $`N=1,D=4`$ superspace, the auxiliary bosonic field also admits a nice geometric interpretation of the Goldstone field for the spontaneously broken $`U(1)`$ automorphism group of $`N=1,D=4`$ superalgebra (“$`\gamma ^5`$ invariance”). The physical field content of $`q`$ coincides with that of $`N=1,D=4`$ supermembrane.
The inverse Higgs conditions in this case do not imply the equation of motion for $`q`$. However, the dynamical equation for $`q`$ turns out to admit, like in the $`D=2`$ particle example, a covariant representation as the vanishing of the covariant $`d\theta `$ projection of the coset Cartan form $`\omega _S^a`$ associated with the spontaneously broken SUSY generator $`S_a`$ . We have no place to present details. Let us explain how to construct the off-shell action for this case.
Like in other PBGS cases, the nonlinear realizations approach on its own provides no clear recipe how to construct such an action. This becomes possible using the trick similar to the one exploited in . Namely, let us start from a linear realization of $`N=1,D=4`$ SUSY in terms of $`N=1,d=3`$ superfields $`\mathrm{\Phi },\xi _aD_a\rho `$ with the following transformation rules under the second SUSY:
$$\delta \rho =\theta ^a\eta _a2D^a\mathrm{\Phi }\eta _a,\delta \mathrm{\Phi }=\frac{1}{2}\eta ^aD_a\rho ,$$
(25)
where $`\eta _a`$ is the transformation parameter and $`D_a`$ is the flat $`N=1,d=3`$ spinor covariant derivative, $`\{D_a,D_b\}=_{ab}`$. It is easy to check that the closure of these transformations and those of manifest $`N=1,d=3`$ SUSY is just the superalgebra (23), with $`Z`$ realized as a pure shift of $`\rho `$. The transformation of $`\xi _a=D_a\rho `$ starts with $`\eta _a`$, suggesting the interpretation of this superfield as the linear realization Goldstone fermion. After some work one finds that $`\mathrm{\Phi }`$ can be covariantly expressed in terms of $`\xi _a`$ as follows
$$\mathrm{\Phi }=\frac{1}{2}\frac{\xi ^2}{1+\sqrt{1+D^2\xi ^2}}.$$
Recalling the transformation law of $`\mathrm{\Phi }`$, one finds that the integral
$$S=d^3xd^2\theta \mathrm{\Phi }\frac{1}{2}d^3xd^2\theta \frac{\xi ^2}{1+\sqrt{1+D^2\xi ^2}},\xi ^a=D^a\rho ,$$
(26)
is invariant under the hidden SUSY transformations as well as $`D=4`$ Poincaré translations and so it can be identified with the sought $`d=3`$ worldvolume superspace action of $`N=1,D=4`$ supermembrane. Indeed, it is easy to find that the bosonic core of this action is just the static gauge membrane NG action:
$$S=d^3x\left(1\sqrt{1\frac{1}{2}qq}\right).$$
(27)
One can find the equivalence field redfinition relating $`\xi _a`$ to the nonlinear realization Goldstone fermion $`\psi _a`$ and $`\rho `$ to $`q`$. Also, the equations of motion following from the action (26) can be shown to be equivalent to those conjectured at the level of the Cartan forms. This implies the presence of hidden $`SO(1,4)`$ Lorentz symmetry in the action. It still remains to prove the precise equivalence of this action to the GS one (e.g., along the lines of ref. ) and the action proposed in the superembedding approach . Note an interesting peculiarity: one can add to the lagrangian in (26) the “cosmological” term $`\rho `$. This term is invariant under the second SUSY (up to surface terms) and the shifts $`\rho \rho +const`$. Adding it changes the equation of motion for the auxiliary field and so can influence the structure of the component action (without this term, the auxiliary field is vanishing on shell). The pure physical boson part of the action is always given by (27).
Besides a scalar multiplet $`\rho `$ (or $`q`$), we can choose a vector $`N=1,d=3`$ multiplet as the Goldstone one (like in ). It is represented by $`N=1`$ spinor superfield strength $`\mu _a`$ obeying the constraint:
$$D^a\mu _a=0.$$
(28)
It leaves in $`\mu _a`$ the first fermionic component together with the divergenceless vector $`F_{ab}D_a\mu _b|_{\theta =0}`$ (just the gauge field strength). Due to the vector-scalar $`d=3`$ duality, the superfield $`\mu _a`$ is expected to describe a D2-brane which is dual to the supermembrane.
The relevant action can be found using the previous trick. One can extend $`\mu ^a`$ to the $`N=1,D=4`$ multiplet $`(\mu ^a,\varphi )`$ with the following transformation rules under the second SUSY
$$\delta \mu _a=\eta _aD^2\varphi \eta _a+_{ab}\varphi \eta ^b,\delta \varphi =\frac{1}{2}\eta ^a\mu _a.$$
(29)
Thus $`\mu ^a`$ can also be interpreted as the linear realization Goldstone fermionic superfield. The following expression for $`\varphi `$
$$\varphi =\frac{1}{2}\frac{\mu ^2}{1+\sqrt{1D^2\mu ^2}}$$
(30)
can be checked to be consistent with (29). Then the action
$$S=d^3xd^2\theta \varphi =\frac{1}{2}d^3xd^2\theta \frac{\mu ^2}{1+\sqrt{1D^2\mu ^2}}$$
(31)
is invariant under the second SUSY in virtue of the transformation rule of $`\varphi `$ (29) and the constraint (28). This nonlinear generalization of the standard $`N=1,d=3`$ abelian vector multiplet action $`\mu ^2`$ involves as its bosonic core the $`d=3`$ BI action
$`S={\displaystyle d^3x\left(\sqrt{1+2F^2}1\right)},`$ (32)
$`^{ab}F_{ab}=0F_{ab}=_{ac}G_b^c+_{bc}G_a^c.`$ (33)
So it is $`N=2`$ extension of the $`d=3`$ BI action with nonlinearly realized second SUSY. It can be regarded as the worldvolume superfield action of the space-time filling D2-superbrane in a flat background. It is easy to prove its dual equivalence to the action (26) by inserting the constraint (28) in it with the Lagrange multiplier $`\rho `$ and integrating out $`\mu ^a`$ .
Finally, let us show how the duality between the bosonic NG and BI actions (27), (32) can be recovered in the nonlinear realizations approach of Sect.1. It is easy to find that in the case of nonlinear realization of $`𝒫_{(4)}`$ in the coset $`𝒫_{(4)}/SO(1,2)`$, which corresponds just to membrane in $`D=4`$, the covariant differentials of $`x^m`$ read
$$\omega ^m=dx^m+2\frac{\lambda ^m}{12\lambda ^2}\left(2\lambda _n+_nq\right)dx^n\omega _n^mdx^n,$$
(34)
where $`\lambda ^m(x)`$ is the Lorentz $`SO(1,3)/SO(1,2)`$ Goldstone field in the appropriate parametrization and $`q(x)`$ is the transverse coordinate of membrane. We could eliminate $`\lambda ^m`$ by the inverse Higgs constraint but in the present case it is instructive to reproduce it as the equation of motion for $`\lambda ^m`$. The minimal invariant action constructed as the covariant worldvolume
$$S_{mem}=d^3x\text{det}\omega _n^m,$$
(35)
up to a normalization factor and constant shift, is
$$S_{mem}=d^3x\frac{1}{12\lambda ^2}\left(2\lambda ^2+\lambda q\right).$$
(36)
Varying $`\lambda ^m`$ yields the inverse Higgs expression for it
$$\lambda _m=\frac{1}{2}\frac{_mq}{1+\sqrt{1\frac{1}{2}(q)^2}}.$$
(37)
After this (36) takes the standard NG form (27). On the other hand, one can treat $`q`$ in (36) as the Lagrange multiplier for the differential constraint on $`\lambda ^m`$:
$$_mF^m=0,F^m2\frac{\lambda ^m}{12\lambda ^2}.$$
(38)
Expressing $`\lambda ^2`$ through $`F^m`$
$$\lambda ^2=\frac{1}{4F^2}\left(1\sqrt{1+2F^2}\right)^2,$$
one reduces (36) just to the BI form (32)
$$S_{mem}d^3x\left(\sqrt{1+2F^2}1\right).$$
(39)
This simple consideration shows that within the nonlinear realization approach the $`d=3`$ Maxwell field strength entering the $`d=3`$ BI action acquires the nice geometric interpretation as the Goldstone field representing the Lorentz coset $`SO(1,3)/SO(1,2)`$, while the BI action itself can be algorithmically derived as the action dual to the static gauge membrane NG action.
5. Concluding remarks. The geometric PBGS approach can be thought of as a useful and viable alternative to the standard GS description of superbranes. Its main merit is that it gives manifestly worldsurface supersymmetric off-shell superfield actions. Leaving aside such important conceptual questions as whether it can be helpful for quantization, etc, we list here a few more modest problems solving which could extend its range of applicability.
First of all, it is desirable to develop convenient general recipes of constructing superfield PBGS actions similar to those provided by the nonlinear realizations for the internal symmetries sigma models. At present, constructing such actions is an art to some extent.
It is important to learn how to construct self-consistent PBGS actions on non-trivial bacgrounds involving the worldvolume and target supergravity and super Yang-Mills fields.
At last, it seems interesting to set up PBGS actions for the systems with $`1/4`$ and other exotic partial breaking options. Some progress in this direction is reported in the contribution by Delduc, Krivonos and myself in this Volume.
Acknowledgements. I thank Jerzy Lukierski for inviting me to give this contribution. This work was supported in part by grants RFBR-CNRS 98-02-22034, RFBR 99-02-18417, INTAS-96-0538, INTAS-96-0308 and Nato Grant No.PST.CLG 974874. |
warning/0002/cond-mat0002461.html | ar5iv | text | # Coulomb Blockade Oscillations in the Thermopower of Open Quantum Dots.
## Abstract
We consider Coulomb blockade oscillations of thermoelectric coefficients of a single electron transistor based on a quantum dot strongly coupled to one of the leads. Analytic expression for the thermopower as a function of temperature $`T`$ and the reflection amplitude $`r`$ in the quantum point contact is obtained. Two regimes can be identified: $`TE_C|r|^2`$ and $`TE_C|r|^2`$, where $`E_C`$ is the charging energy of the dot. The former regime is characterized by weak logarithmic dependence of the thermopower on the reflection coefficient, in the latter the thermopower is linear in the reflection coefficient $`|r|^2`$ but depends on temperature only logarithmically.
Thermoelectric effects in mesoscopic devices have been the subject of extensive experimental and theoretical research . The particle-hole asymmetry required for such effects can be strongly enhanced in these systems as compared to the bulk materials. This and the small size of such devices make them promising candidates for technological applications, such as micro-refrigerators . In particular, many experimental and theoretical studies in the last few years have focused on the thermoelectric properties of quantum dots in the Coulomb blockade (CB) regime . Most of them concentrated on the CB oscillations of the thermopower, $`S=\frac{\mathrm{\Delta }V}{\mathrm{\Delta }T}`$, where $`\mathrm{\Delta }T`$ is the temperature difference across the dot, and $`\mathrm{\Delta }V`$ is the voltage necessary to nullify the current.
The theory of the CB oscillations in the thermopower of quantum dots in the weak tunneling regime was constructed in Ref. . This theory takes into account only the lowest order tunneling processes, i.e. the sequential tunneling, and neglects the cotunneling processes. Its results were in agreement with the experiments of Ref. . Later it became possible to experimentally access the regime of lower temperatures and stronger tunneling where the cotunneling processes become dominant. The theoretical description of this regime was recently given in Ref. .
In very interesting recent experiment the CB oscillations in the thermopower of a nearly open quantum dot were studied as a function of the reflection coefficient $`|r|^2`$ in the contact. The setup of these experiments is schematically represented in Fig. 1. Surprisingly, an initial decrease in the amplitude of CB oscillations of thermopower with decreasing $`|r|^2`$ was followed by a plateau with nearly $`|r|`$-independent CB oscillations of thermopower. This saturation was attributed to the effects of elastic cotunneling .
The theory of thermopower for the weak tunneling regime developed in Refs. does not apply to this case. An additional motivation for studying the thermoelectric phenomena in such devices arises from the fact that due to the increased transparency of the contacts, the open dots are better candidates for micro-refrigerator devices than the closed ones.
In this Letter we present a theory of thermoelectric effects in a quantum dot in the nearly open regime. We consider a quantum dot which is coupled by a tunneling junction to the left lead and by a single channel quantum point contact (QPC) to the right lead, see Fig. 1. The reflection amplitude in the QPC is assumed to be small, $`|r|1`$. The mean level spacing $`\delta `$ in the dot is assumed to be vanishingly small. This is a good assumption since experimentally $`\delta T`$.
The previous studies of such systems were devoted to their thermodynamic and transport properties . A special feature of the thermoelectric power $`S`$ is that it is sensitive to the average energy transported by electrons, which in the tunneling approximation depends on the odd part of the density of states (DoS) as a function of energy. Thus the thermoelectric phenomena represent an independent probe of these systems.
In this Letter we find the thermoelectric coefficient $`G_T`$ of the device in Fig. 1 describing the current response $`I`$ at zero bias, $`\mathrm{\Delta }V=0`$, to the difference of the temperatures $`\mathrm{\Delta }T`$ between the two leads: $`G_T=lim\frac{I}{\mathrm{\Delta }T}|_{\mathrm{\Delta }V=0,\mathrm{\Delta }T0}`$. Our main result is the following expression for $`G_T`$ of the dot:
$`G_T`$ $`=`$ $`{\displaystyle \frac{G_L|r|^2T}{6\pi eE_C}}\mathrm{ln}{\displaystyle \frac{E_C}{T+\mathrm{\Gamma }}}\mathrm{sin}(2\pi N)`$ (2)
$`\times {\displaystyle _{\mathrm{}}^+\mathrm{}}{\displaystyle \frac{x^2(x^2+\pi ^2)}{\left[x^2+(\mathrm{\Gamma }/T)^2\right]\mathrm{cosh}^2(x/2)}}dx.`$
Here $`G_Le^2/h`$ is the conductance of the left contact, $`e`$ is the absolute value of the electron charge, $`E_C`$ is the charging energy. We have also introduced the energy scale $`\mathrm{\Gamma }=(8\gamma /\pi ^2)E_C|r|^2\mathrm{cos}^2(\pi N)`$, which depends on the gate voltage $`N`$; here $`\mathrm{ln}\gamma =𝐂0.5772\mathrm{}`$ is the Euler constant. The result (2) was obtained with logarithmic accuracy assuming that $`E_CT,\mathrm{\Gamma }`$.
The thermopower $`S=G_T/G`$ is then obtained from Eq. (2) using the result of Ref. for the conductance $`G`$ of the device shown in Fig. 1, which we reproduce here for completeness:
$$G=\frac{G_L\mathrm{\Gamma }}{8\gamma E_C}_{\mathrm{}}^+\mathrm{}\frac{x^2+\pi ^2}{\left[x^2+(\mathrm{\Gamma }/T)^2\right]\mathrm{cosh}^2(x/2)}𝑑x.$$
(3)
For the two limiting cases, $`T\mathrm{\Gamma }`$ and $`T\mathrm{\Gamma }`$, we obtain simplified expressions for the thermopower:
$`S`$ $`=`$ $`\{\begin{array}{cc}\frac{64\gamma |r|^2}{9\pi ^2e}\mathrm{ln}\frac{E_C}{T}\mathrm{sin}(2\pi N),& \mathrm{for}T\mathrm{\Gamma },\\ \frac{\pi ^3T}{5eE_C}\mathrm{ln}\frac{E_C}{\mathrm{\Gamma }}\mathrm{tan}(\pi N),& \mathrm{for}T\mathrm{\Gamma }.\end{array}`$ (6)
It is difficult to make a direct comparison of the results (2), (6) with the experiments of Ref. since the experimental data were presented in terms of a fit to the weak-tunneling theory of Ref. . Nevertheless we want to point out that even without taking into account the elastic cotunneling effects , in the regime $`T\mathrm{\Gamma }`$ the thermopower only weakly (logarithmically) depends on the reflection coefficient, which is consistent with the observation in Ref. of the thermopower virtually independent of the reflection coefficient. In the opposite regime, $`\mathrm{\Gamma }T`$, the thermopower is nearly independent of the temperature, but scales linearly with the reflection coefficient $`|r|^2`$ vanishing, as expected, at perfect transmission. Note that even at very low temperatures $`TE_C|r|^2`$, one still has $`T\mathrm{\Gamma }`$ near half-integer values of the gate voltage $`N`$ corresponding to the CB peaks of conductance (3). The width $`\delta N`$ of those regions can be easily found from the condition $`\mathrm{\Gamma }(N)T`$. Upon substitution into Eq. (6) it gives the estimate of the amplitude of the CB oscillations of the thermopower $`S_0e^1|r|\sqrt{T/E_C}\mathrm{ln}(E_C/T)`$.
It is interesting to point out that in the low temperature regime $`T\mathrm{\Gamma }`$, when the conductance (3) shows the temperature dependence $`GT^2`$ characteristic of inelastic cotunneling, the thermopower can be expressed in terms of the logarithmic derivative of the conductance with respect to the gate voltage $`2E_CN`$:
$$S=\frac{\pi ^2T}{10eE_C}\mathrm{ln}\left(\frac{E_C}{\mathrm{\Gamma }}\right)\frac{\mathrm{ln}G}{N}.$$
(7)
This form is analogous to the Cutler-Mott formula for the thermopower of a system of non-interacting electrons in a metal, but with a different coefficient in front of the logarithmic derivative. A similar Cutler-Mott type relation holds in the case of weak inelastic cotunneling ; however, the prefactor of Eq. (7) contains an additional large logarithmic factor $`\mathrm{ln}(E_C/\mathrm{\Gamma })`$. In the opposite case of high temperature $`T\mathrm{\Gamma }`$ no expression similar to the Cutler-Mott formula applies.
Below we present the derivation of the result (2). Following Ref. , the electron transport through the right QPC can be described by a one-dimensional model amenable to bosonization, whereas the left contact can be treated in the tunneling approximation. The Hamiltonian of the dot has the form $`\widehat{H}=\widehat{H}_0+\widehat{H}_R+\widehat{H}_L+\widehat{H}_C`$, where
$`\widehat{H}_0`$ $`=`$ $`{\displaystyle \underset{k\alpha }{}}ϵ_ka_{k\alpha }^{}a_{k\alpha }+{\displaystyle \underset{p\alpha }{}}ϵ_pa_{p\alpha }^{}a_{p\alpha }`$ (9)
$`+`$ $`{\displaystyle \frac{v_F}{2\pi }}{\displaystyle \underset{\alpha }{}}{\displaystyle \left\{[\varphi _\alpha (x)]^2+\pi ^2\mathrm{\Pi }_\alpha ^2(x)\right\}𝑑x},`$ (10)
$`\widehat{H}_L`$ $`=`$ $`{\displaystyle \underset{kp\alpha }{}}\left(v_ta_{k\alpha }^{}a_{p\alpha }F+v_t^{}a_{p\alpha }^{}a_{k\alpha }F^{}\right),`$ (11)
$`\widehat{H}_R`$ $`=`$ $`{\displaystyle \frac{D}{\pi }}|r|{\displaystyle \underset{\alpha }{}}\mathrm{cos}[2\varphi _\alpha (0)],`$ (12)
$`\widehat{H}_C`$ $`=`$ $`E_C\left[\widehat{n}+{\displaystyle \frac{1}{\pi }}{\displaystyle \underset{\alpha }{}}\varphi _\alpha (0)N\right]^2.`$ (13)
The operators $`\widehat{H}_R`$, $`\widehat{H}_L`$, and $`\widehat{H}_C`$ describe the backscattering in the right QPC, tunneling through the left contact, and the charging energy of the dot, respectively. In the equations above $`\alpha =,`$ is the spin label, $`a_{p\alpha }`$ and $`a_{k\alpha }`$ are electron annihilation operators in the dot and the left lead respectively, $`D`$ is the energy cutoff in the bosonization, and $`\varphi _\alpha `$ is the bosonization displacement operator describing the electron transport through the right QPC with $`\mathrm{\Pi }_\alpha `$ being its conjugate momentum, $`[\varphi _\alpha (x),\mathrm{\Pi }_\alpha ^{}(x^{})]=i\delta (xx^{})\delta _{\alpha ,\alpha ^{}}`$ (we have put $`\mathrm{}=1`$). The modified form of the tunneling Hamiltonian in Eq. (11) reflects the fact that the electron tunneling event changes the electron number $`\widehat{n}`$ in the dot. This is achieved through the introduction of the charge-lowering operator $`F`$ which satisfies the commutation relation $`[F,\widehat{n}]=F`$. The current operator through the left contact can be obtained from the equation of motion for the charge operator $`\widehat{I}=e\dot{\widehat{n}}=ie[\widehat{n},\widehat{H}]`$. Only $`\widehat{H}_L`$ contributes to this commutator and gives rise to the following expression for the current operator
$$\widehat{I}=ie\underset{kp\alpha }{}\left(v_t^{}a_{p\alpha }^{}a_{k\alpha }F^{}v_ta_{k\alpha }^{}a_{p\alpha }F\right).$$
(14)
We treat the problem in the lowest order in the tunneling Hamiltonian Eq. (11). We also assume that the conductance of the tunneling contact is much less than the conductance quantum, $`G_Le^2/h`$. In this approximation all of the temperature drop happens across the left contact. We take the temperature of the left lead to be $`T+\mathrm{\Delta }T`$ and that of the dot and the right reservoir to be $`T`$. In the linear approximation in $`\mathrm{\Delta }T`$ the current $`I`$ can be expressed through the tunneling DoS $`\nu (ϵ)`$ in the dot as,
$$G_T=\frac{I}{\mathrm{\Delta }T}=\frac{G_L}{4T^2e\nu _0}_{\mathrm{}}^{\mathrm{}}\frac{\nu (ϵ)ϵdϵ}{\mathrm{cosh}^2\left(\frac{\beta ϵ}{2}\right)}.$$
(15)
Here $`\nu _0`$ is the DoS in the dot in the absence of interaction, Eq. (13).
Thus, technically the problem in the tunneling approximation reduces to the calculation of the energy-dependent tunneling DoS, $`\nu (ϵ)`$. We note that $`G_T`$ depends only on the odd (as a function of energy) component of DoS, whereas the conductance $`G`$ depends only on the even one. Therefore, as was mentioned earlier, thermopower measurements represent an independent test of the theory of Coulomb blockade in nearly open dots developed in Refs. . Moreover, in the leading order in $`\mathrm{max}\{T,\mathrm{\Gamma }\}/E_C`$ the odd component of the tunneling DoS vanishes . The thermoelectric coefficient $`G_T`$ is small in the the ratio of $`\mathrm{max}\{T,\mathrm{\Gamma }\}/E_C`$ in comparison to the conductance $`G`$. Its calculation requires going beyond the previously adopted approximations and retaining sub-leading order in $`ϵ/E_C`$ in the tunneling DoS, $`\nu (ϵ)`$.
The tunneling DoS in the dot can be expressed as
$$\nu (ϵ)=\frac{1}{\pi }\mathrm{cosh}\frac{\beta ϵ}{2}_{\mathrm{}}^{\mathrm{}}𝒢\left(\frac{\beta }{2}+it\right)\mathrm{exp}(iϵt)𝑑t,$$
(16)
where $`𝒢\left(\frac{\beta }{2}+it\right)`$ is the Matsubara Green function,
$$𝒢(\tau )=\underset{pp^{}}{}T_\tau a_{p\alpha }(\tau )F(\tau )a_{p^{}\alpha }^{}(0)F^{}(0),$$
(17)
analytically continued to complex time $`\tau =\frac{\beta }{2}+it`$. The angular brackets $`\mathrm{}`$ in Eq. (17) denote the thermal average.
Because the dynamics of the operators $`a_{p\alpha }`$ and $`F`$ are decoupled, the Green function in Eq. (17) factorizes into $`𝒢(\tau )=G_0(\tau )K(\tau )`$, with $`G_0(\tau )=\nu _0\pi T/\mathrm{sin}(\pi T\tau )`$ being the free electron Green function and $`K(\tau )=T_\tau F(\tau )F^{}(0)`$, .
Since the operator $`F^{}(0)`$ in $`K(\tau )`$ changes the value of $`\widehat{n}`$ from zero to one at $`t=0`$, and $`F(\tau )`$ changes it back to zero at $`t=\tau `$, the correlator $`K(\tau )`$ can be rewritten as
$$K(\tau )=\frac{Z(\tau )}{Z(0)},$$
(18)
where $`Z(\tau )`$ is a functional integral over $`\varphi _\alpha `$’s in the presence of the time-dependent charge $`n_\tau (t)=\theta (t)\theta (\tau t)`$. Introducing the charge and spin mode variables in the right contact $`\varphi _{c,s}(x)=[\varphi _{}(x)\pm \varphi _{}(x)]/\sqrt{2}`$, we can write $`Z(\tau )`$ as
$$Z(\tau )=D[\varphi _c,\varphi _s]\mathrm{exp}[𝒮_C(\tau )𝒮_{0,c}𝒮_{0,s}𝒮_R].$$
(19)
Here $`𝒮_{0,c}+𝒮_{0,s}`$ represents the free electron part of the action in the absence of backscattering in the QPC, $`S_C`$ denotes its charging part, and $`S_R`$ represents the backscattering in the QPC. These terms are given by
$`𝒮_{0,i}={\displaystyle _0^\beta }𝑑t{\displaystyle 𝑑x\frac{v_F}{2\pi }\left([\varphi _i]^2+\frac{\dot{\varphi }_i^2}{v_F^2}\right)},i=c,s`$ (21)
$`𝒮_C(\tau )={\displaystyle _0^\beta }𝑑tE_C\left[n_\tau (t)+{\displaystyle \frac{\sqrt{2}}{\pi }}\varphi _c(0,t)N\right]^2,`$ (22)
$`𝒮_R={\displaystyle _0^\beta }𝑑t{\displaystyle \frac{2D}{\pi }}|r|\mathrm{cos}[\sqrt{2}\varphi _c(0,t)]\mathrm{cos}[\sqrt{2}\varphi _s(0,t)].`$ (23)
At frequencies below $`E_C`$ the fluctuations of the charge mode, $`\varphi _c(0,t)`$ are suppressed by the charging energy term (22) and can be integrated out. Furthermore, we can evaluate the functional integral over $`\varphi _c`$ by the saddle point approximation ignoring the backscattering term, Eq. (23).
The action $`𝒮^{\mathrm{sp}}(\tau )`$ and the value of the charge mode $`\varphi _c^{\mathrm{sp}}(0,t)`$ at the saddle point are found to be
$`𝒮^{\mathrm{sp}}(\tau )`$ $`=`$ $`𝒮_C^{\mathrm{sp}}(\tau )+𝒮_{0,c}^{\mathrm{sp}}=\mathrm{ln}{\displaystyle \frac{2\gamma E_C\mathrm{sin}(\pi T\tau )}{\pi ^2T}},`$ (25)
$`\sqrt{2}\varphi _c^{\mathrm{sp}}(0,t)`$ $`=`$ $`\pi [Nn_\tau (t)]+(t)+(\tau t),`$ (26)
$`(t)`$ $`=`$ $`{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\mathrm{sin}(2\pi nTt)}{n+\frac{E_C}{\pi ^2T}}}.`$ (27)
In Eq. (25) we have assumed that $`\tau E_C^1`$, which is a good approximation since we only need $`\tau =\beta /2+it`$ in Eq. (16).
Averaging the backscattering term (23) over the fluctuations of $`\varphi _c`$ we obtain
$`\stackrel{~}{𝒮}_{R,\tau }=\sqrt{{\displaystyle \frac{8\gamma E_CD}{\pi ^3}}}|r|{\displaystyle _0^\beta }𝑑t\mathrm{cos}[\sqrt{2}\varphi _c^{\mathrm{sp}}(0,t)]\mathrm{cos}[\sqrt{2}\varphi _s(0,t)].`$
Since the charge modes can only be intergated out at frequencies below the charging energy, one has to assume that the energy cutoff in the above action is $`DE_C`$.
Equation (19) can now be written as
$`Z(\tau )`$ $`=`$ $`𝒩e^{𝒮^{\mathrm{sp}}(\tau )}Z_s(\tau ),`$ (29)
$`Z_s(\tau )`$ $`=`$ $`{\displaystyle D[\varphi _s]\mathrm{exp}\left(𝒮_{0,s}\stackrel{~}{𝒮}_{R,\tau }\right)},`$ (30)
where $`𝒩`$ is the $`\tau `$-independent factor which arises from the integration over the fluctuations about the saddle point and drops out of $`K(\tau )`$ in Eq. (18). The correlator $`K(\tau )`$ in Eq. (18) then factorizes into $`K(\tau )=K_\mathrm{\Theta }(\tau )K_F(\tau )`$, where
$`K_\mathrm{\Theta }(\tau )`$ $`=`$ $`e^{𝒮^{\mathrm{sp}}(\tau )}={\displaystyle \frac{\pi ^2T}{2\gamma E_C\mathrm{sin}(\pi T\tau )}},`$ (31)
and $`K_F(\tau )`$ is the spin part of the correlator which can be expressed as
$`K_F(\tau )`$ $`=`$ $`Z_s(\tau )/Z_s(0).`$ (32)
The effective action in Eq. (30) can be re-fermionized following Refs. . The Hamiltonian in this representation has the form
$`\widehat{H}=iv_F{\displaystyle \psi ^{}(x)\psi (x)𝑑x}+\lambda (t)\eta [\psi (0)\psi ^{}(0)],`$ (34)
$`\lambda (t)={\displaystyle \frac{2}{\pi }}\sqrt{\gamma v_FE_C}|r|\mathrm{cos}[\sqrt{2}\varphi _c^{\mathrm{sp}}(t)],`$ (35)
where $`\eta =(c+c^{})`$ is a Majorana fermion.
In the limit $`T/E_C0`$ the functions $`(t)`$ in Eqs. (26,27) tend to zero, and to the leading order in $`T/E_C`$ can be neglected . Then the time-dependent coefficient $`\lambda (t)`$ in Eq. (1) becomes $`\lambda _0(t)=\frac{2}{\pi }\sqrt{\gamma v_FE_C}|r|(1)^{n_\tau (t)}\mathrm{cos}[\pi N]`$. In this approximation the odd component of the tunneling DoS in the dot vanishes, thus nullifying the thermopower. Therefore we expand $`K_F(\tau )`$ in Eq. (32) to first order in $`\delta \lambda (t)=\lambda (t)\lambda _0(t)`$. In the fermion representation (1) we obtain for the linear in $`\delta \lambda (t)`$ correction to $`K_F(\tau )`$
$`\mathrm{\Delta }K_F(\tau )`$ $`=`$ $`{\displaystyle _0^\beta }(1)^{n_\tau (t)}\delta \lambda (t)\mathrm{\Phi }(\tau ,t)𝑑t,`$ (37)
$`\mathrm{\Phi }(\tau ,t)`$ $`=`$ $`T_t\eta (\tau )\eta (0)\eta (t)[\psi (0,t)\psi ^{}(0,t)],`$ (38)
where $`\mathrm{}`$ denotes the thermal average with the Hamiltonian (1) with $`\lambda =\frac{2}{\pi }\sqrt{\gamma v_FE_C}|r|\mathrm{cos}[\pi N]`$ independent of time $`t`$.
The average in Eq. (38) can be evaluated with the aid of Wick theorem. It is not difficult to show that the thermopower is an odd function of $`N`$. We therefore need only to retain the odd in $`N`$ component $`\mathrm{\Delta }_{\mathrm{odd}}K_F(\tau )`$ of Eq. (37). Evaluating the integral in Eq. (37) with logarithmic accuracy in $`E_C/\mathrm{max}\{T,\mathrm{\Gamma }\}`$ we find:
$`\mathrm{\Delta }_{\mathrm{odd}}K_F(\tau )`$ $`=`$ $`{\displaystyle \frac{8}{E_C}}\sqrt{{\displaystyle \frac{\gamma \mathrm{\Gamma }E_C}{\pi }}}|r|\mathrm{sin}\left(\pi N\right)\mathrm{ln}{\displaystyle \frac{E_C}{T+\mathrm{\Gamma }}}`$ (40)
$`\times {\displaystyle }{\displaystyle \frac{\xi d\xi }{\xi ^2+\mathrm{\Gamma }^2}}{\displaystyle \frac{e^{\xi |\tau |}}{e^{\beta \xi }+1}}.`$
The upper energy scale $`E_C`$ in the logarithmic factor originates from the above mentioned energy cutoff $`DE_C`$ of the spin excitations. Using Eq. (40) we obtain our main result, Eq. (2).
In conclusion, we have presented a theory of the Coulomb blockade oscillations of the thermoelectric coefficient $`G_T`$ and the thermopower $`S`$ of quantum dots in the anisotropic nearly open regime in the limit where the single particle mean level spacing is negligible. Two distinct regimes can be identified: the one with $`\mathrm{\Gamma }T`$, and the one with $`\mathrm{\Gamma }T`$. In the former the thermopower is linear in temperature but is nearly independent of the reflection coefficient in the QPC and can be expressed in the form of Eq. (7) analogous to the Cutler-Mott formula . In the latter, the thermopower is linear in the reflection coefficient $`|r|^2`$ but depends on the temperature only logarithmically.
We are grateful to B.L. Altshuler, C.M. Marcus, J.M. Martinis, L.W. Molenkamp and B.Z. Spivak for valuable discussions. It is our pleasure to acknowledge the warm hospitality of the Aspen Center for Physics, the ICTP, Trieste, and the Centre for Advanced Studies in Oslo where part of this work was performed. The authors are A.P. Sloan Research Fellows. A.A. is a Packard Research Fellow. This research was supported by the NSF Grants No. DMR-9984002 and DMR-9974435. |
warning/0002/cs0002001.html | ar5iv | text | # Computing large and small stable models11footnote 1This is a full version of an extended abstract presented at the International Conference on Logic Programming, ICLP-99 and included in the proceedings published by MIT Press.
## 1 Introduction
The stable model semantics by Gelfond and Lifschitz \[Gelfond & Lifschitz, 1988\] is one of the two most widely studied semantics for normal logic programs, the other one being the well-founded semantics by Van Gelder, Ross and Schlipf \[Van Gelder et al., 1991\]. Among 2-valued semantics, the stable model semantics is commonly regarded as the one providing the correct meaning to the negation operator in logic programming. It coincides with the least model semantics on the class of Horn programs, and with the well-founded semantics and the perfect model semantics on the class of stratified programs \[Apt et al., 1988\]. In addition, the stable model semantics is closely related to the notion of a default extension by Reiter \[Marek & Truszczyński, 1989, Bidoit & Froidevaux, 1991\]. Logic programming with stable model semantics has applications in knowledge representation, planning and reasoning about action. It was also recently proposed as a computational paradigm well suited for solving combinatorial optimization and constraint satisfaction problems \[Marek & Truszczyński, 1999, Niemelä, 1999\].
Before we proceed, we will recall the definition of a stable model of a logic program, and some related terminology and properties. The reader is referred to \[Marek & Truszczyński, 1993\] for a more detailed treatment of the subject. In the paper we deal only with the propositional case. A logic program rule is an expression $`r`$ of the form
$$r=ab_1,\mathrm{},b_s,\text{not}(c_1),\mathrm{},\text{not}(c_t),$$
where $`a`$, $`b_i`$s and $`c_i`$s are propositional atoms. The atom $`a`$ is called the head of $`r`$ and is denoted by $`h(r)`$. Atoms $`b_i`$ and $`c_i`$ form the body of $`r`$. The set $`\{b_1,\mathrm{},b_s\}`$ is called the positive body of $`r`$ (denoted by $`b^+(r)`$) and the set $`\{c_1,\mathrm{},c_t\}`$ is called the negative body of $`r`$ (denoted by $`b^{}(r)`$). A logic program is a collection of rules. For a logic program $`P`$, by $`\text{At}(P)`$ we denote the set of atoms occurring in its rules and by $`h(P)`$ — the set of atoms appearing as the heads of rules in $`P`$. We will also denote the size of $`P`$, that is, the total number of occurrences of atoms in $`P`$, by $`size(P)`$. Throughout the paper we use $`n`$ to denote the number of atoms in a logic program $`P`$, and $`m`$ to denote the size of $`P`$.
A set of atoms $`M\text{At}(P)`$ satisfies a rule $`r`$ if $`h(r)M`$, or if $`b^+(r)M\mathrm{}`$, or if $`b^{}(r)M\mathrm{}`$. A set of atoms $`M\text{At}(P)`$ is a model of a program $`P`$ if $`M`$ satisfies all rules of $`P`$.
A logic program rule $`r`$ is called Horn if $`b^{}(r)=\mathrm{}`$. A Horn program is a program whose every rule is a Horn rule. The intersection of two models of a Horn program $`P`$ is a model of $`P`$. Since the set of all atoms is a model of $`P`$, it follows that every Horn program $`P`$ has a unique least model. We will denote this model by $`LM(P)`$. The least model of a Horn program $`P`$ can be constructed by means of the van Emden-Kowalski operator $`T_P`$ \[van Emden & Kowalski, 1976\]. Given a Horn program $`P`$ and a set of atoms $`MP`$, we define
$$T_P(M)=\{a:ab_1,\mathrm{},b_sP,\text{and}\{b_1,\mathrm{},b_s\}M\}.$$
We also define
$$T_P^0(M)=\mathrm{},\text{and}T_P^{i+1}(M)=T_P(T_P^i(M)).$$
Since the operator $`T_P`$ is monotone, the sequence $`T_P^i(\mathrm{})`$ is monotone and its union yields the least model of a Horn program $`P`$. That is,
$$LM(P)=\underset{i=0}{\overset{\mathrm{}}{}}T_P^i(\mathrm{}).$$
If $`P`$ is finite, the sequence stabilizes after finitely many steps.
For a logic program rule $`r`$, by $`horn(r)`$ we denote the rule obtained from $`r`$ by eliminating all negated atoms from the body of $`r`$. If $`P`$ is a logic program, we define $`horn(P)=\{horn(r):rP\}`$.
Let $`P`$ be a logic program (possibly with rules containing negated atoms). For a set of atoms $`M\text{At}(P)`$ we define the reduct of $`P`$ with respect to $`M`$ to be the program obtained by eliminating from $`P`$ each rule $`r`$ such that $`b^{}(r)M\mathrm{}`$ (we call such rules blocked by $`M`$), and by removing negated atoms from all other rules in $`P`$. The resulting program is a Horn program. We will denote it by $`P^M`$. As a Horn program, $`P^M`$ has the least model $`LM(P^M)`$. If $`M=LM(P^M)`$, $`M`$ is a stable model of $`P`$. Clearly, if $`M`$ is a stable model of $`P`$, $`Mh(P)`$. Both the notion of the reduct and of a stable model are due to Gelfond and Lifschitz \[Gelfond & Lifschitz, 1988\].
In the paper we restrict our attention to programs whose rules do not contain multiple positive occurrences of the same atom nor multiple negative occurrences of the same atom in the body. It is clear that adopting this assumption does not limit the generality of our considerations. Repetitive occurrences can be eliminated in linear time (in the size of the program) and doing so does not affect stable models of the program.
If $`M`$ is a stable model of $`P`$, each rule $`r`$ such that $`b^+(r)M`$ and $`b^{}(r)M=\mathrm{}`$ (that is, such that $`M`$ satisfies its body), is called a generating rule for $`M`$. Clearly, if $`M`$ is a stable model of $`P`$, it is also a stable model of the program consisting of all rules in $`P`$ that are generating for $`M`$.
There are several ways to look at the search space of possible stable models of a program $`P`$. The most direct way is to look for stable models by considering all candidate subsets of $`h(P)`$. For each candidate subset $`Mh(P)`$, one can compute the corresponding reduct $`P^M`$, its least model $`LM(P^M)`$, and check the equality $`M=LM(P^M)`$ to decide whether $`M`$ is stable. An alternative way is to observe that stable models are determined by subsets of the set of atoms appearing negated in $`P`$. Indeed, let us denote this set by $`Neg(P)`$ and let us consider sets $`M\text{At}(P)`$ and $`BNeg(P)`$. Let $`B^{}=Neg(P)B`$. Then, $`M`$ is a stable model of $`P`$ if and only if $`M=LM(P^B^{})`$, $`BM=\mathrm{}`$ and $`B^{}M`$. Thus, the existence of stable models can be decided by considering subsets of $`Neg(P)`$. Finally, one can consider the search space of all subsets of $`P`$ itself, and regard each such subset as a candidate for the set of generating rules of a stable model. Indeed, if $`M\text{At}(P)`$ and $`P^{}P`$, then $`M`$ is a stable model of $`P`$ if and only if $`M=h(P^{})`$, $`P^{}`$ is the set of all generating rules for $`M`$ in $`P`$ and $`M=LM(horn(P^{}))`$.
The problem with the stable model semantics is that, even in the propositional case, reasoning with logic programs under the stable model semantics is computationally hard. It is well-known that deciding whether a finite propositional logic program has a stable model is NP-complete \[Marek & Truszczyński, 1991\]. Consequently, it is not at all clear that logic programming with the stable model semantics can serve as a practical computational tool.
This issue can be resolved by implementing systems computing stable models and by experimentally studying the performance of these systems. Several such projects are now under way. Niemelä and Simons \[Niemelä & Simons, 1996\] developed a system, smodels, for computing stable models of finite function symbol-free logic programs and reported very promising performance results. For some classes of programs, smodels decides the existence of a stable model in a matter of seconds even if an input program consists of tens of thousands of clauses. Encouraging results on using smodels to solve planning problems are reported in \[Niemelä, 1999\]. Another well-advanced system is DeReS \[Cholewiński et al., 1996\], designed to compute extensions of arbitrary propositional default theories but being especially effective for default theories encoding propositional logic programs. Finally, systems capable of reasoning with disjunctive logic programs were described in \[Eiter et al., 1997\] and \[Aravindan et al., 1997\].
However, faster implementations will ultimately depend on better understanding of the algorithmic aspects of reasoning with logic programs under the stable model semantics. In this paper, we investigate the complexity of deciding whether a finite propositional logic program has stable models of some restricted sizes. Specifically, we study the following two problems ($`|P|`$ stands for the number of rules in a logic program $`P`$):
(Large stable models) Given a finite propositional logic program $`P`$ and an integer $`k`$, decide whether there is a stable model of $`P`$ of size at least $`|P|k`$.
(Small stable models) Given a finite propositional logic program $`P`$ and an integer $`k`$, decide whether there is a stable model of $`P`$ of size no more than $`k`$.
Inputs to the problems LSM and SSM are pairs $`(P,k)`$, where $`P`$ is a finite propositional logic program and $`k`$ is a non-negative integer. Problems of this type are referred to as parametrized decision problems. By fixing a parameter, a parameterized decision problem gives rise to its fixed-parameter version. In the case of problems LSM and SSM, by fixing $`k`$ we obtain the following two fixed-parameter problems ($`k`$ is now no longer a part of input):
Given a finite propositional logic program $`P`$, decide whether $`P`$ has a stable model of size at least $`|P|k`$.
Given a finite propositional logic program $`P`$, decide whether $`P`$ has a stable model of size at most $`k`$.
The problems LSM and SSM are NP-complete. It follows directly from the NP-completeness of the problem of existence of stable models \[Marek & Truszczyński, 1991\]. But fixing $`k`$ makes a difference! Clearly, the fixed-parameter problems $`\text{SSM}(k)`$ and $`\text{LSM}(k)`$ can be solved in polynomial time (unlike the problems SSM and LSM which, most likely, cannot). Indeed, consider a finite propositional logic program $`P`$. Then, there are $`O(n^k)`$ subsets of $`\text{At}(P)`$ (in fact, as pointed out earlier, it is enough to consider subsets of $`h(P)`$ or $`Neg(P)`$) of cardinality at most $`k`$ (we recall that in the paper $`n`$ stands for the number of atoms in $`P`$). For each such subset $`M`$, it can be checked in time linear in $`m`$ — the size of $`P`$ — whether $`M`$ is a stable model of $`P`$. Thus, one can decide whether $`P`$ has a stable model of size at most $`k`$ in time $`O(mn^k)`$.
Similarly, there are only $`O(|P|^k)`$ subsets of $`P`$ of size at least $`|P|k`$. Each such subset is a candidate for the set of generating rules of a stable model of size at least $`|P|k`$ (and smaller subsets, clearly, are not). Given such a subset $`R`$, one can check in time $`O(m)`$ whether $`R`$ generates a stable model for $`P`$. Thus, it follows that there is an algorithm that decides in time $`O(m|P|^k)`$ whether a logic program $`P`$ has a stable model of size at least $`|P|k`$.
While both algorithms are polynomial in the size of the program, their asymptotic complexity is expressed by the product of the size of a program and a polynomial of order $`k`$ in the number of atoms of the program or in the number of rules of the program. Even for small values of $`k`$, say for $`k4`$, the functions $`mn^k`$ and $`m|P|^k`$ grow very fast with $`m=size(P)`$, $`n=|\text{At}(P)|`$ and $`|P|`$, and render the corresponding algorithms infeasible.
An important question is whether algorithms for problems $`\text{SSM}(k)`$ and $`\text{LSM}(k)`$ exist whose order is significantly lower than $`k`$, preferably, a constant independent of $`k`$. The study of this question is the main goal of our paper. A general framework for such investigations was proposed by Downey and Fellows \[Downey & Fellows, 1997\]. They introduced the concepts of fixed-parameter tractability and fixed-parameter intractability that are defined in terms of a certain hierarchy of complexity classes known as the $`W`$ hierarchy.
In the paper, we show that the problem LSM is fixed-parameter tractable and demonstrate an algorithm that for every fixed $`k`$ decides the problem $`\text{LSM}(k)`$ in linear time — a significant improvement over the straightforward algorithm presented earlier.
On the other hand, we demonstrate that the problem SSM is much harder. We present an algorithm to decide the problems $`\text{SSM}(k)`$, for $`k1`$, that is asymptotically faster than the simple algorithm described above but the improvement is rather insignificant. Our algorithm runs in time $`O(mn^{k1})`$, an improvement only by the factor of $`n`$. The difficulty in finding a substantially better algorithm is not coincidental. We provide evidence that the problem SSM is fixed-parameter intractable. This result implies it is unlikely that there is an algorithm to decide the problems $`\text{SSM}(k)`$ whose running time would be given by a polynomial of order independent of $`k`$.
The study of fixed-parameter tractability of problems occurring in the area of nonmonotonic reasoning is a relatively new research topic. Another paper that pursues this direction is \[Gottlob et al., 1999\]. The authors focus there on parameters describing structural properties of programs and show that in some cases, fixing these parameters leads to polynomial algorithms.
Our paper is organized as follows. In Section 2, we recall basic concepts of the theory of fixed-parameter intractability by Downey and Fellows \[Downey & Fellows, 1997\]. The following two sections present the algorithms to decide the problems LSM and SSM, respectively. The next section focuses on the issue of fixed-parameter intractability of the problem SSM and contains the two main results of the paper. The last section contains conclusions and open problems.
## 2 Fixed-parameter intractability
This section recalls basic ideas of the work of Downey and Fellows on fixed-parameter intractability. The reader is referred to \[Downey & Fellows, 1997\] for a detailed treatment of this subject.
Informally, a parametrized decision problem is a decision problem whose inputs are pairs of items, one of which is referred to as a parameter. The graph colorability problem is an example of a parametrized problem. The inputs are pairs $`(G,k)`$, where $`G`$ is an undirected graph and $`k`$ is a non-negative integer. The problem is to decide whether $`G`$ can be colored with at most $`k`$ colors. Another example is the vertex cover problem in a graph. Again, the inputs are graph-integer pairs $`(G,k)`$ and the question is whether $`G`$ has a vertex cover of cardinality $`k`$ or less. The problems SSM and LSM are also examples of parametrized decision problems. Formally, a parametrized decision problem is a set $`L\mathrm{\Sigma }^{}\times \mathrm{\Sigma }^{}`$, where $`\mathrm{\Sigma }`$ is a fixed alphabet.
By selecting a concrete value $`\alpha \mathrm{\Sigma }^{}`$ of the parameter, a parametrized decision problem $`L`$ gives rise to an associated fixed-parameter problem $`L_\alpha =\{x:(x,\alpha )L\}`$. For instance, by fixing the value of $`k`$ to 3, we get a fixed-parameter version of the colorability problem, known as 3-colorability. Inputs to the 3-colorability problem are graphs and the question is to decide whether an input graph can be colored with 3 colors. Clearly, the problems $`\text{SSM}(k)`$ ($`\text{LSM}(k)`$, respectively) are fixed-parameter versions of the problem SSM (LSM, respectively).
The interest in the fixed-parameter problems stems from the fact that they are often computationally easier than the corresponding parametrized problems. For instance, the problems SSM and LSM are NP-complete yet, as we saw earlier, their parametrized versions $`\text{SSM}(k)`$ and $`\text{LSM}(k)`$ can be solved in polynomial time. Similarly, the vertex cover problem is NP-complete but its fixed-parameter versions are in the class P. To see this, observe that to decide whether a graph has a vertex cover of size at most $`k`$, where $`k`$ is a fixed value and not a part of an input, it is enough to generate all subsets with at most $`k`$ elements of the vertex set of a graph, and then check if any of them is a vertex cover. A word of caution is in order here. It is not always the case that fixed-parameter problems are easier. For instance, the 3-colorability problem is still NP-complete.
As we already pointed out, the fact that a problem admits a polynomial-time solution does not necessarily mean that practical algorithms to solve it exist. An algorithm that runs in time $`O(N^{15})`$, where $`N`$ is the size of the input, is hardly more practical than an algorithm with an exponential running time (and may even be a worse choice in practice). The algorithms we presented so far to argue that the problems $`\text{SSM}(k)`$, $`\text{LSM}(k)`$ and the fixed-parameter versions of the vertex cover problem are in P rely on searching through the space of $`N^k`$ possible solutions (where $`N`$ is the number of atoms of a program, the number of rules of a program, or the number of vertices in a graph, respectively). Thus, these algorithms are not practical, except for the very smallest values of $`k`$. The key question is how fast those polynomial-time solvable fixed-parameter problems can really be solved. Or, in other words, can one significantly improve over the brute-force approach?
A technique to deal with such questions is provided by the fixed-parameter intractability theory of Downey and Fellows \[Downey & Fellows, 1997\]. A parametrized problem $`L\mathrm{\Sigma }^{}\times \mathrm{\Sigma }^{}`$ is fixed-parameter tractable if there exist a constant $`p`$, an integer function $`f`$ and an algorithm $`A`$ such that $`A`$ determines whether $`(x,y)L`$ in time $`f(|y|)|x|^p`$ ($`|z|`$ stands for the length of a string $`z\mathrm{\Sigma }^{}`$). The class of fixed-parameter tractable problems will be denoted by FPT. Clearly, if a parametrized problem $`L`$ is in FPT, each of the associated fixed-parameter problems $`L_y`$ is solvable in polynomial time by an algorithm whose exponent does not depend on the value of the parameter $`y`$. It is known (see \[Downey & Fellows, 1997\]) that the vertex cover problem is in FPT.
There is substantial evidence to support a conjecture that some parametrized problems whose fixed-parameter versions are in P are not fixed-parameter tractable. To study and compare complexity of parametrized problems Downey and Fellows proposed the following notion of reducibility<sup>2</sup><sup>2</sup>2The definition given here is sufficient for the needs of this paper. To obtain structural theorems a subtler definition is needed. This topic goes beyond the scope of the present paper. The reader is referred to \[Downey & Fellows, 1997\] for more details.. A parametrized problem $`L`$ can be reduced to a parametrized problem $`L^{}`$ if there exist a constant $`p`$, an integer function $`q`$ and an algorithm $`A`$ that to each instance $`(x,y)`$ of $`L`$ assigns an instance $`(x^{},y^{})`$ of $`L^{}`$ such that
1. $`x^{}`$ depends upon $`x`$ and $`y`$ and $`y^{}`$ depends upon $`y`$ only,
2. $`A`$ runs in time $`O(q(|y|)|x|^p)`$,
3. $`(x,y)L`$ if and only if $`(x^{},y^{})L^{}`$.
Downey and Fellows also defined a hierarchy of complexity classes called the W hierarchy:
$$\mathrm{F}PT\mathrm{W}[1]\mathrm{W}[2]\mathrm{W}[3]\mathrm{}$$
(1)
The classes W\[t\] can be described in terms of problems that are complete for them (a problem $`D`$ is complete for a complexity class $`E`$ if $`DE`$ and every problem in this class can be reduced to $`D`$). Let us call a boolean formula $`t`$-normalized if it is of the form of product-of-sums-of-products … of literals, with $`t`$ being the number of products-of, sums-of expressions in this definition. For example, 2-normalized formulas are products of sums of literals. Thus, the class of 2-normalized formulas is precisely the class of CNF formulas. We define the weighted $`t`$-normalized satisfiability problem as:
Given a $`t`$-normalized formula $`\phi `$, decide whether there is a model of $`\phi `$ with exactly $`k`$ atoms (or, alternatively, decide whether there is a satisfying valuation for $`\phi `$ which assigns the logical value true to exactly $`k`$ atoms)
Downey and Fellows show that for $`t2`$, the problems $`WS(t)`$ are complete for the class W\[t\]. They also show that a restricted version of the problem $`WS(2)`$:
Given a 3CNF formula $`\phi `$ and an integer $`k`$ (parameter), decide whether there is a model of $`\phi `$ with exactly $`k`$ atoms
is complete for the class $`W[1]`$. Downey and Fellows conjecture that all the implications in (1) are proper<sup>3</sup><sup>3</sup>3If true, this conjecture would imply that in the context of fixed-parameter tractability there is a difference between the complexity of weighted satisfiability for 3CNF and CNF formulas.. In particular, they conjecture that problems in the classes W\[t\], with $`t1`$, are not fixed-parameter tractable.
In the paper, we relate the problem SSM to the problems $`WS(2)`$ and $`WS(3)`$ to place the problem SSM in the W hierarchy, to obtain estimates of its complexity and to argue for its fixed-parameter intractability.
## 3 Large stable models
In this section we will show an algorithm for the parametrized problem LSM that runs in time $`O(2^{k+k^2}m)`$, where $`(P,k)`$ is an input instance and, as in all other places in the paper, $`m=size(P)`$. This result implies that the problem LSM is fixed-parameter tractable and that there is an algorithm that for every fixed $`k`$ solves the problem $`\text{LSM}(k)`$ in linear-time.
Given a logic program $`P`$, denote by $`P^{}`$ the logic program obtained from $`P`$ by eliminating from the bodies of the rules in $`P`$ all literals $`\text{not}(a)`$, where $`a`$ is not the head of any rule from $`P`$. The following well-known result states the key property of the program $`P^{}`$.
###### Lemma 3.1
A set of atoms $`M`$ is a stable model of a logic program $`P`$ if and only if $`M`$ is a stable model of $`P^{}`$.
Lemma 3.1 implies that the problem LSM has a positive answer for $`(P,k)`$ if and only if it has a positive answer for $`(P^{},k)`$. Moreover, it is easy to see that $`P^{}`$ can be constructed from $`P`$ in time linear in the size of $`P`$. Thus, when looking for algorithms to decide the problem LSM we may restrict our attention to programs $`P`$ in which every atom appearing negated in the body of a rule appears also as the head of a rule (that is, to such programs $`P`$ for which we have $`Neg(P)h(P)`$).
By $`P^k`$ let us denote the program consisting of those rules $`r`$ in $`P`$ for which $`|b^{}(r)|k`$. We have the following lemma.
###### Lemma 3.2
Let $`P`$ be a logic program such that $`Neg(P)h(P)`$. Let $`M\text{At}(P)`$ be a set of atoms such that $`|M||P|k`$. Then:
1. $`M`$ is a stable model of $`P`$ if and only if $`M`$ is a stable model of $`P^k`$
2. if $`M`$ is a stable model of $`P^k`$, then $`P^k`$ has no more than $`k+k^2`$ different negated literals appearing in the bodies of its rules.
Proof: (1) Consider a rule $`rPP^k`$. Then $`|b^{}(r)|k+1`$ and, consequently, $`b^{}(r)M\mathrm{}`$. Indeed, if $`b^{}(r)M=\mathrm{}`$, then $`|Mb^{}(r)|=|M|+|b^{}(r)|>|P|`$. Since $`Neg(P)h(P)`$, $`b^{}(r)h(P)`$. In addition, (both if we assume that $`M`$ is a stable model of $`P`$ and if we assume that $`M`$ is a stable model of $`P^k`$), we have $`Mh(P)`$. Thus, $`b^{}(r)Mh(P)`$. Now observe that $`|P||h(P)|`$. Thus, $`|Mb^{}(r)||h(P)||P|`$, a contradiction.
Since for every rule $`rPP^k`$ we have $`b^{}(r)M\mathrm{}`$, it follows that $`(P^k)^M=P^M`$. Hence, $`M=LM(P^M)`$ if and only if $`M=LM((P^k)^M)`$. Consequently, $`M`$ is a stable model of $`P`$ if and only if $`M`$ is a stable model of $`P^k`$.
(2) Let $`P^{}`$ be the set of rules from $`P^k`$ such that $`rP^{}`$ if and only if $`b^{}(r)M=\mathrm{}`$ (the rules in $`P^{}`$ contribute to the reduct $`(P^k)^M`$) and let $`P^{\prime \prime }`$ be the set of the remaining rules in $`P^k`$ (these are the rules that are eliminated when the reduct $`(P^k)^M`$ is computed). Since $`Neg(P)h(P)`$, for every rule $`rP`$, $`b^{}(r)h(P)`$. Thus, $`\{b^{}(r):rP^{}\}h(P)M`$. Since $`Mh(P)`$ (as $`M`$ is a stable model of $`P^k`$) and $`|P||h(P)|`$, we have $`|\{b^{}(r):rP^{}\}|k`$. Further, since $`|P^{}||M||P|k|P^k|k`$, it follows that $`|P^{\prime \prime }|k`$. Consequently, $`|\{b^{}(r):rP^{\prime \prime }\}|k^2`$. Hence, the second part of the assertion follows. $`\mathrm{}`$
Let us now consider the following algorithm for the problem $`\text{LSM}(k)`$ (the input to this algorithm is a logic program $`P`$).
1. Eliminate from the input logic program $`P`$ all literals $`\text{not}(a)`$, where $`a`$ is not the head of any rule from $`P`$. Denote the resulting program by $`Q`$.
2. Compute the set of rules $`Q^k`$ consisting of those rules $`r`$ in $`Q`$ for which $`|b^{}(r)|k`$.
3. Decide whether $`Q^k`$ has a stable model $`M`$ such that $`|M||Q|k`$.
This algorithm reports YES if and only if the program $`Q^k`$ has a stable model $`M`$ such that $`|M||Q|k`$. By Lemma 3.2, that happens precisely if and only if $`Q`$ has a stable model $`M`$ such that $`|M||Q|k`$. This last statement, by Lemma 3.1, is equivalent to the statement that $`P`$ has a stable model $`M`$ such that $`|M||P|k`$. In other words, our algorithm correctly decides the problem $`\text{LSM}(k)`$.
Let us notice that steps 1 and 2 can be implemented in time $`O(m)`$, where the constant hidden by the “big O” notation does not depend on $`k`$. To implement step 3, let us recall that every stable model of a logic program is determined by some subset of the set of atoms that appear negated in the program (each such subset uniquely determines the reduct, as we stated in the introduction; see also \[Bondarenko et al., 1993\]). By Lemma 3.2, the set of such atoms in the program $`Q^k`$ has cardinality at most $`k+k^2`$. Checking for each subset of this set whether it determines a stable model of $`Q^k`$ can be implemented in time $`O(size(Q^k))=O(m)`$. Consequently, our algorithm runs in time $`O(2^{k+k^2}m)`$ (with the constant hidden by the “big O” notation independent of $`k`$).
###### Theorem 3.3
The problem LSM is fixed-parameter tractable. Moreover, for each fixed $`k`$ there is a linear-time algorithm to decide whether a logic program $`P`$ has a stable model of size at least $`|P|k`$.
## 4 Computing stable models of size at most $`k`$
In the introduction we pointed out that there is a straightforward algorithm to decide the problem $`\text{SSM}(k)`$ that runs in time $`O(mn^k)`$, where $`m=size(P)`$ and $`n=|\text{At}(P)|`$. For $`k1`$ (the assumption we adopt in this section), this algorithm can be slightly improved. Namely, we will now describe an algorithm for the problem $`\text{SSM}(k)`$ that runs in time $`O(F(k)mn^{k1})`$, where $`F`$ is some integer function. Thus, if $`k`$ is fixed and not a part of the input, this improved algorithm runs in time $`O(mn^{k1})`$.
We present our algorithm under the assumption that input logic programs are proper. We say that a logic program rule $`r`$ is proper if:
$`h(r)b^+(r)`$, and
$`b^+(r)b^{}(r)=\mathrm{}`$
We say that a logic program $`P`$ is proper if all its rules are proper. Rules that violate at least one of the conditions (P1) and (P2) (that is, rules that are not proper) have no influence on the collection of stable models of a program as we have the following well-known result (see, for instance, \[Brass & Dix, 1997\]).
###### Lemma 4.1
A set of atoms $`M`$ is a stable model of a logic program $`P`$ if and only if $`M`$ is a stable model of the subprogram of $`P`$ consisting of all proper rules in $`P`$.
It is easy to see that rules that violate (P1) or (P2) can be eliminated from a logic program $`P`$ in time $`O(m)`$. Thus, the restriction to proper programs does not affect the generality of our discussion.
For a proper logic program $`P`$ and for a set $`A\text{At}(P)`$ of atoms, we define $`P(A)`$ to be the program consisting of all those rules $`r`$ of $`P`$ that are not blocked by $`A`$ (in other words, those that satisfy $`b^{}(r)A=\mathrm{}`$) and whose positive body is contained in $`A`$ (in other words, such that $`b^+(r)A`$).
Let $`P`$ be a logic program and let $`A\text{At}(P)`$ be a set of atoms. A stable model $`M`$ of $`P`$ is called $`A`$-based if
1. $`M`$ is of the form $`A\{a\}`$, where $`a\text{At}(P)A`$, and
2. $`MLM(P(A)^M)`$ (in other words, when computing $`LM(P^M)`$, the derivation of $`A`$ does not require that $`a`$ be derived first).
We have the following simple lemma.
###### Lemma 4.2
Let $`k`$ be an integer such that $`k1`$. A proper logic program $`P`$ has a stable model of cardinality $`k`$ if and only if for some $`A\text{At}(P)`$, with $`|A|=k1`$, $`P`$ has an $`A`$-based stable model.
It follows from Lemma 4.2 that when deciding the existence of $`k`$-element stable models, $`k1`$, it is enough to focus on the existence of $`A`$-based stable models. This is the approach we take here. In most general terms, our algorithm for the problem $`\text{SSM}(k)`$ consists of generating all subsets $`A\text{At}(P)`$, with $`|A|k1`$, and for each such subset $`A`$, of checking whether $`P`$ has an $`A`$-based stable model. This latter task is the key.
We will now describe an algorithm that, given a logic program $`P`$ and a set $`A\text{At}(P)`$, decides whether $`P`$ has an $`A`$-based stable model. To this end, we define $`P^{}(A)`$ to be the program consisting of all those rules $`r`$ of $`P`$ such that:
1. $`b^{}(r)A=\mathrm{}`$ ($`r`$ is not blocked by $`A`$)
2. $`h(r)A`$
3. $`b^+(r)A`$ consists of exactly one element; we will denote it by $`a_r`$.
Our algorithm is based on the following result allowing us to restrict attention to the program $`P(A)`$ (the statement of the lemma and its proof rely on the terminology introduced above).
###### Lemma 4.3
Let $`A`$ be a set of atoms. A proper logic program $`P`$ has an $`A`$-based stable model if and only if $`P(A)`$ has an $`A`$-based stable model $`M=A\{a\}`$, such that $`a\{a_r:rP^{}(A)\}`$.
Proof: ($``$) Let $`M`$ be an $`A`$-based stable model of $`P`$. Assume that $`M=A\{a\}`$, for some $`aA`$. Since $`P(A)^MP^M`$, $`LM(P(A)^M)LM(P^M)=M`$. Since $`M`$ is $`A`$-based, we have that $`MLM(P(A)^M)`$. It follows that $`M`$ is an $`A`$-based stable model of $`P(A)`$.
Let us assume that there is a rule $`sP^{}(A)`$ such that $`a=a_s`$. The rule $`s`$ is not blocked by $`A`$. Since $`ab^+(s)`$, we have that $`ab^{}(s)`$ (we recall that all rules in $`P`$ are proper). Hence, $`s`$ is not blocked by $`\{a\}`$ either. Consequently, $`horn(s)P^M`$. Since $`sP^{}(A)`$, the body of $`horn(s)`$ (that is, $`b^+(s)`$) is contained in $`M`$. The set $`M`$ is a least model of $`P^M`$. In particular, $`M`$ satisfies $`horn(s)`$. Thus, it follows that $`h(s)M`$. In the same time, $`h(s)a`$ (as $`s`$ is proper). Thus, $`h(s)A`$, a contradiction (we recall that $`sP^{}(A)`$). It follows that $`a\{a_r:rP^{}(A)\}`$.
($``$) We will now assume that $`M=A\{a\}`$ is an $`A`$-based stable model of $`P(A)`$ such that $`a\{a_r:rP^{}(A)\}`$. Similarly as before, we have $`M=LM(P(A)^M)LM(P^M)`$. Let us assume that $`LM(P^M)M\mathrm{}`$. Then there is a rule $`t`$ in $`P^M`$ such that the body of $`t`$ is contained in $`M`$ and $`h(t)M`$. Let $`s`$ be a rule in $`P`$ that gives rise to $`t`$ when constructing the reduct. Assume first that the body of $`t`$ (that is, $`b^+(s)`$) is contained in $`A`$. Then $`sP(A)`$, $`tP(A)^M`$ and, consequently, $`h(t)LM(P(A)^M)=M`$, a contradiction.
Thus, the body of $`t`$ is not contained in $`A`$. Since the body of $`t`$ is contained in $`M`$, it consists of $`a`$ and, possibly, some other elements, all of which are in $`A`$. It follows that $`sP^{}(A)`$. Consequently, $`a=a_s`$ and $`a\{a_r:rP^{}(A)\}`$, a contradiction. Thus, $`LM(P^M)=M`$, that is, $`M`$ is a stable model of $`P`$. Since $`M=LM(P(A)^M)`$, it follows that $`M`$ is an $`A`$-based model of $`P`$. $`\mathrm{}`$
Let $`A`$ be a set of atoms. A logic program with negation, $`P`$, is an $`A`$-program if $`P=P(A)`$, that is if for every rule $`rP`$ we have $`b^+(P)A`$ and $`b^{}(P)A=\mathrm{}`$. Clearly, the program $`P(A)`$, described above, is an $`A`$-program. We will now focus on $`A`$-programs and their $`A`$-based stable models.
Let $`A`$ be a set of atoms. We denote by $`R(A)`$ the set of all proper Horn rules over the set of atoms $`A`$. Clearly, the cardinality of $`R(A)`$ depends on the cardinality of $`A`$ only. Further, we define $`P(A)`$ to be the set of all Horn programs $`QR(A)`$ satisfying the condition $`LM(Q)=A`$. As in the case of $`R(A)`$, the cardinality of $`P(A)`$ also depends on the size of $`A`$ only.
We will now describe conditions that determine whether an $`A`$-program $`P`$ has an $`A`$-based stable model. To this end, with every atom $`a\text{At}(P)A`$, we associate the following values:
* $`F(a)=1`$ if there is a rule $`s`$ in $`P`$ with $`h(s)A\{a\}`$ and $`ab^{}(s)`$; $`F(a)=0`$, otherwise
* $`G(a)=`$ the number of rules $`s`$ in $`P`$ with $`h(s)=a`$ and $`ab^{}(s)`$.
Further, with every proper Horn rule $`rR(A)`$ and every atom $`a\text{At}(P)A`$, we associate the quantity:
* $`H(r,a)=1`$ if there is a rule $`s`$ in $`P`$ with $`horn(s)=r`$ and $`ab^{}(s)`$; $`H(r,a)=0`$, otherwise.
The following lemma characterizes $`A`$-based stable models of an $`A`$-program. Both the statement of the lemma and its proof rely on the terminology introduced above.
###### Lemma 4.4
Let $`A`$ be a set of atoms, let $`P`$ be an $`A`$-program and let $`a`$ be an atom such that $`a\text{At}(P)A`$. Then $`A\{a\}`$ is an $`A`$-based stable model of $`P`$ if and only if $`F(a)=0`$, $`G(a)>0`$, and for some program $`QP(A)`$ and for every rule $`rQ`$, $`H(r,a)>0`$.
Proof: $`()`$ We denote $`M=A\{a\}`$ and assume that $`M`$ is an $`A`$-based stable model for $`P`$. It follows that $`M=LM(P^M)`$. Let $`P_A`$ be the subprogram of $`P`$ consisting of those rules of $`P`$ whose head belongs to $`A`$. Since $`M`$ is an $`A`$-based stable model of $`P`$, we have $`A=LM(P_A^M)`$. Let $`Q`$ be the program obtained from $`P_A^M`$ by removing multiple occurrences of rules. Clearly, $`QP(A)`$. It follows directly from the definition of the reduct that for every rule $`rQ`$, $`H(r,a)=1`$.
Next, we observe that $`aLM(P^M)`$. Thus, $`G(a)>0`$. Let us assume that $`F(a)=1`$. Let $`r`$ be a rule in $`P`$ such that $`h(r)A\{a\}`$ and $`ab^{}(r)`$. Since $`P`$ is an $`A`$-program, $`Ab^{}(r)=\mathrm{}`$. Thus, it follows that $`horn(r)P^M`$. We also have that $`b^+(r)AM`$. Since $`M`$ is a model of $`P^M`$, $`h(r)M`$. However, in the same time we have that $`h(r)A\{a\}(=M)`$, a contradiction. It follows that $`F(a)=0`$.
$`()`$ We now assume that for some $`a\text{At}(P)A`$, $`F(a)=0`$, $`G(a)>0`$ and for some program $`QP(A)`$ and for every rule $`rQ`$, $`H(r,a)=1`$. As before, we set $`M=A\{a\}`$. We will show that $`M=LM(P^M)`$.
First, since $`P`$ is an $`A`$-program and $`H(r,a)=1`$ for every rule $`rQ`$, it follows that $`QP(A)^M`$. Thus, $`ALM(P(A)^M)`$. Second, we have that $`G(a)>0`$. Thus, there is a rule $`rP`$ such that $`h(r)=a`$ and $`ab^{}(r)`$. It follows that $`horn(r)Q`$ and $`horn(r)P^M`$. Since $`QP(A)^M`$, $`A=LM(Q)`$ and $`b^+(r)A`$, we obtain that $`aLM(P(A)^M)`$. Thus, $`MLM(P(A)^M)`$. Finally, since $`F(a)=0`$, we have that for every rule $`sP`$ such that $`ab^{}(s)`$, $`h(s)M`$. Thus, $`LM(P^M)`$ does not contain any atom not in $`M`$. Consequently, $`M=LM(P^M)`$ and $`M`$ is a stable model of $`P`$. Since $`MLM(P(A)^M)`$, $`M`$ is an $`A`$-based stable model of $`P`$. $`\mathrm{}`$
We will discuss now effective ways to compute values $`F(a)`$, $`G(a)`$ and $`H(r,a)`$. Clearly, computing the values $`G(a)`$ can be accomplished in time linear in the size of the program, that is, in time $`O(m)`$. Indeed, we start by initializing all values $`G(a)`$ to 0. Then, for each rule $`sP`$, we set $`G(h(s)):=G(h(s))+1`$ if $`h(s)b^{}(s)`$, and leave $`G(h(s))`$ unchanged, otherwise. To decide which is the case requires that we scan all negated lierals in the body of $`s`$. That takes time $`O(|b^{}(s)|)`$. Thus, the overall time is $`O(m)`$.
Computing values $`F(a)`$ and $`H(r,a)`$ is more complicated. First, we prove the following lemma.
###### Lemma 4.5
Let $`P`$ be an $`A`$-program, let $`a\text{At}(P)A`$ and let $`rR(A)`$. Then
1. $`F(a)=1`$ if and only if $`a\{\{h(s)\}b^{}(s):sP,h(s)A\}`$.
2. $`H(r,a)=1`$ if and only if $`a\{b^{}(s):sP,horn(s)=r\}`$.
Proof: (1) Let us assume first that $`F(a)=1`$. Then there is a rule $`sP`$ such that $`h(s)A\{a\}`$ and $`ab^{}(s)`$. Thus, $`a\{h(s)\}b^{}(s)`$. Consequently, the identity $`a\{\{h(s)\}b^{}(s):sP,h(s)A\}`$ follows. All the implications in this argument can be reversed. Hence, we obtain the assertion (1).
(2) Let us assume that $`H(r,a)=1`$. Then, there is a rule $`sP`$ such that $`horn(s)=r`$ and $`ab^{}(s)`$. Consequently, $`a\{b^{}(s):sP,horn(s)=r\}`$. As in (1), all the implications are in fact equivalences and the assertion (2) follows. $`\mathrm{}`$
Lemma 4.5 shows that to compute all the values $`F(a)`$ one has to compute the set
$$\{\{h(s)\}b^{}(s):sP,h(s)A\}.$$
To this end, for each atom $`a`$ we will compute the number of sets in $`\{\{h(s)\}b^{}(s):sP,h(s)A\}`$ that $`a`$ is a member of. We will denote this number by $`C(a)`$. We first initialize all values $`C(a)`$ to 0. Then, we consider all sets in $`\{\{h(s)\}b^{}(s):sP,h(s)A\}`$ in turn. For each such set and for each atom $`a`$ in this set we set $`C(a):=C(a)+1`$. The set $`\{\{h(s)\}b^{}(s):sP,h(s)A\}`$ is given by all those atoms $`a`$ for which $`C(a)`$ is equal to the number of sets in $`\{\{h(s)\}b^{}(s):sP,h(s)A\}`$. It is clear that the time needed for this computation is linear in the size of the program (assuming appropriate linked-list representation of rules). Thus, all the values $`F(a)`$ can be computed in time linear in the size of the program, that is, in $`O(m)`$ steps.
To compute values $`H(r,a)`$ we proceed similarly. First, we compute all the sets $`\{s:sP,horn(s)=r\}`$, where $`rR(A)`$. To this end, we scan all rules in $`P`$ in order and for each of them we find the rule $`rR(A)`$ such that $`horn(s)=r`$. Then we include $`s`$ in the set $`\{s:sP,horn(s)=r\}`$. Given $`s`$, it takes $`O(g|A|)`$ steps to identify rule $`r`$ (where $`g`$ is some function). Indeed, the size of $`b^+(s)`$ is bound by $`|A|`$ as $`P`$ is an $`A`$-program. Moreover, the number of rules in $`R(A)`$ depends on $`|A|`$ only. Thus, the task of computing all sets $`\{s:sP,horn(s)=r\}`$, for $`rR(A)`$, can be accomplished in $`O(g(|A|)|P|)`$ steps. Next, for each these sets of rules, we proceed as in the case of values $`F(a)`$, to compute their intersections. Each such computation takes time $`O(m)`$, where $`m=size(P)`$). Thus, computing all the values $`H(r,a)`$ can be accomplished in time $`O(g(|A|)|P|+|R(A)|m)=O(f(|A|)m)`$, for some function $`f`$.
We can now put all the pieces together. As a result of our considerations, we obtain the following algorithm for deciding the problem $`\text{SSM}(k)`$.
| | | | | | |
| --- | --- | --- | --- | --- | --- |
| Algorithm to decide the problem $`\text{SSM}(k)`$, $`k1`$ |
| Input: A logic program $`P`$ ($`k`$ is not a part of input) |
| |
| (0) | if $`\mathrm{}`$ is a stable model of $`P`$ then return YES and exit; |
| (1) | $`P:=`$ the set of proper rules in $`P`$; |
| (2) | for every $`A\text{At}(P)`$ with $`|A|k1`$ do |
| (3) | | compute the set of rules $`R(A)`$ and the set of programs $`P(A)`$; |
| (4) | | compute the program $`P(A)`$; |
| (5) | | compute the program $`P^{}(A)`$ and the set $`B=\{a_r:rP^{}(A)\}`$; |
| (6) | | given $`P(A)`$ and $`R(A)`$, compute tables $`F`$, $`G`$ and $`H`$ (as described above); |
| (7) | | for every $`a\text{At}(P(A))AB`$ do |
| (8) | | | if |
| (9) | | | | $`F(a)=0`$, $`G(a)>0`$ and |
| (10) | | | | there is a program $`QP(A)`$ s. t. for every rule $`rQ`$, $`H(r,a)>0`$ |
| (11) | | | then report YES and exit; |
| (12) | report NO and exit. |
The correctness of this algorithm follows from Lemmas 4.2 \- 4.4. We will now analyze the running time of this algorithm. Clearly, line (0) can be executed in $`O(m)`$ steps. As we already observed, rules that are not proper can be eliminated from $`P`$ in time $`O(m)`$. Next, there are $`O(n^{k1})`$ iterations of loop (2). In each of them, line (3) takes time $`O(f_1(k))`$, for some function $`f_1`$ (let us recall that $`|R(A)|`$ and $`|P(A)|`$ depend on $`|A|`$ only). Further, lines (4) and (5) can be executed in time $`O(m)`$. Line (6), as we discussed earlier, can be implemented so that to run in $`O(f(k)m)`$ steps. Loop (7) is executed $`O(n)`$ times and each iteration takes $`O(f_2(k))`$ steps, for some function $`f_2`$ (let us again recall that $`|P(A)|`$ depends on $`k`$ only). Thus, the running time of the whole algorithm is $`O(F(k)mn^{k1})`$, for some integer function $`F`$. Consequently, we get the following result.
###### Theorem 4.6
There is an integer function $`F`$ and an algorithm $`A`$ such that $`A`$ decides the problem $`\text{SSM}(k)`$ and runs in time $`O(F(k)mn^{k1})`$ (the constant hidden in the ”big Oh” notation does not depend on $`k`$).
## 5 Complexity of the problem SSM
The algorithm outlined in the previous section is not quite satisfactory. Its running time is still high. A natural question to ask is: are there significantly better algorithms for the problems $`\text{SSM}(k)`$? In this section we address this question by studying the complexity of the problem SSM. Our goal is to show that the problem is difficult in the sense of the W hierarchy. We will show that the problem SSM is $`W[2]`$-hard and that it is in the class W. To this end, we define the $`(k)`$-weighted $`t`$-normalized satisfiability problem as:
Given a $`t`$-normalized formula $`\phi `$, decide whether there is a model of $`\phi `$ with at most $`k`$ atoms ($`k`$ is a parameter).
The problem $`WS^{}(t)`$ is a slight variation of the problem $`WS(t)`$. It is known to be complete for the class W\[t\], for $`t2`$ (see \[Downey & Fellows, 1997\], page 468). To show W-hardness of SSM, we will reduce the problem $`WS^{}(2)`$ to the problem SSM. Given the overwhelming evidence of fixed-parameter intractability of problems that are $`W[2]`$-hard \[Downey & Fellows, 1997\], it is unlikely that algorithms for problems $`\text{SSM}(k)`$ exist whose asymptotic behavior would be given by a polynomial of order independent of $`k`$. To better delineate the location of the problem SSM in the W hierarchy we also provide an upper bound on its hardness by showing that it can be reduced to the problem $`WS^{}(3)`$, thus proving that the problem SSM belongs to the class $`W[3]`$.
We will start by showing that the problem $`\text{SSM}(k)`$ is reducible (in the sense of the definition from Section 2) to the problem $`WS^{}(3)`$. To this end, we describe an encoding of a logic program $`P`$ by means of a collection of clauses $`T(P)`$ so that $`P`$ has a stable model of size at most $`k`$ if and only if $`T(P)`$ has a model with no more than $`(k+1)(k^2+2k)`$ atoms. In the general setting of the class NP, an explicit encoding of the problem of existence of stable models in terms of propositional satisfiability was described in \[Ben-Eliyahu & Dechter, 1994\]. Our encoding, while different in key details, uses some ideas from that paper.
Let us consider an integer $`k`$ and a logic program $`P`$. For each atom $`q`$ in $`P`$ let us introduce new atoms $`c(q)`$, $`c(q,i)`$, $`1ik+1`$, and $`c^{}(q,i)`$, $`2ik+1`$. Intuitively, atom $`c(q)`$ represents the fact that in the process of computing the least model of the reduct of $`P`$ with respect to some set of atoms, atom $`q`$ is computed no later than during the iteration $`k+1`$ of the van Emden-Kowalski operator. Similarly, atom $`c(q,i)`$ represents the fact that in the same process atom $`q`$ is computed exactly in the iteration $`i`$ of the van Emden-Kowalski operator. Finally, atom $`c^{}(q,i)`$, expresses the fact that $`q`$ is computed before the iteration $`i`$ of the van Emden-Kowalski operator. The formulas $`F_1(q,i)`$, $`2ik+1`$, and $`F_2(q)`$ describe some basic relationships between atoms $`c(q)`$, $`c(q,i)`$ and $`c^{}(q,i)`$ that we will require to hold:
$$F_1(q,i)=c^{}(q,i)c(q,1)\mathrm{}c(q,i1),$$
$$F_2(q)=c(q)c(q,1)\mathrm{}c(q,k+1).$$
Let $`r`$ be a rule in $`P`$ with $`h(r)=q`$, say
$$r=qa_1,\mathrm{},a_s,\text{not}(b_1),\mathrm{},\text{not}(b_t).$$
We define a formula $`F_3(r,i)`$, $`2ik+1`$, by
$$F_3(r,i)=c^{}(a_1,i)\mathrm{}c^{}(a_s,i)\neg c(b_1)\mathrm{}\neg c(b_t)\neg c^{}(q,i).$$
We define $`F_3(r,1)=𝐟alse`$ (false is a distinguished contradictory formula in our propositional language) if $`s1`$. Otherwise, we define
$$F_3(r,1)=\neg c(b_1)\mathrm{}\neg c(b_t).$$
Speaking informally, formula $`F_3(r,i)`$ asserts that $`q`$ is computed by means of rule $`r`$ in the iteration $`i`$ of the least model computation process and that it has not been computed earlier.
Let $`r_1,\mathrm{},r_v`$ be all rules in $`P`$ with atom $`q`$ in the head. We define a formula $`F_4(q,i)`$, $`1ik+1`$, by
$$F_4(q,i)=c(q,i)F_3(r_1,i)\mathrm{}F_3(r_v,i).$$
Intuitively, the formula $`F_4(q,i)`$ asserts that when computing the least model of the reduct of $`P`$, atom $`q`$ is first computed in the iteration $`i`$.
We now define the theory $`T_0(P)`$ that encodes the problem of existence of small stable models:
$`T_0(P)`$ $`=`$ $`\{F_1(q,i):q\text{At}(P),2ik+1\}\{F_2(q):q\text{At}(P)\}`$
$`\{F_4(q,i):q\text{At}(P),1ik+1\}.`$
Next, we establish some useful properties of the theory $`T_0(P)`$. First, we consider a set $`U`$ of atoms that is a model of $`T_0(P)`$ and define
$$M(U)=\{q\text{At}(P):c(q)U\}.$$
###### Lemma 5.1
Let $`U`$ be a model of $`T_0(P)`$ and let $`qM(U)`$. Then there is a unique integer $`i`$, $`1ik+1`$, such that $`c(q,i)U`$.
Proof: Since $`U`$ is a model of a formula $`F_2(q)`$, there is an integer $`i`$, $`1ik+1`$, such that $`c(q,i)U`$. To prove uniqueness of such $`i`$, assume that there are two integers $`j_1`$ and $`j_2`$, $`1j_1<j_2k+1`$, such that $`c(q,j_1)U`$ and $`c(q,j_2)U`$. Since $`UF_4(q,j_2)`$, it follows that there is a rule $`rP`$ with $`h(r)=q`$ and such that $`UF_3(r,j_2)`$. In particular, $`U\neg c^{}(q,j_2)`$. In the same time, since $`c(q,j_1)U`$ and $`UF_1(q,j_2)`$, we have $`c^{}(q,j_2)U`$, a contradiction. $`\mathrm{}`$
For every atom $`qM(U)`$ define $`i_q`$ to be the integer whose existence and uniqueness is guaranteed by Lemma 5.1. Define $`i_U=\mathrm{max}\{i_q:qM(U)\}`$. Next, for each $`i`$, $`1ii_U`$, define
$$[M(U)]_i=\{qM(U):i_q=i\}.$$
###### Lemma 5.2
Let $`U`$ be a model of $`T_0(P)`$. Under the terminology introduced above, for every $`i`$, $`1ii_U`$, $`[M(U)]_i\mathrm{}`$.
Proof: We will proceed by downward induction. By the definition of $`i_U`$, $`[M(U)]_{i_U}\mathrm{}`$. Consider $`i`$, $`2ii_U`$, and assume that $`[M(U)]_i\mathrm{}`$. We will show that $`[M(U)]_{i1}\mathrm{}`$. Let $`q[M(U)]_i`$. Clearly, $`c(q,i)U`$ and, since $`UF_4(q,i)`$, there is a rule $`r=qa_1,\mathrm{},a_s,\text{not}(b_1),\mathrm{},\text{not}(b_t)`$ such that $`UF_3(r,i)`$. Consequently, for every $`j`$, $`1js`$, $`c^{}(a_j,i)U`$. Assume that for every $`j`$, $`1js`$, $`c^{}(a_j,i1)U`$. Since $`Uc^{}(q,i1)c^{}(q,i)`$ and since $`U\neg c^{}(q,i)`$, it follows that $`U\neg c^{}(q,i1)`$. Consequently, $`U`$ satisfies the formula $`F_3(r,i1)`$ and, so, $`UF_4(q,i1)`$. It follows that $`c(q,i1)U`$, a contradiction (we recall that $`i_q=i`$). Hence, there is $`j`$, $`1js`$, such that $`c(a_j,i1)U`$. It follows that $`a_j[M(U)]_{i1}`$ and $`[M(U)]_{i1}\mathrm{}`$. $`\mathrm{}`$
###### Lemma 5.3
Let $`U`$ be a model of $`T_0(P)`$ and let $`|M(U)|k`$. Then
1. $`i_Uk`$, and
2. $`M(U)`$ is a stable model of $`P`$.
Proof: (1) The assertion follows directly from the fact that $`|M(U)|k`$ and from Lemma 5.2.
(2) We need to show that $`M(U)=LM(P^{M(U)})`$. We will first show that $`M(U)LM(P^{M(U)})`$. Since $`M(U)=_{i=1}^{i_U}[M(U)]_i`$, we will show that for every $`i`$, $`1ii_U`$, $`[M(U)]_iLM(P^{M(U)})`$. We will proceed by induction. Let $`q[M(U)]_1`$. It follows that there is a rule $`r`$ such that $`UF_3(r,1)`$. Consequently, $`r`$ is of the form $`r=q\text{not}(b_1),\mathrm{},\text{not}(b_t)`$ and $`U\neg c(b_1)\mathrm{}\neg c(b_t)`$. Hence, for every $`j`$, $`1jt`$, $`b_jM(U)`$. Consequently, the rule $`(q.)`$ is in $`P^{M(U)}`$ and, so, $`qLM(P^{M(U)})`$. The inductive step is based on a similar argument. It relies on the inequality $`i_Uk`$ we proved in (1). We leave the details of the inductive step to the reader.
We will next show that $`LM(P^{M(U)})M(U)`$. We will use the characterization of $`LM(P^{M(U)})`$ as the limit of the sequence of iterations of the van Emden-Kowalski operator $`T_{P^{M(U)}}`$:
$$LM(P^{M(U)})=\underset{i=0}{\overset{\mathrm{}}{}}T_{P^{M(U)}}^i(\mathrm{}).$$
We will first show that for every integer $`i`$, $`0ik+1`$, we have: $`T_{P^{M(U)}}^i(\mathrm{})M(U)`$ and for every $`qT_{P^{M(U)}}^i(\mathrm{})`$, $`i_qi`$.
Clearly, $`T_{P^{M(U)}}^0(\mathrm{})=\mathrm{}M(U)`$. Hence, the basis for the induction is established. Assume that for some $`i`$, $`0ik`$, $`T_{P^{M(U)}}^i(\mathrm{})M(U)`$ and that for every $`qT_{P^{M(U)}}^i(\mathrm{})`$, $`i_qi`$. Consider $`qT_{P^{M(U)}}^{i+1}(\mathrm{})`$. If $`Uc^{}(q,i+1)`$, then $`c(q,v)U`$ for some $`v`$, $`1vi`$. Since $`UF_2(q)`$, $`c(q)U`$ and $`qM(U)`$. By Lemma 5.1, it follows that $`i_q=v`$. Hence, $`i_q<i+1`$.
Thus, assume that $`U\neg c^{}(q,i+1)`$. Since $`qT_{P^{M(U)}}^{i+1}(\mathrm{})`$, there is a rule
$$r=qa_1,\mathrm{},a_s,\text{not}(b_1),\mathrm{},\text{not}(b_t)$$
in $`P`$ such that $`b_jM(U)`$, for every $`j`$, $`1jt`$, and $`a_jT_{P^{M(U)}}^i(\mathrm{})`$, $`1is`$. By the induction hypothesis, for every $`j`$, $`1js`$, we have $`a_jM(U)`$ and $`i_{a_j}i`$. It follows that $`UF_3(r,i+1)`$ and, consequently, that $`c(q,i+1)U`$. Since $`UF_2(q)`$, $`c(q)U`$ and $`qM(U)`$. It also follows (Lemma 5.1) that $`i_q=i+1`$.
Thus, we proved that $`_{i=0}^{k+1}T_{P^{M(U)}}^i(\mathrm{})M(U)`$. Since $`|M(U)|k`$, there is $`j`$, $`0jk`$ such that $`T_{P^{M(U)}}^j(\mathrm{})=T_{P^{M(U)}}^{j+1}(\mathrm{})`$. It follows that for every $`j^{}`$, $`j<j^{}`$, $`T_{P^{M(U)}}^j(\mathrm{})=T_{P^{M(U)}}^j^{}(\mathrm{})`$. Consequently, $`T_{P^{M(U)}}^i(\mathrm{})M(U)`$ for every non-negative integer $`i`$. $`\mathrm{}`$
Consider now a stable model $`M`$ of the program $`P`$ and assume that $`|M|k`$. Clearly, $`M=_{i=1}^{\mathrm{}}T_{P^M}^i(\mathrm{})`$. For each atom $`qM`$ define $`s_q`$ to be the least integer $`s`$ such that $`qT_{P^M}^s(\mathrm{})`$. Clearly, $`s_q1`$. Moreover, since $`|M|k`$, it follows that for each $`qM`$, $`s_qk`$. Now, define
$$U_M=\{c(q),c(q,s_q):qM\}\{c^{}(q,i):qM,s_q<ik+1\}$$
###### Lemma 5.4
Let $`M`$ be a stable model of a logic program $`P`$ such that $`|M|k`$. Under the terminology introduced above, the set of atoms $`U_M`$ is a model of $`T_0(P)`$.
Proof: Clearly, $`U_MF_1(q,i)`$ for $`q\text{At}(P)`$ and $`2ik+1`$, and $`U_MF_2(q)`$ for $`q\text{At}(P)`$.
We will now show that $`U_MF_4(q,i)`$, for $`q\text{At}(P)`$ and $`i=1,2,\mathrm{},k+1`$. First, we will consider the case $`qM`$. There are three subcases here depending on the value of $`i`$.
We start with $`i`$ such that $`s_q<ik+1`$. Then $`U_M\vDash ̸\neg c^{}(q,i)`$. It follows that $`U_M\vDash ̸F_3(r,i)`$ for every rule $`rP`$ such that $`h(r)=q`$. Since $`U_M\vDash ̸c(q,i)`$, $`U_MF_4(q,i)`$.
Next, we assume that $`i=s_q`$. Then, there is a rule $`r=qa_1,\mathrm{},a_s,\text{not}(b_1),\mathrm{},`$ $`\text{not}(b_t)`$ in $`P`$ such that $`b_jM`$, for every $`j`$, $`1jt`$, and $`a_jT_{P^M}^{i1}(\mathrm{})`$, $`1js`$. Clearly, $`U_MF_3(r,i)`$. Since $`U_Mc(q,i)`$, it follows that $`U_MF_4(q,i)`$, for $`i=s_q`$.
Finally, let us consider the case $`1i<s_q`$. Assume that there is rule $`rP`$ such that $`h(r)=q`$ and $`U_MF_3(r,i)`$. Let us assume that $`r=qa_1,\mathrm{},a_s,\text{not}(b_1),\mathrm{},`$ $`\text{not}(b_t)`$. It follows that for every $`j`$, $`1jt`$, $`U_M\neg c(b_j)`$. Consequently, for every $`j`$, $`1jt`$, $`b_jM`$ and the rule $`r^{}=qa_1,\mathrm{},a_s`$ belongs to the reduct $`P^M`$. In addition, for every $`j`$, $`1js`$, $`c^{}(a_j,i)U_M`$. Thus, $`a_jM`$ and $`s_{a_j}i1`$. This latter property is equivalent to $`a_jT_{P^M}^{i1}(\mathrm{})`$. Thus, it follows that $`qT_{P^M}^i(\mathrm{})`$ and $`s_qi`$ — a contradiction with the assumption that $`i<s_q`$. Hence, for every rule $`r`$ with the head $`q`$, $`U_M\vDash ̸F_3(r,i)`$. Since for $`i<s_q`$, $`c(q,i)U_M`$, $`U_MF_4(q,i)`$.
To complete the proof, we still need to consider the case $`qM`$. Clearly, for every $`i`$, $`1ik+1`$, $`U_M\vDash ̸c(q,i)`$. Assume that there is $`i`$, $`1ik+1`$, and a rule $`r`$ such that $`h(r)=q`$ and $`U_MF_3(r,i)`$. Let us assume that $`r`$ is of the form $`qa_1,\mathrm{},a_s,\text{not}(b_1),\mathrm{},\text{not}(b_t)`$. It follows that $`c^{}(a_j,i)U_M`$ and, consequently, $`a_jM`$ for every $`j`$, $`1js`$. In addition, it follows that for every $`j`$, $`1jt`$, $`U_M\neg c(b_j)`$ and, consequently, $`b_jM`$. Thus, $`qa_1,\mathrm{},a_s`$ belongs to the reduct $`P^M`$ and, since $`M`$ is a model of the reduct, $`qM`$, a contradiction. It follows that for every $`i`$, $`1ik+1`$, $`U_MF_4(q,i)`$. $`\mathrm{}`$
For each atom $`q\text{At}(P)`$, let us introduce $`k^2+2k`$ new atoms $`d(q,i)`$, $`1ik^2+2k`$, and define
$$T(P)=T_0(P)\{c(q)d(q,i):1ik^2+2k\}.$$
Lemmas 5.1 \- 5.4 add up to a proof of the following result.
###### Theorem 5.5
Let $`k`$ be a non-negative integer and let $`P`$ be a logic program. The program $`P`$ has a stable model of size at most $`k`$ if and only if the theory $`T(P)`$ has a model of size at most $`(k+1)(k^2+2k)`$.
Proof: $`()`$ Let $`M`$ be a stable model of $`P`$ such that $`|M|k`$. By Lemma 5.4, the set $`U_M`$ is a model of $`T_0(P)`$ Consequently, the set
$$U=U_M\{d(q,i):qM,1ik^2+2k\}$$
is a model of $`T(P)`$. Moreover, it is easy to see that $`|U_M|2k+k^2`$. Hence, $`|U|2k+k^2+k(k^2+2k)=(k+1)(k^2+2k)`$.
Conversely, let us assume that some set $`V`$, consisting of atoms appearing in $`T(P)`$ and such that $`|V|(k+1)(k^2+2k)`$, is a model of $`T(P)`$. Let us define $`U`$ to consist of all atoms of the form $`c(q)`$, $`c(q,i)`$ and $`c^{}(q,i)`$ that appear in $`V`$. Clearly, $`U`$ is a model of $`T_0(P)`$. Let us assume that $`M(U)k+1`$ (we recall that the notation $`M(U)`$ was introduced just before Lemma 5.1 was stated). Then, there are at least $`(k+1)(k^2+2k)`$ atoms of type $`d(q,i)`$ in $`V`$. Consequently, $`V>(k+1)(k^2+2k)`$ as it contains also at least $`k+1`$ atoms $`c(q)`$, where $`qM(U)`$. This is a contradiction. Thus, it follows that $`|M(U)|k`$. Moreover, by Lemma 5.3, $`M(U)`$ is a stable model of $`P`$. $`\mathrm{}`$
Let us now define the following sets of formulas. First, for each atom $`q\text{At}(P)`$ we define
$$C_0(q)=\{\neg c(q)d(q,i):1ik^2+2k\}\{c(q)\neg d(q,i):1ik^2+2k\}.$$
Next, we define
$$C_1(q,i)=\{\neg c^{}(q,i)c(q,1)\mathrm{}c(q,i1)\}\{\neg c(q,j)c^{}(q,i):1ji1\},$$
$$C_2(q)=\{\neg c(q)c(q,1)\mathrm{}c(q,k+1)\}\{\neg c(q,j)c(q):1jk+1\},$$
and
$$C_4(q,i)=\{\neg c(q,i)F_3(r_1,i)\mathrm{}F_3(r_v,i)\}\{\neg F_3(r_j,i)c(q,i):1jv\},$$
where $`\{r_1,\mathrm{},r_v\}`$ is the set of all rules in $`P`$ with $`q`$ in the head.
Clearly, the theory
$`T^c(P)`$ $`=`$ $`\{C_0(q):qAt(P)\}\{C_1(q,i):qAt(P),2ik+1\}`$
$`\{C_2(q):qAt(P)\}\{C_4(q,i):qAt(P),1ik+1\}`$
is equivalent to the theory $`T(P)`$. Moreover, it is a collection of sums of products of literals. Therefore, it is a 3-normalized formula. By Theorem 5.5, it follows that the problem SSM can be reduced to the problem $`WS^{}(3)`$. Thus, we get the following result.
###### Theorem 5.6
The problem $`\text{SSM}(k)W[3]`$.
Next, we will show that the problem $`WS^{}(2)`$ can be reduced to the problem SSM. Let $`C=\{c_1,\mathrm{},c_p\}`$ be a collection of clauses. Let $`A=\{x_1,\mathrm{},x_r\}`$ be the set of atoms appearing in clauses in $`C`$. For each atom $`xA`$, introduce $`k`$ new atoms $`x(i)`$, $`1ik`$. By $`S_i`$, $`1ik`$, we denote the logic program consisting of the following $`n`$ clauses:
$`x_1(i)\text{not}(x_2(i)),\mathrm{},\text{not}(x_r(i))`$
$`\mathrm{}`$
$`x_r(i)\text{not}(x_1(i)),\mathrm{},\text{not}(x_{r1}(i))`$
Define $`S=_{i=1}^kS_i`$. Clearly, each stable model of $`S`$ is of the form $`\{x_{j_1}(1),\mathrm{},x_{j_k}(k)\}`$, where $`1j_pr`$ for $`p=1,\mathrm{},k`$. Sets of this form can be viewed as representations of nonempty subsets of the set $`A`$ that have no more than $`k`$ elements. This representation is not one-to-one, that is, some subsets have multiple representations.
Next, define $`P_1`$ to be the program consisting of the clauses
$$x_jx_j(i),j=1,\mathrm{},r,i=1,2,\mathrm{},k.$$
Stable models of the program $`SP_1`$ are of the form $`\{x_{j_1}(1),\mathrm{},x_{j_k}(k)\}M`$, where $`M`$ is a nonempty subset of $`A`$ such that $`|M|k`$ and $`x_{j_1},\mathrm{},x_{j_k}`$ enumerate (possibly with repetitions) all elements of $`M`$.
Finally, for each clause
$$c=a_1\mathrm{}a_s\neg b_1\mathrm{}\neg b_t$$
from $`C`$ define a logic program clause $`p(c)`$:
$$p(c)=fb_1,\mathrm{},b_t,\text{not}(a_1),\mathrm{},\text{not}(a_s),\text{not}(f)$$
where $`f`$ is yet another new atom. Define $`P_2=\{p(c):cC\}`$ and $`P^C=SP_1P_2`$.
###### Theorem 5.7
A set of clauses $`C`$ has a nonempty model with no more than $`k`$ elements if and only if the program $`P^C`$ has a stable model with no more than $`2k`$ elements.
Proof: Let $`M`$ be a nonempty model of $`C`$ such that $`|M|k`$. Let $`x_{j_1},\mathrm{},x_{j_k}`$ be an enumeration of all elements of $`M`$ (possibly with repetitions). Then the set $`M^{}=\{x_{j_1}(1),\mathrm{},x_{j_k}(k)\}M`$ is a stable model of the program $`SP_1`$. Since $`M`$ is a model of $`C`$, it follows that $`(P^C)^M^{}=(SP_1)^M^{}F`$, where $`F`$ consists of the clauses of the form
$$fb_1,\mathrm{},b_t,$$
such that $`t1`$ and for some $`j`$, $`1jt`$, $`b_jM^{}`$. Since $`M^{}=LM((SP_1)^M^{})`$, it follows that
$$M^{}=LM((SP_1)^M^{}F)=LM((P^C)^M^{}).$$
Thus, $`M^{}`$ is a stable model of $`P^C`$. Since $`|M^{}|2k`$, the “only if” part of the assertion follows.
Conversely, assume that $`M^{}`$ is a stable model of $`P^C`$. Clearly, $`fM^{}`$. Consequently,
$$LM((SP_1)^M^{})=LM((SP_1P_2)^M^{})=LM((P^C)^M^{})=M^{}.$$
That is, $`M^{}`$ is a stable model of $`SP_1`$. As mentioned earlier, it follows that $`M^{}=\{x_{j_1}(1),\mathrm{},x_{j_k}(k)\}M`$, where $`M`$ is a nonempty subset of $`\text{At}(P)`$ such that $`|M|k`$ and $`x_{j_1},\mathrm{},x_{j_k}`$ is an enumeration of all elements of $`M`$.
Consider a clause $`c=a_1\mathrm{}a_s\neg b_1\mathrm{}\neg b_t`$ from $`C`$. Since $`M^{}`$ is a stable model of $`P^C`$, it is a model of $`P^C`$. In particular, $`M^{}`$ is a model of $`p(c)`$. Since $`fM^{}`$, it follows that $`M^{}c`$ and, consequently, $`Mc`$. Hence, $`M`$ is a model of $`C`$. $`\mathrm{}`$
Now the reducibility of the problem $`WS^{}(2)`$ to the problem SSM is evident. Given a collection of clauses $`C`$, to check whether it has a model of size at most $`k`$, we first check whether the empty set of atoms is a model of $`C`$. If so, we return the answer YES and terminate the algorithm. Otherwise, we construct the program $`P^C`$ and check whether it has a stable model of size at most $`2k`$. Consequently, we obtain the following result.
###### Theorem 5.8
The problem SSM is W-hard.
## 6 Open problems and conclusions
The paper established several results pertaining to the problem of computing small and large stable models. It also brings up interesting research questions.
First, we proved that the problem LSM is in the class FPT. For problems that are fixed-parameter tractable, it is often possible to design an algorithm running in time $`O(p(N)+f(k))`$, where $`N`$ is the size of the problem, $`k`$ is a parameter, $`p`$ is a polynomial and $`f`$ is a function \[Downey & Fellows, 1997\]. Such algorithms are often practical for quite large ranges of $`N`$ and $`k`$. The algorithm for the LSM problem presented in this paper runs in time $`O(m2^{k+k^2})`$. It seems plausible it can be improved to run in time $`O(m+f(k))`$, for some function $`f`$. Such an algorithm would most certainly be practical for wide range of values of $`m`$ and $`k`$. We propose as an open problem the challenge of designing an algorithm for computing large stable models with this time complexity.
There is a natural variation on the problem of computing large stable models: given a logic program $`P`$ and an integer $`k`$ (parameter), decide whether $`P`$ has a stable model of size at least $`|\text{At}(P)|k`$. This version of the problem LSM was recently proved by Zbigniew Lonc and the author to be W-hard (and, hence, fixed-parameter intractable) \[Lonc & Truszczyński, 2000\]. The upper bound for the complexity of this problem remains unknown.
In the paper, we described an algorithm that for every fixed $`k`$, decides the existence of stable models of size at most $`k`$ in time $`O(n^{k1}m)`$, where $`n`$ is the number of atoms in the program and $`m`$ is its size. This algorithm offers only a slight improvement over the straightforward “guess-and-check” algorithm. An interesting and, it seems, difficult problem is to significantly improve on this algorithm by lowering the exponent in the complexity estimate to $`\alpha k`$, for some constant $`\alpha <1`$.
We also studied the complexity of the problem SSM and showed that it is fixed-parameter intractable. Our results show that SSM is $`W[2]`$-hard. This result implies that the problem SSM is at least as hard as the problem to determine whether a CNF theory has a model of cardinality at most $`k`$, and strongly suggests that algorithms do not exist that would decide problems $`\text{SSM}(k)`$ and run in time $`O(n^c)`$, where $`c`$ is a constant independent on $`k`$. For the upper bound, we proved in this paper that the problem SSM belongs to class $`W[3]`$. Recently, Zbigniew Lonc and the author \[Lonc & Truszczyński, 2000\] showed that the problem SSM is, in fact, in the class $`W[2]`$.
## Acknowledgments
The author thanks Victor Marek and Jennifer Seitzer for useful discussions and comments. The author is grateful to anonymous referees for very careful reading of the manuscript. Their comments helped eliminate some inaccuracies and improve the presentation of the results. This research was supported by the NSF grants CDA-9502645, IRI-9619233 and EPS-9874764. |
warning/0002/cond-mat0002268.html | ar5iv | text | # Revisiting ”swings” in the crossover features of Ising thin films near 𝑇_𝑐(𝐷)
## Abstract
”Swing” effects at the onset of crossover towards two dimensional behavior in thin Ising films are investigated close to $`T_c(D)`$ by means of Monte Carlo calculations. We find that the effect is extremely large for the specific heat effective critical exponent, in comparison with the ”swing” already noted by Capehart and Fisher for the susceptibility. These effects change considerably the system’s evolution with thickness $`(D)`$ from two-dimensional to three-dimensional behavior, forcing the effective exponents to pass near characteristic Tri Critical Point (TCP) values.
Basic features of phase transitions in systems with thin film geometry have been connected with the problem of the crossover from classical to quantum transitions. The change from classical to quantum character of the transition can be mapped to the evolution with thickness of the phase transition in films (se f.i.).That is the reason why a detailed study of phase transitions in films may be particularly useful for the study of quantum phase transitions apart from the intrinsic usefulness of studying changes in systems with a few layers of thickness. The effective critical exponents and the evolution near the critical point has been extensively studied by means of series expansions , the renormalization group , and Monte Carlo calculations in Ising systems , as well as in the X-Y model . For systems with thin film geometry, the correlation length is much smaller than the film thickness $`(D),`$ sufficiently below and above the critical point (i.e. relatively far from $`T_c(D)`$). Once the correlation length grows sufficiently (i.e. close to $`T_c(D)`$) the system notices that its critical behavior cannot be that of a three-dimensional system and the crossover to the two-dimensional behavior begins. From the point of view of the effective critical exponents this means that the system is initially evolving towards three-dimensional behavior until a crossover to two-dimensional behavior takes place. The film thickness can be characterized by the value of the effective critical exponents just at the onset of this crossover.
The pioneering work of Capehart and Fisher noted that for the case of the effective critical exponent corresponding to the susceptibility ($`\gamma _{eff}`$) an ”under-swing” behavior was apparent just before the crossover. This characteristic behavior means that for a certain thickness $`D^{}`$ the effective critical exponent reaches a minimum with a value $`\gamma _m(D^{})<\gamma ^{3D}<\gamma ^{2D}`$. This kind of behavior was attributed to surface effects, due to the lower value of the thickness ($`D<<L`$). In principle one might expect to find an enhancement of the pehenomenon using free boundary conditions in comparison with periodic boundary conditions, as indeed it was seen the case.
Since that time there has not been much work on the problem considered because research has been devoted strictly to the very close vicinity of the critical point. Monte Carlo simulations have shown the existence of this ”under-swing”, but no attempt has been made to characterize this phenomenon. In principle this ”under-swing”is a small effect, since $`\gamma _m(D^{})`$ is close to $`\gamma ^{3D}`$, but several very interesting questions can be asked concerning this phenomenon: a) We know that there should be a value of the thickness $`D`$ for which this effect should be maximum, $`D^{},`$ because eventually $`\gamma _{eff}`$ must increase again as $`(DL)`$ towards $`\gamma ^{3D}`$: What is the value $`D^{}`$ of this characteristic thickness? b) Is it possible to get more pronounced ”swing” effects in other critical exponents?, c) What are the values of these effective critical exponents for $`D^{}`$ corresponding to the maximum ”swing”? d) Is there a substantial difference between the exponent values obtained using periodic and free boundary conditions?
In the present work we will address these questions studying the thickness dependence of the effective critical exponents ($`\beta _{eff},`$ $`\gamma _{eff},`$ $`\delta _{eff},`$ $`\alpha _{eff}`$) of Ising film ($`L\times L\times D`$), describing the evolution from the pure two-dimensional Ising system ($`D=1`$) towards the three-dimensional system ($`D=L`$) system. In order to obtain the actual behavior of the effective critical exponents we will make use of the fact that the scaling relations hold all the way before and all through the crossover region .
In the present work we report results on phase transitions in Ising plates of equal area ($`L=100`$) and different thickness ($`D=2,3,4,5,7,9,12)`$ by Monte Carlo calculations. In order to reduce the critical slowing down effect near the critical point we use the Wolff single cluster algorithm , with more than 50.000 MCS . To ensure equilibrium we start our calculations with an early thermalization ($`T0K`$ and $`H=0`$) and we increase the temperature in very small (non-constant) steps as we get closer and closer to the critical point. These very small temperature steps give rise to large fluctuations in the numerical derivatives. We did smooth the data taking derivatives including up to the 5<sup>th</sup> nearest neighboring points. We obtain the evolution of the effective critical exponents $`\beta _{eff}`$ and $`\gamma _{eff}`$ by a direct determination of the magnetization $`M(T)`$ and of the susceptibility $`\chi (T)=M^2M^2`$ using the standard relations:
$$\beta _{eff}=\frac{logM(T)}{log[T_c(D)T]},\gamma _{eff}=\frac{log\chi (T)}{log[T_c(D)T]}$$
(1)
The critical temperature $`T_c(D)`$ corresponding to each particular thickness $`D`$ is obtained in the usual way by means of the Binder cumulant method (see f.i. ).
Every calculation has been performed using free and periodic boundary conditions. This comparison is important, because real films, due to substrate effects, are not pure free-surface systems but have a mixture of free and constrained surfaces. A direct comparison has been performed elsewhere for the D dependence of the critical temperature. Here we will carry out this comparison explicitly for the behavior of the effective critical exponents.
We present in Fig.1a $`\beta _{eff}`$ vs. $`log[T_c(D)T]`$ for $`D=3,5,9`$. Three zones are clearly visible: (i) initial evolution towards the three-dimensional value, (ii) crossover zone towards the two-dimensional value, and (iii) finite size effects zone. Note that for $`D=3`$ there is nearly no crossover, since the system is still almost two-dimensional, and as $`D`$ increases, the maximum effective critical exponent, defined just before the crossover starts, $`\beta _m(D)`$, grows tending towards the three-dimensional value ($`\beta ^{3D}0.33`$). We remark the clear difference between exponents with periodic and free boundary conditions. Note how the exponent for free boundary conditions always rises (for the same thickness) to a maximum which is closer to the corresponding three-dimensional value than the same effective critical exponent for periodic boundary conditions. Finally both (free and periodic exponents) collapse together in the crossover. This exponent, as will be seen also for $`1/\delta _{eff}`$ below, does not present any ”swing” effect. It means that we do not find any value of $`D`$ for which $`\beta ^{2D}<\beta ^{3D}<\beta _m(D).`$
The ”under-swing” effect noted by Capehart and Fisher is explicit for the case of the susceptibility. In order to check this effect, we present in Fig 1b results for $`\gamma _{eff}`$ vs. $`log[T_c(D)T]`$ for $`D=3,5,9.`$ Note how the ”under-swing” is clearly detectable for values of $`D`$ close to $`D=9.`$ This ”under-swing” is visible not just for the free boundary conditions, but also for the periodic boundary conditions as was pointed out in Ref. . This is the first time in our knowledge that the ”under-swing” effect is explicitly shown to exist under periodic boundary conditions, where surface effects are reduced.
In order to get a more complete picture of the dependence with thickness of the effective critical exponents we study also the effective critical exponents $`1/\delta _m(D)`$ and $`\alpha _m(D)`$. As usual, $`1/\delta _m(D)`$ should be obtained as the value just before $`1/\delta _{eff}`$ starts the crossover to the two-dimensional value, $`1/\delta _{eff}`$ may be obtained making use of the scaling relation
$$1/\delta _{eff}=\left(\frac{\gamma _{eff}}{\beta _{eff}}+1\right)^1$$
(2)
As mentioned above, this relationship has been proven to hold, not just near the critical point, but also at the crossover region, and before . The results for $`1/\delta _{eff}`$ vs. $`log[T_c(D)T]`$ are presented in Fig.1c. The behavior is very similar to the one observed for $`\beta _{eff}`$. There is no ”over-swing”. Thus we do not have any value of $`D`$ for which $`1/\delta ^{2D}<1/\delta ^{3D}<1/\delta _m(D).`$
The other interesting exponent to study is the effective specific heat critical exponent. We are able to determine explicitly the evolution of this effective critical exponent by means of the relation :
$$\alpha _{eff}=22\beta _{eff}\gamma _{eff}$$
(3)
Fig 1d presents the results for $`\alpha _{eff}`$ vs. $`log[T_c(D)T].`$ We find and extremely enhanced ”over-swing”. We find clearly that $`\alpha ^{2D}<\alpha ^{3D}<\alpha _m(D)`$, not just for $`D=9`$ but also for $`D=5`$ and $`D=3`$. Another interesting result is that the effect is very clearly visible both for free and for periodic boundary conditions.
We may note that in Fig1c and 1d the final data for $`TT_c(D)`$ are not representative, because they correspond to the finite size effects region of $`\beta _{eff}`$ and $`\gamma _{eff}`$.
The best way to show the ”swing” effect is perhaps to plot the values obtained for $`\alpha _m(D)`$ vs. $`D`$ and for $`\gamma _m(D)`$ vs. $`D`$, together with the results obtained for $`\beta _m(D)`$ and $`\delta _m(D).`$ These results are presented in Fig.2. Several points are made clear: a) ”Swing” effects exist only for the specific heat $`(\alpha _{eff})`$ and the susceptibility $`(\gamma _{eff})`$ effective critical exponents b) The ”swing” is enhanced clearly in the effective heat exponent, f.i. we get a ratio $`\alpha _m(D=9)/\alpha ^{3D}`$ $`5`$ while for $`\gamma _m(D)`$ we just find $`\gamma ^{3D}/\gamma _m(D=9)1.24.`$ This means that the effect that is relatively small in the susceptibility can not be ignored in the specific heat c) The maximum ”swing” effect is found for values of $`D`$ close to $`10`$, that is, we should take $`D^{}10.`$ d) ”Swings” appears for both, free and periodic boundary conditions, and are more pronounced for the former.
Now we focus attention on the exponent values obtained for $`D=9D^{}`$. The results for periodic and free boundary conditions are presented in Table I. They are compared with the two-dimensional values, the three-dimensional values and with the Tri Critical Point (TCP) values. As it is known, a Tri Critical Point is at the limit separating continuous (2<sup>nd</sup> order) from discontinuous (1<sup>st</sup> order) transitions . Note that the exponent values corresponding to the Tri Critical Point are close to those for $`D=9D^{},`$ with errors ranging from three to twelve percent. Clearly, the evolution of the effective critical exponents \[$`\beta _m(D),\alpha _m(D),1/\delta _m(D)`$ and $`\gamma _m(D)`$\] from the two-dimensional values ($`D=1`$) to the three-dimensional values ($`D=L)`$ is not monotonous but appears in all cases to come close to the respective Tri Critical Point value for $`DD^{}`$. This effect is made more explicit in plots of $`\alpha _m(D)`$ vs. $`\gamma _m(D)`$ and $`1/\delta _m(D)`$ vs. $`\beta _m(D)`$ (see Fig.3a and 3b).
Fig.3a shows that, in the case of $`\alpha _m(D)`$ vs. $`\gamma _m(D),`$the evolution of the plot from $`D<<D^{}`$ $`[\gamma _m(D)\gamma ^{2D},\alpha _m(D)\alpha ^{2D}]`$ onwards shows an spectacular turn towards the Tri Critical Point pair of values $`(\gamma ^{TCP},\alpha ^{TCP}).`$ Then for $`D>D^{}`$ the evolution towards the pure three-dimensional values begins. Note that the data points get away from the box defined by $`(\gamma ^{2D},\alpha ^{2D})(\gamma ^{3D},\alpha ^{3D}),`$ making explicit the existence of ”swing effects”. An interesting feature of our results is that the general behavior appears to follow the same well defined line independently of the boundary conditions used.
For the case of $`1/\delta _m(D)`$ vs. $`\beta _m(D)`$ the values corresponding to a Tri Critical Point are also closely approached. The non-existence of swing effects in this case is also explicit since the pair of values $`(\beta _m(D),1/\delta _m(D))`$ do not leave the box.
In conclusion we have presented Monte Carlo data for the evolution of effective critical exponents (note that these are ”transient” exponents, not assyntotic, critical exponents) in thin Ising films. In summary we have shown that: a) ”Swing effects” are specially enhanced for the specific heat effective exponent, $`\alpha _m(D)`$, b) They appear very clearly for both free and periodic boundary conditions and c) ”Swing effects” force the effective exponents to pass near exponent values corresponding to a Tri Critical Point (for $`D^{}10`$) well before the evolution towards the three-dimensional values begins.
Our work shows that ”swing” effects must become patent especially in the case of the specific heat for any boundary conditions. It would be very interesting to check this point experimentally. This result rises also the basic question of why Tri Critical Point exponents $`(\beta =1/4,1/\delta =1/5,\gamma =1,\alpha =1/2)`$ describe so well the behavior of thin films at the onset of the crossover, for characteristic thicknesses of $`D^{}10`$.
We thank P.A.Serena for computing facilities and we acknowledge financial support from DGCyT through grant PB96-0037. |
warning/0002/hep-th0002102.html | ar5iv | text | # Contents
## 1 Introduction
One of the outstanding issues in both string theory and phenomenology is the choice of vacuum. Recent dramatic advances in the non-perturbative understanding of strings have demonstrated that all string theories, thought previously to be distinct, are in fact related by various dualities, and can be regarded as different phases of a single underlying theory, called variously $`M`$ and/or $`F`$ theory . This deeper non-perturbative understanding does not alter the fact that many classical string vacua appear equally consistent at the perturbative level. However, the new non-perturbative methods may provide us with new tools to understand transitions between these classical vacua, and perhaps eventually provide a dynamical criterion for deciding which vacuum is preferred physically .
Consistent string vacua are constrained by the principles of quantum mechanics applied to extended objects. At the classical level, these are expressed in the conformal symmetry of the supersymmetric world-sheet field theory. Consistent quantization of the string must confront a possible anomaly in conformal symmetry, as manifested in a net non-zero central charge of the Virasoro algebra. Early studies of the quantum mechanics of extended objects indicated that strings could not survive in the familiar dimension $`D=3+1`$ of our space-time. The way initially used to cancel the conformal anomaly was to choose appropriately the dimension of the ambient space-time, for example, $`D=25+1`$ for bosonic strings and $`D=9+1`$ for the supersymmetric and heterotic strings.
This suggested that the surplus $`n=6`$ real dimensions should be compactified. The simplest possibility is on a Calabi-Yau manifold , which is defined by the following conditions:
* It has a complex structure, with $`N=3`$ complex dimensions required for the $`D=9+13+1`$ case of most direct interest, though all the cases $`N=1,2,3,4,..`$ have some interest.
* It is compact.
* It has a Kähler structure.
* It has holonomy group $`SU(n)`$ or $`Sp(n)`$, e.g., $`SU(3)`$ in the $`N=3`$ case.
It has subsequently been realized that one could compactify on an orbifold , rather than a manifold, and also that generalized heterotic strings could be formulated directly in $`D=3+1`$ dimensions, with extra world-sheet degrees of freedom replacing the surplus space coordinates. More recently, the non-perturbative formulation of the theory in eleven or twelve dimensions, as $`M`$ or $`F`$ theory, has opened up new possibilities . However, Calabi-Yau compactifications continue to play a key role in the search for realistic four-dimensional string models, motivating us to revisit their classification.
One of the most important tools in the investigation of such complex manifolds is the feature that their singularities are connected with the structure of Lie algebras. Kaluza was the first to attempt to understand this circumstance, and used this idea to embark on the unification of all the gauge interactions known at that time, namely electromagnetism and gravitation. These ideas were subsequently extended to non-Abelian gauge theories, and string theory can be regarded as the latest stage in the evolution of this programme.
The three-complex-dimensional CY manifolds can be situated in a sequence of complex spaces of increasing dimensions: two-real- (one-complex-)dimensional tori $`T_2`$, the two-complex-dimensional $`K3`$ spaces, the three-complex-dimensional $`CY_3`$ themselves, four-complex-dimensional $`CY_4`$, etc., whose topological structure and classification become progressively more complicated. Their topologies may be described by the Betti-Hodge numbers which count the numbers of distinct one-, two-, three-dimensional, … cycles (holes,…). The topological data of the different CY manifolds determine their physical properties, such as the different numbers of generations $`N_g`$ (which are related to the Euler characteristics of $`CY_3`$ spaces), etc.. This emphasizes the desirability of approaching systematically the problem of their classification and the relations between, e.g., $`CY_3`$ manifolds with different values of the Euler characteristic and hence the number of generations $`N_g`$. Since some non-perturbative tools now exist for studying transitions between different CY manifolds, one could hope eventually to find some dynamical criterion for determining $`N_g`$.
The topologies and classification of the lower-dimensional spaces in this sequence are better known: although our ultimate objective is deeper understanding of $`CY_3`$ spaces, in this paper we study as a warm-up problem the simpler case of the two-complex-dimensional $`K3`$ hypersurfaces. These are of considerable interest in their own right, since, for example, they may appear as fibrations of higher-dimensional $`CY_n`$ spaces. It is well known that any two $`K3`$ spaces are diffeomorphic to each other. This can be seen, for example, by using the polyhedron techniques of Batyrev discussed in Sections 2 and 3, to calculate the Betti-Hodge invariants for all the $`K3`$ hypersurfaces corresponding to the $`\stackrel{}{k}_4`$ vectors we find. It is easy to check that Batyrev’s results yield the same Euler number 24 for all $`K3`$ manifolds .
The quasi-homogeneous polynomial equations (hereafter called CY equations) whose zeroes define the CY spaces as hypersurfaces in complex projective space are defined (2.6, 2.7, 2.8, 2.9) by projective vectors $`\stackrel{}{k}`$, whose components specify the exponents of the polynomials. The number of CY manifolds is large but finite, as follows from the property of reflexivity introduced in Section 2. The central problem in the understanding of classifying these manifolds may be expressed as that of understanding the set of possible projective vectors $`\stackrel{}{k}=(k_1,\mathrm{},k_{n+1})`$, the corresponding Lie algebras and their representations. More precisely, the classification of all CY manifolds contains the following problems:
$``$ To study the structure of the $`K3,CY_3,\mathrm{}`$ projective vectors $`\stackrel{}{k_n}`$, in particular, to find the links with the projective vectors of lower dimensions: $`D=n1,n2,\mathrm{}`$.
$``$ To establish the web of connections between all the projective vectors $`\stackrel{}{k_n}`$ of the same dimension.
$``$ To find an algebraic description of the geometrical structure for all projective vectors, and calculate the corresponding Betti-Hodge invariants.
$``$ To establish the connections between the projective vectors $`\stackrel{}{k_n}`$, the singularities of the corresponding CY hypersurfaces, the gauge groups and their matter representations, such as the number of generations, $`N_g`$.
$``$ To study the duality symmetries and hypermodular transformations of the projective vectors $`\stackrel{}{k_n}`$.
In addition to the topological properties and gauge symmetries already mentioned, it is now well known that string vacua may be related by duality symmetries. This feature is familiar even from simple compactifications on $`S_1`$ spaces of radius $`R`$, which revealed a symmetry: $`R1/R`$ . In the case of compactifications on tori, there are known to be $`S,T`$, and $`U`$ dualities that interrelate five string theories and play key roles in the formulations of $`M`$ and $`F`$ theories . Compactifications on different types of CY manifolds have also been used extensively in verifying these string dualities . For example, in proving the duality between type-IIA and type-IIB string theories, essential use was made of the very important observation that all CY manifolds have mirror partners . Thus, duality in string theory found its origins in a duality of complex geometry.
Further information about string/$`M`$/$`F`$ theory and its compactifications on CY manifolds can be obtained using the methods of toric geometry. The set of homogeneous polynomials of degree $`d`$ in the complex projective space $`CP^n`$ defined by the vector $`\stackrel{}{k_{n+1}}`$ with $`d=k_1+\mathrm{}k_{n+1}`$ defines a convex reflexive polyhedron <sup>*</sup><sup>*</sup>*The notion of a reflexive polyhedron is introduced and defined in Section 2., whose intersection with the integer lattice corresponds to the polynomials of the CY equation. Therefore, instead of studying the complex hypersurfaces directly, one can study the geometry of polyhedrons. This method was first used to look for the solutions of the algebraic equations of degree five or more in terms of radicals . Thus, the problem of classifying CY hypersurfaces is also connected with the problem of solving high-degree polynomial equations in terms of radicals. The solutions of quintic- and higher-degree algebraic equations in terms of radicals may be expressed using elliptic and hyperelliptic functions, respectively. Specifically, it is known that CY manifolds may be represented using double-periodic elliptic or multi-periodic hyperelliptic functions . These functions have therefore been used to describe the behaviour of strings, and they should also be used to construct the ambient space-time in which strings move.
We embark here on a systematic classification of $`K3`$ manifolds, as a prelude to a subsequent classification of $`CY_3`$ manifolds, based on their construction in the framework of toric geometry. Within this approach, CY manifolds and their mirrors are toric varieties that can be associated with polyhedra in spaces of various dimensions. We propose here an inductive algebraic-geometric construction of the projective vectors $`\stackrel{}{k}`$ that define these polyhedra and the related $`K3`$ and CY spaces. This method has the potential to become exhaustive up to any desired complex dimensionality $`d=1,2,3,4,5,6,\mathrm{}`$ (see Figure 1), limited essentially by the available computer power. As a first step in this programme, we present in this article a construction of $`K3`$ spaces, which is complete for those described by simple polynomial zeroes, and in principle for $`K3`$ spaces obtained as the complete intersections of pairs or triples of such polynomial zero loci. In the construction of projective vectors corresponding to hypersurfaces without an intersection with one internal point, the duality between a complex manifold and its mirror (which does contain an intersection) plays an important role. We discuss here also aspects of the $`CY_3`$ construction that are relevant for the classification of $`K3`$ spaces. We also indicate already how one may generate $`CY_3`$ manifolds with elliptic fibrations or $`K3`$ fibers. More aspects of our $`CY_3`$ construction are left for later work.
To get the flavour of our construction, which is based on the formalism reviewed in Sections 2 and 3 , and is discussed in more detail in Sections 4 et seq., consider first $`CP^1`$ space. Starting from the trivial unit ‘vector’ $`\stackrel{}{k}_1(1)`$, we introduce two singly-extended basic vectors
$$\stackrel{}{k}_1^{ex^{}}=(0,1),\stackrel{}{k}_1^{ex^{\prime \prime }}=(1,0),$$
(1.1)
obtained by combining $`\stackrel{}{k}_1`$ with zero in the two possible ways. The basic vectors (1.1) correspond to the sets of polynomials
$`x^ny`$ $``$ $`\{\stackrel{}{\mu }_1\}=(n,\mathrm{\hspace{0.17em}1}):\stackrel{}{\mu }_1\stackrel{}{k}_1^{ex^{}}=d=\mathrm{\hspace{0.17em}1},`$
$`xy^m`$ $``$ $`\{\stackrel{}{\mu }_2\}=(1,m):\stackrel{}{\mu }_2\stackrel{}{k}_1^{ex^{\prime \prime }}=d=\mathrm{\hspace{0.17em}1},`$ (1.2)
respectively. The only polynomial common to these two sequences is $`xy`$, which may be considered as corresponding to the trivial ‘vector’ $`\stackrel{}{k}_1=(1)`$. Consider now the composite vector $`\stackrel{}{k}_2=(1,1)`$, which can be constructed out of the basic vectors (1.1), and is easily seen to correspond the following three monomials of two complex arguments $`(x,y)`$:
$`\{x^2,xy,y^2\}\stackrel{}{\mu }|_{i=1,2,3}`$ $`=`$ $`\{(2,0),(1,1),(0,2)\}`$
$`\stackrel{}{\mu }^{^{}}|_{i=1,2,3}\stackrel{}{\mu }|_{i=1,2,3}\stackrel{}{1}`$ $`=`$ $`\{(1,1),(0,0),(1,+1)\},`$ (1.3)
where we have used the condition: $`\stackrel{}{\mu }\stackrel{}{k}_2=\mu _11+\mu _21=d=2`$, corresponding to $`\stackrel{}{\mu }{}_{}{}^{}\stackrel{}{k}_2=0`$, and we denote by $`d`$ the dimensionality of the projective vectors. It is convenient to parametrize (1.3) in terms of the new basis vector $`\stackrel{}{e}=(1,1)`$:
$$\stackrel{}{\mu }^{^{}}|_{i=1,2,3}\stackrel{}{(}e)|_{i=1,2,3}=\{(1),(0),(+1)\}\times \stackrel{}{e}$$
(1.4)
The three points $`(2,0),(1,1),(0,2)`$ (or $`1,0,+1`$) corresponding to the composite vector $`k_2=(1,1)`$ may be considered as composing a degenerate linear polyhedron with two integer vertices $`\{(2,0),(0,2)\}`$ $`(\pm 1)`$ and one central interior point $`(1,1)`$ $`(0)`$. As we see in more detail later, this polyhedron is self-dual, or reflexive as defined in Section 2.
To describe $`CY_1`$ spaces in $`CP^2`$ projective space, via the analogous projective vectors $`\stackrel{}{k}_3=(1,1,1),(1,1,2),(1,2,3)`$ that are associated with the corresponding polynomial zero loci, one may introduce the two following types of extended vectors: the doubly-extended basic vectors
$$\stackrel{}{k}_1^{ex}=(0,0,1),(0,1,0),(1,0,0)$$
(1.5)
obtained by adding zero to the two-dimensional basic vectors (1.1) in all possible ways, and the three simple extensions of the composite vector $`\stackrel{}{k}_2=(1,1)`$:
$$\stackrel{}{k}_2^{ex}=(0,1,1),(1,0,1),(1,1,0).$$
(1.6)
Then, out of all the extended vectors (1.5) and (1.6) and the corresponding sets of monomials, one should consider only those pairs (triples) whose common monomials correspond to the composite vector $`\stackrel{}{k}_2=(1,1)`$ (to the unit vector) which produces the reflexive linear polyhedron with three integer points (a single point). The condition of reflexivity restricted to the extended vector pairs (triples), … will also be very important for constructing the closed sets of higher-dimensional projective vectors (again reflexive).
For example, consider one such ‘good’ pair,
$`\stackrel{}{k}_2^{ex}=(0,1,1)\stackrel{}{k}_1^{ex}=(1,0,0),`$ (1.7)
with the corresponding set of monomials,
$`\{x^my^2\}`$ $``$ $`\stackrel{}{\mu }=(m,2,0),`$
$`\{x^nyz\}`$ $``$ $`\stackrel{}{\mu }=(n,1,1),`$
$`\{x^pz^2\}`$ $``$ $`\stackrel{}{\mu }=(p,0,2),`$
$`\stackrel{}{\mu }_i\stackrel{}{k}_2^{ex}`$ $`=`$ $`\mathrm{\hspace{0.17em}2},`$ (1.8)
and
$`\{xy^kz^l\}`$ $``$ $`\stackrel{}{\mu }=(1,k,l),`$
$`\stackrel{}{\mu }\stackrel{}{k}_1^{ex}`$ $`=`$ $`1.`$ (1.9)
The common action of these two extended vectors, (0,1,1) and (1,0,0), gives as results only the following three monomials:
$`\{xy^2,xyz,xz^2\}`$
$`\stackrel{}{\mu }|_{i=1,2,3}`$ $`=`$ $`\{(1,2,0),(1,1,1),(1,0,2)\}`$
$`\stackrel{}{\mu }|_{i=1,2,3}\stackrel{}{1}`$ $`=`$ $`\{(0,1,1),(0,0,0),(0,1,1)\}`$
$`\stackrel{}{e}|_{i=1,2,3}`$ $`=`$ $`\{(1),(0),(1)\}`$ (1.10)
which correspond to the $`CP^1`$ case. Such pairs or triples may be termed ‘reflexive’ pairs or triples, because the vertices $`\stackrel{}{e}|_{i=1,2,3}`$ above generate a (degenerate) reflexive polyhedron.
Such pairs, triples and higher-order sets of projective vectors $`\stackrel{}{k}_1`$ may be used to define chains of integer-linear combinations, as explained in more detail in subsection 4.1:
$$m_1\stackrel{}{k}_1+m_2\stackrel{}{k}_2+\mathrm{}$$
(1.11)
We use the term eldest vector for the leading entry in any such chain, with minimal values of $`m_1,m_2,\mathrm{}`$. In the above case, there are just two distinct types of ‘reflexive’ pairs: $`\{(0,0,1),(1,1,0)\}`$ and $`\{(0,1,1),(1,0,1)\}`$, which give rise to two such chains: $`\{(1,1,1),(1,1,2)\}`$ and $`\{(1,1,2)`$, $`(1,2,3)\}`$. There is only one useful ‘reflexive’ triple: $`\{(0,0,1),(0,1,0),(1,0,0)\}`$ defining a non-trivial three-vector chain. Together, these chains can be used to construct all three projective $`\stackrel{}{k}_2`$ vectors. The second possible ‘reflexive’ triple $`\{(0,1,1),(1,0,1),`$ $`(1,1,0)\}`$ produces a chain that consists of only one projective $`\stackrel{}{k}_3`$ vector: $`(1,1,1)`$.
In addition to the zero loci of single polynomials, CY spaces may be found by higher-level contructions as the intersections of the zero loci of two or more polynomial loci. The higher-level $`CY_1`$ spaces found in this way are given in the last Section of this paper.
In the case of the $`K3`$ hypersurfaces in $`CP^3`$ projective space, our construction starts from the five possible types of extended vectors, with all their possible Galois groups of permutations. These types are the triply-extended basic vectors with the cyclic $`C_4`$ group of permutations,
$$\stackrel{}{k}_1^{ex}=(0,0,0,1):|C_4|=4,$$
(1.12)
the doubly-extended composite vectors with the $`D_3`$ dihedral group of permutations,
$$\stackrel{}{k}_2^{ex}=(0,0,1,1):|D_3|=6,$$
(1.13)
and the following singly-extended composite vectors with the cyclic $`C_4`$, alternating $`A_4`$ and symmetric $`S_4`$ groups of permutations, respectively:
$`\stackrel{}{k}_3^{ex}`$ $`=`$ $`(0,1,1,1):|C_4|=4,`$ (1.14)
$`\stackrel{}{k}_3^{ex}`$ $`=`$ $`(0,1,1,2):|A_4|=12,`$ (1.15)
$`\stackrel{}{k}_3^{ex}`$ $`=`$ $`(0,1,2,3):|S_4|=24.`$ (1.16)
The $`A_4`$ and $`S_4`$ groups of permutations can be identified with the tetrahedral $`T`$ and octahedral $`O`$ rotation groups, respectively. Combining these 50 extended vectors in pairs, we find 22 pairs whose common actions correspond to reflexive polyhedra in the plane. These give rise to 22 chains (lattices parametrized by two positive integers), which together yield 90 vectors $`\stackrel{}{k}_4`$ based on such extended structures, that are discussed in more detail in Section 5. In addition, there exist just four triples constructed from the 10 extended vectors $`(0,0,0,1)+`$ permutations and $`(0,0,1,1)+`$ permutations whose common actions give a unique reflexive polyhedron on the line: $`(1),(0),(+1)`$. The corresponding four triple chains (lattices parametrized by three positive integers) yield 91 $`\stackrel{}{k}_4`$ vectors, as discussed in Section 6. As also discussed there, it turns out that most of the vectors $`\stackrel{}{k}_4`$ obtained from the triple combinations are already included among those found in the double chains, so that the combined number of distinct vectors is just 94. The total number of vectors is, however, 95 (see Table 1), because there exists, in addition to the above enumeration, a single vector $`\stackrel{}{k}_4=(7,8,9,12)`$ which has only a trivial intersection consisting just of the zero point. This can be found within our approach using the non-trivial projection structure of its dual, which is an example of the importance of duality in our classification, as discussed in Section 7.
To find all CY manifolds, and thereby to close their algebra with respect the duality between intersection and projection that is described in more detail in Sections 3 and 4, one must consider how to classify the projective structures of CY manifolds. Some of the 22 chains are dual with respect to the ‘intersection-projection’ structure, but more analysis is required to close the CY algebra. As discussed in Section 7, it is useful for this purpose to look for so-called invariant directions. To find all such invariant directions in the case of $`K3`$ spaces, one should consider all triples selected from the following five extended vectors: $`(0,0,0,1),(0,0,1,1),(0,1,1,1),(0,1,1,2),(0,1,2,3)`$, and their possible permutations, whose intersections give the following five types of invariant directions defined by two monomials:
$`\stackrel{}{\pi }_1^\alpha `$ $`=`$ $`\{(1,1,1,1)(0,1,1,3)\},\alpha =\mathrm{\hspace{0.17em}1},\mathrm{\hspace{0.17em}2},`$
$`\stackrel{}{\pi }_2^\alpha `$ $`=`$ $`\{(1,1,1,1)(0,0,0,3)\},\alpha =\mathrm{\hspace{0.17em}1},\mathrm{\hspace{0.17em}2},\mathrm{\hspace{0.17em}3},\mathrm{\hspace{0.17em}4},`$
$`\stackrel{}{\pi }_3^\alpha `$ $`=`$ $`\{(1,1,1,1)(0,0,1,3)\},\alpha =\mathrm{\hspace{0.17em}1},\mathrm{\hspace{0.17em}2},\mathrm{\hspace{0.17em}3},\mathrm{\hspace{0.17em}4},`$
$`\stackrel{}{\pi }_4^\alpha `$ $`=`$ $`\{(1,1,1,1)(0,0,0,4)\},\alpha =\mathrm{\hspace{0.17em}1},\mathrm{\hspace{0.17em}2},\mathrm{\hspace{0.17em}3},\mathrm{\hspace{0.17em}4},`$
$`\stackrel{}{\pi }_5^\alpha `$ $`=`$ $`\{(1,1,1,1)(0,0,1,4)\},\alpha =\mathrm{\hspace{0.17em}1},`$
and the following three types of invariant directions defined by three monomials:
$`\stackrel{}{\pi }_6^\alpha `$ $`=`$ $`\{(0,2,1,1)(1,1,1,1)(2,0,1,1)\},\alpha =\mathrm{\hspace{0.17em}1},\mathrm{\hspace{0.17em}2},`$
$`\stackrel{}{\pi }_7^\alpha `$ $`=`$ $`\{(0,0,1,2)(1,1,1,1)(2,2,1,0)\},\alpha =\mathrm{\hspace{0.17em}1},\mathrm{\hspace{0.17em}2},\mathrm{\hspace{0.17em}3},\mathrm{\hspace{0.17em}4},`$
$`\stackrel{}{\pi }_8^\alpha `$ $`=`$ $`\{(0,0,0,2)(1,1,1,1)(2,2,2,0)\},\alpha =\mathrm{\hspace{0.17em}1},\mathrm{\hspace{0.17em}2},\mathrm{\hspace{0.17em}3},\mathrm{\hspace{0.17em}4},`$
respectively. Each double intersection of a pair of extended vectors from one of these triples gives the same ‘good’ planar polyhedron whose intersection with the plane integer lattice $`Z_2`$ has just one interior point.
By this method, one can classify the projective vectors by projections, finding 78 projective vectors which can be characterized by their invariant directions. Taking into account the projective vectors with intersection-projection duality that have already been found by the double-intersection method, one can recover all 95 $`K3`$ projective vectors, including the exceptional vector $`(7,8,9,12)`$ that was not found previously among the double and triple chains.
Section 8 of this paper contains a systematic description how various gauge groups emerge associated with singularities in our construction of $`K3`$ spaces . These are interesting because of their possible role in studies of $`F`$ theory. Since this may be regarded as a decompactification of type-$`IIA`$ string, understanding of duality between the heterotic string and type $`IIA`$ string in $`D=6`$ dimensions can be used to help understand the duality between the heterotic string on $`T^2`$ and $`F`$ theory on an elliptically-fibered $`K3`$ hypersurface . The gauge group is directly defined by the $`ADE`$ classification of the quotient singularities of hypersurfaces. The Cartan matrix of the Lie group in this case coincides up to a sign with the intersection matrix of the blown-down divisors. There are two different mechanisms leading to enhanced gauge groups on the $`F`$-theory side and on the heterotic side. On the $`F`$-theory side, the singularities of the CY hypersurface give rise to the gauge groups, but on the heterotic side the singularities can give an enhancement of the gauge group if ‘small’ instantons of the gauge bundle lie on these singularities . This question has been studied in terms of the numbers of instantons placed on a singularity of type $`G`$, where $`G`$ is a simply-laced group. Studies of groups associated with singularities of $`K3`$ spaces are also interesting because elliptic $`CY_n`$ ($`n=3,4`$) manifolds with $`K3`$ fibers can be considered to study $`F`$-theory dual compactifications of the $`E_8\times E_8`$ or $`SO(32)`$ string theory. To do this in toric geometry, it is possible to consider the $`K3`$ polyhedron fiber as a subpolyhedron of the $`CY_n`$ polyhedron, and the Dynkin diagrams of the gauge groups of the type-$`IIA`$ string ($`F`$-theory) compactifications on the corresponding threefold (fourfold) can then be seen precisely in the polyhedron of this $`K3`$ hypersurface. By extension, one could consider the case of an elliptic $`CY_4`$ with $`CY_3`$ fiber, where the last is a CY hypersurface with $`K3`$ fiber. We give in Section 8 several detailed examples of group structures associated with chains of $`K3`$ spaces, which our algebraic approach equips us to study systematically.
Finally, Section 9 provides a brief discussion of $`CY_3`$ manifolds and describes how additional CY spaces can be constructed at higher levels as the intersections of multiple polynomial loci. This discussion is illustrated by the examples of higher-level $`CY_1`$ and $`K3`$ spaces obtained via our construction of lower-level $`K3`$ and $`CY_3`$ spaces. We find, for example, 7 new polyhedra describing $`CY_1`$ spaces given by ‘level-one’ intersections of pairs of polynomial loci, and three new ‘level-two’ polyhedra given by triple intersections of polynomial loci. In looking for higher-level $`K3`$ spaces, we start from 100 types of extended vectors in five dimensions, corresponding to 10 270 distinct vectors when permutations are taken into account. We find that these give rise to 4242 two-vector chains of $`CY_3`$ spaces, 259 triple-vector chains and 6 quadruple-vector chains. Analyzing their internal structures, we we find 730 new $`K3`$ polyhedra at level one, of which 146 can be obtained as intersections of polynomials corresponding to simple polyhedra (points, line segments, triangles and tetrahedra). A complete characterization of higher-level $`K3`$ spaces given by multiple intersections of polynomial loci lies beyond our present computing scope, and we leave their further study to later work.
## 2 Calabi-Yau Spaces as Toric Varieties
We recall that an $`n`$-dimensional complex manifold is a $`2n`$-dimensional Riemannian space with a Hermitean metric
$$ds^2=g_{i\overline{j}}dz^id\overline{z}^{\overline{j}}:g_{ij}=g_{\overline{i}\overline{j}}=0,g_{i\overline{j}}=\overline{g}_{\overline{j}i}.$$
(2.1)
on its $`n`$ complex coordinates $`z_i`$. Such a complex manifold is Kähler if the $`(1,1)`$ differential two-form
$$\mathrm{\Omega }=\frac{1}{2}ig_{i\overline{j}}dz^i\mathrm{\Lambda }d\overline{z}^{\overline{j}},$$
(2.2)
is closed, i.e., $`d\mathrm{\Omega }=0`$. In the case of a Kähler manifold, the metric (2.1) is defined by a Kähler potential:
$$g_{i\overline{j}}=\frac{^2K(z^i,\overline{z}^{\overline{j}})}{z^i\overline{z}^{\overline{j}}}.$$
(2.3)
The Kähler property yields the following constraints on components of the Cristoffel symbols:
$`\mathrm{\Gamma }_{\overline{j}k}^i`$ $`=`$ $`\mathrm{\Gamma }_{j\overline{k}}^{\overline{i}}=\mathrm{\Gamma }_{\overline{j}k}^i=\mathrm{\hspace{0.17em}0},`$
$`\mathrm{\Gamma }_{\overline{j}\overline{k}}^{\overline{i}}`$ $`=`$ $`\overline{\mathrm{\Gamma }_{jk}^i}=g^{\overline{i}s}{\displaystyle \frac{g_{\overline{k}s}}{\overline{z}^{\overline{j}}}}.`$ (2.4)
yielding in turn the following form
$$R_{\overline{i}j}=\frac{\mathrm{\Gamma }_{\overline{i}\overline{k}}^{\overline{k}}}{z^j}.$$
(2.5)
for the Ricci tensor.
Since the only compact submanifold of $`C^n`$ is a point , in order to find non-trivial compact submanifolds, one considers weighted complex projective spaces, $`CP^n(k_1,k_2,\mathrm{},k_{n+1})`$, which are characterized by $`(n+1)`$ quasihomogeneous coordinates $`z_1,\mathrm{},z_{n+1}`$, with the identification:
$$(z_1,\mathrm{},z_{n+1})(\lambda ^{k_1}z_1,\mathrm{},\lambda ^{k_{n+1}}z_{n+1}).$$
(2.6)
The loci of zeroes of quasihomogeneous polynomial equations in such weighted projective spaces yield compact submanifolds, as we explain in more detail in the rest of Section 2, where we introduce and review several of the geometric and algebraic techniques used in our subsequent classification. Other compact submanifolds may be obtained as the complete intersections of such polynomial zero constraints, as we discuss in more detail in Section 9.
### 2.1 The Topology of Calabi-Yau Manifolds in the Polyhedron Method
A CY variety $`X`$ in a weighted projective space $`CP^n(\stackrel{}{k})=CP^n(k_1,\mathrm{},k_{n+1})`$ is given by the locus of zeroes of a transversal quasihomogeneous polynomial $`\mathrm{}`$ of degree $`deg(\mathrm{})=d`$, with $`d=_{j=1}^{n+1}k_j`$ :
$`XX_d(k)\{[x_1,\mathrm{},x_{n+1}]CP^n(k)|\mathrm{}(x_1,\mathrm{},x_{n+1})=0\}.`$ (2.7)
The general polynomial of degree $`d`$ is a linear combination
$$\mathrm{}=\underset{\stackrel{}{\mu }}{}c_\stackrel{}{\mu }x^\stackrel{}{\mu }$$
(2.8)
of monomials $`x^\stackrel{}{\mu }=x_1^{\mu _1}x_2^{\mu _2}\mathrm{}x_{r+1}^{\mu _{r+1}}`$ with the condition:
$$\stackrel{}{\mu }\stackrel{}{k}=d.$$
(2.9)
We recall that the existence of a mirror symmetry, according to which each Calabi-Yau manifold should have a dual partner, was first observed pragmatically in the literature . Subsequently, Batyrev found a very elegant way of describing any Calabi-Yau hypersurface in terms of the corresponding Newton polyhedron, associated with degree-$`d`$ monomials in the CY equation, which is the convex hull of all the vectors $`\stackrel{}{\mu }`$ of degree $`d`$. The Batyrev description provides a systematic approach to duality and mirror symmetry.
To each monomial associated with a vector $`\stackrel{}{\mu }`$ of degree $`d`$, i.e., $`\stackrel{}{\mu }\stackrel{}{k}=d`$, one can associate a vector $`\stackrel{}{\mu }^{^{}}\stackrel{}{\mu }\stackrel{}{e}_0:\stackrel{}{e}_0(1,1,\mathrm{},1)`$, so that $`\stackrel{}{\mu }^{^{}}\stackrel{}{k}=0`$. Using the new vector $`\stackrel{}{\mu }^{^{}}`$, hereafter denoted without the prime $`(^{})`$, it is useful to define the lattice $`\mathrm{\Lambda }`$:
$$\mathrm{\Lambda }=\{\stackrel{}{\mu }Z^{r+1}:\stackrel{}{\mu }\stackrel{}{k}=0\}$$
(2.10)
with basis vectors $`e_i`$, and the dual lattice $`\mathrm{\Lambda }^{}`$ with basis $`e_j^{}`$, where $`e_j^{}e_i=\delta _{ij}`$. Consider the polyhedron $`\mathrm{}`$, defined to be the convex hull of $`\{\stackrel{}{\mu }\mathrm{\Lambda }:\mu _i1,i\}`$. Batyrev showed that to describe a Calabi-Yau hypersurface I.e., with trivial canonical bundle and at worst Gorenstein canonical singularities only., such a polyhedron should satisfy the following conditions:
* the vertices of the polyhedron should correspond to the vectors $`\stackrel{}{\mu }`$ with integer components;
* there should be only one interior integer point, called the center;
* the distance of any face of this polyhedron from the center should be equal to unity.
Such an integral polyhedron $`\mathrm{}`$ is called reflexive, and the only interior point of $`\mathrm{}(k_1+\mathrm{}+k_{r+1}=d)`$ may be taken as the origin $`(0,\mathrm{},0)`$. Batyrev showed that the mirror polyhedron
$`\mathrm{}^{}\{\stackrel{}{\nu }\mathrm{\Lambda }^{}:\stackrel{}{\nu }\stackrel{}{\mu }1,\stackrel{}{\mu }\mathrm{}\}`$ (2.11)
of any reflexive integer polyhedron is also reflexive, i.e., is also integral and contains one interior point only. Thus Batyrev proved the existence of dual pairs of hypersurfaces $`M`$ and $`M^{}`$ with dual Newton polyhedra, $`\mathrm{}`$ and $`\mathrm{}^{}`$.
Following Batyrev , to obtain all the topological invariants of the $`K3,CY_3`$, etc., manifolds, one should study the reflexive regular polyhedra in three, four, etc., dimensions. For this purpose, it is useful to recall the types of polyhedra and their duality properties. In three dimensions, the Descartes-Euler polyhedron formula relates the numbers of vertices, $`N_0`$, the number of edges, $`N_1`$ and numbers of faces, $`N_2`$:
$`1N_0+N_1N_2+1=\mathrm{\hspace{0.17em}0}.`$ (2.12)
This formula yields:
$`1\mathrm{\hspace{0.17em}4}+\mathrm{\hspace{0.17em}6}\mathrm{\hspace{0.17em}4}+1=\mathrm{\hspace{0.17em}0}`$ $``$ $`\{3,3\}:\mathrm{Tetrahedron}`$
$`1\mathrm{\hspace{0.17em}8}+\mathrm{\hspace{0.17em}12}\mathrm{\hspace{0.17em}6}+1=\mathrm{\hspace{0.17em}0}`$ $``$ $`\{3,4\}:\mathrm{Cube}`$
$`1\mathrm{\hspace{0.17em}6}+\mathrm{\hspace{0.17em}12}\mathrm{\hspace{0.17em}8}+1=\mathrm{\hspace{0.17em}0}`$ $``$ $`\{4,3\}:\mathrm{Octahedron}`$
$`1\mathrm{\hspace{0.17em}20}+\mathrm{\hspace{0.17em}30}\mathrm{\hspace{0.17em}12}+1=\mathrm{\hspace{0.17em}0}`$ $``$ $`\{5,3\}:\mathrm{Dodecahedron}`$
$`1\mathrm{\hspace{0.17em}12}+\mathrm{\hspace{0.17em}30}\mathrm{\hspace{0.17em}20}+1=\mathrm{\hspace{0.17em}0}`$ $``$ $`\{3,5\}:\mathrm{Icosahedron}`$ (2.13)
in the particular cases of the five Platonic solids, with the duality relations $`TT,CO,DI`$.
As we shall see later when we consider the $`K3`$ classification, it is interesting to recall the link between the classification of the five ADE simply-laced Cartan-Lie algebras and the finite rotation groups in three dimensions, namely the cyclic and dihedral groups and the groups of the tetrahedron, octahedron (cube) and icosahedron (dodecahedron): $`G_MC_n,D_n,T,O,I`$, corresponding to the $`A_n,D_n`$ series and the exceptional groups $`E_{6,7,8}`$, respectively . Any cyclic group $`C_n`$ of order $`n`$ may be represented as the rotations in a plane around an axis $`0x`$ through angles $`(2m\pi )/n`$ for $`m=0,1,2,\mathrm{},n1`$. This symmetry is realized by the group of symmetries of an oriented regular $`n`$-gon. The dihedral group $`D_n`$ consists of the transformations in $`C_n`$ and in addition $`n`$ rotations through angles $`\pi `$ around axes lying in planes orthogonal to $`0x`$, crossing $`0x`$ and making angles with one another that are multiples of $`(2\pi )/n`$. This group has order $`2n`$. In the case of three-dimensional space, there are three exceptional examples $`T,O,I`$ of finite groups, related to the corresponding regular polyhedra. The order of the corresponding $`G_M`$ is equal to the product of the number of the vertexes of the regular polyhedra with the number of edges leaving the vertex:
$`|T|`$ $`=`$ $`|A_4|=\mathrm{\hspace{0.17em}12};`$
$`|O|`$ $`=`$ $`|S_4|=\mathrm{\hspace{0.17em}24};`$
$`|I|`$ $`=`$ $`|A_5|=\mathrm{\hspace{0.17em}60}.`$
The dual polyhedron, whose vertexes are the midpoints of the faces of the corresponding polyhedron, has the same group of symmetry, $`G_M`$. The finite groups of orthogonal transformations in three-dimensional space do not consist only of rotations. It is remarkable to note that every finite group of rotations of three-space that preserves the sphere centred at the origin can be interpreted as a fractional-linear transformation of the Riemann sphere of a complex variable.
Finally, we recall that all $`K3`$ hypersurfaces have the following common values of the topological invariants: the Hodge number $`h_{1,1}`$ is 20, the Betti number $`b_2=22`$, and we have
$`Pic=h_{1,1}(l(\mathrm{\Delta })4{\displaystyle \underset{\theta \mathrm{\Delta }}{}}l^{}(\theta ))\mathrm{\hspace{0.17em}20}.`$ (2.15)
for the Picard number, where $`l(\mathrm{\Delta })`$ is the number of integer points in the polyhedron and $`l^{}(\theta )`$ is the number of integer interior points on the facets.
In the case of the $`CY_3`$ classification, a corresponding important role will be played by the structure and the duality properties of the four regular polyhedra known in four-dimensional Euclidean space . The Descartes-Euler formulae for these cases become:
$`1\mathrm{\hspace{0.17em}5}+\mathrm{\hspace{0.17em}10}\mathrm{\hspace{0.17em}10}+\mathrm{\hspace{0.17em}5}\mathrm{\hspace{0.17em}1}=\mathrm{\hspace{0.17em}0}`$ $``$ $`\{3,3,3\}:\mathrm{Pentahedroid}`$
$`1\mathrm{\hspace{0.17em}16}+\mathrm{\hspace{0.17em}32}\mathrm{\hspace{0.17em}24}+\mathrm{\hspace{0.17em}8}\mathrm{\hspace{0.17em}1}=\mathrm{\hspace{0.17em}0}`$ $``$ $`\{3,3,4\}:\mathrm{Hypercube}`$
$`1\mathrm{\hspace{0.17em}8}+\mathrm{\hspace{0.17em}24}\mathrm{\hspace{0.17em}32}+\mathrm{\hspace{0.17em}16}\mathrm{\hspace{0.17em}1}=\mathrm{\hspace{0.17em}0}`$ $``$ $`\{4,3,3\}:16\mathrm{hedroid}`$
$`1\mathrm{\hspace{0.17em}24}+\mathrm{\hspace{0.17em}96}\mathrm{\hspace{0.17em}96}+\mathrm{\hspace{0.17em}24}\mathrm{\hspace{0.17em}1}=\mathrm{\hspace{0.17em}0}`$ $``$ $`\{3,4,3\}:24\mathrm{hedroid}`$
$`1\mathrm{\hspace{0.17em}600}+\mathrm{\hspace{0.17em}1200}\mathrm{\hspace{0.17em}720}+\mathrm{\hspace{0.17em}120}\mathrm{\hspace{0.17em}1}=\mathrm{\hspace{0.17em}0}`$ $``$ $`\{3,3,5\}:120\mathrm{hedroid}`$
$`1\mathrm{\hspace{0.17em}120}+\mathrm{\hspace{0.17em}720}\mathrm{\hspace{0.17em}1200}+\mathrm{\hspace{0.17em}600}\mathrm{\hspace{0.17em}1}=\mathrm{\hspace{0.17em}0}`$ $``$ $`\{5,3,3\}:600\mathrm{hedroid}.`$
with the duality relations $`PP`$, $`H16\mathrm{hedroid}`$, $`24\mathrm{hedroid}`$ $``$ $`24\mathrm{hedroid}`$,
$`120\mathrm{hedroid}`$ $``$ $`600\mathrm{hedroid}`$.
We do not discuss these relations further in this paper, but do recall that each mirror pair of CY spaces, $`M_{CY}`$ and $`M_{CY}^{}`$, has Hodge numbers that satisfying the mirror symmetry relation :
$`h_{1,1}(M)`$ $`=`$ $`h_{d1,1}(M^{}),`$
$`h_{d1,1}(M)`$ $`=`$ $`h_{1,1}(M^{})`$ (2.17)
This means that the Hodge diamond of $`M_{CY}^{}`$ is a mirror reflection through a diagonal axis of the Hodge diamond of $`M_{CY}`$. The existence of mirror symmetry is a consequence of the dual properties of CY manifolds. A pair of reflexive polyhedra $`(\mathrm{},\mathrm{}^{})`$ gives a pair of mirror CY manifolds and the following identities for the Hodge numbers for $`n4`$:
$`h_{1,1}(\mathrm{})`$ $`=`$ $`h_{d1,1}(\mathrm{}^{})=`$ (2.18)
$`=`$ $`l(\mathrm{}^{})(d+2){\displaystyle \underset{codim\mathrm{\Theta }^{}=1}{}}l^{^{}}(\mathrm{\Theta }^{})`$
$`+`$ $`{\displaystyle \underset{codim\mathrm{\Theta }^{}=2}{}}l^{}(\mathrm{\Theta }^{})l^{^{}}\mathrm{\Theta },`$
$`h_{1,1}(\mathrm{}^{})`$ $`=`$ $`h_{d1,1}(\mathrm{})=`$ (2.19)
$`=`$ $`l(\mathrm{})(d+2){\displaystyle \underset{codim\mathrm{\Theta }=1}{}}l^{}(\mathrm{\Theta })`$
$`+`$ $`{\displaystyle \underset{codim\mathrm{\Theta }=2}{}}l^{}(\mathrm{\Theta })l^{}(\mathrm{\Theta }^{}),`$
$`h_{p,1}={\displaystyle \underset{codim\mathrm{\Theta }^{}=p+1}{}}l^{^{}}(\mathrm{\Theta })l^{^{}}(\mathrm{\Theta }^{}),\mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}1}<p<d1.`$ (2.20)
Here, the quantities $`l(\mathrm{\Theta })`$ and $`l^{}(\mathrm{\Theta })`$ are the numbers of integer points on a face $`\mathrm{\Theta }`$ of $`\mathrm{}`$ and in its interior, and similarly for $`\mathrm{\Theta }^{}`$ and $`\mathrm{}^{}`$. An $`l`$dimensional face $`\mathrm{\Theta }`$ can be defined by its vertices $`(v_{i_1}=\mathrm{}=v_{i_k})`$, and the dual face defined by $`\mathrm{\Theta }^{}=\{m\mathrm{}^{}:(m,v_{i_1})=,\mathrm{},=(m,v_{i_k})=1\}`$ is an $`(nl1)`$-dimensional face of $`\mathrm{}^{}`$. Thus, we have a duality between the $`l`$-dimensional faces of $`\mathrm{}`$ and the $`(nl1)`$-dimensional faces of $`\mathrm{}^{}`$. The last terms in (2.18, 2.19) correspond to the ‘twisted’ contributions, and the last term corresponds to $`d=4`$. In this case, if the manifold has $`SU(4)`$ group holonomy, then $`h_{2,0}=h_{1,0}=0`$, and the remaining non-trivial Hodge number $`h_{2,2}`$ is determined by:
$$h_{2,2}=2[22+2h_{1,1}+h_{3,1}h_{2,1}].$$
(2.21)
Some further comments about $`CY_3`$ spaces are made in Section 9.
### 2.2 The Web of CY Manifolds in the Holomorphic-Quotient Approach to Toric Geometry
It is well known that weighted projective spaces are examples of toric varieties . The complex weighted projective space $`CP^n`$ can be defined as
$$CP^n\frac{C^{n+1}\stackrel{}{0}}{C^{}},$$
(2.22)
with the action $`C^{}`$:
$$(x_1,\mathrm{},x_{n+1})(\lambda ^{k_1}x_1,\mathrm{}.,\lambda ^{k_{n+1}}x_{n+1}),\lambda C\backslash 0.$$
(2.23)
The generalization of the projective space $`CP^n`$ to a toric variety can be expressed in the following form:
$`\mathrm{}{\displaystyle \frac{C^nZ_\mathrm{\Sigma }}{(C^{})^p}},`$ (2.24)
where, instead of removing the origin, as in the case of a simple projective space, here one removes a point set $`Z_\mathrm{\Sigma }`$, and one takes the quotient by a suitable set of $`C^{}`$ actions. Thus, to understand the structure of certain geometrical spaces in the framework of toric geometry, one must specify the combinatorical properties of the $`Z_\mathrm{\Sigma }`$ and the actions $`C^{}`$.
In the toric-geometry approach, algebraic varieties are described by a dual pair of lattices $`M`$ and $`N`$, each isomorphic to $`Z^n`$, and a fan $`\mathrm{\Sigma }^{}`$ defined on $`N_R`$, the real extension of the lattice $`N`$. In the toric-variety description, the equivalence relations of projective vectors can be considered as diagrams in the lattice $`N`$, in which some vectors $`\stackrel{}{v}_i`$ satisfy linear relations (see later some examples in $`P^2(1,1,1)`$, $`P^2(1,1,2)`$, $`P^2(1,2,3)`$ projective spaces). The complex dimension of the variety coincides with the dimension of the lattice $`N`$. To determine the structure of a toric variety in higher dimensions $`d>2`$, it is useful to introduce the notion of a fan . A fan $`\mathrm{\Sigma }^{}`$ is defined as a collection of $`r`$-dimensional ($`0rd`$) convex polyhedral cones with apex in 0, with the properties that with every cone it contains also a face, and that the intersection of any two cones is a face of each one.
In the holomorphic-quotient approach of Batyrev and Cox , a single homogeneous coordinate is assigned to the system $`\mathrm{}_\mathrm{\Sigma }`$ of varieties, in a way similar to the usual construction of $`P^n`$. This holomorphic-quotient construction gives immediately the usual description in terms of projective spaces, and turns out to be more natural in the descriptions of the elliptic, $`K3`$ and other fibrations of higher-dimensional CY spaces.
One can assign a coordinate $`z_k:k=1,\mathrm{},N`$ to each one-dimensional cone in $`\mathrm{\Sigma }`$. The integer points of $`\mathrm{\Delta }^{}N`$ define these one-dimensional cones
$$(v_1,\mathrm{},v_N)=\mathrm{\Sigma }_{1}^{}{}_{}{}^{}$$
(2.25)
of the fan $`\mathrm{\Sigma }^{}`$. The one-dimensional cones span the vector space $`N_R`$ and satisfy $`(Nn)`$ linear relations with non-negative integer coefficients:
$$\underset{l}{}k_j^lv_l=0,k_j^l0.$$
(2.26)
These linear relations can be used to determine equivalence relations on the space $`C^N\backslash Z_\mathrm{\Sigma }^{}`$. A variety $`\mathrm{}_\mathrm{\Sigma }^{}`$ is the space $`C^N\backslash Z_\mathrm{\Sigma }^{}`$ modulo the action of a group which is the product of a finite Abelian group and the torus $`(C^{})^{(Nn)}`$ :
$$(z_1,\mathrm{},z_N)(\lambda ^{k_j^1}z_1,\mathrm{}.,\lambda ^{k_j^N}z_N),j=1,\mathrm{},Nn.$$
(2.27)
The set $`Z_\mathrm{\Sigma }^{}`$ is defined by the fan in the following way:
$$Z_\mathrm{\Sigma }^{}\underset{I}{}((z_1,\mathrm{},z_N)|z_i=0,iI)$$
(2.28)
where the union is taken over all index sets $`I=(i_1,\mathrm{},i_k)`$ such that $`(v_{i_1},\mathrm{},v_{i_k})`$ do not belong to the same maximal cone in $`\mathrm{\Sigma }^{}`$, or several $`z_i`$ can vanish simultaneously only if the corresponding one-dimensional cones $`v_i`$ are from the same cone.
It is clear from the above definitions that toric varieties can have often singularities, which will be very important for understanding the link between the topological properties of Calabi-Yau hypersurfaces and Cartan-Lie algebras: see the more systematic discussion in Section 8. The method of blowing up (blowing down) these singularities was developed in algebraic geometry: it consists of replacing the singular point or curve by a higher-dimensional (lower-dimensional) variety. The structure of the fan $`\mathrm{\Sigma }^{}`$ determines what kind of singularities will appear in Calabi-Yau hypersurfaces. For example, if the fan $`\mathrm{\Sigma }^{}`$ is simplicial, one can get only orbifold singularities in the corresponding variety .
The elements of $`\mathrm{\Sigma }_1^{}`$ are in one-to-one correspondence with divisors
$$D_{v_i}=\mathrm{}_{\mathrm{\Sigma }_{1i}^{}},$$
(2.29)
which are subvarieties given simply by $`z_i=0`$. This circumstance was used to give a simple graphic explanation of Cartan-Lie algebra (CLA) diagrams, whose Coxeter number could be identified with the intersections of the divisors $`D_{v_i}`$.
Two divisors, $`D_{v_i}`$ and $`D_{v_j}`$, can intersect only when the corresponding one-dimensional cones $`v_i`$ and $`v_j`$ lie in a single higher-dimensional cone of the fan $`\mathrm{\Sigma }^{}`$. The divisors $`D_{v_i}`$ form a free Abelian group $`Div(\mathrm{}_\mathrm{\Sigma }^{})`$. In general, a divisor $`DDiv(\mathrm{}_\mathrm{\Sigma }^{})`$ is a linear combination of some irreducible hypersurfaces with integer coefficients:
$$qD=a_iD_{v_i}.$$
(2.30)
If $`a_i\mathrm{\hspace{0.17em}0}`$ for every $`i`$, one can say that $`D>\mathrm{\hspace{0.17em}0}`$. For a meromorphic function $`f`$ on a toric variety, one can define a principal divisor
$$(f)\underset{i}{}ord_{D_i}(f)D_i,$$
(2.31)
where $`ord_{D_i}(f)`$ is the order of the meromorphic function $`f`$ at $`D_i`$. One can further define the zero divisor $`(f)_0`$ and the polar divisor $`(f)_{inf}`$ of the meromorphic function $`f`$, such that
$$(f)=(f)_0(f)_{inf}.$$
(2.32)
Any two divisors $`D_1,D_2`$ are linearly equivalent: $`D_1D_2`$, if their difference is a principal divisor, $`DD^{}=(f)`$ for some appropriate $`f`$. The quotient of all divisors $`Div(\mathrm{}_\mathrm{\Sigma }^{})`$ by the principal divisors forms the Picard group.
The points of $`\mathrm{\Delta }M`$ are in one-to-one correspondence with the monomials in the homogeneous coordinates $`z_i`$. A general polynomial is given by
$$\mathrm{}=\underset{m\mathrm{\Delta }M}{}c_m\underset{l=1}{\overset{N}{}}z_l^{v_l,m+1}.$$
(2.33)
The equation $`\mathrm{}=0`$ is well defined and $`\mathrm{}`$ is holomorphic if the condition
$$v_l,m1foralll$$
(2.34)
is satisfied. The $`c_m`$ parametrize a family $`M_\mathrm{\Delta }`$ of CY surfaces defined by the zero locus of $`\mathrm{}`$.
### 2.3 Three Examples of $`CY_1`$ Spaces
As discussed in Section 1, three $`CY_1`$ spaces may be obtained as simple loci of polynomial zeroes associated with refelxive polyhedra. For a better understanding of the preceding formalism, we consider as warm-up examples the three elliptic reflexive polyhedron pairs $`\mathrm{\Delta }_i`$ and $`\mathrm{\Delta }_i^{}`$, which define the $`CY_1`$ surfaces $`P^2(1,1,1)[3]`$, $`P^2(1,1,2)[4]`$, and $`P^2(1,2,3)[6]`$ Here and subsequently, we use the conventional notation for such surfaces in n-dimensional projective space: $`P^n(k_1,k_2,\mathrm{})[k_1+k_2+\mathrm{}]`$.. The first polyhedron $`\mathrm{\Delta }_I\mathrm{\Delta }(P^2(1,1,1)[3])`$ consists of the following ten integer points:
$`z^3`$ $``$ $`\mu _1^{(I)}=(1,2),`$
$`xz^2`$ $``$ $`\mu _2^{(I)}=(1,1),`$
$`x^2z`$ $``$ $`\mu _3^{(I)}=(1,0),`$
$`x^3`$ $``$ $`\mu _4^{(I)}=(1,1),`$
$`yz^2`$ $``$ $`\mu _5^{(I)}=(0,1),`$
$`xyz`$ $``$ $`\mu _6^{(I)}=(0,0),`$
$`x^2y`$ $``$ $`\mu _7^{(I)}=(0,1),`$
$`y^2z`$ $``$ $`\mu _8^{(I)}=(1,0),`$
$`xy^2`$ $``$ $`\mu _9^{(I)}=(1,1),`$
$`y^3`$ $``$ $`\mu _{10}^{(I)}=(2,1).`$ (2.35)
and the mirror polyhedron $`\mathrm{\Delta }_I^{}\mathrm{\Delta }^{}(P^2(1,1,1)[3])`$ consists of one interior point and three one-dimensional cones:
$`v_1^{(I)}`$ $`=`$ $`(0,1),`$
$`v_2^{(I)}`$ $`=`$ $`(1,0),`$
$`v_3^{(I)}`$ $`=`$ $`(1,1).`$ (2.36)
We use as a basis the exponents of the following monomials:
$`y^2z`$ $``$ $`\stackrel{}{e}_1=(1,1,0),`$
$`yz^2`$ $``$ $`\stackrel{}{e}_2=(1,0,1).`$ (2.37)
where the determinant of this lattice coincides with the dimension of the projective vector $`\stackrel{}{k}=(1,1,1)`$ (see Figure 2):
$`det\{\stackrel{}{e}_1,\stackrel{}{e}_2,\stackrel{}{e}_0\}=dim(\stackrel{}{k})=\mathrm{\hspace{0.17em}3},`$ (2.38)
where $`\stackrel{}{e}_0`$ is the unit vector (1,1,1).
For this projective vector there exist 27 possibilities of choising two monomials for constructing the basis. Of course, all these bases are equivalent, i.e., they are connected by the $`SL(2,Z)`$ modular transformations:
$$L_{i,j}=\left(\begin{array}{cc}a& b\\ c& d\end{array}\right)$$
where $`a,b,c,dZ`$ and $`adbc=1`$. For the mirror polyhedron obtained from this vector, the basis should correspond to a lattice with determinant three times greater than (2.38), namely 9, for example:
$`\stackrel{}{e}_1=(1,2,1),`$
$`\stackrel{}{e}_2=(1,1,2).`$ (2.39)
with
$`det\{\stackrel{}{e}_1,\stackrel{}{e}_2,\stackrel{}{e}_0\}=dim(\stackrel{}{k})=\mathrm{\hspace{0.17em}9}.`$ (2.40)
where $`\stackrel{}{e}_0`$ is again the unit vector (1,1,1).
To describe this toric curve, one should embed it in the toric variety
$$P^2=(C^3\backslash 0)/(C\backslash 0),$$
(2.41)
where the equivalence relation
$$(x_1,x_2,x_3)(\lambda x_1,\lambda x_2,\lambda x_3)for\lambda C\backslash 0$$
(2.42)
is a consequence of the equation:
$`q_1v_1^{(I)}+q_2v_2^{(I)}+q_3v_3^{(I)}=\mathrm{\hspace{0.17em}0},`$ (2.43)
where the $`q_i=1,i=1,2,3`$ are the exponents of $`\lambda `$. The corresponding general polynomial describing a CY surface is (setting $`z_lx_l`$):
$`\mathrm{}_I`$ $`=`$ $`x_1^3+x_2^3+x_3^3+x_1x_2x_3`$ (2.44)
$`+`$ $`x_1^2x_2+x_1^2x_3+x_2^2x_1+x_2^2x_3+x_3^2x_1+x_3^2x_2.`$
and the Weierstrass equation can be written in the following form:
$$y^2z=x^3+axy^3+bz^3$$
(2.45)
where we have set $`x_1=x,x_2=y,x_3=z`$.
The second dual pair of triangle polyhedra $`\mathrm{\Delta }_{II}\mathrm{\Delta }(P^2(1,1,2)[4])`$ and its mirror $`\mathrm{\Delta }_{II}^{}\mathrm{\Delta }^{}(P^2(1,1,2)[4])`$ have nine points
$`y^4`$ $``$ $`\mu _1^{(II)}=(1,2),`$
$`xy^3`$ $``$ $`\mu _2^{(II)}=(1,1),`$
$`x^2y^2`$ $``$ $`\mu _3^{(II)}=(1,0),`$
$`x^3y`$ $``$ $`\mu _4^{(II)}=(1,1),`$
$`x^4`$ $``$ $`\mu _5^{(II)}=(1,2),`$
$`y^2z`$ $``$ $`\mu _6^{(II)}=(0,1),`$
$`xyz`$ $``$ $`\mu _7^{(II)}=(0,0),`$
$`x^2z`$ $``$ $`\mu _8^{(II)}=(0,1),`$
$`z^2`$ $``$ $`\mu _9^{(II)}=(1,0).`$ (2.46)
and five points, respectively (see Figure 3).
We use as a basis the exponents of the following monomials:
$`z^2`$ $``$ $`\stackrel{}{e}_1=(1,1,1),`$
$`y^2z`$ $``$ $`\stackrel{}{e}_2=(1,1,0).`$ (2.47)
where the determinant of this lattice coincides with the dimension of the projective vector $`\stackrel{}{k}=(1,1,2)`$:
$`det\{\stackrel{}{e}_1,\stackrel{}{e}_2,\stackrel{}{e}_0\}=dim(\stackrel{}{k})=\mathrm{\hspace{0.17em}4},`$ (2.48)
where $`\stackrel{}{e}_0`$ is again the unit vector (1,1,1).
To get the mirror polyhedron with 5 integer points, 4 on the edges and one interior point, one should find a basis with lattice determinant twice (2.48), namely 8, for example:
$`\stackrel{}{e}_1`$ $`=`$ $`(1,1,1),`$
$`\stackrel{}{e}_2`$ $`=`$ $`(2,2,0).`$ (2.49)
The following four points define four one-dimensional cones in $`\mathrm{\Sigma }_1(\mathrm{\Delta }_{II}^{})`$:
$`v_1^{(3)}`$ $`=`$ $`(1,0),`$
$`v_2^{(3)}`$ $`=`$ $`(1,0),`$
$`v_3^{(3)}`$ $`=`$ $`(1,1),`$
$`v_4^{(3)}`$ $`=`$ $`(1,1).`$ (2.50)
Using the linear relations between the four one-dimensional cones, the corresponding $`(C^{})^2`$ is seen to be given by ($`z_l\chi _l`$):
$`(\chi _1,\chi _2,\chi _3,\chi _4)(\lambda \mu ^2\chi _1,\lambda \chi _2,\mu \chi _3,\mu \chi _4).`$ (2.51)
and the general polynomial has the following nine terms:
$`\mathrm{}_{II}`$ $`=`$ $`\chi _2^2\chi _3^4+\chi _2^2\chi _3^3\chi _4+\chi _2^2\chi _3^2\chi _4^2+\chi _2^2\chi _3\chi _4^3+\chi _2^2\chi _4^4`$ (2.52)
$`+`$ $`\chi _1\chi _2\chi _3^2+\chi _1\chi _2\chi _3\chi _4+\chi _1\chi _2\chi _4^2+\chi _1^2`$
in this case.
The vectors $`\stackrel{}{k}=(1,1,1)`$ and $`\stackrel{}{k}=(1,1,2)`$ have three common monomials and a related reflexive segment-polyhedron, corresponding to the projective vector $`\stackrel{}{k}_2=(1,1)`$ of $`CP^1`$. This circumstance can be used further in the construction of the projective algebra in which these two vectors appear in the same chain.
The last $`CY_1`$ example involves the plane of the projective vector $`\stackrel{}{k}=(1,2,3)`$, whose polyhedron $`\mathrm{\Delta }_{III}\mathrm{\Delta }(P^2(1,2,3)[6])`$ and its mirror partner $`\mathrm{\Delta }_{III}^{}\mathrm{\Delta }^{}(P^2(1,2,3)[6])`$ both have seven self-dual points, and one can check the existence of the following six one-dimensional cones (see Figure 4):
$`z^2`$ $``$ $`v_1^{(III)}=(1,0),`$
$`x^2y^2`$ $``$ $`v_2^{(III)}=(1,0),`$
$`x^3z`$ $``$ $`v_3^{(III)}=(0,1),`$
$`y^3`$ $``$ $`v_4^{(III)}=(1,1),`$
$`x^4y`$ $``$ $`v_5^{(III)}=(1,1),`$
$`x^6`$ $``$ $`v_6^{(III)}=(1,2).`$ (2.53)
We use as a basis the exponents of the following monomials:
$`z^2`$ $``$ $`\stackrel{}{e}_1=(1,1,1),`$
$`x^3z`$ $``$ $`\stackrel{}{e}_2=(2,1,0).`$ (2.54)
where the determinant of this lattice coincides with the dimension of the projective vector $`\stackrel{}{k}=(1,2,3)`$:
$`det\{\stackrel{}{e}_1,\stackrel{}{e}_2,\stackrel{}{e}_0\}=dim(\stackrel{}{k})=\mathrm{\hspace{0.17em}6}.`$ (2.55)
As in the case of the two projective vectors $`\stackrel{}{k}=(1,1,1)`$ and $`\stackrel{}{k}=(1,1,2)`$, the vectors $`\stackrel{}{k}=(1,1,2)`$ and $`\stackrel{}{k}=(1,2,3)`$ also have three common monomials, corresponding to the reflexive segment polyhedron described by the vector $`\stackrel{}{k}_2=(1,1)`$ in $`CP^1`$ projective space. Hence these vectors will appear in the second chain of the plane projective algebra.
Thus one can see that, with these three plane projective vectors, $`\stackrel{}{k}=(1,1,1)`$, $`\stackrel{}{k}=(1,1,2)`$, $`\stackrel{}{k}=(1,2,3)`$, one finds only triangle reflexive polyhedra intersecting the integer planar lattice in $`10+4^{}`$, $`9+5^{}`$, $`7+7^{}`$ points. Of course, on the plane one can find other reflexive polyhedra, whose intersection with the integer plane lattice will give new $`CP^1`$ surfaces corresponding to different polygons with more then three vertices, such as a reflexive pair of square and rhombus. These new figures can be obtained using the techniques of extended vectors.
In the following, we will go on to study reflexive polyhedron pairs in three-dimensional space. The corresponding general polynomial can be expressed in terms of six variables, and contains seven monomials:
$`\mathrm{}_{III}`$ $`=`$ $`z_1^2z_3+z_2^2z_3z_4^2z_5^2z_6^2+z_1z_2z_3^2z_5^2z_6^3+z_2^2z_4^3z_5+`$ (2.56)
$`+`$ $`z_2^2z_3^2z_4z_5^3z_6^4+z_2^2z_3^3z_5^4z_6^6+z_1z_2z_3z_4z_5z_6.`$
The $`C_{}^{}{}_{}{}^{4}`$ action is determined by the following linear relations:
$`v_1^{(III)}`$ $`+`$ $`v_2^{(III)}=\mathrm{\hspace{0.17em}0},`$
$`2v_1^{(III)}`$ $`+`$ $`v_4^{(III)}+v_5^{(III)}=\mathrm{\hspace{0.17em}0},`$
$`v_1^{(III)}`$ $`+`$ $`v_3^{(III)}+v_4^{(III)}=\mathrm{\hspace{0.17em}0},`$
$`3v_1^{(III)}`$ $`+`$ $`\mathrm{\hspace{0.17em}2}v_4^{(III)}+v_6^{(III)}=\mathrm{\hspace{0.17em}0}`$ (2.57)
between the elements of $`\mathrm{\Sigma }_1(\mathrm{\Delta }_{III}^{})`$, and is given by
$`(z_1,z_2,z_3,z_4,z_5,z_6)(\lambda \mu ^2\nu \rho ^3z_1,\lambda z_2,\nu z_3,\mu \nu \rho ^2z_4,\mu z_5,\rho z_6).`$ (2.58)
One can introduce the following birational map between $`P^2(1,2,3)[6]`$ and $`\mathrm{}_\mathrm{\Sigma }^{}`$:
$`z_1^2z_3`$ $`=`$ $`y_3^2`$ (2.59)
$`z_2^2z_4^3z_5`$ $`=`$ $`y_2^3`$ (2.60)
$`z_2^2z_3^3z_5^4z_6^6`$ $`=`$ $`y_1^6`$ (2.61)
Then, a dimensionally-reduced example of a CY manifold embedded in a toric variety is described by the weight vector $`k=(1,2,3)`$ and the zero locus of the Weierstrass polynomial
$`\mathrm{}_{III}=y_1^6+y_2^3+y_3^2+y_1y_2y_3+y_1^4y_2+y_1^2y_2^2+y_1^3y_3.`$ (2.62)
The elliptic Weierstrass equation can be written in the weighted projective space $`P^2(1,2,3)[6]`$ as
$$y^2=x^3+axz^4+bz^6$$
(2.63)
with the following equivalence relation
$`(x,y,z)(\lambda ^2x,\lambda ^3y,\lambda z),\lambda C\backslash 0`$ (2.64)
in this case.
These examples illustrate how toric varieties can be defined by the quotient of $`C^k\backslash Z_\mathrm{\Sigma }`$, and not only by a group $`(C\backslash 0)^{kn}`$. One should divide $`C^k\backslash Z_\mathrm{\Sigma }`$ also by a finite Abelian group $`G(v_1,\mathrm{},v_k)`$, which is determined by the relations between the $`D_{v_i}`$ divisors. In this case, the toric varieties can often have orbifold singularities, $`C^k\backslash G`$. For example, the toric variety defined by (2.63) looks near the points $`y=z=0`$ and $`x=z=0`$ locally like $`C^2\backslash Z_2`$ (related to the $`SU(2)`$ algebra) and $`C^2\backslash Z_3`$ (related to the $`SU(3)`$ algebra), respectively, as seen in Figure 5.
## 3 Gauge Group Identifications from Toric Geometry
### 3.1 Calabi-Yau Spaces as Toric Fibrations
As discussed in section 2, any Calabi-Yau manifold can be considered as a hypersurface in a toric variety, with a corresponding reflexive polyhedron $`\mathrm{\Delta }`$ with a positive-integer lattice $`\mathrm{\Lambda }`$, associated with a dual polyhedron $`\mathrm{\Delta }^{}`$ in the dual lattice $`\mathrm{\Lambda }^{}`$. The toric variety is determined by a fan $`\mathrm{\Sigma }^{}`$, consisting of the cones which are given by a triangulation of $`\mathrm{\Delta }^{}`$. A large subset of reflexive polyhedra and their corresponding Calabi-Yau manifolds can be classified in terms of their fibration structures. In this way, it is possible, as we discuss later, to connect the structures of all the projective vectors of the one dimensionality with the projective vectors of other dimensionalities, and thereby construct a new algebra in the set of all ‘reflexive’ projective vectors that gives the full set of $`CY_d`$ hypersurfaces in all dimensions: $`d=1,2,3,\mathrm{}`$.
In order to embark on this programme, it is useful first to review two key operations, intersection and projection, which can give possible fibration structures for reflexive polyhedra :
* There may exist a projection operation $`\pi :\mathrm{\Lambda }\mathrm{\Lambda }_{nk}`$, where $`\mathrm{\Lambda }_{nk}`$ is an $`(nk)`$-dimensional sublattice, and $`\pi (\mathrm{\Delta })`$ is also a reflexive polyhedron, and
* there may exist an intersection projection $`J`$ through the origin of a reflexive polyhedron, such that $`J(\mathrm{\Delta })`$ is again an $`(nl)`$-dimensional reflexive polyhedron, and
* these operations may exhibit the following duality properties:
$`\mathrm{\Pi }(\mathrm{\Delta })`$ $``$ $`J(\mathrm{\Delta }^{})`$
$`J(\mathrm{\Delta })`$ $``$ $`\mathrm{\Pi }(\mathrm{\Delta }^{}).`$ (3.1)
For a reflexive polyhedron $`\mathrm{\Delta }`$ with fan $`\mathrm{\Sigma }`$ over a triangulation of the facets of $`\mathrm{\Delta }^{}`$, the CY hypersurface in variety $`\mathrm{}_\mathrm{\Sigma }`$ is given by the zero locus of the polynomial:
$$\mathrm{}=\underset{\mu \mathrm{\Delta }M}{}c_\stackrel{}{\mu }\underset{i=1}{\overset{N}{}}z_i^{\stackrel{}{v}_i\stackrel{}{\mu }+1}.$$
(3.2)
One can consider the variety $`\mathrm{}_\mathrm{\Sigma }`$ as a fibration over the base $`\mathrm{}_{\mathrm{\Sigma }_{base}}`$ with generic fiber $`\mathrm{}_{\mathrm{\Sigma }_{fiber}}`$. This fibration structure can be written in terms of homogeneous coordinates. The fiber as an algebraic subvariety is determined by the polyhedron $`\mathrm{\Delta }_{fiber}^{}\mathrm{\Delta }_{CY}^{}`$, whereas the base can be seen as a projection of the fibration along the fiber. The set of one-dimensional cones in $`\mathrm{\Sigma }_{base}`$ (the primitive generator of a cone is zero or $`\stackrel{~}{v}_i`$) is the set of images of one-dimensional cones in $`\mathrm{\Sigma }_{CY}`$ (with primitive generator $`v_j`$) that do not lie in $`N_{fiber}`$. The image $`\mathrm{\Sigma }_{base}`$ of $`\mathrm{\Sigma }_{CY}`$ under $`\mathrm{\Pi }:N_{CY}N_{base}`$ gives us the following relation:
$`\mathrm{\Pi }v_i=r_i^j\stackrel{~}{v}_j,`$ (3.3)
if $`\mathrm{\Pi }v_i`$ is in the set of one-dimensional cones determined by $`\stackrel{~}{v}_j`$ $`r_i^jN`$, otherwise $`r_i^j=0`$.
Similarly, the base space is the weighted projective space with the torus transformation:
$`(\stackrel{~}{x}_1,\mathrm{},\stackrel{~}{x}_{\stackrel{~}{N}})(\lambda ^{\stackrel{~}{k}_j^1}\stackrel{~}{x}_1,\mathrm{},\lambda ^{\stackrel{~}{k}_j^{\stackrel{~}{N}}}\stackrel{~}{x}_{\stackrel{~}{N}}),j=1,\mathrm{},\stackrel{~}{N}\stackrel{~}{n},`$ (3.4)
where the $`\stackrel{~}{k}_j^i`$ are integers such that $`_j\stackrel{~}{k}_i^j\stackrel{~}{v}_j=0`$. The projection map from the variety $`\mathrm{}_\mathrm{\Sigma }`$ to the base can be written as
$`\stackrel{~}{x}_i={\displaystyle \underset{j}{}}x_j^{r_j^i},`$ (3.5)
corresponding to the following redefinitions of the torus transformation for $`\stackrel{~}{x}_i`$:
$`\mathrm{\Pi }:\stackrel{~}{x}_i\lambda ^{k_l^jr_j^i}\stackrel{~}{x}_i,{\displaystyle k_l^jr_j^i\stackrel{~}{v}_i}=\mathrm{\hspace{0.17em}0}.`$ (3.6)
In the toric description of $`K3`$ surfaces with elliptic fibers, denoted by $`\mathrm{\Delta }_{}^{}{}_{fiber}{}^{}`$, one can consider the following divisors: $`D_{fiber}`$, $`D_{section}`$, $`D_{v_a}`$ and $`D_{v_b}`$. The last pair of divisors correspond to lattice points of $`\mathrm{\Delta }^{}`$ that are ‘above’ or ‘below’ the fiber, respectively. Let us consider the case when all divisors $`D_{v_a}`$ (or $`D_{v_b}`$) shrink to zero size. In this case, there appears a $`K3`$ hypersurface with two point singularities, which belong to the ADE classification. The process of blowing up these singularities gives the primordial $`K3`$ manifold, and its intersection structure is given by the structure of the edges. The Cartan-Lie algebra (CLA) diagrams of the gauge groups that appear when the exceptional fibers are blown down to points are nothing but the edge diagrams of the upper and lower parts of $`\mathrm{\Delta }^{}`$ without vertices, respectively. A simple well-known example with elliptic fiber and with base $`P^1`$ is given by the following Weierstrass equation for the fiber:
$$y^2=x^3+f(z_1,z_2)xz^4+g(z_1,z_2)z^6,$$
(3.7)
where the coefficients $`f(z_1,z_2),g(z_1,z_2)`$ are functions on the base.
In the following parts of this Section, we discuss some examples of $`K3`$ spaces from our general classification, and explain the identification of their corresponding gauge groups.
### 3.2 Examples of $`K3`$ Toric Fibrations with $`J=\mathrm{\Pi }`$ Weierstrass structure
As a first example, we consider the case of the elliptic $`K3`$ hypersurface with elliptic fiber $`P^2(1,2,3)[6]`$ defined by the integer positive lattice with basis (we explain this lattice basis later in terms of our algebraic description):
$$\left(\begin{array}{c}\stackrel{}{e}_1\\ \stackrel{}{e}_2\\ \stackrel{}{e}_3\end{array}\right)=\left(\begin{array}{cccc}m& n& 0& 0\\ 2& 2& 1& 0\\ 1& 1& 1& 1\end{array}\right),$$
where we consider the following 12 pairs of integer numbers $`(m,n)`$ which are taken from the numbers: $`1,2,3,4,5,6`$,
$$\{(1,1),(1,2),(1,3),(1,4),(1,5),(1,6),(2,3),(2,5)(3,4),(4,5),(5,6)\}.$$
With this choice of the pairs, the basis above determines a self-dual set of 12 projective $`\stackrel{}{k}_4`$-vectors:
$`m=1,n=1`$ $``$ $`\stackrel{}{k}_4=(1,1,4,\mathrm{\hspace{0.17em}\hspace{0.17em}6})[12],(5,6,22,33)`$
$`m=1,n=2`$ $``$ $`\stackrel{}{k}_4=(1,2,6,\mathrm{\hspace{0.17em}\hspace{0.17em}9})[18],(3,5,16,24)`$
$`m=1,n=3`$ $``$ $`\stackrel{}{k}_4=(1,3,8,\mathrm{\hspace{0.17em}\hspace{0.17em}12})[24],(2,5,14,21)`$
$`m=1,n=4`$ $``$ $`\stackrel{}{k}_4=(1,4,10,15)[30],DI^{^{}}`$
$`m=1,n=5`$ $``$ $`\stackrel{}{k}_4=(1,5,12,16)[36],selfdual`$
$`m=1,n=6`$ $``$ $`\stackrel{}{k}_4=(1,6,14,21)[42],selfdual`$
$`m=2,n=3`$ $``$ $`\stackrel{}{k}_4=(2,3,10,15)[30],selfdual`$
$`m=2,n=5`$ $``$ $`\stackrel{}{k}_4=(2,5,14,21)[42],(1,3,8,12)`$
$`m=3,n=4`$ $``$ $`\stackrel{}{k}_4=(3,4,14,21)[42],DI^{^{\prime \prime }}`$
$`m=3,n=5`$ $``$ $`\stackrel{}{k}_4=(3,5,16,24)[48],(1,2,6,9)`$
$`m=4,n=5`$ $``$ $`\stackrel{}{k}_4=(4,5,18,27)[54],DI^{^{\prime \prime \prime }}`$
$`m=5,n=6`$ $``$ $`\stackrel{}{k}_4=(5,6,22,33)[66],(1,1,4,6)`$
Later this set will emerge as the intersection-projection symmetric $`XIX`$ chain ($`J=\mathrm{\Pi }`$) of our algebraic classification. In this example, one can see that the projective vectors corresponding to the tetrahedra produce a self-dual set. We also show in (LABEL:selfdset) the duality relations between 6 other vectors and some of the vectors in Table 1.
However, three of the projective vectors in (LABEL:selfdset), $`\stackrel{}{k}_4=(1,4,10,15)[30]`$, $`(3,4,14,21)[42]`$ and $`(4,5,18,27)[54]`$, correspond to polyhedra with 5 vertices, and their duals can be found among higher-level $`K3`$ spaces. They are found by double intersections (DI) among the five-dimensional extensions of the $`K3`$ vectors shown in Table 1:
$`\stackrel{}{k}_4=(1,4,10,15)[30]`$ $`\stackrel{DI^{^{}}}{}\{\stackrel{}{k}_5^{ex}=(0,1,6,8,15)[30]\}{\displaystyle \{\stackrel{}{k}_5^{ex}=(6,1,0,14,21)[42]\}}`$
$`\stackrel{}{k}_4=(3,4,14,21)[42]`$ $`\stackrel{DI^{^{\prime \prime }}}{}\{\stackrel{}{k}_5^{ex}=(2,1,0,6,9)[18]\}{\displaystyle \{\stackrel{}{k}_5^{ex}=(0,1,2,4,7)[14]\}}`$
$`\stackrel{}{k}_4=(4,5,18,27)[54]`$ $`\stackrel{DI^{^{\prime \prime \prime }}}{}\{\stackrel{}{k}_5^{ex}=(1,0,1,4,6)[12]\}{\displaystyle \{\stackrel{}{k}_5^{ex}=(0,1,1,3,5)[10]\}}`$
as discussed in more detail in Section 6.
The ascending Picard numbers for polyhedra in this chain include:
$`(\mathrm{\Delta }(P^3(1,6,14,21)[42]):\mathrm{}`$ $`=`$ $`24(24^{}),Pic=10(10^{})`$
$`(\mathrm{\Delta }(P^3(1,5,12,18)[36]):\mathrm{}`$ $`=`$ $`24(24^{}),Pic=10(10^{})`$
$`(\mathrm{\Delta }(P^3(1,4,10,15)[30]):\mathrm{}`$ $`=`$ $`25(20^{}),Pic=9(11^{})`$
$`(\mathrm{\Delta }(P^3(1,3,8,12)[24]):\mathrm{}`$ $`=`$ $`27(15^{}),Pic=8(14^{})`$
$`(\mathrm{\Delta }(P^3(1,2,6,9)[18]):\mathrm{}`$ $`=`$ $`30(12^{}),Pic=6(16^{})`$
$`(\mathrm{\Delta }(P^3(1,1,4,6)[12]):\mathrm{}`$ $`=`$ $`39(9),Pic=2(18^{})\mathrm{}\mathrm{}\mathrm{}`$ (3.10)
In the case of the mirror polyhedron chain, there is the inverse property: $`\mathrm{\Delta }^{}(P^3(1,6,14,21)[42])`$ corresponds to the maximal member of the set of mirror polyhedra. These Picard numbers are tabulated in Table 1, together with those of the other $`K3`$ spaces.
In the chain (LABEL:selfdset), the mirror polyhedra, $`\mathrm{\Delta }^{}`$, have an intersection plane $`H_{fiber}^{}`$ through the interior point which defines an elliptic-fiber triangle with seven integer points, $`P^2(1,2,3)[6]`$ (see Figures 6,7):
$$\mathrm{\Delta }_{fiber}^{}=\mathrm{\Delta }^{}H_{fiber}^{}.$$
(3.11)
By mirror symmetry in the polyhedron $`\mathrm{\Delta }`$, a projection operator $`\pi `$ can be defined: $`\pi :MM_{n1}`$, where $`M_{n1}`$ is an $`(n1)`$-dimensional sublattice, such that $`\pi (\mathrm{\Delta })`$ is a reflexive polyhedron in $`M_{n1}`$. This reflexive polyhedron also consists of seven points, so it is self-dual. Also, one can find a planar intersection $`H`$ through $`\mathrm{\Delta }`$ and through the interior point, which also produces the reflexive polyhedron with seven points, namely the fiber $`P^2(1,2,3)[6]`$ (see Figures 6,7):
$$\mathrm{\Delta }_{fiber}=\mathrm{\Delta }H_{fiber}.$$
(3.12)
The dual pair of tetrahedra $`\mathrm{\Delta }(P^3(1,1,4,6)[12]`$ and $`\mathrm{\Delta }(P^3(5,6,22,33)[66]`$ consist of $`39`$ and $`9`$ points, respectively, as seen in Figure 6. They are the biggest and smallest polyhedra in the chain (LABEL:selfdset), and all other tetrahedra in this chain can be found in this Figure. This contains, in particular, the two self-dual polyhedra $`\mathrm{\Delta }(P^3(1,6,14,21)[42]`$ and $`\mathrm{\Delta }(P^3(2,3,10,15)[30]`$ consist of $`24+24^{}`$ and $`18+18^{}`$ points, respectively, as seen in Figure 7:
$`(0,0,1),(0,1.1),(1,2,1),(6,2,1);`$
$`(0,0,1),(0,1.1),(2,2,1),(3,2,1).`$
We now consider the intersection of the three-dimensional polyhedron $`\mathrm{\Delta }(P^3(1,6,14,21)[42])`$ with the two-dimensional plane $`H`$ through the interior point. The intersection of this plane with the polyhedron, $`H\mathrm{\Delta }`$, forms a reflexive polyhedron fiber $`P^2(1,2,3)`$ with seven points. The equation of this plane in canonical coordinates $`\mu _1,\mu _2,\mu _3`$ is: $`m_1=\mathrm{\hspace{0.17em}0}.`$ The fiber consists of the following polyhedron points:
$`v_0`$ $`=`$ $`(0,\underset{¯}{0,0})`$
$`v_1`$ $`=`$ $`(0,\underset{¯}{1,0})`$
$`v_2`$ $`=`$ $`(0,\underset{¯}{0,1})`$
$`v_3`$ $`=`$ $`(0,\underset{¯}{1,1})`$
$`v_4`$ $`=`$ $`(0,\underset{¯}{0,1})`$
$`v_5`$ $`=`$ $`(0,\underset{¯}{1,1})`$
$`v_6`$ $`=`$ $`(0,\underset{¯}{2,1}).`$ (3.14)
Here and subsequently, the components of the vector corresponding to the fiber are underlined.
With respect to this fiber, the base is one-dimensional: $`P^1`$, and its fan $`F_2`$ consists of the divisors corresponding to the interior point and two divisors corresponding to two rays, $`R_1=+\stackrel{}{e}_1`$ and $`R_2=\stackrel{}{e}_1`$, with directions from the point $`(0,2,1)`$ to the point $`(6,2,1)`$ and from the point $`(0,2,1)`$ to $`(1,2,,1)`$, respectively. The points of $`\pi _B^1(R_i)`$ ( i.e., the points projected onto $`R_i`$ by $`\pi _B`$) for the rays $`R_i,(i=1=+,i=2=)`$ are of the form $`(\pm \mathrm{},b,c)`$, where $`(0,b,c)`$ is the point of the fiber.
The 16 points of $`\pi _B^1(R_1)`$ are listed in the Table 3: they correspond to the divisors $`D_{v_i}`$, which produce the $`E_8`$ algebra . Also, from this Table one can easily read the Coxeter numbers/weights. There is only one point in $`\pi _B^1(R_2)`$, namely
$$\stackrel{~}{v}_{1}^{}{}_{}{}^{1}=(1,\underset{¯}{2,1})$$
(3.15)
which therefore does not correspond to any non-trivial group.
### 3.3 Example of Gauge-Group Identification
Consider again the toric variety determined by the dual pair of polyhedra $`\mathrm{\Delta }(P^3(1,1,4,6)[12])`$ and its dual $`\mathrm{\Delta }^{}`$ shown in Figure 6. The mirror polyhedron contains the intersection $`H^{}`$ through the interior point, the elliptic fiber $`P^2(1,2,3)`$. For all integer points of $`\mathrm{\Delta }^{}`$ (apart from the interior point), one can define in a convenient basis the corresponding complex variables:
$`v_1`$ $`=`$ $`(0,\underset{¯}{2,3})z_1`$
$`v_2`$ $`=`$ $`(0,\underset{¯}{1,2})z_2`$
$`v_3`$ $`=`$ $`(0,\underset{¯}{1,1})z_3`$
$`v_4`$ $`=`$ $`(0,\underset{¯}{0,1})z_4`$
$`v_0`$ $`=`$ $`(0,\underset{¯}{0,0})`$
$`v_6`$ $`=`$ $`(0,\underset{¯}{1,0})z_6`$
$`v_7`$ $`=`$ $`(0,\underset{¯}{0,1})z_7`$ (3.16)
and
$`v_8`$ $`=`$ $`(1,\underset{¯}{4,6})z_8`$
$`v_9`$ $`=`$ $`(1,\underset{¯}{0,0})z_9.`$ (3.17)
There are some linear relations between integer points inside the fiber:
$`v_1+\mathrm{\hspace{0.17em}2}v_6+\mathrm{\hspace{0.17em}3}v_7`$ $`=`$ $`\mathrm{\hspace{0.17em}0},`$
$`v_2+v_6+\mathrm{\hspace{0.17em}2}v_7`$ $`=`$ $`\mathrm{\hspace{0.17em}0},`$
$`v_3+v_6+v_7`$ $`=`$ $`\mathrm{\hspace{0.17em}0},`$
$`v_4+v_6+v_7`$ $`=`$ $`\mathrm{\hspace{0.17em}0}`$ (3.18)
and also the following relation between points in $`\mathrm{\Delta }^{}`$:
$`v_8+v_9+\mathrm{\hspace{0.17em}4}v_6+\mathrm{\hspace{0.17em}6}v_7`$ $`=`$ $`\mathrm{\hspace{0.17em}0}`$ (3.19)
The polyhedron $`\mathrm{\Delta }(P^3(1,1,4,6))`$ contains 39 points, which can be subdivided as follows. There are seven points in the fiber $`P^2(1,2,3)`$, determined by the intersection of the plane $`m_1+\mathrm{\hspace{0.17em}2}m_2+\mathrm{\hspace{0.17em}3}m_3=\mathrm{\hspace{0.17em}0}`$ and the positive integer lattice. This plane separates the remaining 32 points in $`16`$ ‘left’ and $`16`$ ‘right’ points.
These ‘left’ and ‘right’ points define singularities of the $`E_{8_L}`$ and $`E_{8_R}`$ types, respectively, which may be illustrated as follows. The plane $`H(\mathrm{\Delta })=m_1+2m_2+3m_3`$ contains the following seven points:
$`t_1`$ $`=`$ $`(5,\underset{¯}{1,1})(z_8^6z_9^6)(\underset{¯}{z_1^6z_2^4z_3^3z_4^2}),`$
$`t_2`$ $`=`$ $`(3,\underset{¯}{0,1})(z_8^4z_9^4)(\underset{¯}{z_1^4z_2^3z_2^3z_4^2z_6}),`$
$`t_3`$ $`=`$ $`(2,\underset{¯}{1,0})(z_8^3z_9^3)(\underset{¯}{z_1^3z_2^2z_3^2z_4z_7}),`$
$`t_4`$ $`=`$ $`(1,\underset{¯}{1,1})(z_8^2z_9^2)(\underset{¯}{z_1^2z_2^2z_3z_4^2z_6^2}),`$
$`t_5`$ $`=`$ $`(0,\underset{¯}{0,0})(z_8z_9)(\underset{¯}{z_1^6z_2^4z_3^3z_4^2}),`$
$`t_6`$ $`=`$ $`(1,\underset{¯}{2,1})(\underset{¯}{z_2z_4^2z_6^3}),`$
$`t_7`$ $`=`$ $`(1,\underset{¯}{1,1})(\underset{¯}{z_3z_7^2}).`$
The Weierstrass equation for the $`E_{8_L}`$ group based on the polyhedron $`\mathrm{\Delta }(P^3(1,1,4,6))`$ can be written in the form:
$`\underset{¯}{z}_6^3+\underset{¯}{z}_6^2(a_2^{(1)}z_8z_9^3+a_2^{(2)}z_9^4)+`$
$`\underset{¯}{z}_1^4\underset{¯}{z}_6(a_4^{(1)}z_8^3z_9^5+a_4^{(2)}z_8^2z_9^6+a_4^{(3)}z_8z_9^7+a_4^{(4)}z_9^8)+`$
$`\underset{¯}{z}_1^6(a_6^{(1)}z_8^5z_9^7+a_6^{(2)}z_8^4z_9^8+a_6^{(3)}z_8^3z_9^9+a_6^{(4)}z_8^2z_9^{(10)}+a_6^{(5)}z_8z_9^{(11)}+a_6^{(6)}z_9^{(12)})`$
$`=\underset{¯}{z}_7^2+a_1\underset{¯}{z}_6\underset{¯}{z}_7z_9^2+\underset{¯}{z}_7(a_3^{(1)}z_8^2z_9^4+a_3^{(2)}z_8z_9^{(5)}).`$ (3.21)
The second Weierstrass equation for the $`E_{8_R}`$ group can be obtained from this equation by interchanging the variables desrcibing the base: $`z_8z_9`$ <sup>§</sup><sup>§</sup>§The coefficients $`a_i`$ correspond to the notations of .. The Weierstrass triangle equation can be presented in the following general form, where we denote $`\underset{¯}{z}_6=x`$, $`\underset{¯}{z}_7=y`$:
$`y^2+a_1xy+a_3y=x^3+a_2x^2+a_4x+a_6,`$ (3.22)
where the $`a_i`$ are polynomial functions on the base. The Weierstrass equation can be written in more simplified form as:
$$y^2=x^3+xf+g,$$
(3.23)
with discriminant
$$\mathrm{\Delta }=\mathrm{\hspace{0.17em}4}f^3+\mathrm{\hspace{0.17em}27}g^2.$$
(3.24)
In the zero locus of the discriminant, some divisors $`D_i`$ define the degeneration of the torus fiber.
In addition to the method described above, there is a somewhat different way to find the singularity type . As we saw in the above example, the polynomials $`f`$ and $`g`$ can be homogeneous of orders 8 and 12, respectively, with a fibration that is degenerate over 24 points of the base. For this form of Weierstrass equation, there exists the ADE classification of degenerations of elliptic fibers. In this approach, the type of degeneration of the fiber is determined by the orders of vanishing of the functions $`f`$, $`g`$ and $`\delta `$. In the case of the general Weierstrass equation, a general algorithm for the ADE classification of elliptic singularities was considered by Tate . For convenience, we repeat in the Table 4 some results of Tate’s algorithm, from which one can recover the $`E_8\times E_8`$ type of Lie-algebra singularity for the (1,1,4,6) polyhedron.
## 4 The Composite Structure of Projective Vectors
We now embark in more detail on our construction of the projective vectors $`\stackrel{}{k}`$ which determine CY hypersurfaces, as previewed briefly in the Introduction and based on the polyhedron technique and the concept of duality reviewed in Section 2. We develop this construction inductively, studying the structure of these vectors initially in low dimensions and then proceeding to higher ones.
### 4.1 Initiation to the Dual Algebra of CY Projective Vectors
Our starting point is the trivial zero-dimensional ‘vector’,
$$\stackrel{}{k}_1=(1).$$
(4.1)
which defines the trivial self-dual polyhedron comprising a single point, with the simplest possible associated monomial:
$$x\mu _1=1\mu _1^{}=(0).$$
The next step is to consider the only polyhedron on the line $`R^1`$ which is also self-dual, and whose intersection with the integer lattice on the line contains three integer points:
$$\mu _1^{}=(1),\mu _1^{}=(0),\mu _1^{}=(+1).$$
(4.2)
The projective vector corresponding to this linear polyhedron is
$$\stackrel{}{k}_2=(1,1),$$
(4.3)
which can be constructed from the $`\stackrel{}{k}_1`$ vector, by the following procedure.
We extend the vector $`\stackrel{}{k}_1`$ to a two-dimensional vector in $`CP_1`$, by inserting a zero component in all possible ways:
$`\stackrel{}{k}_1^{ex^{}}`$ $`=`$ $`(0,1)`$
$`\stackrel{}{k}_1^{ex^{\prime \prime }}`$ $`=`$ $`(1,0).`$ (4.4)
The following monomials correspond to these ‘extended’ vectors:
$`\mu ^{^{}}=(\nu ,1)`$ $``$ $`x^\nu y`$
$`\mu ^{^{\prime \prime }}=(1,\xi )`$ $``$ $`xy^\xi `$ (4.5)
with the arbitrary integer numbers $`\nu ,\xi `$. From all the possible $`\stackrel{}{k}`$ pairs:
$`(\stackrel{}{k}^{ex^{}}\stackrel{}{k}^{ex^{}}),(\stackrel{}{k}^{ex^{\prime \prime }}\stackrel{}{k}^{ex^{\prime \prime }})(\stackrel{}{k}^{ex^{}}\stackrel{}{k}^{ex^{\prime \prime }}),`$ (4.6)
we select only those whose intersections give a reflexive polyhedron. In this simple two-dimensional case, only a single pair is so selected, namely $`\stackrel{}{k}_1^{ex^{}}`$ and $`\stackrel{}{k}_1^{ex^{\prime \prime }}`$:
$`\stackrel{}{k}_1^{ex^{}}{\displaystyle \stackrel{}{k}_1^{ex^{\prime \prime }}}=\mathrm{\hspace{0.17em}1}.`$ (4.7)
and the reflexive polyhedron comprises just a single point. The corresponding monomial is $`xy`$, whose degree is unity for both variables: $`deg_x=1`$ and $`deg_y=1`$.
We now introduce a second operation on these ‘extended’ vectors $`\stackrel{}{k}^{{}_{}{}^{}\mathrm{}}`$, which is ‘dual’ to the intersection, namely the ‘sum’ operation:
$$\stackrel{}{k}_1^{ex^{}}\stackrel{}{k}_1^{ex^{\prime \prime }}=\stackrel{}{k}_2=(0,1)+(1,0)=(1,1).$$
(4.8)
In this simple case, one has three quadratic monomials:
$`x^2`$ $``$ $`\mu _1=(2,0)\mu _1^{}=(1);`$
$`xy`$ $``$ $`\mu _2=(1,1)\mu _2^{}=(0);`$
$`y^2`$ $``$ $`\mu _3=(0,2)\mu _3^{}=(+1).`$
If a projective vector is multiplied by a positive integer number $`mZ^+`$, it still determines the same hypersurface. Therefore, we should also consider sums of such vectors, characterized by two positive integer numbers, $`m,n`$:
$$m\stackrel{}{k}_1^{ex^{}}+n\stackrel{}{k}_1^{ex^{\prime \prime }}.$$
(4.10)
It turns out that, in order to get a reflexive polyhedron with only one interior point, the numbers $`m`$ and $`n`$ have to be lower than certain maximal values: $`m_{max}`$ and $`n_{max}`$, respectively. In our first trivial example, we find that
$$m_{max}=1,n_{max}=\mathrm{\hspace{0.17em}1}.$$
(4.11)
In general, the set of all pairs $`(m,n)`$ with $`mm_{max}`$ and $`nn_{max}`$ generate a ‘chain’ of possible reflexive polyhedra, which happens to be trivial in this simple case.
Following the previous procedure, to construct all possible vectors on the plane we should start from two vectors, $`\stackrel{}{k}_1`$ and $`\stackrel{}{k}_2`$, ‘extended’ to dimension three in $`CP_2`$ space:
$`\stackrel{}{k}_1^{ex^{}}`$ $`=`$ $`(0,0,1),\stackrel{}{k}_1^{ex^{\prime \prime }}=(0,1,0),\stackrel{}{k}_1^{ex^{\prime \prime \prime }}=(1,0,0);`$
$`\stackrel{}{k}_2^{ex^{}}`$ $`=`$ $`(0,1,1),\stackrel{}{k}_2^{ex^{\prime \prime }}=(1,1,0),\stackrel{}{k}_2^{ex^{\prime \prime \prime }}=(1,0,1).`$ (4.12)
The next step consists of finding all possible pairs of these three-dimensional vectors whose intersection gives the only reflexive polyhedron of dimension two, which corresponds to the polyhedron projective vector $`\stackrel{}{k}_2=(1,1)`$. Only two pairs (plus cyclic permutations) satisfy this constraint:
$`[\stackrel{}{k}_1^{ex^{}}(0,0,1)]{\displaystyle [\stackrel{}{k}_2^{ex^{\prime \prime }}(1,1,0)]}=[\stackrel{}{k}_2(1,1)]_J`$ (4.13)
and
$`[\stackrel{}{k}_2^{ex^{}}(0,1,1)]{\displaystyle [\stackrel{}{k}_2^{ex^{\prime \prime }}(1,1,0)]}`$ $`=`$ $`[\stackrel{}{k}_2(1,1)]_J.`$ (4.14)
In these two cases, the corresponding monomials are:
$`x^2z`$ $``$ $`\mu _1=(2,0,1)(1);`$
$`xyz`$ $``$ $`\mu _2=(1,1,1)(0);`$
$`y^2z`$ $``$ $`\mu _3=(0,2,1)(+1);`$
and
$`x^2z^2`$ $``$ $`\mu _1=(2,0,2)(1);`$
$`xyz`$ $``$ $`\mu _2=(1,1,1)(0);`$
$`y^2`$ $``$ $`\mu _3=(0,2,0)(+1);`$
respectively. These lead to the two following chains:
$`I.\stackrel{}{k}_3(1)`$ $`=`$ $`\mathrm{\hspace{0.17em}1}\stackrel{}{k}_1^{ex^{}}+\mathrm{\hspace{0.17em}1}\stackrel{}{k}_2^{ex^{\prime \prime }}=(1,1,1);m=1,n=1`$
$`\stackrel{}{k}_3(2)`$ $`=`$ $`\mathrm{\hspace{0.17em}2}\stackrel{}{k}_2^{ex^{}}+\mathrm{\hspace{0.17em}1}\stackrel{}{k}_2^{ex^{\prime \prime }}=(1,1,2);m=2,n=1`$
$`m_{max}`$ $`=`$ $`dim(\stackrel{}{k}_2^{ex^{\prime \prime }})=\mathrm{\hspace{0.17em}2},n_{max}=dim(\stackrel{}{k}_2^{ex^{}})=\mathrm{\hspace{0.17em}1}`$ (4.17)
and
$`II.\stackrel{}{k}_3(2)`$ $`=`$ $`\mathrm{\hspace{0.17em}1}\stackrel{}{k}_2^{ex^{}}+\mathrm{\hspace{0.17em}1}\stackrel{}{k}_2^{ex^{\prime \prime \prime }}=(1,1,2);m=1,n=1;`$
$`\stackrel{}{k}_3(3)`$ $`=`$ $`\mathrm{\hspace{0.17em}2}\stackrel{}{k}_2^{ex^{}}+\mathrm{\hspace{0.17em}1}\stackrel{}{k}_2^{ex^{\prime \prime \prime }}=(1,2,3);m=2,n=1,`$
$`m_{max}`$ $`=`$ $`dim(\stackrel{}{k}_2^{ex^{\prime \prime }})=\mathrm{\hspace{0.17em}2},n_{max}=dim(\stackrel{}{k}_2^{ex^{}})=\mathrm{\hspace{0.17em}2}.`$ (4.18)
Where the eldest vectors are given on the first lines of the two preceding equations, and we note that the vector $`(1,1,2)`$ is common to both chains.
It turns out that, also in higher dimensions, some $`\stackrel{}{k}`$ vectors are common to more than one chain. Thus it is possible to make a transition from one chain to another by passing through the common vectors. The algebra of projective vectors with the two operations $``$ and $``$ should be closed under duality symmetry:
$`J\mathrm{\Pi },`$ (4.19)
where the symbols $`J`$ and $`\mathrm{\Pi }`$ denotes two dual conjugate operations: intersection and projection, respectively. In this way, all vectors $`\stackrel{}{k}_d`$ can be found. This structure underpins the idea of a web of transitions between all Calabi-Yau manifolds.
### 4.2 General Formulation of Calabi-Yau Algebra
In the spirit of the simple constructions of the previous subsection, we can also construct the corresponding closed $`\stackrel{}{k}_4`$ algebra in the case of $`K3`$ hypersurfaces. However, before giving the results, we first briefly formulate a theorem underlying the construction of a $`\stackrel{}{k}_{d+1}`$ projective vector, determining an associated reflexive $`d+1`$-dimensional polyhedron and $`CY_d`$ hypersurface, starting from $`\stackrel{}{k}_d`$ projective vectors, which determine a $`d`$-dimensional reflexive polyhedron with one interior point and a corresponding $`CY_{d1}`$ hypersurface. This theorem underlies our systematic inductive algebraic construction of CY manifolds.
The theorem is based on two general points:
* First: from the vector $`\stackrel{}{k}_d`$, we construct the ‘extended’ vectors $`\stackrel{}{k}_{d+1}^{ex}`$ using the rule:
$$()\stackrel{}{k}_d=(k_1,\mathrm{},k_2)\stackrel{\pi ^1}{}\stackrel{}{k}_{d+1}^{ex(i)}=(k_1,\mathrm{},0^i,\mathrm{},k_d).$$
(4.20)
* Second: we consider only those pairs of all possible ‘extended’ vectors, $`\stackrel{}{k}_{d+1}^{ex(i)}`$ and $`\stackrel{}{k}_{d+1}^{ex(j)}`$ with $`0i,jd`$, whose intersection gives the reflexive polyhedron of dimension $`d`$ with only one interior point. We denote this operation by:
$$()\stackrel{}{k}_{d+1}^{ex(i)}\stackrel{}{k}_{d+1}^{ex(j)}=[\stackrel{}{k}_d]_J.$$
(4.21)
The statement of the theorem is:
* If by the rule (\*) one can get, from the projective $`\stackrel{}{k}_d`$-vector, a set of ‘extended’ vectors $`\stackrel{}{k}_{d+1}^{ex(i)}`$, $`0id`$, and for any pair of such “extended” $`\stackrel{}{k}_{d+1}^{ex(i)}`$-vectors the conditions (\**) are fulfilled, then the sum of these two ‘extended’ vectors will give an eldest projective vector $`\stackrel{}{k}_{d+1}`$, which determines a reflexive polyhedron with only one interior point.
* Two finite positive integer numbers, $`n_{max},m_{max}Z_+`$, exist such that any linear combination of two vectors $`\stackrel{}{k}_{d+1}^{i,j}(n,m)`$, with integer coefficients $`mm_{max};nn_{max}`$ produce a CY hypersurface. We call ‘chain’ the set of vectors generated by any such pair of ‘extended’ vectors:
$`p\stackrel{}{k}_{d+1}^{i,j}(n,m)`$ $`=`$ $`m\stackrel{}{k}_{d+1}^{ex(i)}+n\stackrel{}{k}_{d+1}^{ex(j)};`$
$`\stackrel{}{k}_{d+1}^{i,j}(1,1)`$ $`=`$ $`\stackrel{}{k}_{d+1}^{i,j}(eld)`$ (4.22)
* The intersection of the vector $`\stackrel{}{k}_{d+1}^{i,j}(m,n)`$ with the vector $`\stackrel{}{k}_{d+1}^{ex(i)}`$ is equal to the intersection of this vector with the vector $`\stackrel{}{k}_{d+1}^{ex(j)}`$:
$`[\stackrel{}{k}_{d+1}^{i,j}(m,n)]{\displaystyle [\stackrel{}{k}_{d+1}^{ex(j)}]}=[\stackrel{}{k}_{d+1}^{i,j}(m,n)]{\displaystyle [\stackrel{}{k}_{d+1}^{ex(i)}]}.`$ (4.23)
We can also formulate a converse theorem:
* If one can decompose a reflexive projective vector $`\stackrel{}{k}_{d+1}`$ as the sum of two reflexive projective vectors $`\stackrel{}{k}_{d+1}^{^{}}`$ and $`\stackrel{}{k}_{d+1}^{\mathrm{`}\mathrm{`}}`$, then there exists the intersection of the vector $`\stackrel{}{k}_{d+1}`$ with either of these two vectors, which defines a projective vector $`\stackrel{}{k}_d`$ and a reflexive polyhedron with only one interior point.
The above theorem provides a description of all $`CY_{d+1}`$ hypersurfaces with $`d`$-dimensional fibers in terms of two positive-integer parameters. Similarly, one can also consider the intersections of three (or more) ‘doubly-extended’ vectors $`\stackrel{}{k}_{d+1}^{ex(^{})}`$, $`\stackrel{}{k}_{d+1}^{ex(^{\prime \prime })}`$, $`\stackrel{}{k}_{d+1}^{ex(^{\prime \prime \prime })}`$ (by ‘doubly-extended’ we mean that they may be obtained by inserting two zero components in $`\stackrel{}{k}_{d1}`$ vectors). One should check that this intersection gives a reflexive polyhedron in the $`d2`$ space:
$`[\stackrel{}{k}_{d1}^{ex(2^{^{}})}]{\displaystyle [\stackrel{}{k}_{d1}^{ex(2^{^{\prime \prime }})}][\stackrel{}{k}_{d1}^{ex(2^{^{\prime \prime \prime }})}]}=[\stackrel{}{k}_{d1}]_J.`$ (4.24)
In this way, one may obtain a $`3,4,\mathrm{},d`$ positive-integer parameter description of the $`(d+1)`$-dimensional polyhedra with $`(d1),(d2),\mathrm{}`$-dimensional fiber sections:
$`p\stackrel{}{k}_{d+1}=m\stackrel{}{k}_{d1}^{ex(2^{^{}})}+n\stackrel{}{k}_{d1}^{ex(2^{^{\prime \prime }})}+l\stackrel{}{k}_{d1}^{ex(2^{^{\prime \prime \prime }})}.`$ (4.25)
Finally, one can obtain additional lists of $`\stackrel{}{k}_{d+1}`$ vectors by using three ‘extended’ vectors, $`\stackrel{}{k}_d^{ex^r}`$, $`\stackrel{}{k}_d^{ex^i}`$ $`\stackrel{}{k}_d^{ex^j}`$ (and similarly using four $`\stackrel{}{k}_{\mathrm{}}^{ex}`$, etc.), and a special algebra of summing these vectors only if the following three conditions are fulfilled:.
$`1.[\stackrel{}{k}_d^{ex^r}]{\displaystyle [\stackrel{}{k}_d^{ex^i}]}`$ $`=`$ $`[\stackrel{}{k}_{d1}]_J^{^{}};`$
$`2.[\stackrel{}{k}_d^{ex^i}]{\displaystyle [\stackrel{}{k}_d^{ex^j}]}`$ $`=`$ $`[\stackrel{}{k}_{d1}]_J^{^{\prime \prime }};`$
$`3.[\stackrel{}{k}_d^{ex^j}]{\displaystyle [\stackrel{}{k}_d^{ex^r}]}`$ $`=`$ $`[\stackrel{}{k}_{d1}]_J^{^{\prime \prime \prime }}.`$
In this way, one may obtain a complete description of the positive-integer lattice which defines all reflexive $`\stackrel{}{k}`$ vectors.
## 5 Two-Vector Chains of $`K3`$ Spaces
We now embark on a parametrization of the $`\stackrel{}{k}_4`$ vectors defining $`K3`$ hypersurfaces with fiber sections. Our description of $`K3`$ hypersurfaces is based on the above understanding of the composite and dual structure of the projective $`\stackrel{}{k}_4`$ vectors. As already mentioned, we find a link between this structure and the finite subgroups of the group of rotations of three-space, namely the cyclic and dihedral groups and the symmetry groups of the Platonic solids: the tetrahedron, the octahedron-cube and the icosahedron-dodecahedron:
* $`C_n:n=1,2,3,\mathrm{}`$, the cyclic group of finite rotations in the plane around an axis ‘1’ by the angles $`\alpha =2\pi /n`$;
* $`D_n:n=2,3,4,\mathrm{}`$, the dihedral group, comprising all these rotations together with the all reflections of a second axis ‘$`n`$’ lying in this plane, which is orthogonal to the axis ‘l’, and producing with respect to each other the angle $`\alpha /2`$;
* T -The finite group of the transformations leaving invariant the regular tetrahedron, with 12 parameters;
* O- The finite group of the transformations leaving invariant the regular cube and octahedron, with 24 parameters;
* I- The finite group of the transformations leaving invariant the regular icosahedron and dodecahedron, with 60 parameters.
We use the polyhedron technique introduced in the previous Section, taking into account all its duality, intersection and projection properties to study the projective-vector classification of $`K3`$ spaces.
### 5.1 Two-Dimensional Integer Chains of $`K3`$ Hypersurfaces
In the $`K3`$ case, as already foreshadowed in the Introduction, the classification can start from a basis of five types of ‘extended’ vectors. We recall that the structure of the three ‘planar’ projective vectors $`\stackrel{}{k}_3=(1,1,1),(1,1,2),(1,2,3)`$ can easily be understood on the basis of the doubly-extended vector $`\stackrel{}{k}_1^{ext}=(0,0,1)`$ and the singly-extended vector $`\stackrel{}{k}_2^{ext}=(0,1,1)`$. The structure of the underlying composite vector $`\stackrel{}{k}_2=(1,1)`$ is also obvious. The full list of $`K3`$ projective vectors is obtainable from the algebra of the following five extended vectors: the maximally-extended vector of the form
$$\stackrel{}{k}_C^{ext}=(0,0,0,1)$$
(5.1)
with its 4 cyclic permutations, the doubly-extended dihedral vector of the form
$$\stackrel{}{k}_D^{ext}=(0,0,1,1)$$
(5.2)
with its 6 dihedral permutations, the singly-extended tetrahedral vector of the form
$$\stackrel{}{k}_T^{ext}=(0,1,1,1)$$
(5.3)
with its 4 cyclic permutations, the extended octahedral vector of the form
$$\stackrel{}{k}_O^{ext}=(0,1,1,2)$$
(5.4)
with its 12 permutations, and finally the extended icosahedral vector of the form
$$\stackrel{}{k}_I^{ext}=(0,1,2,3)$$
(5.5)
with its 24 permutations, for a total of 50 extended vectors.
Using the algebra of combining pairs of these 50 extended $`\stackrel{}{k}^i`$ vectors, we obtain 90 distinct $`\stackrel{}{k}_4`$ vectors in 22 double chains with a regular planar $`k`$-gon intersection: $`k>3`$ with only one interior point, as seen in Table 1. Combining three extended $`\stackrel{}{k}^i`$ vectors, we obtain four triple chains with self-dual line-segment intersection-projections and one interior point, which contain 91 distinct vectors, of which only four $`\stackrel{}{k}_4`$ vectors are different from the 90 vectors found previously, as also seen in Table 1. Of course, there are some vectors which have a regular planar $`k`$-gon in their intersection and no line-segment intersection. Further, as we see later in Section 7, there is just one vector, $`\stackrel{}{k}_4=(7,8,9,12)`$, which has only a single point intersection, i.e., the intersection consists of the zero point alone, and can be determined by the intersection-projection
$`J(\mathrm{\Delta })`$ $``$ $`\mathrm{\Pi }(\mathrm{\Delta }^{})`$
$`\mathrm{\Pi }(\mathrm{\Delta })`$ $``$ $`J(\mathrm{\Delta }^{})`$ (5.6)
duality, where the polyhedra $`\mathrm{\Delta }`$ and $`\mathrm{\Delta }^{}`$ determine a dual pair of $`K3`$ hypersurfaces. We recall that the sum of the integer points in intersection, $`J(\mathrm{\Delta })`$, and in projection, $`\mathrm{\Pi }(\mathrm{\Delta }^{})`$, is equal to $`14=4+10,5+9,6+8,7+7,8+6,9+5,10+4`$ for the plane intersection-projection and $`6=3+3`$ for the line-segment intersection-projection. This duality plays a very important role in our description. Eleven of the 22 two-vector chains found previously satisfy directly the following condition:
$`J(\mathrm{\Delta }^n)`$ $`=`$ $`\mathrm{\Pi }(\mathrm{\Delta }^n)=\mathrm{\Delta }^{n1},`$
$`\mathrm{\Pi }(\mathrm{\Delta }_{}^{n}{}_{}{}^{})`$ $`=`$ $`J(\mathrm{\Delta }_{}^{n}{}_{}{}^{})=\mathrm{\Delta }_{}^{n1}{}_{}{}^{},`$ (5.7)
which means that the number of integer points in the intersection of the polyhedron (mirror polyhedron) forming the reflexive polyhedron of lower dimension is equal to the number of projective lines crossing these integer points of the polyhedron (mirror). The projections of these lines on a plane in the polyhedron and a plane in its mirror polyhedron reproduce, of course, the reflexive polyhedra of lower dimension. Only for self-dual polyhedra can one have
$$J(\mathrm{\Delta }^n)=\mathrm{\Pi }(\mathrm{\Delta }^n)=\mathrm{\Pi }(\mathrm{\Delta }_{}^{n}{}_{}{}^{})=J(\mathrm{\Delta }_{}^{n}{}_{}{}^{})=\mathrm{\Delta }^{n1}=\mathrm{\Delta }_{}^{n1}{}_{}{}^{},$$
(5.8)
namely the most symmetrical form of these relations.
Following the recipe presented as our central Theorem in Section 4, we present Table 5, which lists all the $`\stackrel{}{k}_4`$ projective vectors derived from pairs of extended vectors of lower dimension, which fall into the 22 chains listed. In each case, we list the maximum integers $`m,n`$ in the chains, which are determined by the dimensions of the extended $`\stackrel{}{k}_i`$ vectors. This Table includes all the 90 projective $`\stackrel{}{k}_4`$ vectors found using our construction. All of these $`\stackrel{}{k}_4`$ vectors define $`K3`$ hypersurfaces which could be obtained using the $`Z^n`$ symmetry coset action They may also be used to construct higher-level $`CY_1`$ spaces as the intersections of polynomial loci, as discussed in Section 9..
For illustration, we give in Table 6 the eldest vectors in each chain, i.e., the first members of all 22 chains, which have $`m=1,n=1`$. As one can see, some vectors are common to more than one chain. Using our understanding of the origins of the intersections, and duality, we can classify these 22 chains in five classes, as indicated by the groupings in Table 6, which correspond to the intersections, as indicated.
It should be noted, however, that the above doubly-extended vector structure does not exhaust the full list of possible $`K3`$ projective vectors. The projective vectors
$`(\stackrel{}{k}_4)_{91}`$ $`=`$ $`(4,5,7,9),`$
$`(\stackrel{}{k}_4)_{92}`$ $`=`$ $`(5,6,8,11),`$
$`(\stackrel{}{k}_4)_{93}`$ $`=`$ $`(5,7,8,20),`$
$`(\stackrel{}{k}_4)_{94}`$ $`=`$ $`(7,8,10,25),`$
$`(\stackrel{}{k}_4)_{95}`$ $`=`$ $`(7,8,9,12),`$ (5.9)
have no planar reflexive polyhedron intersections, and therefore were not included in this list. To obtain most of the additional $`\stackrel{}{k}_4`$\- vectors (5.9), one must consider chains constructed from three extended vectors of the type $`\stackrel{}{k}^{ex}=(0,0,1)`$ and $`\stackrel{}{k}^{ex}=(0,1,1)`$, with all possible permutations, having in the intersection the line-segment polyhedron consisting of three integer points. All these chains will be $`J_1\mathrm{\Pi }_1`$ self-dual: $`J_1=\mathrm{\Pi }_1=3`$. It is easy to see that only four different such triple chains can be built, as discussed in Section 6. These chains are much longer than the previous two-vector chains, although their total number, 91, is also less than the full number of all $`K3`$ vectors. The projective vectors
$`(\stackrel{}{k}_4)_{12}`$ $`=`$ $`(3,5,6,7),`$
$`(\stackrel{}{k}_4)_{13}`$ $`=`$ $`(3,6,7,8),`$
$`(\stackrel{}{k}_4)_{14}`$ $`=`$ $`(5,6,7,9),`$
$`(\stackrel{}{k}_4)_{95}`$ $`=`$ $`(7,8,9,12)`$ (5.10)
are not involved in these chains. However, the union of the doubly-extended and triply-extended vector chains gives a total of 94 $`\stackrel{}{k}_4`$ projective vectors. Only the $`\stackrel{}{k}_4=(7,8,9,12)`$ vector has just a point-intersection structure, and is not found by either the double- or triple-vector constructions, as discussed in more detail in Section 7.
To preview how it arises, note that, by $`J\mathrm{\Pi }`$ duality, we know that to $`\stackrel{}{k}=(1,1,1,1)`$, which has three intersection planes (1,1,1) with ten points, there must correspond a $`\stackrel{}{k}`$ which has three different $`\pi `$ projections with four points. Since it should have a non-trivial projection structure, namely a four-point planar polyhedron with one interior point in three independent directions, its external points should satisfy the following condition:
$`{\displaystyle \frac{1}{4}}\{M_1+M_2+M_3+M_4\}=M_0=(1,1,1,1).`$ (5.11)
In three-space, these points can only be taken as:
$`M_1=(4,1,0,0),M_2=(0,3,1,0),M_3=(0,0,4,0),M_4=(0,0,0,3).`$ (5.12)
One can easily check that this polyhedron has three projections: $`\pi _{x_1}`$, $`\pi _{x_2}`$, $`\pi _{x_3}`$, with four points giving the $`(1,1,1)`$ planar polyhedron. The four points $`M_i`$ (5.12) give the exceptional vector $`\stackrel{}{k}=(7,8,9,12)`$. By projection, one can see that the five integer points of this polyhedron produce the $`(1,1,1)`$ planar polyhedron with four points.
### 5.2 Invariant Monomials and the $`J\mathrm{\Pi }`$ Structure of Calabi-Yau Equations
The experience provided by working with $`K3`$ hypersurfaces can aid in the classification of Calabi-Yau manifolds. Also for this more complicated case, one should use the duality conditions: one must be prepared to study the intersection structures of polyhedra and their mirrors and/or to study the projection structures for polyhedra and mirror polyhedra.
This ‘intersection-projection’ structure of the $`\stackrel{}{k}_4`$ vectors from doubly-, triply- and quadruply-extended vectors allows us to introduce the concept of invariant monomials in the CY equations. These invariant monomials are homogeneous under the action of the extended vectors, i.e., if
$`\stackrel{}{z^{}}=\lambda ^{\stackrel{}{k}_j^{ex}}\stackrel{}{z},j=1,2,3,\mathrm{},`$ (5.13)
then
$`\stackrel{}{z^{}}^\stackrel{}{\mu }=\lambda ^{\stackrel{}{k}_j^{ex}\stackrel{}{\mu }}\stackrel{}{z}^\stackrel{}{\mu }=\lambda ^{d_j}\stackrel{}{z}^\stackrel{}{\mu },`$ (5.14)
where $`d_j=dim(\stackrel{}{k}_j^{ex})`$ and $`j=1,2,3,\mathrm{}`$ is the number of extended $`\stackrel{}{k}_j^{ex}`$ vectors. The invariant monomials, $`\mathrm{}_i`$, correspond to the reflexive polyhedra produced by the invariant set $`\mathrm{\Psi }_{inv}`$ which is the same for all the chains.
These extended vectors can be formed from the following five familiar types of projective vectors of lower dimensions:
$`\stackrel{}{k}_1`$ $`=`$ $`(0,0,1),`$
$`\stackrel{}{k}_1`$ $`=`$ $`(0,1,1),`$
$`\stackrel{}{k}_3`$ $`=`$ $`(1,1,1),\stackrel{}{k}_3=(1,1,2),\stackrel{}{k}_3=(1,2,3).`$ (5.15)
A chain of $`\stackrel{}{k}_4`$ projective vectors can be generated from the linear sums of extended vectors, for example, for $`j=1,2`$ one can get:
$`\stackrel{}{k}_4(m,n)`$ $`=`$ $`m\stackrel{}{k}_1^{ex}+n\stackrel{}{k}_2^{ex}`$
$`\mathrm{if}\stackrel{}{k}_1^{ex}{\displaystyle \stackrel{}{k}_2^{ex}}`$ $`=`$ $`\{\mathrm{}_i:\mathrm{}_i\mathrm{\Sigma }_{inv}\}.`$ (5.16)
The invariant monomials are universal for all the $`\stackrel{}{k}_4`$ vectors in this chain.
To construct the $`\stackrel{}{k}_4`$ vectors determining $`K3`$ hypersurfaces, i.e., determining the corresponding polyhedra with the property of reflexivity, one has to give a correct set of invariant monomials. We have constructed the 22 sets of invariant monomials corresponding to the doubly-extended vector structures among the $`\stackrel{}{k}_4`$ projective vectors. In this case, these sets of the invariant monomials give in the intersection reflexive polyhedra of lower dimensions. The number of invariant monomials for this doubly-extended vector structure is given by
$$31=\mathrm{\hspace{0.17em}\hspace{0.17em}1}+\mathrm{\hspace{0.17em}4}\times 2+\mathrm{\hspace{0.17em}22},$$
(5.17)
where the last number corresponds to the Betti number for $`K3`$ hypersurfaces: $`b_2=22`$. The structure of the $`\stackrel{}{k}_4`$ projective vectors obtained from the triply-extended vectors, namely $`\stackrel{}{k}^{ex}=(0,0,0,1)`$ and $`\stackrel{}{k}^{ex}=(0,0,1,1)`$, is given by the following four types of invariant monomials:
$`\mathrm{\Psi }_{I_3}:(2,0,1,1),(0,2,1,1),(1,1,1,1,)`$ $``$ $`x^2zu,y^2zu,xyzu;`$
$`\mathrm{\Psi }_{II_3}:(2,2,1,0),(0,0,1,2),(1,1,1,1,)`$ $``$ $`x^2y^2z,zu^2,xyzu;`$
$`\mathrm{\Psi }_{III_3}:(2,2,2,0),(0,0,0,2),(1,1,1,1,)`$ $``$ $`x^2y^2z^2,u^2,xyzu;`$
$`\mathrm{\Psi }_{IV_3}:(2,0,0,2),(0,2,2,0),(1,1,1,1,)`$ $``$ $`x^2u^2,y^2z^2,xyzu.`$ (5.18)
The four chains corresponding to these sets of invariant monomials are (see Tables 5,6,7 and 8):
$`\stackrel{}{k}_4(\mathrm{\Psi }_{I_3})`$ $`=`$ $`M(1,1,0,0)+N(0,0,1,0)+L(0,0,0,1)=(M,M,N,L),`$
$`\stackrel{}{k}_4(\mathrm{\Psi }_{II_3})`$ $`=`$ $`M(1,0,0,1)+N(0,1,0,1)+L(0,0,1,0)=(M,N,L,M+N),`$
$`\stackrel{}{k}_4(\mathrm{\Psi }_{III_3})`$ $`=`$ $`M(1,0,0,1)+N(0,1,0,1)+L(0,0,1,1)=(M,N,L,M+N+L),`$
$`\stackrel{}{k}_4\mathrm{\Psi }_{(IV_3})`$ $`=`$ $`M(1,0,1,0)+N(0,1,0,1)+L(0,0,1,1)=(M,N,M+L,N+L).`$
In these chains There is in fact another ‘good’ triple intersection, of the extended vectors $`(1,1,0,0),(0,0,1,1),(0,0,0,1)`$, but the chain $`I_3^{}=(M,M,N,N+L)`$ it produces has the same three invariant monomials, $`(0,2,1,1)+(2,0,1,1)+(1,1,1,1)`$ as the $`I_3`$ chain, which includes all its projective vectors., the sets of projective vectors are subject to the following additional projective restrictions:
$`\stackrel{}{k}_4(\mathrm{\Psi }_{I_3})\stackrel{}{e}_I`$ $`=`$ $`\mathrm{\hspace{0.17em}0},\stackrel{}{e}_I=(1,1,0,0)`$
$`\stackrel{}{k}_4(\mathrm{\Psi }_{II_3})\stackrel{}{e}_{II}`$ $`=`$ $`\mathrm{\hspace{0.17em}0},\stackrel{}{e}_{II}=(1,1,0,1)`$
$`\stackrel{}{k}_4(\mathrm{\Psi }_{III_3})\stackrel{}{e}_{III}`$ $`=`$ $`\mathrm{\hspace{0.17em}0},\stackrel{}{e}_{III}=(1,1,1,1)`$
$`\stackrel{}{k}_4(\mathrm{\Psi }_{IV_3})\stackrel{}{e}_{IV}`$ $`=`$ $`\mathrm{\hspace{0.17em}0},\stackrel{}{e}_{IV}=(1,1,1,1)`$
Corresponding to these chains, the following triple intersections
$`\stackrel{}{k}_M^{ex}{\displaystyle \stackrel{}{k}_N^{ex}\stackrel{}{k}_L^{ex}}=\mathrm{\Psi }_{I_3},\mathrm{\Psi }_{II_3},\mathrm{\Psi }_{III_3},\mathrm{\Psi }_{IV_3}.`$ (5.21)
have the above-mentioned invariant monomials.
The $`K3`$ algebra has the interesting consequence that all the $`\{1+4+22\}`$ invariant monomials that give ‘good’ planar reflexive polyhedra in the 22 two-vector chains also can be found by triple constructions. Therefore it is interesting to list now the 22 types of invariant monomials whose origin is also connected with the triple intersections of all types of projective vectors, the triply-extended vectors $`\stackrel{}{k}_1^{ex}=(0,0,0,1)`$, the doubly-extended vectors $`\stackrel{}{k}_2^{ex}=(0,0,1,1)`$, and the singly-extended vectors, $`\stackrel{}{k}_3^{ex}=(0,1,1,1),(0,1,1,2),(0,1,2,3)`$.
These monomials, $`\stackrel{}{z}^\stackrel{}{\mu }`$, are invariant under action of the extended vectors
$`\stackrel{}{k}_i^{ex}\stackrel{}{\mu }`$ $`=`$ $`dim(\stackrel{}{k}_i^{ex}),`$
$`\stackrel{}{k}_j^{ex}\stackrel{}{\mu }`$ $`=`$ $`dim(\stackrel{}{k}_j^{ex}),`$
$`\stackrel{}{k}_l^{ex}\stackrel{}{\mu }`$ $`=`$ $`dim(\stackrel{}{k}_l^{ex}).`$ (5.22)
The directions of the possible projections $`\mathrm{\Pi }`$ are determined <sup>\**</sup><sup>\**</sup>\**Additional constraints on the invariant monomials are given in Section 7, reducing their number to 9 = 1 + 3 + 5. by the degenerate monomial $`(xyzu)\stackrel{}{\mu }=(1,1,1,1)`$ and by the exponents of the following 22 invariant monomials, $`\mu =(\mu _1,\mu _2,\mu _3,\mu _4)`$:
$`\underset{¯}{(3,0,0,0),(3,1,0,0),(3,1,1,0),}(3,2,0,0),`$
$`(3,2,1,0),(3,3,0,0),(3,3,1,0),`$
$`\underset{¯}{(4,0,0,0),(4,1,0,0),}(4,1,1,0),(4,2,0,0),`$
$`(4,2,1,0),(4,3,0,0),(4,3,1,0),(4,4,0,0),(4,4,1,0),`$
$`(6,0,0,0),(6,1,0,0),(6,2,0,0),(6,3,0,0),`$
$`(6,4,0,0),(6,6,0,0).`$ (5.23)
where the underlines pick out those triple intersections where the intersections of pairs of vectors also specify reflexive polyhedra, which will be important later. The four other types of possible projections were already defined above.
The algebraic-geometry sense of $`(J,\mathrm{\Pi })(\mathrm{\Delta })(\mathrm{\Pi },J)(\mathrm{\Delta }^{})`$ duality for $`K3`$ hypersurfaces can be interpreted through the invariant monomials: the list of the invariant monomials for the two-extended-vector classification and the list of all of the three-extended-vector classification are the same, and the total number of them is equal to $`31=1+4\times 2+22`$. The $`J(\mathrm{\Delta },\mathrm{\Delta }^{})\mathrm{\Pi }(\mathrm{\Delta }^{},\mathrm{\Delta })`$ duality can be interpreted at a deeper level for $`J=\mathrm{\Pi }`$ chains: the invariant monomials are identical for corresponding CY equation and for its mirror equation. The projection-projection structure gives additional information about the form of the corresponding CY equation. For example, this structure determines the subset of monomials corresponding to the invariant monomials. As result, the homogeneous CY equation can be written in according in terms of the intersection-projection structure of the projective $`\stackrel{}{k}`$ vectors:
$`\mathrm{}(\stackrel{}{z})={\displaystyle \underset{i}{\overset{J}{}}}\stackrel{}{z}^{\stackrel{}{m}_0^i}\{{\displaystyle \underset{p}{\overset{\mathrm{\Pi }}{}}}a_{\stackrel{}{m}_0^i}^p\stackrel{}{z}^{n_p\stackrel{}{e}^\mathrm{\Pi }}\}=\mathrm{\hspace{0.17em}0}.`$ (5.24)
Here the $`\stackrel{}{z}^{\stackrel{}{m}_0^i}`$ are the invariant monomials which are defined by intersection structure, the vector $`\stackrel{}{e}^\mathrm{\Pi }`$ is the direction of the projection, and the $`n_p`$ are integer numbers.
## 6 Three-Vector Chains of $`K3`$ Spaces
As already mentioned, one can find additional chains of $`K3`$ projective vectors $`\stackrel{}{k}_4`$ if one considers systems of three extended vectors of the type $`\stackrel{}{k}_1^{ex}=(0,0,0,1)`$ and $`\stackrel{}{k}_2^{ex}=(0,0,1,1)`$, which have in their intersections only three integer points or only three invariant monomials. As also already remarked, there are only four different chains, corresponding to the four kinds of invariant monomial triples. We have also commented that these new chains yield only four additional $`K3`$ vectors, whilst the remaining vector, $`\stackrel{}{k}_4=(7,8,9,12)`$, can be constructed out of four extended vectors, as discussed in the following Section. The relationship between the two- and three-vector constructions, and their substantial overlap, is the subject of this Section.
### 6.1 The Three-Vector Chain $`I_3`$: $`\stackrel{}{k}_4=(M,M,N,L)`$
In this chain, the dimension $`(d=2M+N+L)`$ and the eldest vector is $`\stackrel{}{k}_{eld}=(1,1,1,1)`$, whose invariant monomials are $`(2,0,1,1)+(0,2,1,1)`$. The relations between this three-vector chain and the previously-discussed two-vector chains can easily be found. We consider the first three vectors in Table 7, which also form the two-vector chain $`I`$:
$`I:`$ $`m(1,1,1,0)+n(0,0,0,1)=(m,m,m,n)`$ (6.1)
$`M=N=m=[dim]\{(0,0,0,1)\}=\mathrm{\hspace{0.17em}1},`$
$`L=n[dim]\{(1,1,1,0)\}=\mathrm{\hspace{0.17em}3}.`$
Similarly, one can consider four vectors $`(2,2,1,1)`$, $`(3,3,1,2)`$, $`(4,4,1,3)`$ and $`(5,5,2,3)`$, which form the two-vector chain $`II`$:
$`II:`$ $`m(1,1,1,0)+n(1,1,0,1)=(m,n,m+n,m+n)`$
$`N=m[dim]\{(1,1,0,1)\}=\mathrm{\hspace{0.17em}3},`$
$`L=n[dim]\{(1,1,1,0)\}=\mathrm{\hspace{0.17em}3},`$
$`M=m+n<\mathrm{\hspace{0.17em}6}.`$
The four vectors $`(1,1,2,1)`$, $`(1,1,2,2)`$, $`(1,1,2,3)`$ and $`(1,1,2,4)`$ from the two-vector chain $`IV`$ have the following relations with this triple chain:
$`IV:`$ $`m(1,1,2,0)+n(0,0,0,1)=(m,m,2m,n)`$
$`M=m[dim]\{(0,0,0,1)\}=\mathrm{\hspace{0.17em}1},`$
$`N=\mathrm{\hspace{0.17em}2}m=\mathrm{\hspace{0.17em}2},`$
$`L=n[dim]\{(1,1,2,0)\}=\mathrm{\hspace{0.17em}4}.`$
The six vectors $`(1,1,1,3)`$, $`(1,1,2,4)`$, $`(1,1,3,5)`$, $`(1,1,4,6)`$, $`(2,2,1,5)`$ and $`(2,2,3,7)`$ in Table 7 correspond to the two-vector chain $`V`$:
$`V:`$ $`m(1,1,0,2)+n(0,0,1,1)=(m,m,n,2m+n)`$
$`M=m[dim]\{(0,0,1,1)\}=\mathrm{\hspace{0.17em}2},`$
$`N=n[dim]\{(1,1,0,2)\}=\mathrm{\hspace{0.17em}4},`$
$`L=\mathrm{\hspace{0.17em}2}m+n<=\mathrm{\hspace{0.17em}8}.`$
The next three vectors $`(1,1,1,1)`$, $`(3,3,2,4)`$ and $`(2,2,1,3)`$ from the two-vector chain $`VII`$ have the following connection to this triple chain:
$`VII:`$ $`m(1,1,2,0)+n(1,1,0,2)=(m+n,m+n,2m,2n)`$
$`M=m+n<\mathrm{\hspace{0.17em}4},`$
$`N=\mathrm{\hspace{0.17em}2}m<\mathrm{\hspace{0.17em}4},`$
$`L=\mathrm{\hspace{0.17em}2}n<\mathrm{\hspace{0.17em}4}.`$
Two vectors $`(1,1,1,1)`$ and $`(1,1,2,2)`$ correspond to the two-vector chain $`X`$:
$`X:`$ $`m(1,1,0,0)+n(0,0,1,1)=(m,m,n,n)`$
$`M=m\mathrm{\hspace{0.17em}2},`$
$`N=n\mathrm{\hspace{0.17em}2},`$
$`L=n\mathrm{\hspace{0.17em}2}.`$
Finally, the values of $`M,N,L`$ of the five projective vectors $`(1,1,1,2)`$, $`(1,1,2,3)`$, $`(1,1,3,4)`$, $`(2,2,1,3)`$ and $`(2,2,3,5)`$ correspond to the fact that they are also from the two-vector chain $`XI`$:
$`XI:`$ $`m(1,1,0,1)+n(0,0,1,1)=(m,m,n,m+n)`$
$`M=m\mathrm{\hspace{0.17em}2},`$
$`N=n\mathrm{\hspace{0.17em}3},`$
$`L=m+n\mathrm{\hspace{0.17em}5}.`$
### 6.2 The Three-Vector Chain $`II_3`$: $`\stackrel{}{k}_4=(M,N,L,M+N)`$
In this chain, shown in Table 8, the dimension $`d=2M+2N+L`$, there is a symmetry: $`MN`$, the eldest vector $`\stackrel{}{k}_{eld}=(1,1,1,2)`$, and the invariant monomials are $`(2,2,1,0)+(0,0,1,2)`$. Comparing this chain with the previous two-vector chains, one can see clearly the possible values of $`M,N,L`$ for the projective vectors $`(M,N,L,M+N)`$. For example, if one compares the four vectors $`(1,1,2,2)`$, $`(1,2,3,3)`$, $`(1,3,4,4)`$ and $`(2,3,5,5)`$ in this triple chain with their structure in the two-vector chain $`II`$, one finds the following relations:
$`II:`$ $`m(1,0,1,1)+n(0,1,1,1)=(m,n,m+n,m+n)`$
$`M=m[dim]\{(0,1,1,1)\}=\mathrm{\hspace{0.17em}3},`$
$`N=n[dim]\{(1,0,1,1)\}=\mathrm{\hspace{0.17em}3},`$
$`L=m+n<\mathrm{\hspace{0.17em}6}.`$
Similarly, we find the following relations between the values of $`M,N,L`$ in the triple chain and the values of $`m,n`$ for double chains:
$`IV:`$ $`m(1,1,0,2)+n(0,0,1,0)=(m,m,n,2m)`$
$`M=N=m[dim]\{(0,0,1,0)\}=\mathrm{\hspace{0.17em}1},`$
$`L=n[dim]\{(1,1,0,2)\}=\mathrm{\hspace{0.17em}4}.`$
$`VI:`$ $`m(1,0,2,1)+n(0,1,2,1)=(m,n,2m+2n,m+n)`$
$`M=m[dim]\{(0,1,2,1)\}=\mathrm{\hspace{0.17em}4},`$
$`N=n[dim]\{(1,0,2,1)\}=\mathrm{\hspace{0.17em}4},`$
$`L=\mathrm{\hspace{0.17em}2}m+2n<\mathrm{\hspace{0.17em}8}.`$
$`VIII:`$ $`m(1,0,2,1)+n(1,1,0,2)=(m+n,n,2m,m+2n)`$
$`M=m+n\mathrm{\hspace{0.17em}8},`$
$`N=n[dim]\{(1,0,2,1)\}=\mathrm{\hspace{0.17em}4},`$
$`L=\mathrm{\hspace{0.17em}2}m\mathrm{\hspace{0.17em}2}[dim]\{(1,1,0,2)\}=\mathrm{\hspace{0.17em}8}.`$
$`XI:`$ $`m(1,0,1,1)+n(0,1,0,1)=(m,n,m,m+n)`$
$`M=m[dim]\{(0,1,0,1)\}=\mathrm{\hspace{0.17em}2},`$
$`N=n[dim]\{(1,0,1,1)\}=\mathrm{\hspace{0.17em}3},`$
$`L=m.`$
$`XIII:`$ $`m(1,0,2,1)+n(0,1,1,1)=(m,n,m,m+n)`$
$`M=m[dim]\{(0,1,1,1)\}=\mathrm{\hspace{0.17em}3},`$
$`N=n[dim]\{(1,0,2,1)\}=\mathrm{\hspace{0.17em}4},`$
$`L=\mathrm{\hspace{0.17em}2}m+n.`$
$`XIV:`$ $`m(1,0,2,1)+n(1,2,0,3)=(m+n,2n,2m,m+3n)`$
$`M=m+n,`$
$`N=\mathrm{\hspace{0.17em}2}n\mathrm{\hspace{0.17em}2}[dim]\{(1,0,2,1)\}=\mathrm{\hspace{0.17em}8},`$
$`L=\mathrm{\hspace{0.17em}2}m\mathrm{\hspace{0.17em}2}[dim]\{(1,2,0,3)\}=\mathrm{\hspace{0.17em}12}.`$
$`XV:`$ $`m(1,2,0,3)+n(0,0,1,0)=(m,2m,n,3m)`$
$`M=m,[dim]\{(0,0,1,0)\}=\mathrm{\hspace{0.17em}1},`$
$`N=\mathrm{\hspace{0.17em}2}m\mathrm{\hspace{0.17em}2}[dim]\{(0,0,1,0)\}=\mathrm{\hspace{0.17em}2},`$
$`L=n[dim]\{(1,2,0,3)\}=\mathrm{\hspace{0.17em}6}.`$
$`XVII:`$ $`m(1,2,0,3)+n(0,1,1,1)=(m,2m+n,n,3m+n)`$
$`M=m,[dim]\{(0,1,1,1)\}=\mathrm{\hspace{0.17em}3},`$
$`N=\mathrm{\hspace{0.17em}2}m+n`$
$`L=n[dim]\{(1,2,0,3)\}=\mathrm{\hspace{0.17em}6}.`$
$`XXII:`$ $`m(1,0,2,1)+n(0,1,0,1)=(m,n,2m,m+n)`$
$`M=m[dim]\{(0,1,0,1)\}=\mathrm{\hspace{0.17em}2},`$
$`N=n[dim]\{(1,0,2,1)\}=\mathrm{\hspace{0.17em}6},`$
$`L=\mathrm{\hspace{0.17em}2}m.`$
### 6.3 The Three-Vector Chain $`III_3`$: $`\stackrel{}{k}_4=(M,N,L,M+N+L)`$
In this chain, tshown in Table 9, he dimension $`d=2M+2N+2L`$, there is $`MNL`$ symmetry, the eldest vector $`\stackrel{}{k}_{eld}=(1,1,1,3)`$, and the invariant monomials are $`(2,2,2,0)+(0,2,2,2)`$. We see in the Table the appearance of the following two-vector chains
$`V:`$ $`m(1,1,0,2)+n(0,0,1,1)=(m,m,n,2m+n)`$
$`M=N=m[dim]\{(0,0,1,1)\}=\mathrm{\hspace{0.17em}2},`$
$`L=n[dim]\{(0,0,1,1)\}=\mathrm{\hspace{0.17em}4}.`$
$`VI:`$ $`m(1,0,1,2)+n(0,1,1,2)=(m,n,m+n,2m+2n)`$
$`M=m[dim]\{(0,1,1,2)\}=\mathrm{\hspace{0.17em}4},`$
$`N=n[dim]\{(1,0,1,2)\}=\mathrm{\hspace{0.17em}4},`$
$`L=m+n.`$
$`IX:`$ $`m(1,0,1,2)+n(0,2,1,3)=(m,2n,m+n,2m+3n)`$
$`M=m[dim]\{(0,2,1,3)\}=\mathrm{\hspace{0.17em}6},`$
$`N=\mathrm{\hspace{0.17em}2}n\mathrm{\hspace{0.17em}2}[dim]\{(1,0,1,2)\}=\mathrm{\hspace{0.17em}8},`$
$`L=m+n.`$
$`XVI:`$ $`m(1,0,2,3)+n(0,1,0,1)=(m,n,2m,3m+n)`$
$`M=m[dim]\{(0,1,0,1)\}=\mathrm{\hspace{0.17em}2},`$
$`N=n[dim]\{(1,0,2,3)\}=\mathrm{\hspace{0.17em}6},`$
$`L=\mathrm{\hspace{0.17em}2}m.`$
$`XVIII:`$ $`m(1,0,1,2)+n(0,1,2,3)=(m,n,m+2n,2m+3n)`$
$`M=m[dim]\{(0,1,2,3)\}=\mathrm{\hspace{0.17em}6},`$
$`N=n[dim]\{(1,0,1,2)\}=\mathrm{\hspace{0.17em}4},`$
$`L=m+2n.`$
$`XIX:`$ $`m(1,0,2,3)+n(0,1,2,3)=(m,n,2m+2n,3m+3n)`$
$`M=m[dim]\{(0,1,2,3)\}=\mathrm{\hspace{0.17em}6},`$
$`N=n[dim]\{(1,0,2,3)\}=\mathrm{\hspace{0.17em}6},`$
$`L=\mathrm{\hspace{0.17em}2}m+2n.`$
$`XX:`$ $`m(2,0,1,3)+n(0,2,1,3)=(2m,2n,m+n,3m+3n)`$
$`M=\mathrm{\hspace{0.17em}2}m\mathrm{\hspace{0.17em}2}[dim]\{(0,1,2,3)\}=\mathrm{\hspace{0.17em}6},`$
$`N=\mathrm{\hspace{0.17em}2}n\mathrm{\hspace{0.17em}2}[dim]\{(1,0,2,3)\}=\mathrm{\hspace{0.17em}6},`$
$`L=m+n.`$
### 6.4 The Three-Vector Chain $`IV_3`$: $`\stackrel{}{k}_4=(M,N,M+L,N+L)`$
In this case (see Table 10), we have the dimension $`d=2M+2N+2L`$, the eldest vector $`\stackrel{}{k}=(1,1,1,1)`$, and the invariant monomials are $`(2,0,0,2)+(0,2,2,0)`$. This three-vector chain includes the following vectors form the two-vector construction:
$`VII:`$ $`m(2,1,1,0)+n(0,1,1,2)=(2m,m+n,m+n,2n)`$
$`M=\mathrm{\hspace{0.17em}2}m\mathrm{\hspace{0.17em}4},`$
$`N=m+n\mathrm{\hspace{0.17em}4},`$
$`L=nm\mathrm{\hspace{0.17em}0}.`$
$`X:`$ $`m(1,1,0,0)+n(0,0,1,1)=(m,m,n,n)`$
$`M=m\mathrm{\hspace{0.17em}2},`$
$`N=m\mathrm{\hspace{0.17em}2},`$
$`L=nm\mathrm{\hspace{0.17em}0}.`$
$`XII:`$ $`m(3,2,1,0)+n(0,1,2,3)=(3m,2m+n,m+2n,3n)`$
$`M=\mathrm{\hspace{0.17em}3}m,`$
$`N=\mathrm{\hspace{0.17em}2}m+n,`$
$`L=\mathrm{\hspace{0.17em}2}n2m,`$
$`(m,n)=(1,2),(2,1);(1,1),(1,4),(4,1),(2,5),(5,2).`$
$`XXI:`$ $`m(1,2,3,0)+n(1,2,0,3)=(m,2m,3m,3n)`$
$`M=m,`$
$`N=\mathrm{\hspace{0.17em}2}m,`$
$`L=m,`$
$`(m,n)=(1,1),(1,2),(2,1),,(1,5),(5,1),(4,5),(5,4).`$
$`XXII:`$ $`m(1,0,2,1)+n(0,1,0,1)=(m,n,2m,m+n)`$
$`M=m[dim]\{(1,0,2,1)\}=\mathrm{\hspace{0.17em}4},`$
$`N=n[dim]\{(1,0,2,1)\}=\mathrm{\hspace{0.17em}4},`$
$`L=m.`$
## 7 The Dual $`K3`$ Algebra from Four-Dimensional Extended Vectors
As discussed in the Introduction, the enumeration of $`K3`$ reflexive polyhedra obtained at level zero from pairs of projective vectors (Section 5) and triples (Section 6) is not quite complete. The one remaining example, corresponding to $`\stackrel{}{k}_4=(7,8,9,12)`$, can be found using the intersection-projection and duality properties outlined in Section 3, as we now discuss. This method can be used to build projective-vector chains using the rich projective structure of $`K3`$ vectors. For example, one can construct a chain with, as youngest vector, $`\stackrel{}{k}_4=(7,8,10,25)`$, which is dual to the eldest vector $`\stackrel{}{k}_4=(1,1,1,3)`$ contained in the triple chain $`III_3`$. Similarly, one can consider other cases, e.g., building a chain with youngest vector $`\stackrel{}{k}_4=(5,6,8,11)`$, contained in the triple chain $`II_3`$.
### 7.1 The Dual $`\stackrel{}{\pi }`$ Projective-Vector Structure of $`K3`$ Hypersurfaces
We obtained in section 6, as an interesting application of the $`K3`$ algebra, all the $`1+(4\times 2)+22`$ invariant monomials of the 22 double-intersection $`K3`$ chains via the triple intersections of $`K3`$ extended vectors. These invariant monomials correspond to particular directions relative to the reflexive polyhedra, which can be used to find the projection structures of the vectors. In particular, they can be used to find all the projective vectors which have no planar-intersection structure at all. Because of duality, their polyhedra have sufficient invariant directions that the projections on the corresponding perpendicular planes give reflexive planar polyhedra. Examples include youngest vectors which are dual to eldest vectors as well as other relations in the corresponding chain, e.g., as we shall see, the remaining $`K3`$ vector (7,8,9,12) is dual to (1,1,1,1), (7,8,10,25) is dual to (1,1,1,3), etc..
To understand this more deeply, we consider triple chains built using a special subalgebra of the four-dimensional extended vectors: $`\stackrel{}{k}_3^{ex(i)}`$ $`=(0,0,0,1)`$, $`(0,0,1,1)`$, $`(0,1,1,1)`$, $`(0,1,1,2)`$ and $`(0,1,2,3)`$, with all possible permutations. We consider triples $`\stackrel{}{k}_3^{ex(i,j,l)}`$ of these vectors with the property that each pair $`(i,j),(j,l),(l,i)`$ gives a reflexive planar polyhedron:
$`[\stackrel{}{k}_3^{ex(i)}]{\displaystyle [\stackrel{}{k}_3^{ex(j)}]}=[\stackrel{}{k}_3].`$ (7.1)
We note that the triple intersections of these triples of extended vectors always define an invariant direction, $`\stackrel{}{\pi }`$. In some cases, the triple intersection contains just two monomial vectors, and $`\stackrel{}{\pi }`$ is simply defined by their difference:
$`[\stackrel{}{k}_3^{ex(i)}]{\displaystyle [\stackrel{}{k}_3^{ex(j)}][\stackrel{}{k}_3^{ex(r)}]}\stackrel{}{\pi }_N=\{\stackrel{}{\mu }\stackrel{}{\mu }_0\},`$ (7.2)
where $`\stackrel{}{\mu }_0=(1,1,1,1)`$ is the basic monomial $`z^{\stackrel{}{\mu }_0+\stackrel{}{1}}=xyzu`$. These cases are listed in Table 11.
These pairs of invariant monomials correspond to directions $`\stackrel{}{\pi _i}=\stackrel{}{\mu }_i\stackrel{}{\mu }_0`$ in the exponent/monomial hyperspace given by the following vectors $`\stackrel{}{\mu }_N:N=1,2,3,4,5`$:
$`\stackrel{}{\mu }_1`$ $`=`$ $`(3,0,1,1),`$
$`\stackrel{}{\mu }_2`$ $`=`$ $`(0,0,0,3),`$
$`\stackrel{}{\mu }_3`$ $`=`$ $`(0,3,0,1),`$
$`\stackrel{}{\mu }_4`$ $`=`$ $`(0,0,4,0),`$
$`\stackrel{}{\mu }_5`$ $`=`$ $`(4,1,0,0).`$
as can be seen in Table 11.
In the other cases, the triple intersections contain three points which form a degenerate linear polyhedron, which also defines a unique direction $`\stackrel{}{\pi }`$ determined by three points, one of which $`(\stackrel{}{\mu }_0)`$ corresponds to the origin:
$`[\stackrel{}{k}_3^{ex(i)}]{\displaystyle [\stackrel{}{k}_3^{ex(j)}][\stackrel{}{k}_3^{ex(r)}]}\stackrel{}{\pi }_N=\{\stackrel{}{\mu }_+\stackrel{}{\mu }_0\}=\{\stackrel{}{\mu }_0\stackrel{}{\mu }_{}\},`$ (7.4)
as seen in Table 12.
It is easy to see that five of the invariant monomials from Table 11 produce a reflexive three-dimensional polyhedron. For example, from $`\stackrel{}{\mu }_2`$, $`\stackrel{}{\mu }_3`$, $`\stackrel{}{\mu }_4`$ and $`\stackrel{}{\mu }_5`$ one obtains the following exceptional vector whose associated polyhedron has no intersection substructure:
$`\stackrel{}{\mu }_\alpha \stackrel{}{k}_4`$ $`=`$ $`d=k_1+k_2+k_3+k_4,\alpha =0,2,3,4,5`$
$`\stackrel{}{k}_4`$ $`=`$ $`(7,8,9,12)[d=36],`$ (7.5)
where we used the constraint
$`\stackrel{}{\mu }_0={\displaystyle \frac{1}{4}}(\stackrel{}{\mu }_2+\stackrel{}{\mu }_3+\stackrel{}{\mu }_4+\stackrel{}{\mu }_5).`$ (7.6)
Thus duality enables us to identify the missing 95th $`K3`$ vector, which was not generated previously in our systematic study of the two- and three-vector chains. We recall that they contain totals of 90 and 91 vectors, respectively, of which only 94 were distinct.
Similarly, using these invariant monomials, one can find the rest of the exceptional $`\stackrel{}{k}_4`$ vectors, $`(3,5,6,7)`$, $`(3,6,7,8)`$, $`(5,6,7,9)`$ which were not included in the triple chains, together with $`(3,4,5,6)`$. They have intersection polyhedra that are not linear. These other exceptional $`\stackrel{}{k}_4`$ vectors are defined as follows:
$`\stackrel{}{\mu }_\alpha \stackrel{}{k}_4`$ $`=`$ $`d\alpha =0,1,2,3,3^{}`$
$`\stackrel{}{k}_4`$ $`=`$ $`(3,5,6,7)[d=21],`$ (7.7)
where again the following constraint has been used:
$`\stackrel{}{\mu }_0={\displaystyle \frac{1}{4}}(\stackrel{}{\mu }_1+\stackrel{}{\mu }_2+\stackrel{}{\mu }_3+\stackrel{}{\mu }_3^{})=`$
$`(3,1,0,1)+(0,0,0,3)+(1,0,3,0)+(0,3,1,0);`$
and
$`\stackrel{}{\mu }_\alpha \stackrel{}{k}_4`$ $`=`$ $`d\alpha =0,1,2,3,4`$
$`\stackrel{}{k}_4`$ $`=`$ $`(3,6,7,8)[d=24],`$ (7.8)
with the constraint:
$`\stackrel{}{\mu }_0={\displaystyle \frac{1}{4}}(\stackrel{}{\mu }_1+\stackrel{}{\mu }_2+\stackrel{}{\mu }_3+\stackrel{}{\mu }_4)=`$
$`(3,0,1,1)+(0,0,0,3)+(1,0,3,0)+(0,4,0,0).`$
We also find
$`\stackrel{}{\mu }_\alpha \stackrel{}{k}_4`$ $`=`$ $`d\alpha =2,3,3^{},5`$
$`\stackrel{}{k}_4`$ $`=`$ $`(5,6,7,9)[d=27],`$ (7.9)
where the following constraint also has been used:
$`\stackrel{}{\mu }_0={\displaystyle \frac{1}{4}}(\stackrel{}{\mu }_2+\stackrel{}{\mu }_3+\stackrel{}{\mu }_3^{}+\stackrel{}{\mu }_5)=`$
$`(0,0,0,3)+(0,3,0,1)+(0,1,3,0)+(4,0,1,0).`$
and
$`\stackrel{}{\mu }_\alpha \stackrel{}{k}_4`$ $`=`$ $`d\alpha =1,2,3,3^{}`$
$`\stackrel{}{k}_4`$ $`=`$ $`(3,4,5,6)[d=18],`$ (7.10)
where the following constraint also has been used:
$`\stackrel{}{\mu }_0={\displaystyle \frac{1}{4}}(\stackrel{}{\mu }_1+\stackrel{}{\mu }_2+\stackrel{}{\mu }_3+\stackrel{}{\mu }_3^{})=`$
$`(3,1,1,0)+(0,0,0,3)+(1,0,3,0)+(0,3,0,1).`$
### 7.2 Projective Chains of $`K3`$ Spaces Constructed from $`\stackrel{}{\pi }_N`$ Vectors
Using the invariant directions found in the previous Subsection, one can construct new triple chains:
$`p[\stackrel{}{k}_4]_{\stackrel{}{\pi }_N}=m\stackrel{}{k}^{ex(i)}+n\stackrel{}{k}^{ex(j)}+r\stackrel{}{k}^{ex(l)}`$ (7.11)
each corresponding to a direction $`\stackrel{}{\pi }`$ determined by an intersection of invariant monomial pairs. Each good projective vector in such a chain, determined by an invariant direction, contains the monomial/projective direction in its polyhedron. With respect to this direction, the polyhedron is projected onto a ‘good’ planar reflexive polyhedron. If the projective vector appears in several different chains, its polyhedron will have ‘good’ projections corresponding to each of these chains. This property can be used to make a classification by their projections of the projective vectors and their reflexive polyhedra. One finds that 78 projective $`K3`$ vectors out of 95 have such aprojective property. Taking into account the rest of the vectors which already were known from double-intersection $`J=\mathrm{\Pi }`$-symmetric chains, one can recover all 95 projective $`K3`$ vectors.
The distribution of the 3-dimensional set of positive-integer numbers $`m,n,r`$ depends on the dimension of the three extended vectors $`d^{(i)}=_\alpha \{\stackrel{}{k}_3^{ex(i)}\}_\alpha `$, $`i=1,2,3`$, participating in the construction of the chain, can have some ‘blank spots’, corresponding to ‘false vectors’ which do not correspond to any reflexive polyhedron. The origin of this phenomenon is connected with the structure of Calabi-Yau algebra, i.e., some of the projective vectors have different expansions (double-, triple-,…) in terms of the extended vectors. So, for example, if a vector is forbidden in two-vector expansions, it should also be forbidden in triple, etc., expansions, which is what we call a false vector. The self-consistency of the algebra entails the absences of some combinations of integer numbers $`m,n,r`$, even though all of them are below their maximum values. We already have met and discussed this phenomenon in the classification of triple-vector chains.
As seen in Table 11, one can give examples of triple intersections giving just one good vector which has three different projections with $`\mathrm{\Pi }=4`$:
$`[\stackrel{}{k}_4]_{\stackrel{}{\pi }_1^{(2)}}{\displaystyle [\stackrel{}{k}_4]_{\stackrel{}{\pi }_2^{(2)}}[\stackrel{}{k}_4]_{\stackrel{}{\pi }_3^{(4)}}}`$ $``$ $`(3,5,6,7)[21]`$
$`[\stackrel{}{k}_4]_{\stackrel{}{\pi }_1^{(2)}}{\displaystyle [\stackrel{}{k}_4]_{\stackrel{}{\pi }_2^{(2)}}[\stackrel{}{k}_4]_{\stackrel{}{\pi }_4^{(2)}}}`$ $``$ $`(3,6,7,8)[24]`$
$`[\stackrel{}{k}_4]_{\stackrel{}{\pi }_2^{(2)}}{\displaystyle [\stackrel{}{k}_4]_{\stackrel{}{\pi }_3^{(3)}}[\stackrel{}{k}_4]_{\stackrel{}{\pi }_3^{(1)}}}`$ $``$ $`(5,6,7,9)[27].`$
Moreover, the exceptional vector, which has four different projections with $`\mathrm{\Pi }=4`$, is given by the intersection of four such chains, i.e.:
$`[\stackrel{}{k}_4]_{\stackrel{}{\pi }_2^{(1)}}{\displaystyle [\stackrel{}{k}_4]_{\stackrel{}{\pi }_3^{(2)}}[\stackrel{}{k}_4]_{\stackrel{}{\pi }_4^{(2)}}[\stackrel{}{k}_4]_{\stackrel{}{\pi }_5^{(1)}}}`$ $``$ $`(7,8,9,12)[36].`$ (7.13)
To understand this in more detail, we consider one chain with projection $`\mathrm{\Pi }=4`$, which is determined by the invariant direction $`\stackrel{}{\pi }_2^{(1)}`$. The vectors of this chain are represented as linear combinations with positive-integer coefficients, $`M,N,L,`$ of the following three projective vectors, taken from the third line in Table 11:
$`\stackrel{}{k}_4(\stackrel{}{\pi }_2^{(1)})`$ $`=`$ $`M(0,1,1,1)+N(1,0,1,1)+L(1,1,0,1)`$
$`=`$ $`(N+L,M+L,M+N,M+N+L)`$
The basis is constructed out of the exceptional invariant monomials determining the $`\stackrel{}{\pi }`$ directions. Projecting on the perpendicular plane gives us planar reflexive polyhedra, so the third basis vector
$$\stackrel{}{e}_3=(1,1,1,2)(0,0,0,3).$$
(7.15)
is common to all the chains discussed in this Subsection.
Looking at the distribution of allowed integers $`M,N,L`$, we see ‘blank spots’ such as $`M=N=L=1`$, corresponding to the ‘false vector’ $`(2,2,2,3)`$, which is forbidden by the double-vector classification: it would require $`m=2`$ in the chain $`(2,2,2,3)=m(1,1,1,0)+n(0,0,0,1)`$, but actually $`m_{max}=1`$ for this chain. Also, all the polyhedra corresponding to these projective vectors have the other invariant directions $`\stackrel{}{\pi }_3^{(2)}(1,0,3,0)`$ with $`\mathrm{\Pi }=4`$ and should produce the following triple-vector expansion chain:
$`\stackrel{}{k}_4(\stackrel{}{\pi }_3^{(2)})`$ $`=`$ $`M(0,1,1,1)+N(1,0,1,2)+L(3,2,1,0)`$
$`=`$ $`(N+\mathrm{\hspace{0.17em}3}L,M+\mathrm{\hspace{0.17em}2}L,M+N+L,M+\mathrm{\hspace{0.17em}2}N)`$
Projecting on the perpendicular plane to the vector
$$\stackrel{}{e}_3=(0,1,2,1)(1,0,3,0).$$
(7.17)
gives us planar reflexive polyhedron with 4 points. This chain is a little longer and contains other projective vectors. Similarly, one can find using the other projective directions, $`\stackrel{}{\pi }_4^{(\alpha )}`$ and $`\stackrel{}{\pi }_5^{(1)}`$, two new triple expansion chains. Together these four invariant directions, $`\stackrel{}{\pi }_i^{(\alpha )}`$, (i=2,3,4,5), with the constructions of the corresponding triple projective chains contain 40 projective vectors (see Table 1).
One can compare the projection set, $`\stackrel{}{\pi }_2^{(1)}`$ and $`\mathrm{\Pi }=4`$, with the double-vector-intersection chain with $`J=4`$. It is interesting to note that six vectors from the projective chain shown in Table 13 also appear in the $`III`$-intersection chain with $`J=4`$ shown in Table 14. Conversely, the chain shown in this latter Table has just two vectors: $`(3,1,2,3),(1,3,5,6)`$ that are not contained in Table 13. The intersection structure of the $`III`$ chain shown in Table 14 is obtained from the following two vectors:
$`\stackrel{}{k}_4(III)`$ $`=`$ $`m(0,1,1,1)+n(3,0,1,2)`$
$`=`$ $`(n,m,m+n,m+n)`$
$`1m\mathrm{\hspace{0.17em}6},\mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}1}n\mathrm{\hspace{0.17em}2}.`$
The corresponding four invariant monomials are:
$`\stackrel{}{\mu }_0^1`$ $`=`$ $`(0,0,0,3)u^3`$
$`\stackrel{}{\mu }_0^2`$ $`=`$ $`(1,0,3,0)xz^3`$
$`\stackrel{}{\mu }_0^3`$ $`=`$ $`(2,3,0,0)x^2y^3`$
$`\stackrel{}{\mu }_0^4`$ $`=`$ $`(1,1,1,1)xyzu.`$ (7.19)
and the corresponding basis can be chosen in the form:
$`\stackrel{}{e}_1`$ $`=`$ $`(0,mn,m,0)`$
$`\stackrel{}{e}_2`$ $`=`$ $`(0,1,2,1)`$
$`\stackrel{}{e}_3`$ $`=`$ $`(1,1,1,2)`$
The canonical expression for the determinant of this lattice is
$`det(\stackrel{}{e}_1,\stackrel{}{e}_2,\stackrel{}{e}_3,\stackrel{}{e}_0)=\mathrm{\hspace{0.17em}3}m+\mathrm{\hspace{0.17em}6}n=d,`$ (7.21)
where $`\stackrel{}{e}_0(1,1,1,1)`$.
### 7.3 Example of a $`J,\mathrm{\Pi }=10`$ Double-Intersection Chain
To see another aspect of mirror symmetry and duality, consider the $`II`$ chain with intersection $`J(\mathrm{\Delta })=\mathrm{\Pi }(\mathrm{\Delta })=10`$ and $`J(\mathrm{\Delta }^{})=\mathrm{\Pi }(\mathrm{\Delta }^{})=4`$ shown in Table 15. The decomposition of this chain is in terms of the following two vectors:
$`\stackrel{}{k}_4`$ $`=`$ $`m(0,1,1,1)+n(1,0,1,1)`$
$`=`$ $`(n,m,m+n,m+n)`$
$`1m\mathrm{\hspace{0.17em}3},\mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}1}n\mathrm{\hspace{0.17em}3}.`$
The basis of the lattice in which the polyhedral intersection with the set of positive-integer points corresponds to Table 15 is the following:
$`\stackrel{}{e}_1`$ $`=`$ $`(m,n,0,0)`$
$`\stackrel{}{e}_2`$ $`=`$ $`(1,1,1,0)`$
$`\stackrel{}{e}_3`$ $`=`$ $`(1,1,0,1)`$ (7.23)
and the corresponding determinant is
$`det(\stackrel{}{e}_1,\stackrel{}{e}_2,\stackrel{}{e}_3,\stackrel{}{e}_0)=\mathrm{\hspace{0.17em}3}m+\mathrm{\hspace{0.17em}3}n=d,`$ (7.24)
where $`\stackrel{}{e}_0=(1,1,1,1)`$ again. The ten corresponding invariant monomials are:
$`\stackrel{}{\mu }_0^1`$ $`=`$ $`(3,3,0,0)x^3y^3`$
$`\stackrel{}{\mu }_0^2`$ $`=`$ $`(2,2,1,0)x^2y^2z`$
$`\stackrel{}{\mu }_0^3`$ $`=`$ $`(1,1,2,0)xy^{}z^2`$
$`\stackrel{}{\mu }_0^4`$ $`=`$ $`(0,0,3,0)z^3`$
$`\stackrel{}{\mu }_0^5`$ $`=`$ $`(2,2,0,1)x^2y^2u,`$
$`\stackrel{}{\mu }_0^6`$ $`=`$ $`(1,1,1,1)xyzu,`$
$`\stackrel{}{\mu }_0^7`$ $`=`$ $`(0,0,2,1)z^2u,`$
$`\stackrel{}{\mu }_0^8`$ $`=`$ $`(1,1,0,2)xyu^2,`$
$`\stackrel{}{\mu }_0^9`$ $`=`$ $`(0,0,1,2)zu^2,`$
$`\stackrel{}{\mu }_0^{10}`$ $`=`$ $`(0,0,0,3)u^3.`$ (7.25)
For the vector $`\stackrel{}{k}_4=(1,1,2,2)`$, one can consider the basis
$`\stackrel{}{e}_1`$ $`=`$ $`(3,3,0,0)`$
$`\stackrel{}{e}_2`$ $`=`$ $`(1,1,1,0)`$
$`\stackrel{}{e}_3`$ $`=`$ $`(1,1,0,1)`$ (7.26)
with determinant 18, in which the dual pair of polyhedra:
$`1_L+\mathrm{\hspace{0.17em}10}_J+\mathrm{\hspace{0.17em}1}_R=\mathrm{\hspace{0.17em}12},`$
$`4_L^{}+\mathrm{\hspace{0.17em}4}_J^{}+\mathrm{\hspace{0.17em}4}_R^{}=\mathrm{\hspace{0.17em}12}^{}.`$ (7.27)
both contain 12 points and 12 mirror points, respectively.
### 7.4 Example of a Chain with $`\mathrm{\Pi }=5`$ and Eldest Vector $`\stackrel{}{k}_4=(7,8,10,25)`$
Now we present in Table 16 a projective chain with $`\mathrm{\Pi }=5`$, constructed from the invariant direction $`\stackrel{}{\pi }_8^{(1)}`$ with the invariant monomials $`(0,0,0,2)+(2,2,2,0)`$. The 14 projective vectors of this chain are represented as linear combinations with positive-integer coefficients, $`M,N,L,Q`$: $`Q=2,1`$ of the following three vectors: Projecting on the perpendicular plane gives us planar reflexive polyhedra, so the third basis vector
$$\stackrel{}{e}_3=(1,1,1,1)(0,0,0,2).$$
(7.28)
is common to all the chains discussed in this Subsection.
The vectors of this chain are represented as linear combinations with positive-integer coefficients, $`M,N,L,Q`$: $`Q=2,1`$ of the following three vectors:
$`Q\stackrel{}{k}_4(\stackrel{}{\pi }_8^{(1)})`$ $`=`$ $`M(0,1,1,1)+N(1,0,1,2)+L(1,1,0,2)`$
$`=`$ $`(N+L,M+L,M+N,\mathrm{\hspace{0.17em}2}M+\mathrm{\hspace{0.17em}2}N+\mathrm{\hspace{0.17em}2}L)`$
where the third basis vector,
$$\stackrel{}{e}_3=(1,1,1,1)(0,0,0,2).$$
(7.30)
is common to all the chain.
There can be constructed additional three chains, $`\stackrel{}{\pi }_8^{(2,3,4)}`$, with the same invariant direction, (0,0,0,2)-(2,2,2,0), and the same youngest vector, but with the different triple intersections and therefore with the different projective chains. Together one can find inside all of four projective chains, $`\stackrel{}{\pi }_8^{(\alpha )},\alpha =1,2,3,4`$, a total of 33 projective vectors (see Table 1).
It is interesting to note that the chain $`\stackrel{}{\pi }_8^{(1)}`$ has 11 $`\stackrel{}{k}_4`$ vectors with $`\mathrm{\Pi }=5`$ in common with the $`IX_J`$ chain where $`J=5`$, whose structure is obtained from the following two vectors:
$`\stackrel{}{k}_4(IX)`$ $`=`$ $`m(0,1,1,2)+n(2,1,0,3)`$
$`=`$ $`(\mathrm{\hspace{0.17em}2}n,m+n,m,\mathrm{\hspace{0.17em}2}m+\mathrm{\hspace{0.17em}3}n)`$
$`1m\mathrm{\hspace{0.17em}6},\mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}1}n\mathrm{\hspace{0.17em}4}.`$
The chain $`IX`$ of $`\stackrel{}{k}_4`$ projective vectors with the structure $`5_{J=\mathrm{}}9_{\mathrm{\Pi }=\mathrm{}}`$ is presented in Table 17.
The lattice determinant and the basis are given by the following expressions:
$`\stackrel{}{e}_1`$ $`=`$ $`(0,m,m+n,0)`$
$`\stackrel{}{e}_2`$ $`=`$ $`(1,2,2,0)`$
$`\stackrel{}{e}_3`$ $`=`$ $`(1,1,1,1),`$ (7.32)
and
$`det(\stackrel{}{e}_1,\stackrel{}{e}_2,\stackrel{}{e}_3,\stackrel{}{e}_0)=\mathrm{\hspace{0.17em}4}m+\mathrm{\hspace{0.17em}6}n=d,`$ (7.33)
where $`\stackrel{}{e}_0=(1,1,1,1)`$.
The possible values of $`m`$ and $`n`$ for this chain are also determined by the dimensions of the extended vectors, $`d(\stackrel{}{k}^{ex(i)})=6`$ and $`d(\stackrel{}{k}^{ex(j)})=4`$, with the additional constraint $`n_{max}=3<dim(0,1,1,2)`$ (see Table 17):
$`p\stackrel{}{k}_4(IX)`$ $`=`$ $`m(0,1,1,2)+n(2,0,1,3)`$
$`p=\mathrm{\hspace{0.17em}1}`$ $``$ $`\mathrm{\hspace{0.17em}1}m\mathrm{\hspace{0.17em}6};,\mathrm{\hspace{0.17em}\hspace{0.17em}1}n\mathrm{\hspace{0.17em}3}`$ (7.34)
The 5 invariant monomials for this chain are the following:
$`\stackrel{}{\mu }_0^1=(1,4,0,0)`$ $``$ $`xy^4`$
$`\stackrel{}{\mu }_0^2=(2,2,2,0)`$ $``$ $`x^2y^2z^2`$
$`\stackrel{}{\mu }_0^3=(3,0,4,0)`$ $``$ $`x^3z^4`$
$`\stackrel{}{\mu }_0^4=(1,1,1,1)`$ $``$ $`xyzu`$
$`\stackrel{}{\mu }_0^5=(0,0,0,2)`$ $``$ $`u^2.`$ (7.35)
### 7.5 Example of a $`J=\mathrm{\Pi }=9`$ Chain
To see another aspect of mirror symmetry and duality, we now consider the chain $`VI`$ with intersection $`J(\mathrm{\Delta })=\mathrm{\Pi }(\mathrm{\Delta })=9`$ and $`J(\mathrm{\Delta }^{})=\mathrm{\Pi }(\mathrm{\Delta }^{})=5`$ shown in Table 15, which is constructed from the extended vectors $`\stackrel{}{k}^i=(0,1,1,2)`$ and $`\stackrel{}{k}^j=(1,0,1,2)`$. In this case, duality gives very simple connections between the numbers of integer points in the dual polyhedron pair, as seen in Table 18.
The canonical basis for chain $`VI`$ is:
$`\stackrel{}{e}_1`$ $`=`$ $`(m,n,0,0)`$
$`\stackrel{}{e}_2`$ $`=`$ $`(1,1,1,0)`$
$`\stackrel{}{e}_3`$ $`=`$ $`(1,1,1,1)`$ (7.36)
with the following restriction on the determinant
$`det(\stackrel{}{e}_1,\stackrel{}{e}_2,\stackrel{}{e}_3,\stackrel{}{e}_0)=\mathrm{\hspace{0.17em}4}m+\mathrm{\hspace{0.17em}4}n=d,`$ (7.37)
where $`\stackrel{}{e}_0=(1,1,1,1)`$.
The possible values of $`m`$ and $`n`$ for this chain are determined by the dimensions of the extended vectors, without any unexpected puzzles:
$`p\stackrel{}{k}_4(VI)`$ $`=`$ $`m(0,1,1,2)+n(1,0,1,2)`$
$`p=\mathrm{\hspace{0.17em}1}`$ $``$ $`\mathrm{\hspace{0.17em}1}m\mathrm{\hspace{0.17em}4};\mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}1}n\mathrm{\hspace{0.17em}4}.`$ (7.38)
and the following:
$`\stackrel{}{\mu }_0^1=(4,4,0,0)`$ $``$ $`x^4y^4`$
$`\stackrel{}{\mu }_0^2=(3,3,1,0)`$ $``$ $`x^3y^3z`$
$`\stackrel{}{\mu }_0^3=(2,2,2,0)`$ $``$ $`x^2y^2z^2`$
$`\stackrel{}{\mu }_0^4=(1,1,3,0)`$ $``$ $`xyz^3`$
$`\stackrel{}{\mu }_0^5=(0,0,4,0)`$ $``$ $`z^4`$
$`\stackrel{}{\mu }_0^6=(2,2,0,1)`$ $``$ $`x^2y^2u`$
$`\stackrel{}{\mu }_0^7=(1,1,1,1)`$ $``$ $`xyzu`$
$`\stackrel{}{\mu }_0^8=(0,0,2,1)`$ $``$ $`z^2u`$
$`\stackrel{}{\mu }_0^9=(0,0,0,2)`$ $``$ $`u^2.`$ (7.39)
are the 9 invariant monomials $`\mathrm{\Psi }_{inv}`$ for this chain.
Analogously, one can consider the projective chain $`\stackrel{}{\pi }_7^{(\alpha )}`$ ($`\mathrm{\Pi }=5`$) with youngest vector (5,6,8,11), and compare it with the double-intersection chain $`VIII`$, constructed from the extended vectors $`\stackrel{}{k}(VIII)=m(0,1,1,2)+n(1,1,2,0)`$; $`(d=4m+4n`$. $`m_{max}=3,n_{max}=4)`$, $`\stackrel{}{k}_{eld}=(1,2,3,2)[8]`$. Among the 95 $`K3`$ projective vectors, 26 have such an invariant-direction structure, and therefore can be found in corresponding projective chains (see Table 1).
## 8 $`K3`$ Hypersurfaces and Cartan-Lie Algebra Graphs
We discuss in this Section more details of the emergence of Cartan-Lie algebra graphs in our construction of CY spaces.
### 8.1 Cartan-Lie Algebra Graphs and the Classification of Chains of Projective Vectors
As we commented already in the Introduction and in Section 2, the structure of the projective $`\stackrel{}{k}_4`$ vectors in 22 chains leads to interesting relations with the five classical regular dual polyhedron pairs in three-dimensional space: the one-dimensional point, two-dimensional line segment and three-dimensional tetrahedron, octahedron-cube and icosahedron-dodecahedron. There are also interesting correspondences with the Cartan-Lie algebra $`CLA`$ graphs for the five types of groups in the $`ADE_{6,7,8}`$ series: see Figure 8. The $`CLA_{J,\mathrm{\Pi }}`$ graphs, which can be seen in the polyhedra of the corresponding $`\stackrel{}{k}_4`$ projective vectors, follow completely the structure of the five possible extended vectors:
$`\stackrel{}{k}_C^{ext}=(0,0,0,1)`$ $``$ $`A_r;`$
$`\stackrel{}{k}_D^{ext}=(0,0,1,1)`$ $``$ $`D_r;`$
$`\stackrel{}{k}_T^{ext}=(0,1,1,1)`$ $``$ $`E_6;`$
$`\stackrel{}{k}_O^{ext}=(0,1,1,2)`$ $``$ $`E_7;`$
$`\stackrel{}{k}_I^{ext}=(0,1,2,3)`$ $``$ $`E_8.`$ (8.1)
We give in Table 19 the ADE structures and the $`CD_J`$ diagrams of all the eldest $`K3`$ projective vectors from the 22 double chains. An illustration is given in Figure 8, and the rest of this Section discusses the examples of chains $`XV`$ to $`XIX`$, illustrating the power of our systematic approach.
### 8.2 The $`K3`$ Chain $`XV`$ with Graphs in the $`E_8^{(1)}A_r^{(1)}`$ Series
Here we give the list of $`\stackrel{}{k}_4`$ vectors which can be constructed from the Weierstrass vectors $`\stackrel{}{k}_3(1,2,3)`$ and $`\stackrel{}{k}_1=(1)`$, shown as chain $`XV`$ in Table 20. The number of $`\stackrel{}{k}_4`$ vectors in this chain is determined by the positive-integer numbers: $`m=1,n6`$, according to the dimensions of the corresponding component $`\stackrel{}{k}^i`$.
The basis for this chain, see Figure 9, can be written in the the following form:
$`\stackrel{}{e}_1`$ $`=`$ $`(n,0,0,m)`$
$`\stackrel{}{e}_2`$ $`=`$ $`(2,1,0.0)`$
$`\stackrel{}{e}_3`$ $`=`$ $`(3,0,1,0)`$ (8.2)
The determinant of this canonical basis coincides, of course, with the dimensions of the $`\stackrel{}{k}_4`$ vectors:
$`det(\stackrel{}{e}_1,\stackrel{}{e}_2,\stackrel{}{e}_3,\stackrel{}{e}_0)=\mathrm{\hspace{0.17em}6}m+\mathrm{\hspace{0.17em}1}n=d,`$ (8.3)
where $`\stackrel{}{e}_0=(1,1,1,1)`$. The decomposition of this chain is again determined by the dimension of the extended vectors $`d(\stackrel{}{k}^{ex(i)})=k_1^{ex(i)}+k_2^{ex(i)}+k_3^{ex(i)}+k_4^{ex(i)}`$, as seen in Table 20:
$`\stackrel{}{k}_4(XV)=m(1,2,3,0)`$ $`+`$ $`n(0,0,0,1)`$
$`m=\mathrm{\hspace{0.17em}1}`$ $`\mathrm{\hspace{0.17em}\hspace{0.17em}1}n\mathrm{\hspace{0.17em}6}.`$ (8.4)
The seven invariant monomials corresponding to this chain are:
$`\stackrel{}{\mu }_0^1=(6,0,0,1)`$ $``$ $`x^6u`$
$`\stackrel{}{\mu }_0^2=(4,1,0,1)`$ $``$ $`x^4yu`$
$`\stackrel{}{\mu }_0^3=(2,2,0,1)`$ $``$ $`x^2y^2u`$
$`\stackrel{}{\mu }_0^4=(0,3,0,1)`$ $``$ $`y^3u`$
$`\stackrel{}{\mu }_0^5=(3,0,1,1)`$ $``$ $`x^3zu`$
$`\stackrel{}{\mu }_0^6=(1,1,1,1)`$ $``$ $`xyzu`$
$`\stackrel{}{\mu }_0^7=(0,0,2,1)`$ $``$ $`z^2u.`$ (8.5)
Considering the dual pairs for these vectors, one can see that the singularities of the eldest vector $`\stackrel{}{k}_4=(1,2,3,1)`$ correspond to some graphs of the $`A_{6}^{(1)}{}_{L}{}^{}E_{8}^{(1)}{}_{R}{}^{}`$ series, as seen in Figure 9. For instance, if one looks at the integer points in the edges of the polyhedron on the left (right) side of the intersection by the hyperplane $`\stackrel{}{k}^i=(0,1,2,3)`$, one sees graphs with $`A_{6}^{(1)}{}_{L}{}^{}`$ and $`E_{8}^{(1)}{}_{R}{}^{}`$ Lie algebras. Going to the last minimal $`\stackrel{}{k}=(1,2,3,6)`$ of this chain, we find that the right graph degenerates and left points reproduce $`A_{11}^{(1)}`$ with the maximum possible rank in this chain. Thus, the six $`\stackrel{}{k}`$ vectors in this chain produce the following graphs in the $`A`$ series: $`A_6^{(1)},A_7^{(1)},A_8^{(1)},A_9^{(1)},A_{10}^{(1)},A_{11}^{(1)}`$.
### 8.3 The $`K3`$ Chain $`XVI`$ with Graphs in the $`E_8^{(1)}D_r`$ Series
The basis for the chain shown in Table 21 is
$`\stackrel{}{e}_1`$ $`=`$ $`(m,n,0,0)`$
$`\stackrel{}{e}_2`$ $`=`$ $`(0,2,1,0)`$
$`\stackrel{}{e}_3`$ $`=`$ $`(1,1,1,1),`$ (8.6)
with
$`det(\stackrel{}{e}_1,\stackrel{}{e}_2,\stackrel{}{e}_3,\stackrel{}{e}_0)=\mathrm{\hspace{0.17em}6}m+\mathrm{\hspace{0.17em}2}n=d,`$ (8.7)
where $`\stackrel{}{e}_0=(1,1,1,1)`$ again. The decomposition of this chain is completely determined by the dimensions of the vectors shown in Table 21:
$`p\stackrel{}{k}_4(XVI)`$ $`=`$ $`m(0,1,2,3)+n(1,0,0,1)`$
$`p=\mathrm{\hspace{0.17em}1}^{}`$ $``$ $`\mathrm{\hspace{0.17em}1}m\mathrm{\hspace{0.17em}2};\mathrm{\hspace{0.17em}\hspace{0.17em}1}n\mathrm{\hspace{0.17em}6}`$
$`p=\mathrm{\hspace{0.17em}2}`$ $``$ $`m=n=2.`$ (8.8)
The seven invariant monomials corresponding to this chain are the following:
$`\stackrel{}{\mu }_0^1=(2,6,0,0)`$ $``$ $`x^2y^6`$
$`\stackrel{}{\mu }_0^2=(2,4,1,0)`$ $``$ $`x^2y^4z`$
$`\stackrel{}{\mu }_0^3=(2,2,2,0)`$ $``$ $`x^2y^2z^2`$
$`\stackrel{}{\mu }_0^4=(2,0,3,0)`$ $``$ $`x^2z^3`$
$`\stackrel{}{\mu }_0^5=(1,3,0,1)`$ $``$ $`xy^3u`$
$`\stackrel{}{\mu }_0^6=(1,1,1,1)`$ $``$ $`xyzu`$
$`\stackrel{}{\mu }_0^7=(0,0,0,2)`$ $``$ $`u^2.`$ (8.9)
The example of the $`E_{8}^{(1)}{}_{L}{}^{}D_{8R}`$ graph associated with the eldest $`(1,1,2,4))[8]`$ polyhedron in Table 21 is shown in Figure 10.
### 8.4 The $`J=\mathrm{\Pi }`$ Symmetric Chain $`XVII`$ with Exceptional Graph $`E_6\times E_8`$
We show in Table 22 the projective vectors constructed from $`\stackrel{}{k}_3^{ex}=(0,1,1,1)`$ and $`\stackrel{}{k}_3^{ex}=(1,0,2,3)`$. In this case, the number of points in the maximal polyhedron with $`m=n=1`$ can easily be calculated: $`33=(10)_L+(7)_{int}+(16)_R`$. The ‘right’ $`15_R+1_R`$ points form the graph for the affine $`E_8^{(1)}`$ Lie algebra, as shown in Figure 11:
$`6=\mathrm{\hspace{0.17em}1}+\mathrm{\hspace{0.17em}1}+\mathrm{\hspace{0.17em}1}+\mathrm{\hspace{0.17em}1}+1+1`$ $``$ $`\{(P_{x_0})_1+(P_{x_1})_2+(P_{x_2})_3+(P_{x_3})_4+(P_{x_4})_5+(P_{x_5})_6\}`$
$`3=\mathrm{\hspace{0.17em}3}`$ $``$ $`\{(P_{x_6,x_6^{^{}},x_6^{^{\prime \prime }}})_3\}`$
$`6=\mathrm{\hspace{0.17em}4}+\mathrm{\hspace{0.17em}2}`$ $``$ $`\{(P_{x_7,x_7^{^{}},x_7^{^{\prime \prime }},x_7^{^{\prime \prime \prime }}})_4+(P_{x_8,x_8^{^{}}})_2\}.`$ (8.10)
The ‘left’ points in this polyhedron, $`9_L+1_L`$, correspond to the $`E_6^{(1)}`$ affine series with the Coxeter numbers:
$`3=\mathrm{\hspace{0.17em}1}+\mathrm{\hspace{0.17em}1}+\mathrm{\hspace{0.17em}1}`$ $``$ $`\{(P_{x_1})_1+(P_{x_2})_2+(P_{x_3})_3\}`$
$`3=\mathrm{\hspace{0.17em}2}+\mathrm{\hspace{0.17em}1}`$ $``$ $`\{(P_{x_4,x_4^{^{}}})_2+(P_{x_0})_1\}`$
$`3=\mathrm{\hspace{0.17em}2}+\mathrm{\hspace{0.17em}1}`$ $``$ $`\{(P_{x_5,x_5^{^{}}})_2+(P_{x_6})_1\}.`$ (8.11)
For $`m_{max}=d(\stackrel{}{k}(1,2,3))=6`$ and $`n_{min}=1`$, the corresponding polyhedron contains 18 points: $`18=(10)_L+(7)_{int}+(1)_R`$. Conversely, for $`m_{min}=1`$ and $`n_{max}=3=dim(\stackrel{}{k}(1,1,1))`$, the self-dual vector $`\stackrel{}{k}=(3,1,7,10)`$ has 24 integer points: $`24=(1)_L+(7)_{int}+(16)_R`$. Finally, the polyhedron with $`m=5`$ and $`n=3`$ contains the minimal possible number of integer points, namely $`9=(1)_L+(7)_{int}+(1)_R`$. This minimal vector $`(3,5,11,14)[33]`$ is the dual conjugate of the vector $`\stackrel{}{k}=(1,1,4,6)[12]`$.
The canonical basis of the chain shown in Table 22 is:
$`\stackrel{}{e}_1`$ $`=`$ $`(m,n,0,0)`$
$`\stackrel{}{e}_2`$ $`=`$ $`(2,1,1,0)`$
$`\stackrel{}{e}_3`$ $`=`$ $`(1,0,1,1),`$ (8.12)
with
$`det(\stackrel{}{e}_1,\stackrel{}{e}_2,\stackrel{}{e}_3,\stackrel{}{e}_1)=\mathrm{\hspace{0.17em}3}m+\mathrm{\hspace{0.17em}6}n=d,`$ (8.13)
where $`\stackrel{}{e_1}=(1,1,1,1)`$. The possible values of $`m`$ and $`n`$ for this chain are determined in the standard way from the dimensions of the extended vectors, $`d(\stackrel{}{k}^{ex(j)})=6`$ and $`d(\stackrel{}{k}^{ex(i)})=3`$, as seen in Table 22:
$`p\stackrel{}{k}_4(XVII)`$ $`=`$ $`m(0,1,1,1)+n(1,0,2,3)`$
$`p=\mathrm{\hspace{0.17em}1}^{}`$ $``$ $`1m\mathrm{\hspace{0.17em}6};;\mathrm{\hspace{0.17em}\hspace{0.17em}1}n\mathrm{\hspace{0.17em}3};`$
$`p=\mathrm{\hspace{0.17em}2}`$ $``$ $`m=n=\mathrm{\hspace{0.17em}2};`$
$`p=\mathrm{\hspace{0.17em}3}`$ $``$ $`m=n=\mathrm{\hspace{0.17em}3}.`$ (8.14)
The seven invariant monomials corresponding to this chain are the following:
$`\stackrel{}{\mu }_0^1=(6,3,0,0,)`$ $``$ $`x^6y^3`$
$`\stackrel{}{\mu }_0^2=(4,2,1,0,)`$ $``$ $`x^4y^2z`$
$`\stackrel{}{\mu }_0^3=(2,1,2,0,)`$ $``$ $`x^2yz`$
$`\stackrel{}{\mu }_0^4=(0,0,3,0,)`$ $``$ $`z^3`$
$`\stackrel{}{\mu }_0^5=(3,2,0,1,)`$ $``$ $`x^3y^2u`$
$`\stackrel{}{\mu }_0^6=(1,1,1,1,)`$ $``$ $`xyzu`$
$`\stackrel{}{\mu }_0^7=(0,1,0,2,)`$ $``$ $`yu^2`$ (8.15)
and the corresponding $`E_{6}^{(1)}{}_{L}{}^{}E_{8}^{(1)}{}_{R}{}^{}`$ graph associated with the eldest $`(1,1,3,4))[9]`$ polyhedron in chain $`XVII`$ is shown in Table 23 and Figure 11.
### 8.5 The $`J=\mathrm{\Pi }`$ Symmetric Chain $`XVIII`$ with Exceptional Graph $`E_7\times E_8`$
This chain can be built from the vectors $`\stackrel{}{k}_4^{ex^i}=(0,1,1,2)`$ and $`\stackrel{}{k}_4^{ex^j}=(1,0,2,3)`$, with positive integers $`m6`$ and $`n4`$. The maximal ($`m=n=1`$) polyhedron in this chain is again completely determined by the dimensions 4 and 6 of the projective vectors $`\stackrel{}{k}_4^{ex^i}`$ and $`\stackrel{}{k}_4^{ex^j}`$, respectively:
$$36=(13)_L+(7)_{J=\mathrm{\Pi }}+(16)_R.$$
(8.16)
The ‘right’ $`15_R+1_R`$ and ‘left’ $`12_L+1_L`$ points produce the graphs for the affine $`E_8^{(1)}`$ and $`E_7^{(1)}`$ Lie algebras, respectively, as seen in Figure 12. The vector $`\stackrel{}{k}=(3,4,9,14)[28]`$ is self-dual with $`E_8^{(1)}`$ graphs for the dual polyhedron pair. The ‘minimal’ vector $`\stackrel{}{k}`$ gives the following set of integer lattice points in the polyhedron:
$$(1)_L+(7)_{int}+(1)_R=9.$$
(8.17)
The canonical basis for the chain shown in Table 24 is:
$`\stackrel{}{e}_1`$ $`=`$ $`(m,n,0,0)`$
$`\stackrel{}{e}_2`$ $`=`$ $`(2,1,1,0)`$
$`\stackrel{}{e}_3`$ $`=`$ $`(1,1,1,1),`$ (8.18)
with
$`det(\stackrel{}{e}_1,\stackrel{}{e}_2,\stackrel{}{e}_3,\stackrel{}{1})=\mathrm{\hspace{0.17em}4}m+\mathrm{\hspace{0.17em}6}n=d,`$ (8.19)
The possible values of $`m`$ and $`n`$ for this chain fill up the dimensions of the extended vectors $`d(\stackrel{}{k}^{ex(j)})=6`$ and $`d(\stackrel{}{k}^{ex(i)})=4`$, as seen in Table 24:
$`p\stackrel{}{k}_4(XVIII)`$ $`=`$ $`m(0,1,1,2)+n(1,0,2,3)`$
$`p=\mathrm{\hspace{0.17em}1}^{}`$ $``$ $`1m\mathrm{\hspace{0.17em}6};\mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}1}n\mathrm{\hspace{0.17em}4};`$
$`p=\mathrm{\hspace{0.17em}2}`$ $``$ $`m=n=\mathrm{\hspace{0.17em}2};`$
$`p=\mathrm{\hspace{0.17em}3}`$ $``$ $`m=n=\mathrm{\hspace{0.17em}3};`$
$`p=\mathrm{\hspace{0.17em}4}`$ $``$ $`m=n=\mathrm{\hspace{0.17em}4}.`$ (8.20)
The seven invariant monomials corresponding to this chain are the following:
$`\stackrel{}{\mu }_0^1=(6,4,0,0,)`$ $``$ $`x^6y^4`$
$`\stackrel{}{\mu }_0^2=(4,3,1,0,)`$ $``$ $`x^4y^3z`$
$`\stackrel{}{\mu }_0^3=(2,2,2,0,)`$ $``$ $`x^2y^2z^2`$
$`\stackrel{}{\mu }_0^4=(0,1,3,0,)`$ $``$ $`yz^3`$
$`\stackrel{}{\mu }_0^5=(3,2,0,1,)`$ $``$ $`x^3y^2u`$
$`\stackrel{}{\mu }_0^6=(1,1,1,1,)`$ $``$ $`xyzu`$
$`\stackrel{}{\mu }_0^7=(0,1,0,2,)`$ $``$ $`yu^2`$ (8.21)
The $`E_{7}^{(1)}{}_{L}{}^{}E_{8}^{(1)}{}_{R}{}^{}`$ graph associated with the eldest $`(1,1,3,5))[10]`$ polyhedron in chain $`XVIII`$ can be seen in Table 25 and Figure 12.
### 8.6 Chain $`XIX`$ with $`(7_J,7_\mathrm{\Pi })`$ Weierstrass Triangle Fibrations
We now consider the chain $`XIX`$ of $`\stackrel{}{k}_4`$ projective vectors with $`E_{8}^{}{}_{L}{}^{}`$ and $`E_{8}^{}{}_{R}{}^{}`$ graphs. This chain starts from the $`m=n=1`$ polyhedron, which is left-right symmetric with respect to the intersection $`P^2(1,2,3)`$. This polyhedron $`P^3(1,1,4,6)`$ contains $`39=16_L+(7)_{J=\mathrm{\Pi }}+16_R`$ integer points: see Table 26 and Figure 13. The minimal vector $`\stackrel{}{k}=(5,6,22,33)[66]`$ is the dual conjugate of the eldest vector $`\stackrel{}{k}=(1,1,4,6)[12]`$, the vector $`\stackrel{}{k}=(1,6,14,21)[42]`$ is self-dual, and its dual pair of $`K3`$ polyhedra yield the self-dual $`E_8^{(1)}`$ graph.
The basis of the chain shown in Table 26 is the following:
$`\stackrel{}{e}_1`$ $`=`$ $`(m,n,0,0)`$
$`\stackrel{}{e}_2`$ $`=`$ $`(2,2,1,0)`$
$`\stackrel{}{e}_3`$ $`=`$ $`(1,1,1,1),`$ (8.22)
with
$`det(\stackrel{}{e}_1,\stackrel{}{e}_2,\stackrel{}{e}_3,\stackrel{}{e}_0)=\mathrm{\hspace{0.17em}6}m+\mathrm{\hspace{0.17em}6}n=d,`$ (8.23)
where $`\stackrel{}{e}_0=(1,1,1,1)`$. The possible values of $`m`$ and $`n`$ for this chain are completely determined by the dimensions of the vectors $`d(\stackrel{}{k}^{ex(j)})=6`$ and $`d(\stackrel{}{k}^{ex(i)})=6`$ (see Table 26):
$`p\stackrel{}{k}_4(XIX)`$ $`=`$ $`m(0,1,2,3)+n(1,0,2,3)`$
$`p=\mathrm{\hspace{0.17em}1}^{}`$ $``$ $`1m\mathrm{\hspace{0.17em}6};\mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}1}n\mathrm{\hspace{0.17em}6};`$
$`p=\mathrm{\hspace{0.17em}2}`$ $``$ $`m=n=\mathrm{\hspace{0.17em}2};`$
$`p=\mathrm{\hspace{0.17em}3}`$ $``$ $`m=n=\mathrm{\hspace{0.17em}3};`$
$`p=\mathrm{\hspace{0.17em}4}`$ $``$ $`m=n=\mathrm{\hspace{0.17em}4};`$
$`p=\mathrm{\hspace{0.17em}6}`$ $``$ $`m=n=\mathrm{\hspace{0.17em}6}.`$ (8.24)
The seven invariant monomials corresponding to this chain are the following:
$`\stackrel{}{\mu }_0^1=(6,6,0,0)`$ $``$ $`x^6y^6`$
$`\stackrel{}{\mu }_0^2=(4,4,1,0)`$ $``$ $`x^4y^4z`$
$`\stackrel{}{\mu }_0^3=(2,2,2,0)`$ $``$ $`x^2y^2z^2`$
$`\stackrel{}{\mu }_0^4=(0,1,3,0)`$ $``$ $`yz^3`$
$`\stackrel{}{\mu }_0^5=(3,3,0,1)`$ $``$ $`x^3y^3u`$
$`\stackrel{}{\mu }_0^6=(1,1,1,1)`$ $``$ $`xyzu`$
$`\stackrel{}{\mu }_0^7=(0,0,0,2)`$ $``$ $`u^2.`$ (8.25)
Using these invariant monomials and basis the CY equations for all the $`\stackrel{}{k}(l=m+n)`$ projective vectors of this chain can be written in the following form:
$`F(\stackrel{}{z})_{m,n}={\displaystyle \underset{j=1}{\overset{j=7}{}}}\stackrel{}{z}^{\stackrel{}{\mu }_0^j}\{{\displaystyle \underset{p=1}{\overset{p=\mathrm{\Pi }_{jL}}{}}}a_{\stackrel{}{\mu }_0^j}^{pL}\stackrel{}{z}^{n_{pL}(\stackrel{}{e}_1)}+{\displaystyle \underset{p=1}{\overset{p=\mathrm{\Pi }_{jR}}{}}}a_{\stackrel{}{\mu }_0^j}^{pR}\stackrel{}{z}^{n_{pR}(\stackrel{}{e}_1)}\},`$ (8.26)
where the basis vector $`\stackrel{}{e}_1=(m,n,0,0)`$. The $`E_{8}^{(1)}{}_{L}{}^{}E_{8}^{(1)}{}_{R}{}^{}`$ graph obtained from the eldest $`(1,1,4,6)[12]`$ polyhedron in chain $`XIX`$ is shown in Table 27 and Figure 13.
## 9 Perspectives on the Further Classification of $`CY_3`$ and $`K3`$ Spaces
Although a fuller study of $`CY_3`$ spaces lies outside the scope of this paper, a preliminary study is of interest here, for the following reason. In addition to the 95 $`K3`$ spaces (Table 1) related to the zeroes of single polynomials discussed in previous Sections, others can be found by ‘higher-level’ constructions as the intersections of the loci of zeroes of quasihomogeneous polynomials, which are obtainable from $`CY_3`$ spaces, as we now discuss.
When going on to consider the general construction of $`\stackrel{}{k}_5`$ projective vectors in $`CP^4`$ that describe $`CY_3`$ hypersurfaces, we start from the 95 simple extensions of these $`K3`$ vectors as well as 5 multiple extensions of lower-dimensional vectors, together with all their possible permutations. Corresponding to the previous five forms of extended vectors, one finds the following sets and permutations: quadruply-extended basic vectors with the cyclic $`C_5`$ group of permutations:
$$\stackrel{}{k}_1^{ex}=(0,0,0,0,1):|C_5|=5;$$
(9.1)
triply-extended composite vectors with the dihedral $`D_5`$ group of permutations:
$$\stackrel{}{k}_2^{ex}=(0,0,0,1,1):|D_5|=10;$$
(9.2)
the following doubly-extended composite vectors with the $`D_5^{},A_5^{}`$ and $`A_5`$ groups of permutations:
$`\stackrel{}{k}_3^{ex}`$ $`=`$ $`(0,0,1,1,1):|D_5^{}|=10,`$ (9.3)
$`\stackrel{}{k}_3^{ex}`$ $`=`$ $`(0,0,1,1,2):|A_5^{}|=30,`$ (9.4)
$`\stackrel{}{k}_3^{ex}`$ $`=`$ $`(0,0,1,2,3);|A_5|=60.`$ (9.5)
we recall that the alternating group of permutations $`A_5`$ can be identified with the icosahedral symmetry group $`I`$. All the other extended $`\stackrel{}{k}_5`$ vectors can be obtained similarly from 95 $`K3`$ vectors, utilising the symmetric group $`S_5`$ or some subgroups. The full set of extended $`\stackrel{}{k}_5`$ vectors is displayed in Table 28. As noted in its caption, the total number of extended vectors is 10 270.
As an illustration how our method may be used to classify $`CY_3`$ manifolds, we now describe briefly how to obtain the complete list of $`\stackrel{}{k}_5`$ vectors with $`K3`$ intersections, which we find to be distributed in 4242 chains. To build the chains for $`CY_3`$ which have a double-vector structure, each of which is parametrized by a pair of positive integers, should find the ‘good’ pairs of ‘extended’ vectors (i.e., those whose intersection gives a reflexive $`K3`$ hypersurface), which involves checking all the $`10270\times 10271/2=52731315`$ possible pairs of vectors from Table 28. It was just such a search by computer that led to the 4242 double chains mentioned above, together with their eldest vectors. For more complete information about these chains, see .
These chains give many $`CY_3`$ projective vectors, but not all. The complete list also includes the ‘good’ triples which have elliptic fibres. This requires looking for good triples among the following five types of five-dimensional extended vectors:
$`1.`$ $`(0,0,0,0,1)\mathrm{\hspace{0.17em}5},`$ (9.6)
$`2.`$ $`(0,0,0,1,1)\mathrm{\hspace{0.17em}10},`$
$`3.`$ $`(0,0,1,1,1)\mathrm{\hspace{0.17em}10},`$
$`4.`$ $`(0,0,1,1,2)\mathrm{\hspace{0.17em}30},`$
$`5.`$ $`(0,0,1,2,3)\mathrm{\hspace{0.17em}60},`$
where the number after the arrow on each line of (9.6) corresponds to the number of permutations in each case. We have found 259 such good triples, together with their eldest vectors, corresponding to 259 elliptic chains. The union of the $`K3`$ and elliptic projective vectors still does not yield the full dual set of $`\stackrel{}{k_5}`$ projective vectors. We must also construct another set of chains using quadruples from among the following multiply-extended vectors:
$`1.`$ $`(0,0,0,0,1)\mathrm{\hspace{0.17em}5},`$ (9.7)
$`2.`$ $`(0,0,0,1,1)\mathrm{\hspace{0.17em}10},`$
The number of $`CY_3`$ chains found in this way is just six.
In addition to these $`4242`$ double, $`259`$ triple and $`6`$ quadruple $`CY_3`$ chains (to be compared with the $`22`$ double and $`4`$ triple $`K3`$ chains found previously), one must find all the vectors whose intersection contains only one central interior point (to be compared with the exceptional $`K3`$ vector $`(7,8,9,12)`$). We have found just two such examples in the case of $`CY_3`$, namely $`(41,48,51,52,64)`$ and $`(51,60,64,65,80)`$, again using the intersection-projection duality technique. The eldest vectors for all the $`CY_3`$ projective vector chains we have found can be obtained from .
In the cases of dimension higher than three, the concept of intersection-projection duality is richer, and leads to one important and by now well-known consequence , namely the isomorphism between different homology groups for dual pairs of $`CY_d`$ manifolds $`M,M^{}`$, and specifically the following relation:
$$H^{p,q}(M)H^{dp,q}(M^{})$$
(9.8)
for $`0p,qd`$. We leave a more complete discussion of duality of $`CY_3`$ spaces to future work, limiting our discussion here of their ramifications for the classification of $`K3`$
Our construction based on the 10 270 extended vectors obtained from the $`100(=95+5)`$ types of projective vectors in lower dimensions $`n=1,2,3,4`$ shown in Table 28 yielded all the $`4242(259,\mathrm{})`$ eldest vectors representing $`CY_3`$ spaces with $`K3`$ (elliptic, …) fibers. However, this method of construction simultaneously provides a new higher-level list of $`K3`$ spaces defined by planar polyhedra. To explain this, let us first assign to all $`K3`$ spaces defined by $`n`$-dimensional projective vectors level zero, and denote them by $`\mathrm{\Pi }_0`$. Then, level one $`K3`$ spaces consist of all the ‘good’ intersections<sup>††</sup><sup>††</sup>††In the sense that they give $`n`$-dimensional reflexive polyhedra. of two $`(n+1)`$-dimensional extended vectors, denoted by $`\mathrm{\Pi }_1`$. This yields a list of reflexive polyhedra that is more complete than the previous list of polyhedra obtained from $`n`$-dimensional projective vectors, i.e., $`\mathrm{\Pi }_0\mathrm{\Pi }_1`$. Continuing, one may define the set of all ‘good’ intersections of level two, $`\mathrm{\Pi }_2`$, by considering the intersections of three $`(n+2)`$-dimensional extended vectors, and similarly for the higher levels $`3,4,\mathrm{}`$:
$`\mathrm{\Pi }_0\mathrm{\Pi }_1\mathrm{\Pi }_2\mathrm{}\mathrm{\Pi }_{last}`$ (9.9)
until this process gives us no new reflexive polyhedra. Since the number of distinct reflexive polyhedra in any dimension is finite, e.g., the maximal number of integer points for planar polyherdra is 10, for $`K3`$ polyhedra it is 39, etc., there exists a maximum last level, after which one cannot find any new types of polyhedra.
Following this approach in the simple case of $`CY_1`$ spaces, we recall that we found three planar polyhedra (triangles) at level zero, determined by the three projective vectors $`(1,1,1)`$, $`(1,1,2)`$ and $`(1,2,3)`$. At level one, constructing the 22 chains of $`K3`$ projective vectors via the 22 ‘good’ intersections of the five types of four-dimensional extended vectors, we now find 7 new planar polyhedra in 9 of the 22 two-vector $`K3`$ chains, differing from the previous three triangles by the numbers of vertices $`(V,V^{})`$ and/or by the numbers of integer points $`(N,N^{})`$ and/or by the areas of these planar polyhedra, as shown in Table 29. To look for further new polyhedra at level 2, one should consider the five following types of vectors: $`(1),(1,1),(1,1,1),(1,1,2)`$ and $`(1,2,3)`$, extended to five dimensions. Taking into account all the 50 possible permutations, and looking for the ‘good’ triple intersections, we find among the 259 ‘good’ planar reflexive polyhedra mentioned above just three distinct new polyhedra, which are exhibited in Table 30.
Extending this procedure, we found among the 4242 chains of $`CY_3`$ spaces with ‘good’ intersections 730 new $`K3`$ polyhedra at level one, many with multiple realizations as in Tables 29 and 30. As an example how such new $`K3`$ spaces emerge, consider the following two-vector $`CY_3`$ chain: $`m(0,1,1,4,6)+n(1,0,1,4,6)`$. The maximum values of $`m`$ and $`n`$ are determined by the dimensions of these extended vectors, namely $`d=12`$. This chain contains 46 different $`\stackrel{}{k}_5`$ projective vectors. The four-dimensional pentahedroid corresponding to the eldest vector in this chain is shown in Figure 14. As can be seen there, in addition to its 5 vertices, the pentahedroid has 10 one-dimensional edges, 10 two-dimensional triangular faces, and 5 three-dimensional tetrahedral facets. This pentahedroid contains two realizations of the tetrahedron corresponding to $`\stackrel{}{k}_4=(1,1,4,6)`$, whose intersection contains an elliptic fibre corresponding to $`\stackrel{}{k}_3=(1,2,3)`$.
A snapshot of the complete $`m(0,1,1,4,6)+n(1,0,1,4,6)`$ chain is shown in Figure 15, where the number of points $`N`$ in each member of the chain is plotted as a function of $`d=k_1+k_2+k_3+k_4+k_5`$ for each of the allowed values of $`m`$. We note a systematic tendency for $`N`$ to decrease as $`d`$ increases. (The structure of the chain is, of course, symmetric under the interchange: $`nm`$). The corresponding plot for the dual polyhedra is shown in Figure 16: here we see that the number of points $`N^{}`$ increases as $`d`$ increases.
To get another impression of the rich new structures emerging at levels one and above, we consider a ‘tetrahedron subalgebra’ of our $`K3`$ algebra, i.e., we consider only those projective vectors corresponding to point- and segment-polyhedra, triangles and tetrahedra. With this restriction, we start from only 32 $`K3`$ projective vectors, corresponding to four-vertex tetrahedra and five of our previous extended vectors. In this way, the number of reflexive polyhedra at level one is reduced to just 632, consisting of 460 tetrahedra and 172 reflexive polyhedra with numbers of vertices between 5 and 10. In this list of 632 polyhedra, there are actually only 146 distinct new types of polyhedra, as shown in Table 31. More information about them can be obtained from : we leave their more detailed study to later work.
The method described here has a very simple geometrical interpretation. According to the chain structure, each $`CY_3`$ can have a complex internal structure, and correspondingly its vector can be extended as a sum of two $`K3`$, three elliptic, four two-component or five single-component extended vectors. Another nice feature of this chain structure is that it gives us complete information about the integer lattice which determines all the CY equations. Moreover, it also gives us the possibility of summarizing the singularity structure of $`CY_3`$ spaces. As we discussed in Section 8, the $`K3`$ polyhedron structure gives us a systematic way of classifying the corresponding Cartan-Lie algebra graphs. It will be interesting to make a full corresponding analysis for $`CY_3`$ hypersurfaces, taking duality into account. This method could also provide the full classification of Betti-Hodge topological numbers for $`CY_3`$ manifolds. Moreover, this algebraic method enables us to ‘walk’ between different dimensions, e.g., to classify $`CY_4`$, $`\mathrm{}_5`$,… manifolds (Figure 1). The greatest limitations may be our abilities to analyze this algebra and/or the available computer resources.
A fuller analysis of our structural classification of the $`\stackrel{}{k}_5`$ vectors for $`CY_3`$ manifolds will be given in later work. An important aspect of this procedure is that we can study the structures of the positive-integer lattices which correspond to the $`\stackrel{}{k}`$ vectors, introducing the corresponding modular (for two-dimensional sublattices) and hypermodular (for 3-, 4- or higher-dimensional lattices) transformations. These yield duality groups that are more general than the well-known $`S,T`$ and $`U`$ dualities, including them as subgroups. Moreover, the study of the geometric properties of the one-dimensional complex torus, two-dimensional $`K3`$ hypersurfaces and Calabi-Yau manifolds with dimensions $`d=3,4,\mathrm{}`$ gives insight into the possible rank and dimensions of the Lie algebras which may be important for the understanding of the nature of the symmetries used in high-energy physics.
Acknowledgements
G.V. thanks H. Dahmen, G. Harigel, L. Fellin, V. Petrov, A. Zichichi and the CERN Theory Division for support. The work of D.V.N. was supported in part by DOE grant no. DE-FG-0395ER40917. |
warning/0002/math-ph0002028.html | ar5iv | text | # Untitled Document
Existence of algebraic decay in non-Abelian ferromagnets
A. Patrascioiu
Physics Department and Center for the Study of Complex Systems
University of Arizona, Tucson, AZ 85721
(Received December 1, 1991)
The low temperature regime of non-Abelian two dimensional ferromagnets is investigated. The method involves mapping such models into certain site-bond percolation processes and using ergodicity in a novel fashion. It is concluded that all ferromagnets possessing a continuous symmetry (Abelian or not) exhibit algebraic decay of correlations at sufficiently low temperatures.
PACS: 05.07.Fh.05.50+q.11.15Ha.64.60t
In a recent letter Seiler and I proposed studying the phase structure of the 2D $`O(N)`$ models by mapping them into a correlated site-bond percolation problem. This approach was applied to certain discrete spin modes and to the $`O(2)`$ model, for which we rederived the Froehlich and Spencer result regarding the existence of a massless phase at sufficiently low temperatures $`1/\beta `$. In this paper I report an extension of the percolation approach to $`O(N)`$ $`N3`$. It leads to the conclusion that a massless phase exists in all $`O(N)`$ models.
For completeness I will repeat the main points of Ref. (see also Ref. for a more complete discussion). With any $`O(N)`$ spin configuration one can associate an Ising spin configuration by dividing the sphere $`S(N1)`$ into two hemispheres and introducing an Ising variable $`\sigma =\pm 1`$, which specifies in which hemisphere the spin points. In this manner the standard nearest neighbor action (s.n.n.a.) for the $`O(N)`$ model allows rewriting the partition function as
$$Z=\underset{\{\sigma \}}{}\left(\underset{i\mathrm{\Lambda }}{}𝑑s_i𝑑\stackrel{}{s}_{pi}\right)\mathrm{exp}\left[\beta \underset{i,j}{}(s_is_j\sigma _i\sigma _j+\stackrel{}{s}_{pi}\stackrel{}{s}_{pj})\right]$$
$`(1)`$
Here $`\stackrel{}{u}`$ is the unit vector chosen for specifying the hemispherical decomposition, $`s_{}=|\stackrel{}{s}\stackrel{}{u}|`$ and $`\stackrel{}{s}_p\stackrel{}{u}=0`$. With respect to the Ising variables the action is ferromagnetic, hence amenable to the Fortuin-Kasteleyn transformation . This procedure associates to the Ising problem a correlated site-bond percolation process defined as follows:
FK1-identify clusters of like-$`\sigma `$ spins (H-clusters)
FK2-within each H-cluster occupy bonds randomly with probability $`1exp(2\beta _{ij})`$ (obtain FK-cluster)
FK3-assign to every site within a given FK-cluster the same $`\sigma `$ value, obtained by choosing randomly + or - with probability 1/2.
Here $`\beta _{ij}`$ is the space dependent inverse temperature, which for the s.n.n.a. would be $`\beta s_is_j`$. Fortuin and Kasteleyn proved that the mean FK-cluster size (expected size of the cluster attached to the origin) equals the magnetic susceptibility of the Ising variable
$$\chi _{_{Is}}\frac{1}{|\mathrm{\Lambda }|}\underset{x,y\mathrm{\Lambda }}{}\sigma _x\sigma _y.$$
$`(2)`$
In particular the latter diverges when the mean FK-cluster size diverges.
To apply the F-K procedure to the $`O(N)`$ models, Seiler and I considered a modified model called ‘cut’ action: the Gibbs factor is s.n.n.a. only if $`|\stackrel{}{s}_i\stackrel{}{s}_j|<ϵ,0<ϵ<2`$ and 0 otherwise. We then formulated the following three conjectures:
C1: The Mermin-Wagner theorem applies to the ‘cut’ model.
C2: The $`O(N)`$ models (‘cut’ or not) are ergodic.
C3: On a triangular lattice $`T`$ a percolation process produced by a measure enjoying the symmetries of the lattice can contain at most one percolating cluster.
I refer the reader to Refs. and for a thorough discussion of the motivations behind these three conjectures and of the comparison of the ‘cut’ and the s.n.n.a. models. I will elaborate only on C2, which is central to the arguments presented in this paper. Imagine a very large lattice on which one has used the Monte Carlo procedure to simulate the $`O(N)`$ model. If one has achieved thermalization, then this configuration is ‘typical.’ In the infinite volume limit a typical configuration has two important properties:
P1: spacial averages equal ensemble averages (Birkoff’s theorem)
P2: the configuration is (statistically) invariant under additional Monte Carlo steps.
I will briefly sketch the argument used in Ref. to prove that the ‘cut’ $`O(2)`$ model must exhibit algebraic decay of its correlation functions for $`ϵ`$ sufficiently small. In Eq. (2) let $``$ stand for expectation value measured with the full Gibbs measure. By P1 and P2, $`\chi _{_{Is}}`$ can be computed as a quenched expectation value provided the spins $`s_{i||}`$ are assigned the values of a typical configuration. Since the Gibbs measure is invariant under lattice translations and (discrete) rotations, by C1 and C3 a typical configuration cannot contain a percolating H-cluster. An interesting theorem by Russo states that if a translational invariant percolation process on a $`T`$ lattice is such that neither clusters of the set $`E`$ nor of its complement $`\overline{E}`$ percolate, then the mean cluster size of both $`E`$ and $`\overline{E}`$ must diverge. Taking $`E`$ to stand for $`\sigma =+1`$ and $`\overline{E}`$ for $`\sigma =1`$ shows that the mean size of the H-clusters must diverge. (This statement is not surprising since at $`\beta =0`$ and $`ϵ=2`$ the $`O(N)`$ model is equivalent to the Bernoulli site-percolation process with $`p=1/2`$ and for the latter the critical density on a $`T`$ lattice is indeed 1/2.) The FK-clusters are subclusters of the H-clusters obtained via rule FK2. In the ‘cut’ model, this rule must be amended. Indeed because of rule FK3, the constraint could be violated unless bonds are occupied at all sites having $`s_{}>dϵ/2`$. Therefore, in a ‘cut’ $`O(N)`$ model, the FK-clusters must contain D-clusters defined by the condition $`s_{}>d`$. In Ref. we showed that for the ‘cut’ $`O(2)`$ model simple applications of C1 and C3 required that neither $`D`$-clusters nor $`\overline{D}`$-clusters $`(s_{||}<d)`$ can percolate and then, by Russo’s theorem, both must have divergent mean size, QED. From the discussion presented thus far it follows that in any ‘cut’ $`O(N)`$ $`N>1`$ model on a $`T`$ lattice, if neither clusters of $`D`$ nor of $`\overline{D}`$ percolate, the mean FK-cluster size must diverge and hence correlations must decay algebraically. In fact in the ‘cut’ model $`D`$-clusters can never percolate. Indeed the set $`D`$ consists of two disconnected pieces, both of which are contained in H-clusters and I have already argued that H-clusters cannot percolate. Therefore, the only question is whether $`\overline{D}`$-clusters could percolate for $`ϵ`$ sufficiently small? The reason for which a topological answer to this question exists in $`O(2)`$ is that in that case the set $`\overline{D}`$ consists also of two disconnected pieces, which, for $`ϵ<\sqrt{2}`$, cannot communicate. Obviously in $`O(N)N3`$, $`\overline{D}`$ is a connected set and a new strategy must be employed. In the sequel I will state three independent arguments, that in the ‘cut’ $`O(N)`$ model $`\overline{D}`$-clusters cannot percolate for $`ϵ`$ sufficiently small or $`\beta `$ sufficiently large. Each argument requires a new conjecture and I will address their merits too.
Argument 1
This is a proof by contradiction. For simplicity I will discuss the s.n.n.a. $`O(3)`$ model $`(ϵ=2)`$ at $`\beta `$ large and choose $`\stackrel{}{u}=\widehat{z}`$. I will take $`d`$ small but independent of $`\beta `$ \- so that by FK2, when $`\beta `$ is large,the bond occupation probability for sites in $`D`$ goes to 1. I will assume that a cluster of $`\overline{D}`$ percolates and show that that assumption suggests that a certain magnetic susceptibility (Eq. (4)) diverges. To that end I introduce spherical coordinates and rewrite the partition function as
$$Z=\left(\underset{i\mathrm{\Lambda }}{}_0^\pi d\theta _i_0^{2\pi }d\phi _i\right)\mathrm{exp}\{\beta \underset{i,j}{}[\mathrm{cos}\theta _i\mathrm{cos}\theta _j+\mathrm{sin}\theta _i\mathrm{sin}\theta _j\mathrm{cos}(\phi _i\phi _j]\}$$
$`(3)`$
Consider the following susceptibility
$$\chi _\phi \frac{1}{|\mathrm{\Lambda }|}\underset{x_iy\mathrm{\Lambda }}{}\mathrm{cos}(\phi _x\phi _y)$$
$`(4)`$
By P1 and P2 $`\chi _\phi `$ could be measured by quenching the $`\theta `$ variables to the values $`\overline{\theta }`$ they would take in a typical configuration. That is
$$\chi _\phi =\frac{1}{|\mathrm{\Lambda }|}\underset{x_iy\mathrm{\Lambda }}{}\mathrm{cos}(\phi _x\phi _y)_q$$
$`(5)`$
where $`_q`$ means expectation value computed with the measure
$$\left(\underset{i\mathrm{\Lambda }}{}_0^{2\pi }𝑑\phi _i\right)\mathrm{exp}\left[\beta \underset{i,j}{}\mathrm{sin}\overline{\theta }_i\mathrm{sin}\overline{\theta }_j\mathrm{cos}(\phi _i\phi _j)\right].$$
$`(6)`$
Since the quenched model is an $`O(2)`$ model (albeit with space dependent couplings), one can employ Ginibre’s inequality to bound $`\chi _\phi `$ from below by the value it would take if in the measure (6) one replaced $`\beta \mathrm{sin}\theta _i\mathrm{sin}\theta _j`$ by 0 at all sites where $`\sqrt{\beta \mathrm{sin}\theta _i}<c`$ for some $`c>0`$. Under the assumption that $`\overline{D}`$ percolates, by C3, these sites could not possibly percolate, but would form islands. The average size of these islands relative to the average distance between them would decrease with beta. Indeed by the Mermin-Wagner theorem, the probability of finding the spin at a site taking values in some subset of the sphere $`A`$ of volume $`V(A)`$ is equal to $`V(A)/4\pi `$. (For $`\beta `$ large, one can use perturbation theory to estimate the average size of these islands, which becomes actually independent of $`\beta `$.) Thus the assumption that the equatorial strip $`\overline{D}`$ percolates implies that $`\chi _\phi `$ is bounded from below by the susceptibility of an $`O(2)`$ model at large inverse temperature, but on a lattice having some small, randomly distributed holes. Although I am not aware of any rigorous result proving that, the following conjecture seems eminently reasonable. C4: Consider a $`T`$ lattice and dilute bonds randomly with a probability smaller than the percolation probability for unoccupied bonds. Then there exists a $`\beta _{kt}<\mathrm{}`$ such that for any $`\beta >\beta _{kt}`$ the susceptibility diverges.
Before motivating this conjecture, let me say that there is no reason to expect that if in the $`O(3)`$ model $`\overline{D}`$ percolated, the polar caps would be distributed as the holes produced by a Bernoulli process. Their actual distribution would be controlled by the full $`O(3)`$ measure. However, if $`\overline{D}`$ percolated and especially if the model had a mass gap, by some central limit theorem, one would expect the polar caps to form islands and their distribution to be random at distances much larger than the correlation length.
The intuition for C4 comes from the following rigorous results:
a) Georgii proved that if one randomly dilutes sites or bonds on a regular lattice with $`D2`$, then provided a remaining cluster percolates, there exists an inverse temperature $`\beta _c<\mathrm{}`$ such that for $`\beta >\beta _c`$ there exists long range order (l.r.o.).
b) De Massi et al. , proved that under the same conditions as above, the Laplacian retains its continuous spectrum.
In the language of the Coulomb gas, my conjecture is that if one introduces in the gas perfect conductors, randomly distributed, if the perfectly conducting regions do not percolate, at sufficiently low temperatures, the Coulomb gas does not exhibit Debye screening (the introduction of the perfect conductors will only affect the dielectric constant).
To conclude this argument, the contradiction is this: if one assumes that for the $`O(3)`$ model $`\overline{D}`$ percolates and $`\chi _{_{Is}}`$ is finite, then clearly so is the $`s_z`$-susceptibility (since $`s_z1`$). On the other hand C4 strongly suggests that the $`s_xs_y`$ susceptibility would diverge when $`\beta `$ is large. This is a clear violation of $`O(3)`$ invariance, hence the assumption that $`\overline{D}`$ percolates must be false. Although not transparent, the topology of $`O(3)`$ is crucial for this argument. Indeed one may wonder if a similar reasoning could not be used to relate the $`O(2)`$ model to the Ising model and thus prove that the latter must exhibit l.r.o. at large $`\beta `$, in violation of the Mermin-Wagner theorem? The answer is no, precisely because $`\overline{D}`$ is no longer a connected set and thus it could not possibly percolate.
Argument 2
This is again a proof by contradiction. For simplicity I consider the ‘cut’ $`O(3)`$ model and choose $`\stackrel{}{u}=\widehat{z}`$. I would like to argue that if the equatorial strip $`\overline{D}`$ percolated, then $`O(3)`$ invariance would be broken.
Next let me consider the realistic case of a $`T`$ lattice and an $`ϵ`$ small, yet $`ϵ>0`$. Suppose that in fact the equatorial strip $`\overline{D}`$ does percolate and hence its complement $`D`$ forms islands. In the ‘cut’ model, the lines $`s_z=c>d`$ will have to form closed loops, nested inside these islands. Consider now a c-tilted equator, namely the great circle passing thru $`s_z=c`$ and $`s_x=0`$. Since neither the hemisphere $`s_x>0`$ nor $`s_x<0`$ can percolate, any site of the lattice must be surrounded by an infinite sequence $`X(k)kZ`$ of concentric closed loops $`s_x=0`$. (By the line $`s_x=0`$ I mean a line on the dual lattice such that $`s_{x_i}s_{x_j}0`$; same type of qualifications apply to all other lines appearing in this discussion.) $`O(3)`$ invariance requires that the average number of intersections of the $`X`$ lines with the c-tilted equators is independent of c. However if $`\overline{D}`$ percolates, then in any typical configuration there exists a $`k_0<\mathrm{}`$ such that any $`X(k)`$ line with $`k>k_0`$ intersects the percolating cluster. That means that infinitely many $`X`$ lines cross the $`c=0`$ tilted equator, while they may or may not cross the c-tilted equators with $`c>0`$. In other words if $`\overline{D}`$ percolates, then one would expect the average number of crossings of the $`X`$ lines with the c-tilted equators to decrease with c, in violation of $`O(3)`$ invariance. If on the other hand $`\overline{D}`$ does not percolate, then both $`D`$ and $`\overline{D}`$ form rings and no a priori asymmetry in the average number of crossings of the $`X`$ lines with the c-tilted equators exists. (An example where $`\overline{D}`$ percolates is the Richard model , which is a modified $`O(3)`$ model in which $`|s_z|<1b`$ for some $`b>0`$, hence this model is only $`O(2)`$ invariant. The percolation approach used in Ref. can be employed to prove rigorously that this model has to be massless for $`ϵ`$ sufficiently small - see Ref. ; $`\chi __\phi `$ diverges, yet $`\chi _{_{Is}}<\mathrm{}`$.)
In the discussion above I used the word ‘expect’ because one could say that even though if $`\overline{D}`$ percolates the regions with $`s_z>c`$ are hidden inside regions of smaller $`s_z`$ values, they are larger and thus restore $`O(3)`$ invariance. However $`O(3)`$ invariance requires that any typical configuration has the following two properties:
T1: The area is preserved.
T2: The gradient is preserved.
Property T1 means that the density of sites where the spin points in some region A is proportional to the volume $`V(A)`$. Property T2 says that if one selects two points on the sphere $`p_1`$ and $`p_2`$, separated by a distance $`L`$, the average distance between sites where the spin points in the neighborhood of $`p_1`$ respectively $`p_2`$ depends only on $`L`$ (it is independent of which $`p_1`$ and $`p_2`$ are chosen, provided they are at distance $`L`$). Obviously both properties are required by C1. C5: If in the ‘cut’ $`O(3)`$ model $`\overline{D}`$ percolated, then the typical configuration would violate $`T_1`$ or $`T_2`$ (or both) and, hence, $`O(3)`$ invariance.
The motivation for C5 is this: if $`\overline{D}`$ percolated, then, as already argued, $`D`$ would form islands - as opposed to rings, which are formed when neither $`\overline{D}`$ nor $`D`$ percolates on a $`T`$ lattice. The basic difference between a system forming islands and one forming rings is that islands are basically of finite size - the probability to find an island of diameter $`L`$ decreases exponentially with $`L`$; on the contrary, if the system forms rings, there exists an infinite sequence of clusters surrounding each other and hence no exponential suppression of large clusters. Thus if the system forms islands the typical configuration will contain mostly mappings of a hemisphere over some finite region of $`T`$. It is easy to check that such maps cannot preserve both T1 and T2. No such difficulty exists if one considers rings - arbitrarily large regions of $`T`$.
Argument 3
As I have already noted, if $`\overline{D}`$ percolates, then $`D`$ forms islands. Moreover, $`D`$ consists of two disconnected pieces $`D_u`$ and $`D_l`$. When $`ϵ`$ is sufficiently small, the volume of $`D_u`$, $`V(D_u)`$ is much larger than that of $`\overline{D}`$, $`V(\overline{D})`$. On the other hand the area of the boundary of $`D_u`$, $`S(D_u)`$ is half $`S(\overline{D})`$. Is it reasonable to expect that under these circumstances, the mean cluster size of $`D_u`$ is finite while that of $`\overline{D}`$ infinite? The answer is provided by the following conjecture: C6: In the ‘cut’ $`O(N)`$ model, if two sets $`A`$ and $`B`$ have V(A)=V(B) and $`S(A)<S(B)`$, then there exists $`ϵ_0(A,B)>0`$ such that for any $`ϵ<ϵ_0,`$ $`AB`$, where $``$ represents the mean cluster size.
The conjecture says that at given volume, the larger the surface of a set, the smaller its average cluster size. The reason for adding the qualifier that $`ϵ<ϵ_0`$ is that for $`ϵ>0`$ the surface of the clusters of a set A need not consist of points on the surface of A. I believe that this conjecture is intuitively clear. It can be proved in 1D. In 2D it was verified numerically for $`O(3)`$ as follows: A was the Northern polar cap of area $`4\pi /3`$, B the equatorial strip of the same area and $`ϵ`$ was such that the Northern and Southern polar caps could barely communicate. The data indicated that the mean cluster size of both A and B increased as $`L^{2\eta }`$ ($`L`$-linear size of the lattice) and that $`\eta _A<\eta _B`$.
If C6 is true, it cannot be true that in the ‘cut’ $`O(N)`$ models the equatorial strip $`\overline{D}`$ percolates. Indeed if $`\overline{D}`$ percolates, its mean cluster size is divergent. By C6, for $`ϵ`$ sufficiently small, so is the mean cluster size of $`D`$. By Russo’s theorem, on a $`T`$ lattice, that can occur only if neither $`D`$ nor $`\overline{D}`$ percolates. QED.
Discussion
The arguments presented above indicate that all 2D $`O(N)`$ models possess a massless phase. (This situation contradicts common wisdom. Evidence in favor of the latter is analyzed separately and found wanting.) The arguments moreover suggest that although at large $`\beta `$ extended topological defects - instantons - may exist in non-Abelian models, they are supressed entropically with respect to spin waves. This situation, already conjectured by the author in 1986 , suggests that for
$$N3$$
the 2-point function may behave as
$$\stackrel{}{s}_{}\stackrel{}{s}_xa(\beta )\frac{e^{m(\beta )x}}{\sqrt{x}}+b(\beta )\frac{1}{x^{\eta (\beta )}}.$$
$`(7)`$
I have no basis at the present time to estimate $`a(\beta )`$ and $`b(\beta )`$, nor whether $`\eta `$ depends on $`\beta `$ in any given model. However, it could be that $`a`$ and $`b`$ are such that at intermediate distances the decay is exponential to a very good approximation (a similar effect governs the time evolution of a metastable state in nonrelativistic quantum mechanics ).
Finally a word about perturbation theory. The fact that the 2D $`O(N)`$ models possess a massless phase for $`\beta `$ sufficiently large does not imply that in 2D perturbation theory fails to produce the correct asymptotic expansion at fixed distances (as it does in 1D for $`N3`$). However if one defines the Callan-Symanzik $`\beta `$-function by requiring that say $`\stackrel{}{s}(0)\stackrel{}{s}(x)/\stackrel{}{s}(0)\stackrel{}{s}(y)`$ is a renormalization group invariant for $`x,y1`$, then clearly an algebraic decay for $`\beta >\beta _{kt}(N)`$ implies that the Callan-Symanzik $`\beta `$-function could be chosen to be vanishing. If my conjecture about Eq. (7) proved to be correct, one could also define the $`\beta `$-function as $`d\beta /dln(m)`$, in which case one may find the famous asymptotic freedom answer. However I find it hard to believe that if that were the case, the continuum limit constructed by letting $`\beta \mathrm{}`$ would not contain (coupled) massless excitations (of course a continuum limit could also be constructed for any $`\mathrm{}>\beta >\beta _{kt}(N)`$ \- that field theory would be a massless theory).
Many of the ideas expressed in this paper stem from my long time collaboration with Erhard Seiler. I am also gratefulfor the hospitality extended to me by the Max Planck Institut fur Pysik und Astrophysik - Munich.
References
A. Patrascioiu and E. Seiler, Phys. Rev. Lett. 68, 1395 (1992).
J. Froehlich and T. Spencer, Comm. Math. Phys. 81, 455 (1981).
A. Patrascioiu and E. Seiler, J.Stat.Phys. 69, 55 (1992).
C. M. Fortuin and P. W. Kasteleyn, J. Phys. Soc. JPN (suppl.) 24, 86 (1969).
L. Russo, Z. Wahrsch. Verw. Gebiete 42, 39 (1978).
J. Ginibre, Comm. Math. Phys. 16, 310 (1970).
H. O. Georgii, Comm. Math. Phys. 81, 527 (1981).
A. DeMassi,P. A. Ferrari, S. Goldstein and W. D. Wick, J. Stat. Phys. 55, 787 (1989).
J.-L. Richard, Phys. Lett. B 134, 75 (1987).
A. Patrascioiu and E. Seiler, The Difference between Abelian and Non-Abelian Models: Facts and Fancy, MPI preprint, 1991, math-ph/9903038.
A. Patrascioiu, Phys. Rev. Lett. 58, 2285 (1987).
A. Patrascioiu, Phys. Rev. D 24, 496 (1981). |
warning/0002/hep-ph0002157.html | ar5iv | text | # The Proton Spin Puzzle: A Status Report∗
## 1 Introduction
Experimentally, the polarized structure functions $`g_1`$ and $`g_2`$ are determined by measuring two asymmetries:
$`A_{}={\displaystyle \frac{d\sigma ^{}d\sigma ^{}}{d\sigma ^{}+d\sigma ^{}}},A_{}={\displaystyle \frac{d\sigma ^{}d\sigma ^{}}{d\sigma ^{}+d\sigma ^{}}},`$ (1.1)
where $`d\sigma ^{}`$ ($`d\sigma ^{}`$) is the differential cross section for the longitudinal lepton spin parallel (antiparallel) to the longitudinal nucleon spin, and $`d\sigma ^{}`$ ($`d\sigma ^{}`$) is the differential cross section for the lepton spin antiparallel (parallel) to the lepton momentum and nucleon spin direction transverse to the lepton momentum and towards the direction of the scattered lepton. From the parton-model or from the OPE approach, the first moment of the polarized proton structure function
$`\mathrm{\Gamma }_1^p(Q^2){\displaystyle _0^1}g_1^p(x,Q^2)𝑑x,`$ (1.2)
can be related to the combinations of the quark spin components via
$`\mathrm{\Gamma }_1^p={\displaystyle \frac{1}{2}}{\displaystyle \underset{q}{}}e_q^2\mathrm{\Delta }q(Q^2)={\displaystyle \frac{1}{2}}{\displaystyle \underset{q}{}}e_q^2p,s|\overline{q}\gamma _\mu \gamma _5q|p,ss^\mu ,`$ (1.3)
where $`\mathrm{\Delta }q`$ represents the net helicity of the quark flavor $`q`$ along the direction of the proton spin in the infinite momentum frame:
$`\mathrm{\Delta }q={\displaystyle _0^1}\mathrm{\Delta }q(x)𝑑x{\displaystyle _0^1}\left[q^{}(x)+\overline{q}^{}(x)q^{}(x)\overline{q}^{}(x)\right]𝑑x.`$ (1.4)
At energies $`Q^210\mathrm{GeV}^2`$ or smaller, only three light flavors are relevant:
$`\mathrm{\Gamma }_1^p(Q^2)={\displaystyle \frac{1}{2}}\left({\displaystyle \frac{4}{9}}\mathrm{\Delta }u(Q^2)+{\displaystyle \frac{1}{9}}\mathrm{\Delta }d(Q^2)+{\displaystyle \frac{1}{9}}\mathrm{\Delta }s(Q^2)\right).`$ (1.5)
Other information on the quark polarization is available from the low-energy nucleon axial coupling constants $`g_A^3`$ and $`g_A^8`$:
$`g_A^3(Q^2)`$ $``$ $`p,s|\overline{u}\gamma _\mu \gamma _5u\overline{d}\gamma _\mu \gamma _5d|p,ss^\mu =\mathrm{\Delta }u(Q^2)\mathrm{\Delta }d(Q^2),`$ (1.6)
$`g_A^8(Q^2)`$ $``$ $`p,s|\overline{u}\gamma _\mu \gamma _5u+\overline{d}\gamma _\mu \gamma _5d2\overline{s}\gamma _\mu \gamma _5s|p,ss^\mu =\mathrm{\Delta }u(Q^2)+\mathrm{\Delta }d(Q^2)2\mathrm{\Delta }s(Q^2).`$
Since there is no anomalous dimension associated with the axial-vector currents $`A_\mu ^3`$ and $`A_\mu ^8`$, the non-singlet couplings $`g_A^3`$ and $`g_A^8`$ do not evolve with $`Q^2`$ and hence can be determined at $`q^2=0`$ from low-energy neutron and hyperon beta decays. Under SU(3)-flavor symmetry, the non-singlet couplings are related to the SU(3) parameters $`F`$ and $`D`$ by
$`g_A^3=F+D,g_A^8=\mathrm{\hspace{0.17em}3}FD.`$ (1.7)
We use the updated coupling $`g_A^3=1.2670\pm 0.0035`$ and the values
$`F=0.463\pm 0.008,D=\mathrm{\hspace{0.17em}0.804}\pm 0.008,F/D=\mathrm{\hspace{0.17em}0.576}\pm 0.016`$ (1.8)
to obtain $`g_A^8=0.585\pm 0.025`$ .
Prior to the EMC measurement of polarized structure functions, a prediction for $`\mathrm{\Gamma }_1^p`$ was made based on the assumption that the strange sea in the nucleon is unpolarized, i.e., $`\mathrm{\Delta }s=0`$. It follows from (1.5) and (1.6) that
$`\mathrm{\Gamma }_1^p(Q^2)={\displaystyle \frac{1}{12}}g_A^3+{\displaystyle \frac{5}{36}}g_A^8.`$ (1.9)
This is the Ellis-Jaffe sum rule : $`\mathrm{\Gamma }_1^p=0.185\pm 0.003`$ in the absence of QCD corrections and equals to $`0.171\pm 0.006`$ at $`Q^2=10\mathrm{GeV}^2`$ to leading-order corrections. The 1987 EMC experiment then came to a surprise. The result published in 1988 and later indicated that $`\mathrm{\Gamma }_1^p=0.126\pm 0.018`$, substantially lower than the expectation from the Ellis-Jaffe conjecture. From the EMC measurement of $`\mathrm{\Gamma }_1^p`$, we obtain
$`\mathrm{\Delta }u=\mathrm{\hspace{0.17em}0.77}\pm 0.06,\mathrm{\Delta }d=0.49\pm 0.06,\mathrm{\Delta }s=0.15\pm 0.06,`$ (1.10)
and
$`g_A^0(Q^2)`$ $``$ $`p,s|\overline{u}\gamma _\mu \gamma _5u+\overline{d}\gamma _\mu \gamma _5d+\overline{s}\gamma _\mu \gamma _5s|p,ss^\mu `$ (1.11)
$`=`$ $`\mathrm{\Delta }u(Q^2)+\mathrm{\Delta }d(Q^2)+\mathrm{\Delta }s(Q^2)\mathrm{\Delta }\mathrm{\Sigma }(Q^2)=\mathrm{\hspace{0.17em}0.14}\pm 0.18`$
at $`Q^2=10.7\mathrm{GeV}^2`$. The EMC results exhibit two surprising features: The strange-sea polarization is sizeable and negative, and the total contribution of quark helicities to the proton spin is small and consistent with zero. This is sometimes referred to as the “proton spin crisis”.
The so-called “proton spin crisis” is not pertinent since the proton helicity content explored in the DIS experiment is, strictly speaking, defined in the infinite momentum frame in terms of QCD current quarks and gluons, whereas the spin structure of the proton in the proton rest frame is referred to the constituent quarks. That is, the quark helicity $`\mathrm{\Delta }q`$ defined in the infinite momentum frame is generally not the same as the constituent quark spin component in the proton rest frame, just like that it is not sensible to compare apple with orange. What trigged by the EMC experiment is the “proton helicity decomposition puzzle” rather than the “proton spin crisis” (for a review, see ). The non-relativistic SU(6) constituent quark model predicts $`\mathrm{\Delta }\mathrm{\Sigma }^{}\mathrm{\Delta }U+\mathrm{\Delta }D=1`$ ($`\mathrm{\Delta }U`$ denoting the constituent up quark spin and likewise for $`\mathrm{\Delta }D`$), but its prediction for the axial-vector coupling constant $`g_A^3=\frac{5}{3}`$ is too large compared to the measured value of $`1.2670\pm 0.0035`$ . In the relativistic quark model, the proton is no longer a low-lying $`S`$-wave state since the quark orbital angular momentum is nonvanishing due to the presence of quark transverse momentum in the lower component of the Dirac spinor. Realistic models e.g. the cloudy bag model , predict $`\mathrm{\Delta }\mathrm{\Sigma }^{}0.60`$; that is, about 40% of the proton spin is carried by the orbital angular momentum of the constituent quarks. In the parton picture, the naive expectation of $`\mathrm{\Delta }\mathrm{\Sigma }`$, which is equal to $`g_A^8=0.59`$ under the assumption of vanishing sea polarization, is very close to the relativistic quark model’s prediction of $`\mathrm{\Delta }\mathrm{\Sigma }^{}`$. One of the main theoretical problems is that hard gluons cannot induce sea polarization perturbatively for massless quarks due to helicity conservation. Hence, it is difficult to accommodate a large strange-sea polarization in the naive parton model.
## 2 Experimental Progress
Before 1993 it took 5 years on the average to carry out a new polarized DIS experiment (see Table I). This situation was dramatically changed after 1993. Many new experiments measuring the nucleon and deuteron spin-dependent structure functions became available. The experimental progress is certainly quite remarkable in the past years.
Since experimental measurements only cover a limited kinematic range, an extrapolation to unmeasured $`x0`$ and $`x1`$ regions is necessary. At small $`x`$, a Regge behavior $`g_1(x)x^{\alpha (0)}`$ is conventionally assumed in earlier experimental analyses. The improvement by the new measurements is two-folds: First, the small $`x`$ region has been pushed down to the order $`10^3`$ in SMC experiments (see Table I). Second, an extrapolation to the unmeasured small $`x`$ region is done by performing a NLO QCD fit rather than by assuming Regge behaviour. The uncertainties of $`\mathrm{\Gamma }_1`$ which mostly arise from the small $`x`$ extrapolation are substantially reduced. From Table I it is also clear that the EMC experiment has been confirmed by all successive polarized DIS measurements.
Comparing to the original EMC measurement, the statistic and systematic errors of the combined world average for $`\mathrm{\Gamma }_1^p`$ are substantially reduced. The result is $`\mathrm{\Gamma }_1^p=(0.120.13)\pm 0.007`$ at $`Q^2=(510)\mathrm{GeV}^2`$. Consequently, $`\mathrm{\Delta }\mathrm{\Sigma }=(0.200.30)\pm 0.04`$. For example, $`\mathrm{\Delta }\mathrm{\Sigma }=0.25\pm 0.04`$ leads to
$`\mathrm{\Delta }u=\mathrm{\hspace{0.17em}0.81}\pm 0.01,\mathrm{\Delta }d=0.45\pm 0.01,\mathrm{\Delta }s=0.11\pm 0.01.`$ (2.1)
We will employ (2.1) as the benchmarked values for $`\mathrm{\Delta }q`$ in ensuing discussions.
The Bjorken sum rule evaluated up to $`\alpha _s^3`$ for three light flavors is
$`\mathrm{\Gamma }_1^p(Q^2)\mathrm{\Gamma }_1^n(Q^2)={\displaystyle \frac{1}{6}}{\displaystyle \frac{g_A}{g_V}}\left[1{\displaystyle \frac{\alpha _s(Q^2)}{\pi }}{\displaystyle \frac{43}{12}}\left({\displaystyle \frac{\alpha _s(Q^2)}{\pi }}\right)^220.22\left({\displaystyle \frac{\alpha _s(Q^2)}{\pi }}\right)^3\right].`$ (2.2)
A serious test of the Bjorken sum rule, which is a rigorous consequence of QCD, became possible since 1993. The current experimental results are
$`\mathrm{E143}\text{[17]}:`$ $`\mathrm{\Gamma }_1^p\mathrm{\Gamma }_1^n=\mathrm{\hspace{0.17em}0.164}\pm 0.021,`$
$`\mathrm{SMC}\text{[20]}:`$ $`\mathrm{\Gamma }_1^p\mathrm{\Gamma }_1^n=\mathrm{\hspace{0.17em}0.174}_{0.012}^{+0.024},`$ (2.3)
at $`Q^2=5\mathrm{GeV}^2`$, to be compared with the prediction
$`\mathrm{\Gamma }_1^p\mathrm{\Gamma }_1^n`$ $`=`$ $`0.181\pm 0.003`$ (2.4)
at the same energies. Therefore, the Bjorken sum rule has been confirmed by data to an accuracy of 10% level.
The quark polarization $`\mathrm{\Delta }q`$ is usually determined from the inclusive data of $`g_1`$ by assuming SU(3) flavor symmetry. Moreover, inclusive DIS determines the sum of polarized quark and antiquark distributions, but not the valence and sea quark spin distributions. Semi-inclusive polarized experiments in principle allow the determination of $`\mathrm{\Delta }q`$ for each flavor and disentangle valence and sea polarizations separately . Hence, SU(3) flavor symmetry can be tested by comparing the measured first moments of the flavor distributions to the SU(3) predictions . Semi-inclusive data are available from SMC and HERMES (see Table II). In order to present the experimental results, we digress for a moment to adopt a different definition for quark spin densities here: $`\mathrm{\Delta }q_s(x)=q_s^{}(x)q_s^{}(x)`$ and $`\mathrm{\Delta }\overline{q}(x)=\overline{q}^{}(x)\overline{q}^{}(x)`$, where $`\mathrm{\Delta }q_s(x)=\mathrm{\Delta }q(x)\mathrm{\Delta }q_v(x)`$ is the sea spin distribution for the quark flavor $`q`$. The SMC analysis is based on the assumption of SU(3) flavor symmetric sea: $`\mathrm{\Delta }\overline{u}(x)=\mathrm{\Delta }\overline{d}(x)=\mathrm{\Delta }\overline{s}(x)=\mathrm{\Delta }u_s(x)=\mathrm{\Delta }d_s(x)=\mathrm{\Delta }s(x)`$, while the HERMES results shown in Table II rely on the assumption of flavor independent polarization:
$`{\displaystyle \frac{\mathrm{\Delta }u_s(x)}{u_s(x)}}={\displaystyle \frac{\mathrm{\Delta }d_s(x)}{d_s(x)}}={\displaystyle \frac{\mathrm{\Delta }s(x)}{s(x)}}={\displaystyle \frac{\mathrm{\Delta }\overline{u}(x)}{\overline{u}(x)}}={\displaystyle \frac{\mathrm{\Delta }\overline{d}(x)}{\overline{d}(x)}}={\displaystyle \frac{\mathrm{\Delta }\overline{s}(x)}{\overline{s}(x)}}.`$ (2.5)
Note that the ansatz<sup>1</sup><sup>1</sup>1This ansatz is generally not fulfilled by the assumption of flavor independent polarization made by HERMES . of $`\mathrm{\Delta }\overline{q}(x)=\mathrm{\Delta }q_s(x)`$ has to be made in both experiments in order to extract the valence quark polarization $`\mathrm{\Delta }q_v`$ from the measurement of $`\mathrm{\Delta }q`$ and $`\mathrm{\Delta }\overline{q}`$, i.e., $`\mathrm{\Delta }q_v(x)=\mathrm{\Delta }q(x)\mathrm{\Delta }\overline{q}(x)`$. However, one caveat has to be mentioned: A priori the antiqaurk spin $`\mathrm{\Delta }\overline{q}`$ can be different from the sea polarization $`\mathrm{\Delta }q_s`$ if they are not produced from gluons. For example, antiquarks are not polarized in the model of . Under the assumption of SU(3) symmetric sea polarization, we are led to the predictions $`g_A^3=\mathrm{\Delta }u_v\mathrm{\Delta }d_v`$, $`g_A^8=\mathrm{\Delta }u_v+\mathrm{\Delta }d_v`$ and hence $`\mathrm{\Delta }u_v=2F=0.93\pm 0.02`$ and $`\mathrm{\Delta }d_v=FD=0.34\pm 0.02`$. Note that the valence polarization $`\mathrm{\Delta }q_v`$ should be scale independent.
The measurement of the gluon spin $`\mathrm{\Delta }G`$ by all possible means is very important both theoretically and experimentally (see for various processes sensitive to the gluon spin distributions). A global fit to the present inclusive DIS data of $`g_1(x)`$ cannot even fix the sign of $`\mathrm{\Delta }G`$ decisively (see Sec. 3.3), not mentioning its magnitude. One way of measuring $`\mathrm{\Delta }G(x)`$ directly is via the photon gluon fusion process occurred in the semi-inclusive DIS reaction. A recent HERMES measurement of the longitudinal spin asymmetry in photoproduction of pairs of hadrons with high transverse momentum indicates that $`\mathrm{\Delta }G(x)/G(x)=0.41\pm 0.18\pm 0.03`$ at $`x=0.17`$ to LO QCD. Hence $`\mathrm{\Delta }G(x)`$ is found to be positive in the intermediate $`x`$ region.
## 3 Theoretical Progress
In my opinion there are four important progresses in theory:
* The role played by the gluon to the first moment of $`g_1`$ is clarified. There are two extreme factorization schemes of interest.
* First-principles calculations of the quark spin and orbital angular momentum by lattice QCD became available.
* A complete and consistent NLO analysis of $`g_1`$ data became possible.
* Evoluation and gauge dependence of the quark orbital angular momentum are explored.
### 3.1 Anomalous gluon and sea quark interpretations
#### 3.1.1 Anomalous gluon interpretation
We see from Sec. II that the polarized DIS data indicate that the fraction of the proton spin carried by the light quarks inside the proton is $`\mathrm{\Delta }\mathrm{\Sigma }=(0.200.30)`$ and the strange-quark polarization is $`\mathrm{\Delta }s0.10`$ at $`Q^2=(510)\mathrm{GeV}^2`$. The question is what kind of mechanism can generate a sizeable and negative sea polarization. It is difficult, if not impossible, to accommodate a large $`\mathrm{\Delta }s`$ in the naive parton model because massless quarks and antiquarks created from gluons have opposite helicities owing to helicity conservation. This implies that sea polarization for massless quarks cannot be induced perturbatively from hard gluons, irrespective of gluon polarization. It is also unlikely that the observed $`\mathrm{\Delta }s`$ comes solely from nonperturbative effects or from chiral-symmetry breaking due to nonvanishing quark masses.
As an attempt to understand the polarized DIS data, we consider QCD corrections to the polarized proton structure function $`g_1^p(x)`$. To the next-to-leading order (NLO) of $`\alpha _s`$, the expression for $`g_1^p(x)`$ is
$`g_1^p(x,Q^2)`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle e_q^2\left[\mathrm{\Delta }C_q(x,\alpha _s)\mathrm{\Delta }q(x,Q^2)+\mathrm{\Delta }C_G(x,\alpha _s)\mathrm{\Delta }G(x,Q^2)\right]}`$
$`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle }e_q^2[\mathrm{\Delta }q^{(0)}(x,Q^2)+\mathrm{\Delta }q^{(1)}(x,Q^2)+\mathrm{\Delta }q_s^G(x,Q^2)`$
$`+\mathrm{\Delta }C_q^{(1)}(x,\alpha _s)\mathrm{\Delta }q^{(0)}(x,Q^2)+\mathrm{\Delta }C_G^{(1)}(x,\alpha _s)\mathrm{\Delta }G(x,Q^2)+\mathrm{}],`$
where uses of $`\mathrm{\Delta }C_q^{(0)}(x)=\delta (1x)`$ and $`\mathrm{\Delta }C_G^{(0)}(x)=0`$ have been made, $``$ denotes the convolution
$`f(x)g(x)={\displaystyle _x^1}{\displaystyle \frac{dy}{y}}f\left({\displaystyle \frac{x}{y}}\right)g(y),`$ (3.2)
and $`\mathrm{\Delta }C_q`$, $`\mathrm{\Delta }C_G`$ are short-distance quark and gluon coefficient functions, respectively. More specifically, $`\mathrm{\Delta }C_G^{(1)}`$ arises from the hard part of the polarized photon-gluon cross section, while $`\mathrm{\Delta }C_q^{(1)}`$ from the short-distance part of the photon-quark cross section. Contrary to the coefficient functions, $`\mathrm{\Delta }q_s^G(x)`$ and $`\mathrm{\Delta }q^{(1)}(x)`$ come from the soft part of polarized photon-gluon and photon-quark scatterings, respectively. Explicitly, they are given by
$`\mathrm{\Delta }q^{(1)}(x,Q^2)=\mathrm{\Delta }\varphi _{q/q}^{(1)}(x)\mathrm{\Delta }q^{(0)}(x,Q^2),\mathrm{\Delta }q_s^G(x,Q^2)=\mathrm{\Delta }\varphi _{q/G}^{(1)}(x)\mathrm{\Delta }G(x,Q^2),`$ (3.3)
where $`\mathrm{\Delta }\varphi _{j/i}(x)`$ is the polarized distribution of parton $`j`$ in parton $`i`$. Diagrammatically, $`\mathrm{\Delta }\varphi _{q/q}^{(1)}`$ and $`\mathrm{\Delta }\varphi _{q/G}^{(1)}`$ are depicted in Fig. 1.
The photon-gluon scattering box diagram is ultraviolet finite but it depends on the choice of infrared and collinear regulators. However, the hard part of the box diagram should be soft-regulator independent. In the improved parton model, this is done by introducing a factorization scale $`\mu _{\mathrm{fact}}`$ so that the region k2>
μfact2subscriptsuperscript𝑘2perpendicular-to>
superscriptsubscript𝜇fact2k^{2}_{\perp}{\ \lower-1.2pt\vbox{\hbox{\hbox to0.0pt{$>$\hss}\lower 5.0pt\vbox{\hbox{$\sim$}
}}}\ }\mu_{\rm fact}^{2} contributes to the hard photon-gluon cross section. The results are (see e.g. )
$`\mathrm{\Delta }C_G^{(1)}(x,Q^2,\mu _{\mathrm{fact}}^2)_{\mathrm{CI}}`$ $`=`$ $`(2x1)\left(\mathrm{ln}{\displaystyle \frac{Q^2}{\mu _{\mathrm{fact}}^2}}+\mathrm{ln}{\displaystyle \frac{1x}{x}}1\right),`$
$`{\displaystyle _0^1}𝑑x\mathrm{\Delta }C_G^{(1)}(x)_{\mathrm{CI}}`$ $`=`$ $`{\displaystyle \frac{\alpha _s}{2\pi }},`$ (3.4)
where the reason for introducing the subscript “CI” will become clear below. Hence,
$`\mathrm{\Gamma }_1^p(Q^2){\displaystyle _0^1}𝑑xg_1^p(x,Q^2)={\displaystyle \frac{1}{2}}\left(1{\displaystyle \frac{\alpha _s}{\pi }}\right){\displaystyle \underset{i}{}}\left[\mathrm{\Delta }q_i(Q^2)_{\mathrm{CI}}{\displaystyle \frac{\alpha _s(Q^2)}{2\pi }}\mathrm{\Delta }G(Q^2)\right].`$ (3.5)
The $`(1\frac{\alpha _s}{\pi })`$ term in Eq. (3.5) comes from the QCD loop correction, while the $`\alpha _s\mathrm{\Delta }G`$ term arises from the box diagram of photon-gluon scattering. If the gluon polarization inside the proton is positive, a partial cancellation between $`\mathrm{\Delta }q_{\mathrm{CI}}`$ and $`\frac{\alpha _s}{2\pi }\mathrm{\Delta }G`$ will explain why the observed $`\mathrm{\Gamma }_1^p`$ is smaller than what naively expected from the Ellis-Jaffe sum rule. Note that unlike the usual QCD corrections, the QCD effect due to photon-gluon scattering is very special: The term $`\alpha _s\mathrm{\Delta }G`$ is conserved to the leading-order QCD evolution; that is, $`\mathrm{\Delta }G`$ grows with $`\mathrm{ln}Q^2`$, whereas $`\alpha _s`$ is inversely proportional to $`\mathrm{ln}Q^2`$.
As a consequence, Eq. (2.1) is modified to
$`\mathrm{\Delta }u_{\mathrm{CI}}{\displaystyle \frac{\alpha _s}{2\pi }}\mathrm{\Delta }G`$ $`=`$ $`0.81\pm 0.01,`$
$`\mathrm{\Delta }d_{\mathrm{CI}}{\displaystyle \frac{\alpha _s}{2\pi }}\mathrm{\Delta }G`$ $`=`$ $`0.45\pm 0.01,`$ (3.6)
$`\mathrm{\Delta }s_{\mathrm{CI}}{\displaystyle \frac{\alpha _s}{2\pi }}\mathrm{\Delta }G`$ $`=`$ $`0.11\pm 0.01,`$
and
$`g_A^0=(\mathrm{\Delta }u+\mathrm{\Delta }d+\mathrm{\Delta }s)_{\mathrm{CI}}{\displaystyle \frac{3\alpha _s}{2\pi }}\mathrm{\Delta }G=\mathrm{\hspace{0.17em}0.25}\pm 0.04`$ (3.7)
at $`Q^2=510\mathrm{GeV}^2`$. Eqs. (3.1.1) and (3.7) imply that in the presence of anomalous gluon contributions, $`\mathrm{\Delta }\mathrm{\Sigma }_{\mathrm{CI}}`$ is not necessarily small while $`\mathrm{\Delta }s_{\mathrm{CI}}`$ is not necessarily large. In the absence of sea polarization and in the framework of perturbative QCD, it is easily seen that $`\mathrm{\Delta }G\mathrm{\Delta }s(2\pi /\alpha _s)2.5`$ at $`Q^2=10\mathrm{GeV}^2`$ and $`\mathrm{\Delta }\mathrm{\Sigma }_{\mathrm{CI}}0.58.`$ It thus provides a nice and simple solution to the proton spin puzzle: This improved parton picture is reconciled, to a large degree, with the constituent quark model and yet explains the suppression of $`\mathrm{\Gamma }_1^p`$, provided that $`\mathrm{\Delta }G`$ is positive and large enough. This anomalous gluon interpretation of the observed $`\mathrm{\Gamma }_1^p`$, as first proposed in (see also ), looks appealing and has become a popular explanation since 1988.
Some historical remarks are in order:
* Long before the EMC experiment, there already existed three theoretical calculations of the photon-gluon box diagram. Kodaira was the first one to compute the moments of the structure functions $`g_{1,2}`$. Since he worked in the OPE framework, there is no decomposition of $`g_A^0`$ in terms of quark and gluon spin components. The anomalous gluonic contribution to $`\mathrm{\Gamma }_1^p`$ was first put forward by Lam and Li in 1982. A calculation of the gluonic coefficient function using the dimensional regularization was first made by Ratcliffe .
* The original results for the photon-gluon scattering cross section obtained by are not perturbative QCD reliable as they depend on the choice of soft regulators. The first moment of $`\mathrm{\Delta }C_G^{(1)}(x)`$ is equal to $`\alpha _s/(2\pi )`$ in but vanishes in the work of . After the soft part below the factorization scale $`\mu _{\mathrm{fact}}`$ is removed, the gluon coefficient function is given by Eq. (3.1.1) which is soft-cutoff independent.
#### 3.1.2 Sea quark interpretation
According to the operator product expansion (OPE), the moments of structure functions can be expressed in terms of hard coefficients which are calculable by perturbative QCD and forward matrix elements of local gauge-invariant operators which are nonperturbative in nature:
$`{\displaystyle _0^1}𝑑xx^{n1}g(x)={\displaystyle \underset{i}{}}C_i^n(Q^2)N|O_i^n(0)|N.`$ (3.8)
It turns out that there is no gauge-invariant twist-2, spin-1 local gluonic operator for the first moment of $`g_1(x)`$ (see e.g. ). Here we face a dilemma here: On the one hand, the anomalous gluon interpretation sounds more attractive and is able to reconcile to a large degree with the conventional quark model; on the other hand, the sea-quark interpretation of $`\mathrm{\Gamma }_1^p`$ relies on a more solid theory of the OPE. In fact, these two popular explanations for the $`g_1^p`$ data have been under hot debate over many years before 1995. Though the OPE approach is model independent, it faces the questions of what is the deep reason for the absence of gluonic contributions to $`\mathrm{\Gamma }_1^p`$ and how are we going to understand a large and negative strange-quark polarization ?
#### 3.1.3 Factorization scheme dependence
In spite of much controversy on the aforementioned issue, this dispute was actually resolved almost a decade ago . The key point is that a different interpretation for $`\mathrm{\Gamma }_1^p`$ corresponds to a different factorization definition for the quark spin density and the hard photon-gluon cross section. The choice of the “ultraviolet” cutoff for soft contributions specifies the factorization convention.<sup>2</sup><sup>2</sup>2It is misleading to identify the regularization scheme for soft divergences, e.g. the off-shell scheme, with the factorization scheme; the former is merely employed to get rid of infrared and collinear divergences appearing in the calculation of partonic cross sections. However, the hard gluon coefficient functions are soft-cutoff independent, and they depend on the ultraviolet regulator on the triangle diagram for defining the polarized quark distribution inside the gluon. In other words, it is the choice of ultraviolet regulator rather than the soft cutoff that specifies the factorization scheme. More specifically, since $`\mathrm{\Delta }\varphi ^{(1)}`$ in Eq. (3.1) is ultraviolet divergent, it is clear that, just like the case of unpolarized deep inelastic scattering, the coefficient functions $`\mathrm{\Delta }C_q`$ and $`\mathrm{\Delta }C_G`$ depend on how the parton spin distributions $`\mathrm{\Delta }\varphi _{j/i}^{(1)}`$ are defined, or how the ultraviolet regulator is specified on $`\mathrm{\Delta }\varphi ^{(1)}`$. That is, the ambiguities in defining $`\mathrm{\Delta }\varphi _{q/q}^{(1)}`$ and $`\mathrm{\Delta }\varphi _{q/G}^{(1)}`$ are reflected on the ambiguities in extracting $`\mathrm{\Delta }C_q^{(1)}`$ and $`\mathrm{\Delta }C_G^{(1)}`$. Consequently, the decomposition of the photon-gluon and photon-quark cross sections into the hard and soft parts depends on the choice of the factorization scheme and the factorization scale $`\mu `$. Of course, the physical quantity $`g_1^p(x)`$ is independent of the factorization prescription (for a review on the issue of factorization, see ).
However, the situation for the polarized DIS case is more complicated: In addition to all the ambiguities that spin-averaged parton distributions have, the parton spin densities are subject to two extra ambiguities, namely, the axial anomaly and the definition of $`\gamma _5`$ in $`n`$ dimension. It is well known that the polarized triangle diagram for $`\mathrm{\Delta }\varphi _{q/G}^{(1)}`$ (see Fig. 1) has an axial anomaly. There are two extreme ultraviolet regulators of interest. One of them, which we refer to as the chiral-invariant (CI) factorization scheme (sometimes called as the “jet scheme” , the “parton model scheme” or the “$`k_{}`$ cut-off scheme”), respects chiral symmetry and gauge invariance but not the axial anomaly. This corresponds to a direct brute-force cutoff $`\mu `$ on the $`k_{}`$ integration in the triangle diagram ( i.e. k2<
μ2subscriptsuperscript𝑘2perpendicular-to<
superscript𝜇2k^{2}_{\perp}{\ \lower-1.2pt\vbox{\hbox{\hbox to0.0pt{$<$\hss}\lower 5.0pt\vbox{\hbox{$\sim$}
}}}\ }\mu^{2}) with $`k_{}`$ being the quark transverse momentum perpendicular to the virtual photon direction. Since the gluonic anomaly is manifested at $`k_{}^2\mathrm{}`$, it is evident that the spin-dependent quark distribution \[i.e. $`\mathrm{\Delta }q^{(1)}(x)`$\] in the CI factorization scheme is anomaly-free. Note that this is the $`k_{}`$-factorization scheme employed in the usual improved parton model .
The other ultraviolet cutoff on the triangle diagram of Fig. 1, as employed in the approach of the OPE, is chosen to satisfy gauge symmetry and the gluonic anomaly. As a result, chiral symmetry is broken in this gauge-invariant (GI) factorization scheme and the sea polarization is perturbatively induced from hard gluons via the anomaly. This perturbative mechanism for sea quark polarization is independent of the light quark masses. A straightforward calculation gives
$`\mathrm{\Delta }\varphi _{q/G}^{(1)}(x)_{\mathrm{GI}}=\mathrm{\Delta }\varphi _{q/G}^{(1)}(x)_{\mathrm{CI}}{\displaystyle \frac{\alpha _s}{\pi }}(1x),`$ (3.9)
where the term $`\frac{\alpha _s}{\pi }(1x)`$ originates from the QCD anomaly arising from the region where $`k_{}^2\mathrm{}`$. Two remarks are in order. First, this term has been erroneously claimed<sup>3</sup><sup>3</sup>3Some misconceptions in the literature about the $`\overline{\mathrm{MS}}`$ scheme have to be clarified. It has been argued that the GI scheme is pathologic and inappropriate since $`\mathrm{\Delta }C_G^{(1)}(x)_{\mathrm{GI}}`$, which is “hard” by definition, contains an unwanted “soft” term proportional to $`(1x)`$ \[see Eq. (3.12)\]. The cross section of the photon-gluon box diagram contains a term proportional to $`(1x)`$ if the infrared and collinear divergences are regulated by the quark mass or by the dimensional regulator. If the ultraviolet regulator for the triangle diagram is chirality-preserving, the $`(1x)`$ term, which arises from the region where $`k_T^2m_q^2`$, does not contribute to the hard gluon coefficient, as it should be. However, if the ultraviolet regulator preserves the axial anomaly and gauge invariance, for example, the $`\overline{\mathrm{MS}}`$ regulator, chirality will be broken and the axial anomaly is absorbed in the quark spin density. It turns out that the effect of the axial anomaly, which is manifested at $`k_T^2\mu _{\mathrm{fact}}^2`$ (the factorization scale), has the $`x`$ dependence of the $`(1x)`$ form. This explains why $`\mathrm{\Delta }C_G`$ in the $`\overline{\mathrm{MS}}`$ scheme has a term proportional to $`(1x)`$ because the axial-anomaly effect must be subtracted from the previous gluon coefficient function computed in the chiral-invariant scheme. As pointed out in and again in , the $`(1x)`$ term in the gluonic coefficient function in the $`\overline{\mathrm{MS}}`$ scheme is purely “hard”, contrary to what has been claimed previously. in some literature to be a soft term coming from $`k_{}^2m_q^2`$. Second, although the quark spin distribution inside the gluon $`\mathrm{\Delta }\varphi _{q/G}^{(1)}(x)`$ cannot be reliably calculated by perturbative QCD, its difference in GI and CI schemes is trustworthy in QCD. Since the polarized valence quark distributions are $`k_{}`$-factorization scheme independent, the total quark spin distributions in GI and CI schemes are related via Eqs. (3.3) and (3.9) to be
$`\mathrm{\Delta }q(x,Q^2)_{\mathrm{GI}}=\mathrm{\Delta }q(x,Q^2)_{\mathrm{CI}}{\displaystyle \frac{\alpha _s(Q^2)}{\pi }}(1x)\mathrm{\Delta }G(x,Q^2).`$ (3.10)
For a derivation of this important result based on a different approach, namely, the nonlocal light-ray operator technique, see Müller and Teryaev .
The axial anomaly in the box diagram for polarized photon-gluon scattering also occurs at $`k_{}^2\mathrm{}`$, more precisely, at $`k_{}^2=[(1x)/4x]Q^2`$ with $`x0`$. It is natural to expect that the axial anomaly resides in the gluon coefficient function $`\mathrm{\Delta }C_G^{(1)}`$ in the CI scheme, whereas its effect in the GI scheme is shifted to the quark spin density. Since $`\mathrm{\Delta }C_G^{(1)}(x)`$ is the hard part of the polarized photon-gluon cross section, which is sometimes denoted by $`g_1^G(x)`$, the polarized structure function of the gluon target, we have
$`\mathrm{\Delta }C_G^{(1)}(x)=g_1^G(x)\mathrm{\Delta }\varphi _{q/G}^{(1)}(x).`$ (3.11)
It follows from Eqs. (3.9) and (3.10) that
$`\mathrm{\Delta }C_G^{(1)}(x)_{\mathrm{GI}}=\mathrm{\Delta }C_G^{(1)}(x)_{\mathrm{CI}}+{\displaystyle \frac{\alpha _s}{\pi }}(1x).`$ (3.12)
The first moments of $`\mathrm{\Delta }C_G(x)`$, $`_q\mathrm{\Delta }q(x)`$ and $`g_1^p(x)`$ are given by
$`{\displaystyle _0^1}𝑑x\mathrm{\Delta }C_G^{(1)}(x)_{\mathrm{GI}}=0,{\displaystyle _0^1}𝑑x\mathrm{\Delta }C_G^{(1)}(x)_{\mathrm{CI}}={\displaystyle \frac{\alpha _s}{2\pi }},`$
$`\mathrm{\Delta }\mathrm{\Sigma }_{\mathrm{GI}}(Q^2)=\mathrm{\Delta }\mathrm{\Sigma }_{\mathrm{CI}}(Q^2){\displaystyle \frac{n_f\alpha _s(Q^2)}{2\pi }}\mathrm{\Delta }G(Q^2),`$ (3.13)
and
$`\mathrm{\Gamma }_1^p`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle e_q^2\left(\mathrm{\Delta }q_{\mathrm{CI}}(Q^2)\frac{\alpha _s(Q^2)}{2\pi }\mathrm{\Delta }G(Q^2)\right)}`$ (3.14)
$`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle e_q^2\mathrm{\Delta }q_{\mathrm{GI}}(Q^2)},`$
where we have neglected contributions to $`g_1^p`$ from $`\mathrm{\Delta }\varphi _{q/q}^{(1)}`$ and $`\mathrm{\Delta }C_q^{(1)}`$. Note that $`\mathrm{\Delta }\mathrm{\Sigma }_{\mathrm{GI}}(Q^2)`$ is equivalent to the singlet axial charge $`p,s|J_\mu ^5|p,s`$. The well-known results (3.10-3.14) indicate that $`\mathrm{\Gamma }_1^p`$ receives anomalous gluon contributions in the CI factorization scheme (e.g. the improved parton model), whereas hard gluons do not play any role in $`\mathrm{\Gamma }_1^p`$ in the GI scheme such as the OPE approach. From (3.14) it is evident that the sea quark or anomalous gluon interpretation for the suppression of $`\mathrm{\Gamma }_1^p`$ observed experimentally is simply a matter of convention .
The $`\overline{\mathrm{MS}}`$ scheme is the most common one chosen in the GI factorization convention. The so-called Adler-Bardeen (AB) factorization scheme often adopted in the literature is obtained from the GI scheme by adding the $`x`$-independent term $`\alpha _s/(2\pi )`$ to $`\mathrm{\Delta }C_G^{\mathrm{GI}}`$ via
$`\mathrm{\Delta }C_G(x)_{\mathrm{AB}}=\mathrm{\Delta }C_G(x)_{\mathrm{GI}}{\displaystyle \frac{\alpha _s}{2\pi }},`$ (3.15)
while
$`\mathrm{\Delta }q(x,Q^2)_{\mathrm{AB}}=\mathrm{\Delta }q(x,Q^2)_{\mathrm{GI}}+{\displaystyle \frac{\alpha _s(Q^2)}{2\pi }}{\displaystyle _x^1}{\displaystyle \frac{dy}{y}}\mathrm{\Delta }G(y,Q^2).`$ (3.16)
In general, one can define a family of schemes labelled by the parameter $`a`$ :
$`\mathrm{\Delta }q(x)_a=\mathrm{\Delta }q_{\mathrm{GI}}(x)+{\displaystyle \frac{\alpha _s}{2\pi }}\left[(2x1)(a1)+2(1x)\right]\mathrm{\Delta }G(x),`$ (3.17)
which satisfy the relation
$`\mathrm{\Delta }\mathrm{\Sigma }_a=\mathrm{\Delta }\mathrm{\Sigma }_{\mathrm{GI}}+{\displaystyle \frac{3\alpha _s}{2\pi }}\mathrm{\Delta }G,`$ (3.18)
to the first moment, but differ in their higher moments. The AB scheme corresponds to $`a=2`$ and the CI scheme to $`a=1`$. Since the $`x`$ dependence of the axial-anomaly contribution in the quark sector is fixed to be $`(1x)`$, it is obvious that all the schemes in this family including the AB scheme cannot consistently put all hard anomaly effects into gluonic coefficient functions unless $`a=1`$, contrary to the original claim made in .
Finally, it should be stressed that the quark coefficient function $`\mathrm{\Delta }C_q^{(1)}(x)`$ in the dimensional regularization scheme is subject to another ambiguity, namely, the definition of $`\gamma _5`$ in $`n`$ dimension used to specify the ultraviolet cutoff on $`\mathrm{\Delta }\varphi _{q/q}^{(1)}`$ (see Fig. 1). For example, $`\mathrm{\Delta }C_q^{(1)}(x)`$ calculated in the $`\gamma _5`$ prescription of ’t Hooft and Veltman, Breitenlohner and Maison (HVBM) is different from that computed in the dimension reduction scheme. The result
$`\mathrm{\Delta }C_q^{(1)}(x)=C_q^{(1)}(x){\displaystyle \frac{2\alpha _s}{3}}(1+x)`$ (3.19)
usually seen in the literature is obtained in the HVBM scheme, where $`C_q(x)`$ is the unpolarized quark coefficient function. Of course, the quantity $`\mathrm{\Delta }q^{(1)}(x)+\mathrm{\Delta }C_q(x)\mathrm{\Delta }q^{(0)}(x)`$ and hence $`g_1^p(x)`$ is independent of the definition of $`\gamma _5`$ in dimensional regularization.
#### 3.1.4 A brief summary
In addition to all the ambiguities that spin-averaged parton distributions have, the parton spin densities are subject to two extra ambiguities, namely, the axial anomaly and the definition of $`\gamma _5`$ in $`n`$ dimension. There are two extreme cases of interest: Either the hard axial anomaly is manifested in the matrix elements of the quark current (GI scheme) or it is absorbed in the gluonic coefficient function so that the quark matrix element is anomaly-free (CI scheme). It should be stressed that in the so-called AB scheme, not all hard anomaly effects are lumped into $`\mathrm{\Delta }C_G`$ and hence the corresponding quark matrix element is not anomaly-free. Of course, it appears that the CI scheme is close to the intuitive parton picture as the quark spin distribution which does not contain hard contributions from the anomaly is scale independent. (The price to be paid is that $`\mathrm{\Delta }q_{\mathrm{CI}}`$ cannot be expressed as the matrix element of a local gauge-invariant operator and hence it is difficult to compute by lattice QCD. Also, a priori there is no reason that $`\mathrm{\Delta }\mathrm{\Sigma }`$ should have a simple quark interpretation .) Nevertheless, physically and mathematically GI and CI schemes are equivalent.
Two remarks are in order. (i) It is worth emphasizing that although the suppression of $`\mathrm{\Gamma }_1^p`$ can be accommodated by anomalous gluon and/or sea quark contributions, no quantitative prediction of $`\mathrm{\Gamma }_1^p`$ can be made. An attempt of explaining the smallness of $`g_A^0`$ has been made in the large-$`N_c`$ approach . (ii) So far we have focused on the perturbative part of the axial anomaly. The perturbative QCD result (3.10) indicates that the difference $`\mathrm{\Delta }q_s^{\mathrm{GI}}\mathrm{\Delta }q_s^{\mathrm{CI}}`$ ($`\mathrm{\Delta }q=\mathrm{\Delta }q_v+\mathrm{\Delta }q_s`$ with the valence polaization $`\mathrm{\Delta }q_v`$ being scheme independent) is induced perturbatively from hard gluons via the anomaly mechanism and its sign is predicted to be negative. By contrast, $`\mathrm{\Delta }q_s^{\mathrm{CI}}(x)`$ can be regarded as an intrinsic sea-quark spin density produced nonperturbatively. The well-known solution to the $`U_A(1)`$ problem in QCD involves two important ingredients: the QCD anomaly and the QCD vacuum with a nontrivial topological structure, namely the $`\theta `$-vacuum constructed from instantons which are nonperturbative gluon configurations. Since the instanton-induced interactions can flip quark helicity, in analog to the baryon-number nonconservation induced by the ’t Hooft mechanism, the quark-antiquark pair created from the QCD vacuum via instantons can have a net helicity. It has been suggested that this mechanism of quark helicity nonconservation provides a natural and nonperturbative way of generating negative sea-quark polarization .
In retrospect, the dispute among the anomalous gluon and sea-quark explanations of the suppression of $`\mathrm{\Gamma }_1`$ mostly before 1996 is considerably unfortunate and annoying since the fact that $`g_1(x)`$ is independent of the definition of the quark spin density and hence the choice of the factorization scheme due to the axial-anomaly ambiguity is presumably well known to all the practitioners in the field, especially to those QCD experts working in the area.
### 3.2 Lattice calculation of the proton spin content
The present lattice calculation is starting to shed light on the proton spin contents. An evaluation of $`\mathrm{\Delta }q_{\mathrm{GI}}`$, the gauge-invariant quark spin component defined by $`\mathrm{\Delta }q_{\mathrm{GI}}=p,s|\overline{q}\gamma ^\mu \gamma _5q|p,ss_\mu `$, involves a disconnected insertion in addition to the connected insertion (see Fig. 2; the infinitely many possible gluon lines and additional quark loops are implicit). The sea-quark spin contribution comes from the disconnected insertion.
There are four lattice calculations done by Dong et al. , Fukugita et al. , Göckeler et al. in the quenched approximation and Güsken et al. in full lattice QCD. Note that the disconnected contribution is not evaluated in . From Table III it is clear that the lattice results for $`g_A^0`$ and $`\mathrm{\Delta }s`$ are in agreement with experiment, while the full lattice QCD calculations for $`g_A^3,\mathrm{\Delta }u`$ and $`\mathrm{\Delta }d`$ are too small compared to experiment. In particular, there is a 30% discrepancy between the lattice QCD estimate of $`g_A^3`$ and the experimental value. This points to the presence of sizeable higher order or even nonperturbative contributions to the renormalized factor $`Z_A^{\mathrm{NS}}`$ on the non-singlet current.
As for the chiral-invariant quantity $`\mathrm{\Delta }\mathrm{\Sigma }_{\mathrm{CI}}`$, it involves the matrix element of $`\stackrel{~}{J}_5^+`$ in light-front gauge where $`\stackrel{~}{J}_{\mu 5}`$ is an anomaly-free singlet axial vector current and hence sizeable gauge configurations are needed in lattice calculations for $`\mathrm{\Delta }\mathrm{\Sigma }_{\mathrm{CI}}`$. Nevertheless, it is conceivable to have lattice results for $`\mathrm{\Delta }G`$ and $`\mathrm{\Delta }q_{\mathrm{CI}}`$ soon in the near future. The lattice calculation of the quark total angular momentum was also available recently, see Sec. 3.4.
### 3.3 NLO analysis of polarization data
The experimental data of $`g_1(x,Q^2)`$ taken at different $`x`$-bin correspond to different ranges of $`Q^2`$; that is, $`Q^2`$ of the data is $`x`$-bin dependent. Hence, it is desirable to evolve the data to a common value of $`Q^2`$ in order to determine the moments of $`g_1`$ and test the Bjorken sum rule. Because of the availability of the two-loop polarized splitting functions $`\mathrm{\Delta }P_{ij}^{(1)}(x)`$ , it became possible to embark on a full next-to-leading order (NLO) analysis of the experimental data of polarized structure functions by taking into account the measured $`x`$ dependence of $`Q^2`$ at each $`x`$ bin. The NLO analyses have been performed in the $`\overline{\mathrm{MS}}`$ scheme (a family of the GI scheme) , the Adler-Bardeen (AB) scheme and the CI scheme . Of course, physical quantities such as the polarized structure function $`g_1(x)`$ are independent of choice of the factorization convention. Physically, the spin-dependent valence quark and gluon distributions should be the same in all factorization schemes.
The $`Q^2`$ dependence of parton spin densities is determined by the spin-dependent DGLAP equations:
$`{\displaystyle \frac{d}{dt}}\mathrm{\Delta }q_{\mathrm{NS}}(x,t)={\displaystyle \frac{\alpha _s(t)}{2\pi }}\mathrm{\Delta }P_{qq}^{\mathrm{NS}}(x)\mathrm{\Delta }q_{\mathrm{NS}}(x,t),`$
$`{\displaystyle \frac{d}{dt}}\left(\begin{array}{c}\mathrm{\Delta }q_\mathrm{S}(x,t)\\ \mathrm{\Delta }G(x,t)\end{array}\right)={\displaystyle \frac{\alpha _s(t)}{2\pi }}\left(\begin{array}{cc}\mathrm{\Delta }P_{qq}^\mathrm{S}(x)& 2n_f\mathrm{\Delta }P_{qG}(x)\\ \mathrm{\Delta }P_{Gq}(x)& \mathrm{\Delta }P_{GG}(x)\end{array}\right)\left(\begin{array}{c}\mathrm{\Delta }q_\mathrm{S}(x,t)\\ \mathrm{\Delta }G(x,t)\end{array}\right),`$ (3.20)
with $`t=\mathrm{ln}(Q^2/\mathrm{\Lambda }_{_{\mathrm{QCD}}}^2)`$,
$`\mathrm{\Delta }q_{\mathrm{NS}}(x)=\mathrm{\Delta }q_i(x)\mathrm{\Delta }q_j(x),\mathrm{\Delta }q_\mathrm{S}(x)={\displaystyle \underset{i}{}}\mathrm{\Delta }q_i(x),`$ (3.21)
and
$`\mathrm{\Delta }P_{ij}(x)=\mathrm{\Delta }P_{ij}^{(0)}(x)+{\displaystyle \frac{\alpha _s}{2\pi }}\mathrm{\Delta }P_{ij}^{(1)}(x)+\mathrm{}.`$ (3.22)
The spin-dependent anomalous dimensions are defined as
$`\mathrm{\Delta }\gamma _{ij}^n={\displaystyle _0^1}\mathrm{\Delta }P_{ij}(x)x^{n1}𝑑x=\mathrm{\Delta }\gamma _{ij}^{(0),n}+{\displaystyle \frac{\alpha _s}{2\pi }}\mathrm{\Delta }\gamma _{ij}^{(1),n}+\mathrm{}.`$ (3.23)
To the NLO, $`\mathrm{\Delta }P_{qq}^{(1)}`$ and $`\mathrm{\Delta }P_{qG}^{(1)}`$ were calculated in the $`\overline{\mathrm{MS}}`$ scheme by Zijlstra and van Neerven . However, the other two polarized splitting functions $`\mathrm{\Delta }P_{Gq}^{(1)}`$ and $`\mathrm{\Delta }P_{GG}^{(1)}`$ were not available until 1995 . The recent analysis shows that the NLO $`\chi ^2`$ is significantly smaller than that of LO, indicating the necessity of NLO fit of data in practice.
Since the unpolarized parton densities are mostly parameterized in the $`\overline{\mathrm{MS}}`$ scheme and the two-loop splitting functions are available in the same scheme, it is quite natural to perform the NLO analysis in the $`\overline{\mathrm{MS}}`$ scheme. In principle one can also work in any other factorization scheme. The splitting functions $`\mathrm{\Delta }P_{ij}^{(1)}`$ in the CI scheme, for example, is obtained by applying Eq. (3.10) to the spin-dependent DGLAP equation, see . It is worth accentuating again that though it is perfectly all right to analyze the data in the AB scheme, one has to keep in mind that not all hard effects are absorbed in the gluonic coefficient functions in such a scheme.
The sea-quark and gluon spin distributions cannot be separately determined from current experimental data. In other words, while the shapes of the spin-dependent valence quark distributions are fairly constrained by the data, sea-quark and gluon spin densities are almost completely undetermined. In particular, $`\mathrm{\Delta }G`$ is rather weakly constrained by the data. This is understandable because the gluon polarization contributes to $`g_1`$ only at the NLO level through the convolution $`\mathrm{\Delta }C_G\mathrm{\Delta }G`$. In principle, measurements of scaling violation in $`g_1(x,Q^2)`$ via, for example, the derivative of $`g_1(x,Q^2)`$ with respect to $`Q^2`$, in next-generation experiments will allow an estimate of the gluon spin density and the overall size of gluon polarization. Of course, the data should be sufficiently accurate in order to study the gluon spin density. Meanwhile, it is even more important to probe $`\mathrm{\Delta }G(x)`$ independently in those hadron-hadron collision processes where gluons play a dominant role.
Though the polarized structure function is factorization scheme independent, it is important to perform NLO analyses in different schemes to test the reliability and consistency of the theory. It is found in that the polarized parton distributions obtained in $`\overline{\mathrm{MS}}`$, AB and CI schemes agree well for the non-singlet spin densities, while the first moment of $`\mathrm{\Delta }G(x)`$ obtained in the AB or CI scheme is different from that in the $`\overline{\mathrm{MS}}`$ scheme, reflecting that the present data can hardly constrain the gluon spin distribution. Typically, the extracted value of the gluon spin at $`Q^2=1\mathrm{GeV}^2`$ lies in the range 0<
ΔG<
20<
Δ𝐺<
20{\ \lower-1.2pt\vbox{\hbox{\hbox to0.0pt{$<$\hss}\lower 5.0pt\vbox{\hbox{$\sim$}
}}}\ }\Delta G{\ \lower-1.2pt\vbox{\hbox{\hbox to0.0pt{$<$\hss}\lower 5.0pt\vbox{\hbox{$\sim$}
}}}\ }2 . However, the recent analysis shows that the LO fit cannot decide on the sign of $`\mathrm{\Delta }G`$, while the NLO analysis yields a negative first moment of the gluon density.<sup>4</sup><sup>4</sup>4It has been found by Jaffe that the gluon spin component is negative $`\mathrm{\Delta }G0.4`$ in the MIT bag model and even more negative in the non-relativistic quark model. However, it was explained in that the negative $`\mathrm{\Delta }G`$ obtained by Jaffe is a consequence of neglecting some self-interaction effects. This illustrates again that it is difficult to pin down the gluon spin distribution from present polarized DIS data. Note that a recent NLO analysis in shows that the polarized strange quark density is significantly different from zero independently of the factorization schemes used in the analysis:
$`\mathrm{\Delta }s_{\overline{\mathrm{MS}}}=0.102\pm 0.012,\mathrm{\Delta }s_{\mathrm{CI}}=0.064\pm 0.010,\mathrm{\Delta }s_{\mathrm{AB}}=0.058\pm 0.012,`$ (3.24)
at $`Q^2=1\mathrm{GeV}^2`$. Note that the sea polarization $`\mathrm{\Delta }s`$ is scheme dependent:
$`\mathrm{\Delta }s_{\mathrm{CI},\mathrm{AB}}=\mathrm{\Delta }s_{\mathrm{GI}}+{\displaystyle \frac{\alpha _s}{2\pi }}\mathrm{\Delta }G.`$ (3.25)
The presence of the sea polarization in the CI or AB scheme implies that the gluon polarization is not at its maximal value given by $`\mathrm{\Delta }G(2\pi /\alpha _s)\mathrm{\Delta }s_{\mathrm{GI}}`$.
### 3.4 Orbital angular momentum
The orbital angular momentum plays a role in the proton spin structure. For example, the growth of large $`\mathrm{\Delta }G`$ with $`Q^2`$ is balanced by the large negative orbital angular momentum of the quark-gluon pair. It is also known that the reduction of the total spin component $`\mathrm{\Delta }\mathrm{\Sigma }`$ due to the presence of the quark transverse momentum in the lower component of the Dirac spinor is traded with the quark orbital angular momentum.
In the proton spin sum rule:
$`{\displaystyle \frac{1}{2}}=J_q+J_G,`$ (3.26)
the total angular momenta $`J_q`$ and $`J_G`$ of quarks and gluons respectively are gauge invariant. However, the decomposition into spin and orbital components, $`J_q=\frac{1}{2}\mathrm{\Delta }\mathrm{\Sigma }+L_q`$ and $`J_G=\mathrm{\Delta }G+L_G`$, is gauge dependent. This leads to difficulties in defining the partonic spin densities: There exist no local gauge invariant operators that could represent the densities of gluon spin and the orbital angular momenta of quarks and gluons. It is known that the spin and orbital angular momenta of quarks and gluons appearing in the decomposition
$`J^z`$ $`=`$ $`J_q^z+J_G^z=S_q^z+L_q^z+S_G^z+L_G^z`$ (3.27)
$`=`$ $`{\displaystyle d^3x\left[\frac{1}{2}\overline{\psi }\gamma ^3\gamma ^5\psi +\psi ^{}(\stackrel{}{x}\times i\stackrel{}{})^3\psi +(\stackrel{}{E}\times \stackrel{}{A})^3E^k(\stackrel{}{x}\times \stackrel{}{})^3A^k\right]},`$
obtained in the the light-cone gauge $`A^+=0`$ and in the infinite momentum frame, are separately gauge variant except for the quark spin operator $`S_q^z`$. However, the gluon spin and its distribution, for example, are physical, gauge invariant quantities and can be measured experimentally. The point is that $`\mathrm{\Delta }G`$ can be expressed as the matrix element of a string-like nonlocal gauge-invariant operator $`O_G`$ . As stressed in , what one measures experimentally is the matrix element of the gauge-invariant operator $`O_G`$. But in the light-cone gauge, the above operator reduces to the gluon spin operator $`S_G^z=d^3x(\stackrel{}{E}\times \stackrel{}{A})^3`$; that is, the gauge-invariant extension of the gluon spin operator in the light-cone gauge is measurable. Likewise, a gauge-invariant operator that reduces to the quark orbital angular momentum in the light-cone gauge has been discussed in (see however a different discussion in ).
In principle, the total angular momenta of quarks and gluons, $`J_q`$ and $`J_G`$ respectively, can be measured in deeply virtual Compton scattering in a special kinematic region where single quark scattering dominates . It has also been suggested that the orbital angular momentum might be deduced from an azimuthal asymmetry in hadron production with a transversely polarized target.
The evolution of the quark and gluon orbital angular momenta was first discussed by Ratcliffe . Ji, Tang and Hoodbhoy have derived a complete leading-log evolution equation in the light-cone gauge:
$`{\displaystyle \frac{d}{dt}}\left(\begin{array}{c}L_q\\ L_G\end{array}\right)={\displaystyle \frac{\alpha _s(t)}{2\pi }}\left(\begin{array}{cc}\frac{4}{3}C_F& \frac{n_f}{3}\\ \frac{4}{3}C_F& \frac{n_f}{3}\end{array}\right)\left(\begin{array}{c}L_q\\ L_G\end{array}\right)+{\displaystyle \frac{\alpha _s(t)}{2\pi }}\left(\begin{array}{cc}\frac{2}{3}C_F& \frac{n_f}{3}\\ \frac{5}{6}C_F& \frac{11}{2}\end{array}\right)\left(\begin{array}{c}\mathrm{\Delta }\mathrm{\Sigma }\\ \mathrm{\Delta }G\end{array}\right),`$ (3.28)
with the solutions
$`L_q(Q^2)={\displaystyle \frac{1}{2}}\mathrm{\Delta }\mathrm{\Sigma }+{\displaystyle \frac{1}{2}}{\displaystyle \frac{3n_f}{16+3n_f}}+f(Q^2)\left(L_q(Q_0^2)+{\displaystyle \frac{1}{2}}\mathrm{\Delta }\mathrm{\Sigma }{\displaystyle \frac{1}{2}}{\displaystyle \frac{3n_f}{16+3n_f}}\right),`$
$`L_G(Q^2)=\mathrm{\Delta }G(Q^2)+{\displaystyle \frac{1}{2}}{\displaystyle \frac{16}{16+3n_f}}+f(Q^2)\left(L_G(Q_0^2)+\mathrm{\Delta }G(Q_0^2){\displaystyle \frac{1}{2}}{\displaystyle \frac{16}{16+3n_f}}\right),`$
where
$`f(Q^2)=\left({\displaystyle \frac{\mathrm{ln}Q_0^2/\mathrm{\Lambda }_{_{\mathrm{QCD}}}^2}{\mathrm{ln}Q^2/\mathrm{\Lambda }_{_{\mathrm{QCD}}}^2}}\right)^{\frac{32+6n_f}{332n_f}}`$ (3.30)
and $`\mathrm{\Delta }\mathrm{\Sigma }`$ is $`Q^2`$ independent to the leading-log approximation. We see that the growth of $`\mathrm{\Delta }G`$ with $`Q^2`$ is compensated by the gluon orbital angular momentum, which also increases like $`\mathrm{ln}Q^2`$ but with opposite sign. The solution (3.4) has an interesting implication in the asymptotic limit $`Q^2\mathrm{}`$, namely
$`J_q(Q^2)={\displaystyle \frac{1}{2}}\mathrm{\Delta }\mathrm{\Sigma }+L_q(Q^2){\displaystyle \frac{1}{2}}{\displaystyle \frac{3n_f}{16+3n_f}},`$
$`J_G(Q^2)=\mathrm{\Delta }G(Q^2)+L_G(Q^2){\displaystyle \frac{1}{2}}{\displaystyle \frac{16}{16+3n_f}}.`$ (3.31)
Thus, history repeats herself: The partition of the nucleon spin between quarks and gluons follows the well-known partition of the nucleon momentum. Taking $`n_f=6`$, we see that $`J_q:J_G=0.53:\mathrm{\hspace{0.17em}0.47}`$. If the evolution of $`J_q`$ and $`J_G`$ is very slow, which is empirically known to be true for the momentum sum rule that half of the proton’s momentum is carried by gluons even at a moderate $`Q^2`$, then $`\mathrm{\Delta }\mathrm{\Sigma }0.25`$ at $`Q^2=10\mathrm{GeV}^2`$ implies that $`L_q0.13`$ at the same $`Q^2`$, recalling that the quark orbital angular momentum is expected to be of order 0.20 in the relativistic quark model.
Recently the quark orbital angular momentum of the nucleon was calculated from lattice QCD by considering the form factor of the quark energy-momentum tensor $`T_{\mu \nu }`$ . The total angular momentum of the quarks is found to be
$`J_q=0.30\pm 0.07.`$ (3.32)
That is about 60% of the proton spin is attributable to the quarks. Hence, the quark orbital angular momentum is $`L_q=0.17\pm 0.06`$ when combining with the previous quark spin content $`\frac{1}{2}\mathrm{\Delta }\mathrm{\Sigma }=0.13\pm 0.06`$ . Therefore, about 25% of the proton spin originates from the quark spin and about 35% comes from the quark orbital angular momentum.
It must be stressed that $`J_q`$ should be factorization scheme independent. This means that a replacement of $`\mathrm{\Delta }\mathrm{\Sigma }_{\mathrm{GI}}`$ by $`\mathrm{\Delta }\mathrm{\Sigma }_{\mathrm{CI}}`$ in the spin sum rule (3.26) requires that the difference $`\mathrm{\Delta }\mathrm{\Sigma }_{\mathrm{CI}}\mathrm{\Delta }\mathrm{\Sigma }_{\mathrm{GI}}=(n_f\alpha _s/2\pi )\mathrm{\Delta }G`$ be compensated by a counterpart in the gluon orbital angular momentum:
$`L_q^{\mathrm{CI}}=L_q^{\mathrm{GI}}{\displaystyle \frac{n_f\alpha _s}{4\pi }}\mathrm{\Delta }G.`$ (3.33)
It is interesting to note that if $`\mathrm{\Delta }G`$ is of order 2.5 , one will have $`\mathrm{\Delta }\mathrm{\Sigma }_{\mathrm{CI}}0.58`$ and $`L_q^{\mathrm{CI}}0`$ . In other words, while $`\mathrm{\Delta }\mathrm{\Sigma }_{\mathrm{CI}}`$ is close to the relativistic quark model value of $`\mathrm{\Delta }\mathrm{\Sigma }`$, $`L_q^{\mathrm{CI}}`$ deviates farther from the quark model expectation $`L_q^{}0.20`$ (see Sec. I).
## 4 Conclusions
The spin sum rule of the proton in the infinite momentum frame reads
$`{\displaystyle \frac{1}{2}}=J_q+J_G={\displaystyle \frac{1}{2}}\mathrm{\Delta }\mathrm{\Sigma }+L_q+\mathrm{\Delta }G+L_G.`$ (4.1)
The quark spin can be inferred from polarized DIS measurements of $`g_1(x)`$ and its first moment. Due to the ambiguity arising from the axial anomaly, the definition of the sea polarization $`\mathrm{\Delta }q_s`$ and hence $`\mathrm{\Delta }q=\mathrm{\Delta }q_v+\mathrm{\Delta }q_s`$ is $`k_{}`$ factorization dependent, but $`J_q,g_1(x)`$ and $`\mathrm{\Gamma }_1`$ are not. The only spin content which is for sure at present is the observed value $`\mathrm{\Delta }\mathrm{\Sigma }0.200.30`$ in the GI scheme (e.g. the $`\overline{\mathrm{MS}}`$ scheme). The recent lattice calculation yields $`J_q=0.30\pm 0.07`$. Therefore, we have $`L_q^{\mathrm{GI}}0.10\pm 0.06`$ for $`\mathrm{\Delta }\mathrm{\Sigma }_{\mathrm{GI}}0.25`$, to be compared with the quark model prediction $`L_q^{}0.20`$ . The values of $`\mathrm{\Delta }\mathrm{\Sigma }`$ and $`L_q`$ in the CI scheme (e.g. the improved QCD parton model) or the AB scheme depend on the gluon spin. Since $`L_q^{\mathrm{CI}}L_q^{\mathrm{GI}}=n_f\alpha _s\mathrm{\Delta }G/(4\pi )`$, it is clear that if the gluon polarization is positive and large enough, then $`L_q^{\mathrm{CI}}`$ will deviate even farther from the quark picture although $`\mathrm{\Delta }\mathrm{\Sigma }_{\mathrm{CI}}`$ can be made to be close to the constituent relativistic quark model. In the asymptotic limit, $`J_q(\mathrm{})=\frac{1}{2}\mathrm{\Delta }\mathrm{\Sigma }(\mathrm{})+L_q(\mathrm{})\frac{1}{4}`$ and $`J_G(\mathrm{})=\mathrm{\Delta }G(\mathrm{})+L_G(\mathrm{})\frac{1}{4}`$. The recent lattice result $`J_q=0.30\pm 0.07`$ at a moderate $`Q^2`$ seems to suggest that the evolution of $`J_q`$ and hence $`J_G`$ is slow enough.
ACKNOWLEDGMENTS
This work was supported in part by the National Science Council of ROC under Contract No. NSC89-2112-M001-016. |
warning/0002/gr-qc0002023.html | ar5iv | text | # Introduction
## Introduction
Describing evolution of the degrees of freedom of generally covariant theories is an unsolved puzzle, and constitutes one of the challenges of the human thinking of our time. The study of generally covariant theories has been motivated by general relativity, which has this peculiar property (see for instance ). In gravity, general covariance means the theory is diffeomorphism invariant, and this symmetry of gravity implies at the Hamiltonian level that the theory has not a genuine Hamiltonian for describing the evolution of the degrees of freedom of the gravitational field, rather, dynamics is gauge; generated by the first class constraints of the theory. This is the so-called problem of time in classical general relativity .
On the other hand, if the classical regime of general relativity is only a limit, which emerges in a suitable way from its fundamental quantum behavior , then the theory is in trouble. Our standard methods of quantization crash and do not apply to the particular physical situation raised by general relativity. Standard quantization methods in field theory are background-dependent, quantum gravity needs a background independent procedure in its quantization. So, how to make compatible the symmetry of general relativity with the quantum theory? what is the meaning of evolution in quantum gravity? Loop quantum gravity answers the first question, because the quantization of the gravitational field is carried out in a background independent way . With respect to the second, it remains as an open question. Among the several proposals for describing the evolution of the degrees of freedom of generally covariant theories, I find Rovelli’s proposal as one of the most creative ones .
Here, following the spirit of relationism, which is the heart in Rovelli’s point of view, we analyse the ‘problem of time’ in generally covariant theories with vanishing Hamiltonian and with a finite number $`D`$ of degrees of freedom. To obtain the relations involving the coordinates and momenta of the unreduced phase space $`\gamma _{ex}`$ with the physical states that label the points of the physical phase space $`\gamma _{ph}`$, we need to start from the embedding equations which give the dependence of the coordinates and momenta with respect to the $`M`$ time variables $`t^m`$ as well as the $`2D`$ physical states $`(\stackrel{~}{q}^a,\stackrel{~}{p}_a)`$. These equations constitute the classical version of the generalized Heisenberg picture, which arises when these equations are promoted to quantum operators in the reduced Hilbert space of the theory. By plugging the expressions of the time variables $`t^m`$ in terms the original canonical variables into the expressions of coordinates and momenta, we get the full relational evolution of the phase space degrees of freedom for any physical state of $`\gamma _{ph}`$. This way of expressing the full solution of the dynamics of generally covariant theories constitutes the full set of evolving constants of motion required in their dynamics, and is displayed in Sect. I. In addition, an alternative mechanism which generates also the evolving constants of motion is proposed. Of course, we study also the quantum version of the evolving constants of motion. In sect. II, we analyze the parameterized harmonic oscillator (as an example of parameterized systems). In Sect. III, we continue the study of the $`SL(2,R)`$ model which constraint algebra mimics the algebra structure of general relativity. Due to the fact a Schrödinger equation emerges in its quantum dynamics, we compare the generalized Heisenberg picture (related with the evolving constants of motion) which needs $`M`$ time variables with the Schrödinger picture which singles out one time variable only. We also especulate on the classical limit generally covariant theories and its possible relation with the full set of evolving constants of motion. Our conclusions are summarized in Sect. IV.
## I Relational classical and quantum dynamics
Classical dynamics. The classical dynamics of a constrained theory with a finite number of degrees of freedom characterized by first class constraints is as follows . The theory is obtained from the Hamiltonian form of the action
$`S[q^i,p_i,\lambda ^m]`$ $`=`$ $`{\displaystyle _{\tau _1}^{\tau _2}}𝑑\tau \left\{{\displaystyle \frac{dq^i}{d\tau }}p_i\lambda ^mC_m(q^i,p_i)\right\},`$ (1)
which is invariant under arbitrary reparameterizations of the parameter $`\tau `$. Hence, $`\tau `$ is a non physical coordinate time. The unreduced phase space $`\gamma _{ex}`$ is coordinated by the canonical pairs $`(q^i,p_i)`$ ;$`i=1,2,\mathrm{},N`$. The canonical 2-form on $`\gamma _{ex}`$ is $`\omega _{ex}=dp_idq^i`$. Thus, $`(\gamma ,\omega )`$ is a symplectic space. The variation of the action $`S[q^i,p_i,\lambda ^m]`$ with respect to the canonical coordinates $`q^i,p_i`$ gives the equations of motion
$`{\displaystyle \frac{dq^i}{d\tau }}`$ $`=`$ $`\lambda ^m{\displaystyle \frac{C_m(q^i,p_i)}{p_i}},`$ (2)
$`{\displaystyle \frac{dp_i}{d\tau }}`$ $`=`$ $`\lambda ^m{\displaystyle \frac{C_m(q^i,p_i)}{q^i}},`$ (3)
while the variation of the action with respect to the Lagrange multipliers $`\lambda ^m`$ gives the constraint equations
$`C_m`$ $`=`$ $`C_m(q^i,p_i)=0,m=1,2\mathrm{},M.`$ (4)
The variation of the action has been done under the standard boundary conditions $`q^i(\tau _s)=q_s^i`$; $`s=1,2`$, namely, the allowed paths are those with fixed values for the configuration variables at the boundary points $`\tau _s`$. The boundary conditions can be changed and thus to modify the action by suitable boundary terms to allow the gauge symmetry generated by the constraints . The constraints generate Hamiltonian vector fields $`X_{dC_m}`$, which are tangent vectors to the constraint surface, given by
$`X_{dC_m}`$ $`=`$ $`{\displaystyle \frac{C_m}{p_i}}{\displaystyle \frac{}{q^i}}+{\displaystyle \frac{C_m}{q^i}}{\displaystyle \frac{}{p_i}}.`$ (5)
More important, the integral curves of these Hamiltonian vectors fields constitute the gauge submanifold or the orbits of the constraint surface, and the dynamics of the system with respect to $`\tau `$ is the unfolding of this gauge symmetry, i.e., dynamics is gauge.
The first class constraints satisfy, in general, a non Lie algebra
$`\{C_m,C_n\}`$ $`=`$ $`C_{mn}^l(q^i,p_i)C_l,`$ (6)
and the number of independent physical degrees of freedom of the theory is $`D=NM`$. The constraint surface defined by the constraint equations (4) is a $`(2D+M)`$-dimensional manifold. The constraint surface can be parameterized by the set of independent coordinates $`(\stackrel{~}{q}^a,\stackrel{~}{p}_a,t^m)`$, where $`(\stackrel{~}{q}^a,\stackrel{~}{p}_a),a=1,2,\mathrm{},D`$ are (local) canonical variables which coordinatize the open sets of the physical phase space $`\gamma _{ph}`$ of the theory, and $`t^m,m=1,2,\mathrm{},M`$ coordinatize the orbits, i.e., the gauge submanifold of the constraint surface generated by the first class constraints. Notice that the canonical coordinates $`q_0^\alpha `$ , and $`p_{0\alpha }`$ are the physical observables of the system. Of course, they satisfy $`\{\stackrel{~}{q}^a,\stackrel{~}{p}_b\}=\delta _b^a`$ on the physical phase space, and the symplectic form on $`\gamma _{ph}`$ in these coordinates is $`\omega _{ph}=d\stackrel{~}{p}_ad\stackrel{~}{q}^a`$. Therefore, the general solution of the dynamics of the constrained theory can be expressed as
$`q^i`$ $`=`$ $`q^i(t^m;\stackrel{~}{q}^a,\stackrel{~}{p}_a),`$ (7)
$`p_i`$ $`=`$ $`p_i(t^m;\stackrel{~}{q}^a,\stackrel{~}{p}_a).`$ (8)
It is important to emphasize that that such dependence is local. For instance, in the case in which the physical phase space $`\gamma _{ph}`$ is compact, a finite set of physical observables $`(\stackrel{~}{q}^a,\stackrel{~}{p}_a)`$ is needed to coordinate the open sets of $`\gamma _{ph}`$ due to its compactness. So, the constraint surface looks like a ‘fibre bundle’ $`P(\gamma _{ph},\text{Orbits})`$, the constraint surface being the total space $`P`$, the physical phase space $`\gamma _{ph}`$ being the base space, and the orbits being the fibers of the bundle. In the generic case, $`P(\gamma _{ph},\text{Orbits})`$ is locally trivial. This means that non global gauge condition is allowed in general, and that local gauge conditions associated with each open set of the physical phase space can be specified only.
At the same time, the full solution of the theory implies to give the dependence of the physical observables $`\stackrel{~}{q}^a`$, and $`\stackrel{~}{p}_a`$ of the system in terms of the coordinates of the unreduced phase space
$`\stackrel{~}{q}^a`$ $`=`$ $`\stackrel{~}{q}^a(q^i,p_i),`$ (9)
$`\stackrel{~}{p}_a`$ $`=`$ $`\stackrel{~}{p}_a(q^i,p_i),`$ (10)
as well as the orbit coordinates $`t^m`$
$`t^m`$ $`=`$ $`t^m(q^i,p_i).`$ (11)
What these equations tell us is that one single internal time variable is not enough to describe the evolution of the system, rather, $`M`$ internal time variables are needed. In the way the full solution has been expressed in (7), and (8), these $`M`$ time variables are $`t^m`$, $`m=1,2,\mathrm{},M`$. Notice also that these $`M`$ time variables are internal clocks, given by (11). One of the properties of these internal clocks $`t^m`$ is that they do not run taking increasing values of $`t^m`$ when time goes on. In fact, they can run ‘forward’ and ‘backward’ depending on the values of the coordinates and momenta the system is reaching through Eq. (11). Other property is that these clocks can run with different ‘speeds’ for the same reason. So, the meaning of time that arise in generally covariant theories is completely different with respect to the monotonous function we are familiarized with. In was showed that Eqs. (7) and (8) can be obtained by a combination of a canonical transformation plus Hamilton-Jacobi techniques. That approach implies the modification of the original set of first class constraints. This can be done, but it is not necessary in principle. Moreover, the full solution requires (11) (missing in Ref. ) as we have seen, and more important, it is the combination of Eqs. (7), and (8) with Eq. (11) which leads to the relational description of the dynamics of the system, as we will see it later on.
It is worth to mention the relationship between the time variables $`t^m`$ and the full gauge transformation generated by the first class constraints $`C_m`$. Assuming that the full gauge transformation of the original canonical variables is given by
$`q_{}^{}{}_{}{}^{i}`$ $`=`$ $`q_{}^{}{}_{}{}^{i}(q^i,p_i,\alpha ^m(\tau )),`$ (12)
$`p_{}^{}{}_{i}{}^{}`$ $`=`$ $`p_{}^{}{}_{i}{}^{}(q^i,p_i,\alpha ^m(\tau )),`$ (13)
with $`\alpha ^m`$ the $`m`$ gauge parameters involved in the gauge transformation. Then, by plugging (7) and (8) into the right hand side of (13), we get
$`q_{}^{}{}_{}{}^{i}`$ $`=`$ $`q^i(t_{}^{}{}_{}{}^{m};\stackrel{~}{q}^a,\stackrel{~}{p}_a),`$ (14)
$`p_{}^{}{}_{i}{}^{}`$ $`=`$ $`p_i(t_{}^{}{}_{}{}^{m};\stackrel{~}{q}^a,\stackrel{~}{p}_a),`$ (15)
where the functional dependence in the right-hand side of the above equations is exactly the same as that given by (7)-(8) but with
$`t_{}^{}{}_{}{}^{m}`$ $`=`$ $`t_{}^{}{}_{}{}^{m}(t^m,\alpha ^m),`$ (16)
which relates the time variables $`t^m`$ after and before of any finite gauge transformation of the canonical variables (13).
The map $`\varphi :P(\gamma _{ph},\text{Orbits})\gamma _{ex},\varphi (\stackrel{~}{q}^a,\stackrel{~}{p}_a,t^m)(q^i(t^m;\stackrel{~}{q}^a,\stackrel{~}{p}_a),p_i(t^m;\stackrel{~}{q}^a,\stackrel{~}{p}_a))`$ allows us to define on the constraint surface $`P(\gamma _{ph},\text{Orbits})`$ the pull back $`\mathrm{\Omega }=\varphi ^{}\omega _{ex}`$ of the canonical form $`\omega _{ex}=dp_idq^i`$ on $`\gamma _{ex}`$, which is degenerate. Thus, the geometry of constrained systems involves three spaces: the unconstrained phase space $`(\gamma _{ex},\omega _{ex})`$, the constraint surface $`(P(\gamma _{ph},\text{Orbits}),\mathrm{\Omega }=\varphi ^{}\omega _{ex})`$, and the physical phase space $`(\gamma _{ph},\omega _{ph})`$. The map that connects the constraint surface and the physical phase space is the projection $`\pi :P(\gamma _{ph},\text{Orbits})\gamma _{ph}`$, $`\pi (\stackrel{~}{q}^a,\stackrel{~}{p}_a,t^m)(\stackrel{~}{q}^a,\stackrel{~}{p}_a)`$. Due to the fact that the Hamiltonian vector fields (5) are tangent vectors to the orbits, they can be expressed in terms of the local coordinates $`\stackrel{~}{q}^a,\stackrel{~}{p}_a,t^m`$ of the constraint surface as
$`X_{dC_m}`$ $`=`$ $`\{C_m,t^n\}(\stackrel{~}{q}^a,\stackrel{~}{p}_a,t^m){\displaystyle \frac{}{t^n}}.`$ (17)
The observables (9), (10), and the orbits (11) also generate Hamiltonian vectors fields, their restriction on the constraint surface are
$`X_{d\stackrel{~}{q}^a}`$ $`=`$ $`{\displaystyle \frac{}{\stackrel{~}{p}_a}}+\{\stackrel{~}{q}^a,t^n\}{\displaystyle \frac{}{t^n}},`$ (18)
$`X_{d\stackrel{~}{p}_a}`$ $`=`$ $`{\displaystyle \frac{}{\stackrel{~}{q}^a}}+\{\stackrel{~}{p}_a,t^n\}{\displaystyle \frac{}{t^n}},`$ (19)
$`X_{dt^m}`$ $`=`$ $`\{t^m,\stackrel{~}{q}^a\}{\displaystyle \frac{}{\stackrel{~}{q}^a}}+\{t^m,\stackrel{~}{p}_a\}{\displaystyle \frac{}{\stackrel{~}{p}_a}}+\{t^m,t^n\}{\displaystyle \frac{}{t^n}},`$ (20)
where it is understood that all the quantities are evaluated in the point $`(\stackrel{~}{q}^a,\stackrel{~}{p}_a,t^m)`$. Thus, $`\{X_{dC_m},X_{d\stackrel{~}{q}^a},X_{d\stackrel{~}{p}_a}\}`$ is a basis, naturally adapted to the involved geometry, of the tangent space of the constraint surface. The vectors $`X_{dC_m}`$ play the role of vertical vectors because they have vanishing projection on the tangent space of $`\gamma _{ph}`$, $`d\pi \left(X_{dC_m}\right)=0`$. $`X_{d\stackrel{~}{q}^a}`$ , and $`X_{d\stackrel{~}{p}_a}`$ are the horizontal lifts on the constraint surface of the coordinate basis on $`\gamma _{ph}`$, $`d\pi \left(X_{d\stackrel{~}{q}^a}\right)=\frac{}{\stackrel{~}{p}_a}`$ , $`d\pi \left(X_{d\stackrel{~}{p}_a}\right)=\frac{}{\stackrel{~}{q}^a}`$. In summary, the solution of the dynamics of the constrained system means to specify Eqs. (7)-(11). This fact, rises a new problem: the problem of the meaning of physical time of generally covariant theories, i.e. the specification of an internal clock in the framework of the theory with respect to which to describe the evolution of the degrees of freedom of the theory in a gauge invariant way. Let us explain, the dynamics with respect to $`\tau `$ is given by
$`q^i(\tau )`$ $`=`$ $`q^i(t^m(\tau );\stackrel{~}{q}^a,\stackrel{~}{p}_a),`$ (21)
$`p_i(\tau )`$ $`=`$ $`p_i(t^m(\tau );\stackrel{~}{q}^a,\stackrel{~}{p}_a),`$ (22)
for any physical state $`(\stackrel{~}{q}^a,\stackrel{~}{p}_a)`$ of the system. So, this dynamics is non gauge-invariant, i.e., it depends on $`\tau `$. The question is, can we describe evolution of the system in a gauge invariant way? The answer is yes. At first sight, this sounds strange or impossible in a system with gauge freedom. To see how this can be achieved, let us plug the time variables (11) into the full solution (7), and (8)
$`q^i`$ $`=`$ $`q^i(t^m(q^i,p_i);\stackrel{~}{q}^a,\stackrel{~}{p}_a),`$ (23)
$`p_i`$ $`=`$ $`p_i(t^m(q^i,p_i);\stackrel{~}{q}^a,\stackrel{~}{p}_a).`$ (24)
Last equations are very important, they relate the original phase space variables $`q^i`$, and $`p_i`$ with the physical states of the physical phase space $`(\stackrel{~}{q}^a,\stackrel{~}{p}_a)`$. These equations admite two, related, interpretations. First, they give the relational evolution of the coordinates $`q^i`$ and the momenta $`p_i`$ for any fixed point $`(\stackrel{~}{q}^a,\stackrel{~}{p}_a)`$ of the physical phase space, i.e., it is possible to choose M coordinates denoted by $`q^m`$ (or momenta $`p_m`$; or a combination of both) as ‘clocks’ and describe the evolution of the remaining set of coordinates and momenta as functions of the $`q^m`$ for any physical state $`(\stackrel{~}{q}^a,\stackrel{~}{p}_a)`$ of the system. Second, if we fix the values of this $`M`$ coordinates, say $`q^m=q^m`$ then, the before mentioned expressions of coordinates and momenta give $`M`$-parameter families of physical observables defined on $`\gamma _{ph}`$, $`q^m`$ being the parameters. Eqs. (23), and (24) are evolving constants of motion in the sense of Rovelli . This concept captures the essence that the before mentioned observables (defined on $`\gamma _{ph}`$) describe the relational evolution of the coordinates $`q^i`$ and momenta $`p_i`$, and at the same time they are physical observables.
Let us consider particular cases of (11), say
$`t^m`$ $`=`$ $`q^m,m=1,2,\mathrm{},M,`$ (25)
then (23), and (24) acquire the form
$`q^m`$ $`=`$ $`q^m,m=1,2,\mathrm{},M,`$ (26)
$`q^i`$ $`=`$ $`q^i(q^m;\stackrel{~}{q}^a,\stackrel{~}{p}_a),i=M+1,\mathrm{},N,`$ (27)
$`p_i`$ $`=`$ $`p_i(q^m;\stackrel{~}{q}^a,\stackrel{~}{p}_a),i=1,\mathrm{},N.`$ (28)
Thus, the ‘clocks’ are given by $`q^m`$ and last two pairs of equations are the evolving constants of motion involved. Other particular case is given by
$`t^m`$ $`=`$ $`p_m,m=1,2,\mathrm{},M,`$ (29)
and (23), and (24) acquire the form
$`q^i`$ $`=`$ $`q^i(p_m;\stackrel{~}{q}^a,\stackrel{~}{p}_a),i=1,2,\mathrm{},N,`$ (30)
$`p_m`$ $`=`$ $`p_m,m=1,2,\mathrm{},M,`$ (31)
$`p_i`$ $`=`$ $`p_i(p_m;\stackrel{~}{q}^a,\stackrel{~}{p}_a),i=M+1,\mathrm{},N.`$ (32)
In this case, the ‘clocks’ are $`p_m`$ and last two pairs of equations are the evolving constants of motion required.
As we have seen, the general relations involving the coordinates $`q^i`$ and momenta $`p_i`$ with the physical states $`(\stackrel{~}{q}^a,\stackrel{~}{p}_a)`$ is given by (23), and (24). The explicit form of (23), and (24) could be complicated for particular theories, but this fact would rise technical rather than conceptual difficulties (see for the opposite viewpoint where the authors rise questions on interpretation, consistency, and the degree to which the resulting quantum theory emerging from the before classical dynamics coincide with, or generalizes, the usual non-relativistic theory). Thus, Eqs. (23), and (24) constitute the full set of evolving constants needed in the relational description of the dynamics of generally covariant theories with a finite number $`D`$ of degrees of freedom. The solution (23), and (24) sits in the spirit that in covariant theories there is non privileged observable with respect to which to describe evolution, and that only relational evolution makes sense. From this point of view, general covariance forces us to use relational evolution, namely, to describe the change of some variables of the system with respect to the others. This is the essence of relationism, which appears to be the natural language for describing the evolution of the degrees of freedom of generally covariant theories .
In addition, in this paper, we propose an alternative mechanism to generate the evolving constants. This mechanism is essentially to compute the action of the Hamiltonian vector fields $`X_{dC_m}`$ on some evolving constant $`E^1`$
$`X_{dC_m}(E)`$ $`=:`$ $`E^m.`$ (33)
The evolving function $`E^1`$ depends on the canonical coordinates of the unconstrained phase space $`q^i`$, and $`p_i`$ as well as on the canonical coordinates $`\stackrel{~}{q}^a`$, and $`\stackrel{~}{p}_a`$ of the physical phase space. Therefore, in the computation of the action of the Hamiltonian vector fields (5) on the evolving function we can proceed in two ways. First, taking the observables $`\stackrel{~}{q}^a`$, and $`\stackrel{~}{p}_a`$ constants in the dependence of the evolving function $`E`$. This can be done because $`\stackrel{~}{q}^a`$, and $`\stackrel{~}{p}_a`$ are constant along the orbits. b) Taking the explicit dependence of the physical observables in terms of the canonical variables of the unconstrained phase space given by (9), and (10). Of course, both approaches lead to the same results. The repeated application of the Hamiltonian vector fields on the new evolving constants $`E^2`$, $`E^3`$,.., gives another evolving constants, and so on until no new evolving constants are obtained, and the process ends. From the knowledge of the evolving constants and the expressions of the physical observables, the full solution of the dynamics of the system encoded in Eqs. (7)-(11) is obtained.
Quantum dynamics. Let us begin with the quantum description of the system. We use the Dirac method. In the same way as in the classical dynamics we have three spaces $`(\gamma _{ex},\omega _{ex})`$, $`(P(\gamma _{ph},\text{Orbits}),\mathrm{\Omega })`$, and $`(\gamma _{ph},\omega _{ph})`$. In the quantum theory we have three Hilbert spaces; the unconstrained Hilbert space $``$ or a suitable extension of it if the constraints have continuum spectrum, the physical Hilbert space $`_{phys}`$, and the reduced Hilbert space $`_r`$ obtained by projecting $`_{phys}`$. Suppose we have solved the quantum theory in a full way, i.e., we have the physical Hilbert space $`_{phys}`$ of the theory. A general physical state $`\varphi `$ of the system is killed by all the constraints of the theory $`\widehat{C}_m\psi =0`$, and it is given by
$`\psi `$ $`=`$ $`{\displaystyle \underset{n_1,n_2,\mathrm{},n_D}{}}c_{n_1,n_2,\mathrm{},n_D}n_1,n_2,\mathrm{},n_D.`$ (34)
in Dirac notation. Here, the physical states are labelled by the quantum numbers $`n_a`$, $`a=1,2,\mathrm{},D`$ which come from a complete set of commuting physical observables $`\widehat{O}_a`$, $`a=1,2,\mathrm{},D`$ of the system
$`\widehat{O}_an_1,n_2,\mathrm{},n_D`$ $`=`$ $`O(n_a)n_1,n_2,\mathrm{},n_D.`$ (35)
Of course, these quantum observables are combinations of the physical observables $`\widehat{\stackrel{~}{q}}^a`$, and $`\widehat{\stackrel{~}{p}}_a`$. We have come to the heart of the problem, how to describe relational evolution in the quantum theory.
Quantum evolving constants. Let us see how the quantum version of the classical evolving constants looks. The idea is to search for a representation of the physical states (34) in the reduced Hilbert space associated with the physical phase space of the system. Explicitly
$`\psi (\stackrel{~}{q}^a)`$ $`=`$ $`\stackrel{~}{q}^a\psi ={\displaystyle \underset{n_1,n_2,\mathrm{},n_D}{}}c_{n_1,n_2,\mathrm{},n_D}\stackrel{~}{q}^an_1,n_2,\mathrm{},n_D.`$ (36)
The inner product in the Hilbert space
$`\psi \varphi `$ $`=`$ $`{\displaystyle 𝑑\mu (\stackrel{~}{q}^a)\psi ^{}(\stackrel{~}{q}^a)\varphi (\stackrel{~}{q}^a)},`$ (37)
can be determined with the condition that the operators $`\widehat{\stackrel{~}{q}}^a`$, and $`\widehat{\stackrel{~}{p}}_a`$ be hermitian operators and with the implementation of the action of the operators $`\widehat{\stackrel{~}{q}}^a`$, $`\widehat{\stackrel{~}{p}}_a`$ on this Hilbert space. Notice also that is always possible to build creation and annihilation operators $`\widehat{a}_a=\widehat{\stackrel{~}{q}}^a+i\widehat{\stackrel{~}{p}}_a`$, $`\widehat{a}_a^{}=\widehat{\stackrel{~}{q}}^ai\widehat{\stackrel{~}{p}}_a`$ for each pair of canonical operators $`\widehat{\stackrel{~}{q}}^a`$, and $`\widehat{\stackrel{~}{p}}_a`$ because the number of these operators is even. $`\widehat{a}_a`$, $`\widehat{a}_a^{}`$ can help in the construction of $`_r`$.
With the before machinery, the quantum version of the evolving constants is
$`\widehat{q}^i`$ $`=`$ $`q^i(t^m(q^i,p_i);\widehat{\stackrel{~}{q}}^a,\widehat{\stackrel{~}{p}}_a),`$ (38)
$`\widehat{p}_i`$ $`=`$ $`p_i(t^m(q^i,p_i);\widehat{\stackrel{~}{q}}^a,\widehat{\stackrel{~}{p}}_a),`$ (39)
or, equivalently,
$`\psi \widehat{q}^i\psi `$ $`=`$ $`\psi q^i(t^m(q^i,p_i);\widehat{\stackrel{~}{q}}^a,\widehat{\stackrel{~}{p}}_a)\psi ,`$ (40)
$`\psi \widehat{p}_i\psi `$ $`=`$ $`\psi p_i(t^m(q^i,p_i);\widehat{\stackrel{~}{q}}^a,\widehat{\stackrel{~}{p}}_a)\psi ,`$ (41)
where the mean values are computed with the inner product (37). In the case of parameterized systems, last equations reduce to the standard ones which describe the evolution of the position and momenta operators as well as the evolution of the mean values of the position and momenta operators in the Heisenberg picture. Of course, the well-known ordering problems for the operators might appear here too.
## II parameterized harmonic oscillator
### A Classical dynamics
In order to make these ideas concrete, let us consider a familiar example: the parameterized harmonic oscillator, which action is
$`S={\displaystyle 𝑑\tau \left[\frac{dx}{d\tau }p+\frac{dt}{d\tau }p_t\lambda \left(p_t+\frac{p^2}{2m}+\frac{1}{2}m\omega ^2x^2\right)\right]}.`$ (42)
The unconstrained classical space $`\mathrm{\Gamma }`$ is coordinatized by the canonical pairs $`(x,p)`$, and $`(t,p_t)`$. By doing the variation of the action with respect to $`x`$, $`p`$, $`t`$, and $`p_t`$ we find the equations of motion
$`{\displaystyle \frac{dp}{d\tau }}=\lambda m\omega ^2x,{\displaystyle \frac{dx}{d\tau }}=\lambda {\displaystyle \frac{p}{m}},{\displaystyle \frac{dp_t}{d\tau }}=0,{\displaystyle \frac{dt}{d\tau }}=\lambda .`$ (43)
The variation of the action with respect to the Lagrange multiplier $`\lambda `$ gives the first class constraint
$`C=p_t+{\displaystyle \frac{p^2}{2m}}+{\displaystyle \frac{1}{2}}m\omega ^2x^2.`$ (44)
The classical dynamics is the unfolding of the gauge symmetry of the system generated by the first class constraint $`C`$. The gauge orbit on the constrained surface $`C=0`$ is the integral curve of the Hamiltonian vector field
$`X_{dC}`$ $`=`$ $`{\displaystyle \frac{}{t}}{\displaystyle \frac{p}{m}}{\displaystyle \frac{}{x}}+m\omega ^2x{\displaystyle \frac{}{p}}.`$ (45)
If we have a solution $`x(\tau )`$, $`p(\tau )`$, $`t(\tau )`$, and $`p_t(\tau )`$ of the equations of motion (43), any other solution $`x^{}(\tau )`$, $`p^{}(\tau )`$, $`t^{}(\tau )`$, and $`p_{t}^{}{}_{}{}^{}(\tau )`$ can be found through the relations
$`x^{}(\tau )`$ $`=`$ $`\mathrm{cos}(\theta (\tau ))x(\tau )+{\displaystyle \frac{1}{m\omega }}\mathrm{sin}(\theta (\tau ))p(\tau ),`$ (46)
$`p^{}(\tau )`$ $`=`$ $`m\omega \mathrm{sin}(\theta (\tau ))x(\tau )+\mathrm{cos}(\theta (\tau ))p(\tau ),`$ (47)
$`t^{}(\tau )`$ $`=`$ $`{\displaystyle \frac{\theta (\tau )}{\omega }}+t(\tau ),`$ (48)
$`p_{t}^{}{}_{}{}^{}`$ $`=`$ $`p_t,`$ (49)
that connect all the solutions, while the Lagrange multiplier transforms as
$`\lambda ^{}(\tau )`$ $`=`$ $`\lambda (\tau )+{\displaystyle \frac{1}{\omega }}\dot{\theta }(\tau ),`$ (50)
in order to leave the action invariant , here $`\dot{\theta }(\tau )=\frac{d\theta (\tau )}{d\tau }`$.
Let us construct the general solution in a given gauge. We choose the gauge $`\lambda =1`$. We still have one gauge fixing to impose at $`\tau =0`$. We choose $`t(0)=0`$. Using the constraint equation and the solution of (43), we obtain
$`x(\tau )`$ $`=`$ $`A\mathrm{cos}(\omega \tau )+B\mathrm{sin}(\omega \tau ),`$ (51)
$`t(\tau )`$ $`=`$ $`\tau ,`$ (52)
$`p(\tau )`$ $`=`$ $`m\omega A\mathrm{sin}(\omega \tau )+m\omega B\mathrm{cos}(\omega \tau ),`$ (53)
$`p_t(\tau )`$ $`=`$ $`{\displaystyle \frac{1}{2}}m\omega ^2(A^2+B^2),`$ (54)
where $`(A,B)`$ are the physical observables that coordinatize the physical phase space of the system, which is $`R^2`$. It is clear that $`x`$, $`t`$, $`p`$ are non-observables (they depend on $`\tau `$).
The two physical observables $`(A,B)`$ can be expressed in terms of the phase space variables as
$`A`$ $`=`$ $`\mathrm{cos}(\omega t)x{\displaystyle \frac{1}{m\omega }}\mathrm{sin}(\omega t)p,`$ (55)
$`B`$ $`=`$ $`{\displaystyle \frac{1}{m\omega }}\mathrm{cos}(\omega t)p+\mathrm{sin}(\omega t)x`$ (56)
Notice that $`A=x(t=0)x_0`$, and $`B=\frac{p(t=0)}{m\omega }\frac{p_0}{m\omega }`$, i.e., the position $`x_0`$ of the harmonic oscillator when the internal clock measures $`t=0`$, and the momentum $`p_0`$ when the internal clock measures $`t=0`$ are (physical) observables. Moreover, Eq. (56) means that the precise combination of the position $`x=X`$ and the momentum $`p=P`$ of the harmonic oscillator, when the internal clock indicates $`t=T`$ in the form expressed by the formula (56) is an observable of the (composed) system: harmonic oscillator + internal clock. These observables have vanishing Poisson brackets with the first class constraint $`C`$ as required by the formalism of constrained systems. Actually, the Dirac method requires observables to have weakly vanishing Poisson brackets with the first class constraints. Here, the observables $`A`$, $`B`$ have strong vanishing Poisson brackets with the constraint $`C`$. The Poisson brackets between $`A`$ and $`B`$ in the physical phase space reads
$`\{A,B\}`$ $`=`$ $`{\displaystyle \frac{1}{m\omega }}.`$ (57)
Classical evolving constants. From (54), we obtain the evolving constant
$`x`$ $`=`$ $`x_0\mathrm{cos}(\omega t)+{\displaystyle \frac{p_0}{m\omega }}\mathrm{sin}(\omega t),`$ (58)
of the system. As before mentioned, last equation admits two, related, interpretations. First, for any fixed point $`(A,B)`$ (equivalently $`(x_0,p_0)`$) of the physical phase space, (58) gives the relative evolution of the configuration variables $`x`$, and $`t`$ of the system
$`x`$ $`=`$ $`X(t;x_0,p_0)=x_0\mathrm{cos}(\omega t)+{\displaystyle \frac{p_0}{m\omega }}\mathrm{sin}(\omega t).`$ (59)
Second, for any fixed $`t`$, it gives a one-parameter family of physical observables, $`t`$ being the parameter, on the physical phase space.
Generation of evolving constants. We define the function $`E^1`$ on $`\mathrm{\Gamma }`$
$`E^1(x,t,p,p_t)`$ $`:=`$ $`xx_0\mathrm{cos}(\omega t){\displaystyle \frac{p_0}{m\omega }}\mathrm{sin}(\omega t).`$ (60)
The restriction of this function on the constraint surface is $`E^1_C=0`$. The action of the Hamiltonian vector field $`X_{dC}`$ on $`E^1`$ is
$`X_{dC}(E^1)`$ $`=:`$ $`E^2=\omega x_0\mathrm{sin}(\omega t)+{\displaystyle \frac{p_0}{m}}\mathrm{cos}(\omega t){\displaystyle \frac{p}{m}},`$ (61)
and the restriction of $`E^2`$ on the constraint surface is $`E^2_C=0`$, and more important, the equation $`E^2_C=0`$ is precisely an evolving constant
$`p`$ $`=`$ $`m\omega x_0\mathrm{sin}(\omega t)+p_0\mathrm{cos}(\omega t).`$ (62)
Note also that the action of $`X_{dC}`$ on $`E^2`$ gives again $`E^1`$, and the process ends. In other words, the evolving constant (62) was obtained from the application of the Hamiltonian vector field $`X_{dC}`$ on $`E^1`$, and viceversa.
The full solution. In the present case the constraint surface is coordinatized by the coordinates of the physical phase space $`(x_0,p_0)`$ and by the internal time $`t`$. Therefore, Eqs. (7), and (8) acquire the form
$`x`$ $`=`$ $`X(t;x_0,p_0)=x_0\mathrm{cos}(\omega t)+{\displaystyle \frac{p_0}{m\omega }}\mathrm{sin}(\omega t),`$ (63)
$`t`$ $`=`$ $`T(t;x_0,p_0)=t`$ (64)
$`p`$ $`=`$ $`P(t;x_0,p_0)=m\omega x_0\mathrm{sin}(\omega t)+p_0\mathrm{cos}(\omega t),`$ (65)
$`p_t`$ $`=`$ $`P_T(t;x_0,p_0)={\displaystyle \frac{p_{0}^{}{}_{}{}^{2}}{2m}}{\displaystyle \frac{1}{2}}m\omega ^2x_{0}^{}{}_{}{}^{2},`$ (66)
Of course, last equations are also (23), and (24). Notice that Eqs. (9), and (10) acquire the form
$`x_0`$ $`=`$ $`\mathrm{cos}(\omega t)x{\displaystyle \frac{1}{m\omega }}\mathrm{sin}(\omega t)p,`$ (67)
$`p_0`$ $`=`$ $`\mathrm{cos}(\omega t)p+m\omega \mathrm{sin}(\omega t)x,`$ (68)
and the dependence of the orbit coordinate $`x^1=t`$, see (11), is
$`t=T(x,t,p,p_t)`$ $`=`$ $`t.`$ (69)
Last equations constitute the full solution of the classical dynamics of the system.
Notice that the internal time variable $`x^1=t=T(x,t,p,p_t)=t`$ is not a physical observable because the Poisson bracket with the first class constraint does not vanish. Nevertheless, when we take the full solution into account we can express $`t=\stackrel{~}{T}(x,p,p_0,x_0)`$, given by
$`\mathrm{cos}\omega t`$ $`=`$ $`{\displaystyle \frac{\left(\frac{1}{2}m\omega xx_0+\frac{1}{2m}pp_0\right)}{H_0}},`$ (70)
with $`H_0=\frac{1}{2m}p_{0}^{}{}_{}{}^{2}+\frac{1}{2}m\omega x_{0}^{}{}_{}{}^{2}`$. The above expression is an evolving constant. From this point of view, the internal clock $`t`$ defines a two-parameter family of physical observables on the physical phase space; $`x`$, and $`p`$ being the parameters. So, the internal clock $`t`$ becomes a physical clock $`t(x,p)`$, namely, a physical observable when the full solution is considered. We restrict the analysis to a branch of the above multivalued function to compute the time $`t(x,p)`$ at which the particle reaches the position $`x`$ and the momentum $`p`$ evolving from an initial position $`x_0`$ and momentum $`p_0`$
$`t(x,p)`$ $`=`$ $`\stackrel{~}{T}(x,p;x_0,p_0)={\displaystyle \frac{1}{\omega }}\text{arc cos}\left({\displaystyle \frac{\frac{1}{2}m\omega xx_0+\frac{1}{2m}pp_0}{H_0}}\right).`$ (71)
Or in terms of $`x`$ only
$`t_\pm (x)={\displaystyle \frac{1}{\omega }}\text{arc cos}\left({\displaystyle \frac{\frac{1}{2}m\omega xx_0\pm \sqrt{\frac{1}{2m}\left(H_0\frac{1}{2}m\omega x^2\right)}p_0}{H_0}}\right).`$ (72)
These classical expressions have a quantum version as we will see later.
### B Quantum dynamics
At quantum level, as Dirac showed, the physical states are those killed by the first class constraint. We associate abstract operators with the classical coordinates and momenta, given by
$`x\widehat{X},t\widehat{T},p\widehat{P},p_t\widehat{P}_T,`$ (74)
which satisfy the Dirac rule
$`[\widehat{X},\widehat{P}]=i\mathrm{},[\widehat{T},\widehat{P}_T]=i\mathrm{},`$ (75)
and by inserting these operators in the quantum constraint $`\widehat{C}\psi `$, this equation becomes
$`\left(\widehat{P}_T+{\displaystyle \frac{\widehat{P}^2}{2m}}+{\displaystyle \frac{1}{2}}m\omega ^2\widehat{X}^2\right)\psi `$ $`=`$ $`0.`$ (76)
Any physical state can be expressed in terms of the single quantum number of the harmonic oscillator, in abstract Dirac notation
$`\psi `$ $`=`$ $`{\displaystyle \underset{n}{}}C_nn,`$ (77)
$`\widehat{I}`$ $`=`$ $`{\displaystyle \underset{n}{}}nn.`$ (78)
In last expression, the physical states $`\psi `$ are ‘frozen’ (i.e. they are abstract vectors), the complex coefficients $`C_n`$ are constants. Notice that we have not choosen the coordinate basis yet. Taking a ‘coordinate representation’ $`x,t`$ where the operators acquire the form
$`x,t\widehat{X}\psi `$ $`=`$ $`xx,t\psi ,`$ (79)
$`x,t\widehat{T}\psi `$ $`=`$ $`tx,t\psi ,`$ (80)
$`x,t\widehat{P}\psi `$ $`=`$ $`{\displaystyle \frac{\mathrm{}}{i}}{\displaystyle \frac{}{x}}x,t\psi ,`$ (81)
$`x,t\widehat{P}_T\psi `$ $`=`$ $`{\displaystyle \frac{\mathrm{}}{i}}{\displaystyle \frac{}{t}}x,t\psi ,`$ (82)
any physical state vector $`\psi `$ is expanded in the coordinate basis $`x,t`$ as
$`x,t\psi `$ $`=`$ $`\psi (x,t)={\displaystyle \underset{n}{}}C_nx,tn,`$ (83)
$`x^{},t^{}x,t`$ $`=`$ $`{\displaystyle \underset{n}{}}x^{},t^{}nnx,t,`$ (84)
with $`x,tn=e^{\frac{i}{\mathrm{}}E_nt}f_n(x)`$, $`E_n=\mathrm{}\omega \left(n+\frac{1}{2}\right)`$. Thus, in the Dirac framework, the coordinate representation is nothing but the ‘Heisenberg picture’ of the standard quantum mechanics, where the coordinate basis $`x,t`$ is ‘rotating’ and the physical state $`\psi `$ is fixed (see Eq. (78) where the coefficients $`C_n`$ are constant complex numbers).
Schrödinger equation. In addition, we can build a ‘Schrödinger basis’ from the ‘Heisenberg basis’ $`x,t`$. In this ‘Schrödinger basis’, which we denote by $`x`$, the state vector is ‘moving around’ the ‘fixed basis’ $`x`$. Explicitly,
$`\psi (x,t)`$ $`=`$ $`x\psi (t),\text{Schrödinger basis }x,`$ (85)
with
$`\psi (t)`$ $`=`$ $`{\displaystyle \underset{n}{}}\stackrel{~}{C}_n(t)\stackrel{~}{n}={\displaystyle \underset{n}{}}C_ne^{\frac{i}{\mathrm{}}E_nt}\stackrel{~}{n},`$ (86)
$`x\stackrel{~}{n}`$ $`=`$ $`f_n(x).`$ (87)
Taking the derivative with respect to the coordinate $`t`$ of $`\psi (t)`$, the familiar Schrödinger equation emerges in the formalism
$`i\mathrm{}{\displaystyle \frac{d}{dt}}\psi (t)`$ $`=`$ $`\widehat{H}\psi (t),`$ (88)
with $`\widehat{H}=\frac{1}{2m}\widehat{P}^2+\frac{1}{2}m\omega ^2\widehat{X}^2`$. As usual, the physical vector $`\psi (t)`$ evolves in $`t`$ while the coordinate basis $`x`$ is fixed. In other words, if we consider the system composed of the harmonic oscillator plus the clock together, we are describing the evolution of the degrees of freedom of the harmonic oscillator with respect to the internal clock itself, that is to say, the evolution of one part of the system with respect to the rest of it. In the next section, we will carry out the same procedure we applied here in order to analyze the meaning of evolution in generally covariant quantum theories.
Quantum evolving constants. Let us now go to the quantum version of the evolving constants. The Hilbert space is built with the implementation of the physical state vectors $`\psi =_nC_nn`$ in the reduced Hilbert space $`_r`$ associated with the physical phase space of the harmonic oscillator. In the present case
$`\psi (x_0)`$ $`=`$ $`x_0\psi ={\displaystyle \underset{n}{}}C_nf_n(x_0).`$ (89)
The inner product in $`_r`$
$`\psi \varphi `$ $`=`$ $`{\displaystyle 𝑑\mu (x_0)\psi ^{}(x_0)\varphi (x_0)}`$ (90)
can be determined with the condition that the operators $`\widehat{x}_0`$, and $`\widehat{p}_0`$ be hermitian operators.
Thus, in the classical expression
$`x`$ $`=`$ $`X(t;x_0,p_0)=x_0\mathrm{cos}(\omega t)+{\displaystyle \frac{p_0}{m\omega }}\mathrm{sin}(\omega t),`$ (91)
$`x_0`$, and $`p_0`$ are physical observables given by the Eq. (56) and they become operators acting on $`_r`$, so the quantum version of the classical evolving constant is
$`\widehat{x}(t)`$ $`=`$ $`x(t;\widehat{x}_0,\widehat{p}_0)=\widehat{x}_0\mathrm{cos}(\omega t)+{\displaystyle \frac{\widehat{p}_0}{m\omega }}\mathrm{sin}(\omega t),`$ (92)
which is the well-known evolution equation for the position operator $`\widehat{X}`$ in the Heisenberg picture.
In addition, the classical expression
$`p`$ $`=`$ $`P(t;x_0,p_0)=m\omega x_0\mathrm{sin}(\omega t)+p_0\mathrm{cos}(\omega t),`$ (93)
has its quantum analog
$`\widehat{p}(t)`$ $`=`$ $`p(t;\widehat{x}_0,\widehat{p}_0)=m\omega \widehat{x}_0\mathrm{sin}(\omega t)+\widehat{p}_0\mathrm{cos}(\omega t),`$ (94)
and finally
$`\widehat{p}_t`$ $`=`$ $`p_t(t;\widehat{x}_0,\widehat{p}_0)={\displaystyle \frac{\widehat{p}_{0}^{}{}_{}{}^{2}}{2m}}{\displaystyle \frac{1}{2}}m\omega ^2\widehat{x}_{0}^{}{}_{}{}^{2}.`$ (95)
In summary, for the case of parameterized systems, the quantum version of the evolving constants equations constitutes the Heisenberg equations for the physical operators involved in each particular theory. In the case of the harmonic oscillator, Eqs. (92) and (94).
Time operator. The classical expression (71) becomes an operator $`\widehat{T}(X,P)=t(X,P;\widehat{x}_0,\widehat{p}_0)`$ which is defined on the reduced Hilbert space $`_r`$. Taken arbitrarily the order of the operators, we have
$`\widehat{T}(X,P)`$ $`=`$ $`{\displaystyle \frac{1}{\omega }}\text{arc cos}\left({\displaystyle \frac{\frac{1}{2}m\omega X\widehat{x}_0+\frac{1}{2m}P\widehat{p}_0}{\widehat{H}_0}}\right),`$ (96)
with $`\widehat{H}_0=\frac{1}{2m}\widehat{p}_{0}^{}{}_{}{}^{2}+\frac{1}{2}m\omega \widehat{x}_{0}^{}{}_{}{}^{2}`$. From this operator, we can compute the ‘time of arrival’ operator $`\widehat{T}(X)`$
$`\widehat{T}_\pm (X)={\displaystyle \frac{1}{\omega }}\text{arc cos}\left({\displaystyle \frac{\frac{1}{2}m\omega X\widehat{x}_0\pm \sqrt{\frac{1}{2m}\left(\widehat{H}_0\frac{1}{2}m\omega X^2\right)}\widehat{p}_0}{\widehat{H}_0}}\right),`$ (97)
associated with the time at which the harmonic oscillator is detected with an apparatus located in $`x=X`$. The ‘time of arrival’ operator for a free particle has been studied in . The analysis of the ‘time of arrival’ operator for the harmonic oscillator deserves to be studied.
## III $`SL(2,R)`$ model with two Hamiltonian constraints
### A Classical dynamics
Let us see how the relative evolution looks in a non familiar generally covariant model. A nonlinear generally covariant system with two Hamiltonian constraints and with one physical degree of freedom was introduced in . This model mimics the constraint structure of general relativity. Here, we continue the study of this model. In particular, we display the full set of evolving constants required in its classical and quantum dynamics. Moreover, for a Schrödinger-like equation of motion arises in its quantum dynamics, we compare the meaning of time (evolution) in both, evolving constants and Schrödinger-like equation, viewpoints.
First, a brief summary of its classical dynamics, for more details and its physical interpretation see Ref.. The model is defined by the action
$`S[\stackrel{}{u},\stackrel{}{v},N,M,\lambda ]={\displaystyle \frac{1}{2}}{\displaystyle 𝑑t\left[N(𝒟\stackrel{}{u}^2+\stackrel{}{v}^2)+M(𝒟\stackrel{}{v}^2+\stackrel{}{u}^2)\right]},`$ (99)
where
$$𝒟\stackrel{}{u}=\frac{1}{N}(\dot{\stackrel{}{u}}\lambda \stackrel{}{u}),𝒟\stackrel{}{v}=\frac{1}{M}(\dot{\stackrel{}{v}}+\lambda \stackrel{}{v});$$
(100)
the two Lagrangian dynamical variables $`\stackrel{}{u}=(u^1,u^2)`$ and $`\stackrel{}{v}=(v^1,v^2)`$ are two-dimensional real vectors; $`N`$, $`M`$ and $`\lambda `$ are Lagrange multipliers. The squares are taken in $`R^2`$: $`\stackrel{}{u}^2=\stackrel{}{u}\stackrel{}{u}=(u^1)^2+(u^2)^2`$. The action can be put in the Hamiltonian form
$`S[\stackrel{}{u},\stackrel{}{v},\stackrel{}{p},\stackrel{}{\pi },\lambda ^m]={\displaystyle 𝑑\tau \left[\dot{\stackrel{}{u}}\stackrel{}{p}+\dot{\stackrel{}{v}}\stackrel{}{\pi }\lambda ^mC_m\right]}.`$ (101)
The canonical pairs that coordinatize the unconstrained classical phase space are $`(u^1,p^1)`$, $`(u^2,p^2)`$, $`(v^1,\pi ^1)`$, and $`(v^2,\pi ^2)`$. Also $`\lambda ^1=N`$, $`\lambda ^2=M`$, and $`\lambda ^3=\lambda `$. The first class constraints have the form
$`C_1`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(\stackrel{}{p}^2\stackrel{}{v}^2\right),`$ (102)
$`C_2`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(\stackrel{}{\pi }^2\stackrel{}{u}^2\right),`$ (103)
$`C_3`$ $`=`$ $`\stackrel{}{u}\stackrel{}{p}\stackrel{}{v}\stackrel{}{\pi },`$ (104)
which algebra is isomorphic to the $`sl(2,R)`$ Lie algebra
$`\{C_1,C_2\}`$ $`=`$ $`C_3`$ (105)
$`\{C_1,C_3\}`$ $`=`$ $`2C_1`$ (106)
$`\{C_2,C_3\}`$ $`=`$ $`2C_2.`$ (107)
The classical dynamics is the unfolding of the gauge symmetry generated by the Hamiltonian vector fields
$`X_{dC_1}`$ $`=`$ $`\stackrel{}{p}\stackrel{}{}_u\stackrel{}{v}\stackrel{}{}_\pi ,`$ (108)
$`X_{dC_2}`$ $`=`$ $`\stackrel{}{\pi }\stackrel{}{}_v\stackrel{}{u}\stackrel{}{}_p,`$ (109)
$`X_{dC_3}`$ $`=`$ $`\stackrel{}{u}\stackrel{}{}_u+\stackrel{}{v}\stackrel{}{}_v+\stackrel{}{p}\stackrel{}{}_p\stackrel{}{\pi }\stackrel{}{}_\pi ,`$ (110)
associated with the first class constraints of the model.
The physical phase space can be coordinated by the points $`(J,\varphi ,ϵ,ϵ^{})`$, and these physical observables have the following form
$`ϵ`$ $`=`$ $`{\displaystyle \frac{u^1p^2p^1u^2}{|u^1p^2p^1u^2|}},`$ (111)
$`ϵ^{}`$ $`=`$ $`{\displaystyle \frac{\pi ^1v^2v^1\pi ^2}{|\pi ^1v^2v^1\pi ^2|}},`$ (112)
$`J`$ $`=`$ $`|u^1p^2p^1u^2|,`$ (113)
$`\varphi `$ $`=`$ $`\mathrm{arctan}{\displaystyle \frac{u^1v^2p^1\pi ^2}{u^1v^1p^1\pi ^1}}.`$ (114)
The Poisson brackets between $`J`$ and $`\varphi `$ in the reduced phase space reads
$`\{J,\varphi \}=ϵϵ^{}.`$ (115)
Classical evolving constants. Finally, the relation between the Lagrangian variables $`(\stackrel{}{u},\stackrel{}{v})`$ and the physical states $`(J,\varphi ,ϵ,ϵ^{})`$
$`\left[u^1v^1+ϵϵ^{}u^2v^2\right]\mathrm{cos}\varphi +\left[u^1v^2ϵϵ^{}u^2v^1\right]\mathrm{sin}\varphi `$ $`=`$ $`J,`$ (116)
which leads to the notion of evolving constants of the system . The evolving constants give the evolution of the Lagrangian variables of the system in a gauge invariant way, i.e., for any fixed physical state of the system $`(J,\varphi ,ϵ,ϵ^{})`$, Eq. (116) gives the change of one of the four coordinates as a function of the other three coordinates, say
$`U^1(x,y,z;J,\varphi ,ϵ,ϵ^{})={\displaystyle \frac{ϵ^{}x(z\mathrm{cos}\varphi y\mathrm{sin}\varphi )+ϵJ}{ϵ(y\mathrm{cos}\varphi +z\mathrm{sin}\varphi )}}.`$ (117)
This relative evolution among the coordinates is gauge invariant. In addition, for any fixed $`x,y,z`$ last equation gives a three-parameter family of physical observables, the three parameters are the three coordinates $`x,y,z`$, on the physical phase space.
Generation of evolving constants. We start with the evolving constant (116), and define the evolving function $`E^1`$
$`E^1(u,v,p,\pi )`$ $`:=`$ $`\left[u^1v^1+ϵϵ^{}u^2v^2\right]\mathrm{cos}\varphi +\left[u^1v^2ϵϵ^{}u^2v^1\right]\mathrm{sin}\varphi J.`$ (118)
The restriction of $`E^1`$ on the constraint surface vanishes, $`E^1_C=0`$. The action of the Hamiltonian vector field $`X_{dH_1}`$ on $`E^1`$ is
$`X_{dH_1}(E^1)=:E^2`$ $`=`$ $`\left[p_1v^1+ϵϵ^{}p_2v^2\right]\mathrm{cos}\varphi \left[p_1v^2ϵϵ^{}p_2v^1\right]\mathrm{sin}\varphi ,`$ (119)
and the restriction of $`E^2`$ on the constraint surface vanishes, so $`E^2_C=0`$ gives the evolving constant
$`\left[p_1v^1+ϵϵ^{}p_2v^2\right]\mathrm{cos}\varphi +\left[p_1v^2ϵϵ^{}p_2v^1\right]\mathrm{sin}\varphi `$ $`=`$ $`0.`$ (120)
The action of $`X_{dH_1}`$ on $`E^2`$ gives zero, so the process ends. Now, we compute the action of the Hamiltonian vector field $`X_{dH_2}`$ on $`E^1`$
$`X_{dH_2}(E^1)=:E^3`$ $`=`$ $`\left[u^1\pi _1+ϵϵ^{}u^2\pi _2\right]\mathrm{cos}\varphi +\left[u^1\pi _2ϵϵ^{}u^2\pi _1\right]\mathrm{sin}\varphi ,`$ (121)
and the restriction of $`E^3`$ on the constraint surfaces vanishes, so $`E^3_C=0`$ gives the evolving constant
$`\left[u^1\pi _1+ϵϵ^{}u^2\pi _2\right]\mathrm{cos}\varphi +\left[u^1\pi _2ϵϵ^{}u^2\pi _1\right]\mathrm{sin}\varphi =0.`$ (122)
The action of $`X_{dH_2}`$ on $`E^3`$ gives zero, so the process ends. Finally, the computation of the action of the Hamiltonian vector field $`X_{dD}`$ on $`E^1`$
$`X_{dD}(E^1)=:E^4`$ $`=`$ $`E^1J,`$ (123)
so we recover the original evolving constant we start with, and no more evolving can be obtained from (116).
The full solution. Eqs. (23), and (24) acquire the form
$`u^1`$ $`=`$ $`U^1(u^2,v^1,v^2;J,\varphi ,ϵ,ϵ^{})={\displaystyle \frac{ϵ^{}u^2(v^2\mathrm{cos}\varphi v^2\mathrm{sin}\varphi )+ϵJ}{ϵ(v^1\mathrm{cos}\varphi +v^2\mathrm{sin}\varphi )}},`$ (124)
$`u^2`$ $`=`$ $`U^2(u^2,v^1,v^2;J,\varphi ,ϵ,ϵ^{})=u^2,`$ (125)
$`v^1`$ $`=`$ $`V^1(u^2,v^1,v^2;J,\varphi ,ϵ,ϵ^{})=v^1,`$ (126)
$`v^2`$ $`=`$ $`V^2(u^2,v^1,v^2;J,\varphi ,ϵ,ϵ^{})=v^2,`$ (127)
$`p_1`$ $`=`$ $`P_1(u^2,v^1,v^2;J,\varphi ,ϵ,ϵ^{})=ϵ^{}\left(v^1\mathrm{sin}\varphi v^2\mathrm{cos}\varphi \right),`$ (128)
$`p_2`$ $`=`$ $`P_2(u^2,v^1,v^2;J,\varphi ,ϵ,ϵ^{})=ϵ\left(v^1\mathrm{cos}\varphi +v^2\mathrm{sin}\varphi \right),`$ (129)
$`\pi _1`$ $`=`$ $`\mathrm{\Pi }_1(u^2,v^1,v^2;J,\varphi ,ϵ,ϵ^{})={\displaystyle \frac{ϵu^2v^1+ϵ^{}J\mathrm{sin}\varphi }{(v^1\mathrm{cos}\varphi +v^2\mathrm{sin}\varphi )}},`$ (130)
$`\pi _2`$ $`=`$ $`\mathrm{\Pi }_2(u^2,v^1,v^2;J,\varphi ,ϵ,ϵ^{})={\displaystyle \frac{ϵu^2v^2ϵ^{}J\mathrm{cos}\varphi }{(v^1\mathrm{cos}\varphi +v^2\mathrm{sin}\varphi )}},`$ (131)
and the Eqs. (9), and (10) are precisely the Eqs. (114) while the Eqs. (11) acquire the form
$`u^2`$ $`=`$ $`U^2(u^i,v^i,p_i,\pi _i)=u^2,`$ (132)
$`v^1`$ $`=`$ $`V^1(u^i,v^i,p_i,\pi _i)=v^1,`$ (133)
$`v^2`$ $`=`$ $`V^2(u^i,v^i,p_i,\pi _i)=v^2.`$ (134)
So, the dynamics of this model can be described in a relational fashion way. The difference with respect to parameterized systems, as the example of the harmonic oscillator previously analyzed, is that in the present case a single internal time variable is not enough, rather, we need three internal time variables. In the way we have expressed the full solution (131), $`u^2,v^1,v^2`$ are clocks, i.e., once the component $`u^2`$ of the position of the first particle, and the position $`(v^1,v^2)`$ of the second particle are known, the change of the component $`u^1`$ of the first particle and the change of the momenta of both particles $`\stackrel{}{p}`$, $`\stackrel{}{\pi }`$ are also known when the system is an particular physical state $`(J,\varphi ,ϵ,ϵ^{})`$. Therefore, the full relational evolution of the system is expressed in terms of three internal clocks $`u^2`$, $`v^1`$, and $`v^2`$.
### B Quantum dynamics
At quantum level, the model is characterized by the following set of observables
$`\widehat{J}m,ϵ,ϵ^{}`$ $`=`$ $`m\mathrm{}m,ϵ,ϵ^{},`$ (135)
$`\widehat{ϵ}m,ϵ,ϵ^{}`$ $`=`$ $`ϵm,ϵ,ϵ^{},`$ (136)
$`\widehat{ϵ}^{}m,ϵ,ϵ^{}`$ $`=`$ $`ϵ^{}m,ϵ,ϵ^{},`$ (137)
and the physical states are given by
$`\psi ={\displaystyle \underset{mϵ,ϵ^{}}{}}C_{m,ϵ,ϵ^{}}m,ϵ,ϵ^{},`$ (138)
in abstract Dirac notation. In the ‘coordinate representation’ $`u,v,\alpha ,\beta `$, which is nothing but the Heisenberg picture in standard quantum mechanics because all the coordinates $`(u,v,\alpha ,\beta )`$ are put at the same level, the state reads
$`\psi (u,v,\alpha ,\beta )`$ $`=`$ $`u,v,\alpha ,\beta \psi ={\displaystyle \underset{mϵ,ϵ^{}}{}}C_{m,ϵ,ϵ^{}}u,v,\alpha ,\beta m,ϵ,ϵ^{},`$ (139)
with $`u,v,\alpha ,\beta m,ϵ,ϵ^{}=e^{im(ϵ\alpha ϵ^{}\beta )}J_m\left(\frac{uv}{\mathrm{}}\right)`$. Thus the basis $`u,v,\alpha ,\beta `$ is ‘rotating’ and the state $`\psi `$ is fixed, i.e., the coefficients $`C_{mϵϵ^{}}`$ are constant complex numbers. The ‘coordinate representation’ appears as the most ‘democratic’ basis because it does not prefer one coordinate more than the others.
Schrödinger equation. In the same sense that in parameterized systems we were able to build a ‘Schrödinger basis’ from the Heisenberg basis, we can do the same here, and rewrite the physical state (138). In the present example, we can build two Schrödinger bases $`u,v,\beta `$, and $`u,v,\alpha `$. In the first one, the physical state vector (138) is expressed as
$`\psi (u,v,\alpha ,\beta )`$ $`=`$ $`u,v,\beta \psi (\alpha ),`$ (140)
with
$`\psi (\alpha )`$ $`=`$ $`{\displaystyle \underset{mϵ,ϵ^{}}{}}\stackrel{~}{C}_{m,ϵ,ϵ^{}}(\alpha )\stackrel{~}{m,ϵ,ϵ^{}}={\displaystyle \underset{mϵ,ϵ^{}}{}}C_{m,ϵ,ϵ^{}}e^{imϵ\alpha }\stackrel{~}{m,ϵ,ϵ^{}},`$ (141)
$`u,v,\beta \stackrel{~}{m,ϵ,ϵ^{}}`$ $`=`$ $`e^{imϵ^{}\beta }J_m\left({\displaystyle \frac{uv}{\mathrm{}}}\right).`$ (142)
Taking the derivative with respect to the coordinate $`\alpha `$ of $`\psi (\alpha )`$, a Schrödinger equation emerges in the formalism
$`i\mathrm{}{\displaystyle \frac{d}{d\alpha }}\psi (\alpha )`$ $`=`$ $`{\displaystyle \frac{ϵ}{ϵ^{}}}\widehat{O}_{34}\psi (\alpha ),`$ (143)
and the physical observable $`\widehat{O}_{34}`$ has the form
$`u,v,\beta \widehat{O}_{34}\psi `$ $`=`$ $`{\displaystyle \frac{\mathrm{}}{i}}{\displaystyle \frac{}{\beta }}u,v,\beta \psi ,`$ (144)
in the ‘Schrödinger basis’ $`u,v,\beta `$. As expected, in the Schrödinger basis $`u,v,\beta `$ , the state $`\psi (\alpha )`$ evolves while the basis $`u,v,\beta `$ is fixed with respect to $`\alpha `$. This is not a matter of terminology, in fact the evolution equation (143) is well defined, and we are really able of describing evolution under this picture, namely, to describe the change of the some part of the whole state with respect to rest of it, in complete agreement with the spirit of relationism.
In the second Schrödinger basis $`u,v,\alpha `$, the physical state vector (138) is expressed as
$`\psi (u,v,\alpha ,\beta )`$ $`=`$ $`u,v,\alpha \psi (\beta ),`$ (145)
with
$`\psi (\beta )`$ $`=`$ $`{\displaystyle \underset{mϵ,ϵ^{}}{}}\stackrel{~}{\stackrel{~}{C}}_{m,ϵ,ϵ^{}}(\beta )\stackrel{~}{\stackrel{~}{m,ϵ,ϵ^{}}}={\displaystyle \underset{mϵ,ϵ^{}}{}}C_{m,ϵ,ϵ^{}}e^{imϵ^{}\beta }\stackrel{~}{\stackrel{~}{m,ϵ,ϵ^{}}},`$ (146)
$`u,v,\alpha \stackrel{~}{\stackrel{~}{m,ϵ,ϵ^{}}}`$ $`=`$ $`e^{imϵ\alpha }J_m\left({\displaystyle \frac{uv}{\mathrm{}}}\right).`$ (147)
Taking the derivative with respect to the coordinate $`\beta `$ of $`\psi (\beta )`$, a Schrödinger equation emerges in the formalism
$`i\mathrm{}{\displaystyle \frac{d}{d\beta }}\psi (\beta )`$ $`=`$ $`{\displaystyle \frac{ϵ^{}}{ϵ}}\widehat{O}_{12}\psi (\beta ),`$ (148)
and the physical observable $`\widehat{O}_{12}`$ has the form
$`u,v,\alpha \widehat{O}_{12}\psi `$ $`=`$ $`{\displaystyle \frac{\mathrm{}}{i}}{\displaystyle \frac{}{\alpha }}u,v,\alpha \psi ,`$ (149)
in the ‘Schrödinger basis’ $`u,v,\alpha `$.
Quantum evolving constants. The quantum version of the evolving constants is as follows. More precisely, the quantum version of the full classical solution (131) is expressed as
$`\widehat{u}^1`$ $`=`$ $`u^1(u^2,v^1,v^2;\widehat{J},\widehat{\mathrm{sin}\varphi },\widehat{\mathrm{cos}\varphi },\widehat{ϵ},\widehat{ϵ}^{})={\displaystyle \frac{\widehat{ϵ}^{}\widehat{ϵ}(v^2\widehat{\mathrm{cos}\varphi }v^1\widehat{\mathrm{sin}\varphi })+\widehat{J}}{v^1\widehat{\mathrm{cos}\varphi }v^2\widehat{\mathrm{sin}\varphi })}},`$ (150)
$`\widehat{p}_1`$ $`=`$ $`p_1(u^2,v^1,v^2;\widehat{J},\widehat{\mathrm{sin}\varphi },\widehat{\mathrm{cos}\varphi },\widehat{ϵ},\widehat{ϵ}^{})=\widehat{ϵ}^{}\left(v^1\widehat{\mathrm{sin}\varphi }v^2\widehat{\mathrm{cos}\varphi }\right),`$ (151)
$`\widehat{p}_2`$ $`=`$ $`p_2(u^2,v^1,v^2;\widehat{J},\widehat{\mathrm{sin}\varphi },\widehat{\mathrm{cos}\varphi },\widehat{ϵ},\widehat{ϵ}^{})=\widehat{ϵ}\left(v^1\widehat{\mathrm{cos}\varphi }+v^2\widehat{\mathrm{sin}\varphi }\right),`$ (152)
$`\widehat{\pi }_1`$ $`=`$ $`\pi _1(u^2,v^1,v^2;\widehat{J},\widehat{\mathrm{sin}\varphi },\widehat{\mathrm{cos}\varphi },\widehat{ϵ},\widehat{ϵ}^{})={\displaystyle \frac{\widehat{ϵ}u^2v^1+\widehat{ϵ}^{}\widehat{J}\widehat{\mathrm{sin}\varphi }}{v^1\widehat{\mathrm{cos}\varphi }+v^2\widehat{\mathrm{sin}\varphi }}},`$ (153)
$`\widehat{\pi }_2`$ $`=`$ $`\pi _2(u^2,v^1,v^2;\widehat{J},\widehat{\mathrm{sin}\varphi },\widehat{\mathrm{cos}\varphi },\widehat{ϵ},\widehat{ϵ}^{})={\displaystyle \frac{\widehat{ϵ}u^2v^2\widehat{ϵ}^{}\widehat{J}\widehat{\mathrm{cos}\varphi }}{v^1\widehat{\mathrm{cos}\varphi }+v^2\widehat{\mathrm{sin}\varphi }}}.`$ (154)
The meaning of the first equation in (154) is the following: we have to take the mean value of the operator $`\widehat{u}^1`$ with respect to generic states $`\psi `$ of the reduced Hilbert space $`_r`$ of the model .
In summary, the quantum dynamics of parameterized systems can be described in terms of a ‘Schrödinger equation’ or in terms of the ‘Heisenberg picture’. The Schrödinger equation arises as a consequence of the Dirac quantization, as we have seen for the case of the harmonic oscillator. On the other hand, in the $`SL(2,R)`$ model, we were able to build two (dependent) Schrödinger equations, and thus to identify two (dependent) internal time variables $`\alpha `$ and $`\beta `$ with respect to which the physical states of the $`SL(2,R)`$ model evolve. This does not mean that is always possible to single out in general an internal time variable, given by a Schrödinger equation, in generally covariant theories once the Dirac quantization has been performed. Therefore, in general, a Schrödinger equation does not arise in the formalism. The Schrödinger picture, when this picture emerges in the formalism as a consequence of the Dirac quantization, singles out one internal clock only. More important, the quantization of generally covariant theories based on the reduced Hilbert space (generalized Heisenberg picture) need $`M`$ internal clocks, where $`M`$ is the number of first class constraints. In the case of the $`SL(2,R)`$ model, the clocks are $`u^2`$, $`v^1`$, and $`v^2`$ in the generalized Heisenberg picture. In the Schrödinger picture, the internal clock is given by $`\alpha `$ (or $`\beta `$).
Classical limit. Now, we compare the quantum evolving constants of the $`SL(2,R)`$ model with those of the harmonic oscillator in order to get insights on the classical limit of generally covariant theories, and in particular of the $`SL(2,R)`$ model. We expect that the classical limit of generally covariant theories should be attached to the concept of coherent states as it happens in standard quantum mechanics (parameterized systems). In the case of the harmonic oscillator, the coherent states are roughly those states $`\psi `$ in the reduced Hilbert space $`_r`$ such that the mean values $`\psi \widehat{x}(t)\psi `$, and $`\psi \widehat{p}(t)\psi `$ reproduce the classical behavior of the system. Of course this condition is not enough to single out the coherent states of the system. In addition, those states have also to minimize the uncertainty relations of position and momentum. Of course, these two conditions are still not enough to identify the coherent states due to the fact that both conditions are satisfied by both squeezed and coherent states. In the case of the parameterized harmonic oscillator a mechanism that identifies the coherent states is available following standard methods. It is natural to expect that a combination of the coherent states approach to the quantization of generally covariant theories with the full set of evolving constants of motion required in their quantum dynamics displayed here could bring the classical limit of constrained systems.
## IV concluding remarks
We have displayed the full solution of the relational evolution of the degrees of freedom of fully constrained theories with a finite number of degrees of freedom (see Eqs. (23), and (24)). Our procedure follows from the embedding equations of the coordinates and momenta in the unconstrained phase space (see Eqs. (7), and (8)) plus the expressions of the $`M`$ internal time variables (see Eq. (11)). The form of the solution containts all the evolving constants of motion needed in the description of the classical dynamics of fully constrained theories, i.e., we have given the full mathematical solution to the Rovelli’s point of view on the ‘problem of time’ pioneered in Refs. . Of course, the physical (and phylosophical) interpretation is due to Rovelli. Also, we have explored a method to generate those evolving constants. This method consists in the repeated application of the Hamiltonian vector fields associated with the first class constraints on some initial evolving constant. Combining the expressions of this evolving constants with the expressions of the physical observables the full relational evolution of the coordinates and momenta is obtained. Finally, we have also analysed on a general setting the quantum version of the relational evolution of the degrees of freedom of fully constrained theories.
To find the full solution of the relational evolution of the degrees of freedom for gravity, matter fields coupled to gravity (see for the first steps), topological quantum field theories, or for a background-independent string theory constitutes one of the challenges of the new millenium.
## acknowledgments
I am indebted to Carlo Rovelli for conversations about the ‘problem of time’. I also thank financial support provided by the Sistema Nacional de Investigadores (SNI) of the Secretaría de Educación Pública (SEP) of Mexico. |
warning/0002/gr-qc0002044.html | ar5iv | text | # Gravitational Collapse of Dust with a Cosmological Constant
## I Introduction
Recently, Markovic and Shapiro , motivated by some observational suggestions of a positive cosmological constant , have reexamined the effect of this constant on the evolution of a homogeneous dust ball embedded in vacuum. This paper extends their analysis so as to include the inhomogeneous and degenerate cases. The qualitative behavior of the boundary histories are shown by way of effective potential and Penrose-Carter diagrams. The case $`\mathrm{\Lambda }<0`$ is included as it provides for an interesting contrast. The well known case $`\mathrm{\Lambda }=0`$ is not included.
## II Dust
The study of spherically symmetric distributions of matter without pressure in the general theory of relativity has a long history. It is fair to say that the dynamics of this “Lemaître - Tolman - Bondi” solution are well understood, even with a non-vanishing cosmological constant . Whereas the discovery of “shell-focusing” singularities in dust added a new dimension to the dynamics , these singularities are now well studied and are not considered here.
We review the dynamics to set the notation. First, recall that the flow lines of all dust distributions are geodesic. As a consequence, with spherical symmetry we can choose synchronous comoving coordinates $`(\text{r},\theta ,\varphi ,\tau )`$ so that the line element associated with the dust takes the form
$$ds^2=e^{\alpha (\text{r},\tau )}d\text{r}^2+R(\text{r},\tau )^2(d\theta ^2+\mathrm{sin}^2\theta d\varphi ^2)d\tau ^2.$$
(1)
As long as $`R^{^{}}0(^{^{}}\frac{}{\text{r}})`$ we obtain
$$e^{\alpha (\text{r},\tau )}=\frac{R_{}^{^{}}{}_{}{}^{2}}{1+2E(\text{r})}.$$
(2)
A further integration gives one more independent function of r
$$R_{}^{}{}_{}{}^{2}2E(\text{r})\frac{\mathrm{\Lambda }R^2}{3}=\frac{2M(\text{r})}{R}$$
(3)
where $`{}_{}{}^{}\frac{}{\tau }`$. The energy density follows as
$$4\pi \rho (\text{r},\tau )=\frac{M^{^{}}}{R^2R^{^{}}}.$$
(4)
Many explicit forms of $`R(\text{r},\tau )`$ are known, but these are not of interest here.
## III Vacuum
The $`\mathrm{\Lambda }`$ generalization of the Schwarzschild vacuum is well known. In terms of familiar curvature coordinates $`(r,\theta ,\varphi ,t)`$ the line element is given by
$$ds^2=\frac{dr^2}{f(r)}+r^2(d\theta ^2+\mathrm{sin}^2\theta d\varphi ^2)f(r)dt^2,$$
(5)
where
$$f(r)=1\frac{2m}{r}\frac{\mathrm{\Lambda }r^2}{3}.$$
(6)
The associated generalization of the Birkhoff theorem is well known . It is interesting to note that the $`\mathrm{\Lambda }`$ generalization of the Israel theorem is not known. Geodesically complete forms of the metric (5) along with Penrose - Carter diagrams are now well know .
The coordinates $`(r,\theta ,\varphi ,t)`$ are adapted to two Killing vectors and so geodesics of the metric (5) have two constants of motion. The orbits are stably planar and we choose the plane to be $`\theta =\pi /2`$. The momentum conjugate to $`\varphi `$ is the orbital angular momentum l,
$$r^2\dot{\varphi }=\text{l},$$
(7)
and the momentum conjugate to $`t`$ is the energy $`\gamma `$,
$$f(r)\dot{t}=\gamma .$$
(8)
For timelike geodesics we can take $`{}_{}{}^{.}=\frac{d}{d\lambda }`$ where $`\lambda `$ is the proper time. In what follows we are interested in radial motion so that $`\text{l}=0`$. $`\gamma `$, however, plays a central role. The timelike geodesic equations reduce to
$$\gamma ^2\dot{r}^2=f(r).$$
(9)
We can write $`P(r)f(r)`$ and treat $`P`$ as the effective potential of elementary mechanics.
## IV Junction
The junction of dust and vacuum in spherical symmetry by way of the Darmois - Israel conditions is well understood . To summarize, the continuity of the first fundamental form associated with the boundary ($`\mathrm{\Sigma }`$) ensures that the continuity of $`\theta `$ and $`\varphi `$ in metrics (1) and (5) is allowed and that the history of the boundary is given by
$$R(\text{r}_\mathrm{\Sigma },t)=r_\mathrm{\Sigma }.$$
(10)
The continuity of the second fundamental form guarantees that the flow lines of the boundary particles are simultaneously geodesic of both enveloping 4-geometries. The junction conditions demand that
$$M(\text{r}_\mathrm{\Sigma })=m,$$
(11)
and that for $`R^{^{}}0`$
$$E(\text{r}_\mathrm{\Sigma })=\frac{\gamma ^21}{2}.$$
(12)
The case $`R^{^{}}=0`$ gives $`\gamma =0`$.
## V Discussion
The qualitative history of the geodesics of (5), and via (12) therefore of the dust boundary $`\mathrm{\Sigma }`$, can be obtained from a sketch of $`P`$ (and in particular the requirement that $`\gamma ^2P`$). These are shown in Figure 1. (The roots $`(r_0,r_2,r_3)`$ are given explicitly in .) Note that for $`\mathrm{\Lambda }<0`$ all orbits are closed, in contrast to $`\mathrm{\Lambda }0`$. The case $`\gamma =0`$ is unique in the sense that $`\mathrm{\Sigma }`$ traverses the bifurcation of the Killing horizons (in the non-degenerate cases).
The Penrose - Carter diagrams are shown in Figure 2. The possible histories of $`\mathrm{\Sigma }`$ are shown. The dust can be matched to the left or to the right. The degenerate case $`3m=1/\sqrt{\mathrm{\Lambda }}`$ requires a special coordinate construction . Note that here the case $`\gamma =0`$ is associated with unstable equilibrium at the points of internal infinity.
### Acknowledgments
This work was supported by a grant from the Natural Sciences and Engineering Research Council of Canada. I would like to thank Sean Hayward for pointing out the work by Nakao and José Lemos for reminding me of his work on Oppenheimer-Snyder collapse. |
warning/0002/hep-ph0002206.html | ar5iv | text | # Squark-Chargino Production in Polarized Gamma-Proton Collisions at TeV Energy Scale
## I introduction
Although the Standard Model (SM) of elementary particles has been successful with high precision up to the scale of 100 GeV, there are many theoretical reasons that new physics beyond the SM should exist at TeV scale. Among the models of new physics, the supersymmetry (SUSY) seems to be one of the most promising candidates for TeV scale. A number of TeV energy machines have been proposed or constructing, such as the Large Hadron Collider (LHC), Next Linear $`e^+e^{}`$ Collider (NLC) and Linac-Ring type ep and $`\gamma `$p machines. The latter one can be realized by using the beam of high energy photons produced through the Compton backscattering of laser photons off a TeV energy linear $`e^{}`$ (or $`e^+`$) beam. We have already proposed some Linac-Ring type ep and $`\gamma `$p machines . Here we concentrate ourselves only on three of them, i.e., the HERA+ LC, LHC+Linac 1 and LHC+TESLA. Their calculated center of mass energies and luminosities are given in Table 1.
A supersymmetric SM has a new spectrum of particles called SUSY particles which are the partners of all the known particles with the spins differing by $`\frac{1}{2}`$. Some of the SUSY partners are scalar leptons (sleptons, $`\stackrel{~}{l}`$), scalar quarks (squarks, $`\stackrel{~}{q}`$), wino ($`\stackrel{~}{w}^\pm `$), Higgsino ($`\stackrel{~}{H}_{1,2}^\pm `$) (or mixing of the latter ones, charginos, $`\chi _{1,2}^\pm `$), photino ($`\stackrel{~}{\gamma }`$), zino ($`\stackrel{~}{z}^0`$), Higgsino ($`\stackrel{~}{H}_{1,2}^0`$) (or their mixed states, neutralinos, $`\chi _{1,2,3,4}^0`$ ), gluino ($`\stackrel{~}{g}`$) etc. It is commonly believed that these SUSY particles should have masses below 1 TeV. Experiments at the existing colliders have already put the lower mass limits as $`m_{\stackrel{~}{q}}>176`$ GeV and $`m_{\chi ^\pm }>99`$ GeV.
In this paper we study the associated production of the squarks and charginos at TeV energy $`\gamma `$p colliders with polarized beams. We have already discussed this process with unpolarized beams . Many other processes such as $`\gamma p\stackrel{~}{q}\stackrel{~}{q}^{}`$X, $`\gamma p\stackrel{~}{q}\stackrel{~}{g}`$X, $`\gamma p\stackrel{~}{q}\stackrel{~}{\gamma }(or\stackrel{~}{q}\stackrel{~}{z})`$X, have already been discussed .
## II polarized high energy gamma beam
A beam of laser photons ($`\omega _01.26`$ eV, for example) with high intensity, about $`10^{20}`$ photons per pulse, is Compton-backscattered off high energy electrons ($`E_e`$=250 GeV, for example) from a linear accelerator and turns into hard photons with a conversion coefficient close to unity. The energy of the backscattered photons, $`E_\gamma `$, is restricted by the kinematic condition $`y_{max}=0.83`$ (where $`y=E_\gamma /E_e`$) in order to get rid of background effects, in particular $`e^+e^{}`$ pair production in the collision of a laser photon with a backscattered photon in the conversion region.
The details of the Compton kinematics and calculations of the cross section can be found in ref . The energy spectrum of the high energy real (backscattered) photons, $`f_{\gamma /e}(y)`$, is given by
$`f_{\gamma /e}(y)={\displaystyle \frac{1}{D(\kappa )}}[1y+{\displaystyle \frac{1}{1y}}4r(1r)\lambda _e\lambda _0r\kappa (2r1)(2y)]`$ (1)
where $`\kappa =4E_e\omega _0/m_e^2`$ and $`r=y/\kappa (1y)`$. Here $`\lambda _0`$ and $`\lambda _e`$ are the laser photon and the electron helicities respectively, and $`D(\kappa )`$ is
$`D(\kappa )`$ $`=`$ $`(1{\displaystyle \frac{4}{\kappa }}{\displaystyle \frac{8}{\kappa ^2}})ln(1+\kappa )+{\displaystyle \frac{1}{2}}+{\displaystyle \frac{8}{\kappa }}{\displaystyle \frac{1}{2(1+\kappa )^2}}`$ (2)
$`+`$ $`\lambda _e\lambda _0[(1+{\displaystyle \frac{2}{\kappa }})ln(1+\kappa ){\displaystyle \frac{5}{2}}+{\displaystyle \frac{1}{1+\kappa }}{\displaystyle \frac{1}{2(1+\kappa )^2}}]`$ (3)
In our numerical calculations, we assume $`E_e\omega _0=0.3`$ $`MeV^2`$ or equivalently $`\kappa =4.8`$ which corresponds to the optimum value of $`y_{max}=0.83`$, as mentioned above.
The energy spectrum, $`f_{\gamma /e}(y)`$, does essentially depend on the value $`\lambda _e\lambda _0`$. In the case of opposite helicities ($`\lambda _e\lambda _0=1`$) the spectrum has a very sharp peak at the high energy part of the photons. This allows us to get a highly monochromatic high energy gamma beam by eliminating low energy part of the spectrum . On the contrary, for the same helicities ($`\lambda _e\lambda _0=+1`$) the spectrum is flat.
The average degree of linear polarization of the photon is proportional to the degree of linear polarization of the laser. In our calculations, we assume that the degree of linear polarization of the laser is zero so that the final photons have only the degree of circular polarization $`(\lambda (y)=<\xi _2>0`$ and $`<\xi _1>`$$`=<\xi _3>=0`$ $`)`$. The circular polarization of the backscattered photon is given as follows
$`<\xi _2>=\lambda (y)={\displaystyle \frac{(12r)(\frac{1}{1y}+1y)\lambda _0+\lambda _er\kappa (1+(12r)^2(1y))}{\frac{1}{1y}+1y4r(1r)\lambda _0\lambda _er\kappa (2r1)(2y)}}`$ (4)
For the same initial polarizations ($`\lambda _0\lambda _e=+1`$) , it is seen that $`\lambda (y)+1`$, as nearly independent of $`y`$; while for the case of the opposite polarizations ($`\lambda _0\lambda _e=1`$), the curve $`\lambda (y)`$ smoothly changes from $`1`$ to $`+1`$ as $`y`$ increases from zero to $`0.83`$ .
## III polarized cross-sections for the reaction
The subprocess contributing to our physical process $`\gamma p\stackrel{~}{w}\stackrel{~}{q}`$X is $`\gamma q\stackrel{~}{w}\stackrel{~}{q}`$. The invariant amplitude for the specific subprocess $`\gamma u\stackrel{~}{w}^+\stackrel{~}{d}`$ is the sum of the three terms corresponding to the s-channel $`u`$ quark exchange, the t-channel $`\stackrel{~}{w}`$ wino exchange and the u-channel $`\stackrel{~}{d}`$ squark exchange interactions:
$`_a`$ $`=`$ $`{\displaystyle \frac{iee_qg}{2\widehat{s}}}\overline{u}(p^{})(1\gamma _5)(\overline{)}p+\overline{)}k)\overline{)}ϵu(p)`$ (5)
$`_b`$ $`=`$ $`{\displaystyle \frac{iege_{\stackrel{~}{w}}}{2(\widehat{t}m_{\stackrel{~}{w}}^2)}}\overline{u}(p^{})\overline{)}ϵ(\overline{)}p\overline{)}k+m_{\stackrel{~}{w}})(1\gamma _5)u(p)`$ (6)
$`_c`$ $`=`$ $`{\displaystyle \frac{iee_{\stackrel{~}{q}}g}{2(\widehat{u}m_{\stackrel{~}{q}}^2)}}\overline{u}(p^{})(1\gamma _5)u(p)(pp^{}+k^{}).ϵ`$ (7)
where $`e_q`$, $`e_{\stackrel{~}{q}}`$ and $`e_{\stackrel{~}{w}}`$ are the quark, squark and wino charges, and $`g=e/sin\theta _w`$ is the weak coupling constant. Note that we ignore the quark masses.
Since we are interested in the polarized cross-section, we use the following density matrices for the initial photon and quark:
$`\rho ^{(\gamma )}`$ $`=`$ $`{\displaystyle \frac{1}{2}}(1+\stackrel{}{\xi }.\stackrel{}{\sigma })`$ (8)
$`\rho ^{(q)}`$ $`=`$ $`u\overline{u}=\overline{)}p[1+\gamma _5(\lambda _q+\stackrel{}{\xi }_{}.\stackrel{}{\gamma }_{})]`$ (9)
where $`\xi _1`$, $`\xi _2`$ and $`\xi _3`$ are Stokes parameters. We take into account only circular polarization for the photon which is defined by $`\xi _2`$, as has been already mentioned in the previous section. $`\lambda _q`$ stands for the helicity of the parton-quark that is $`+1(1)`$ for the spin directions parallel (anti-parallel) to its momentum. The last term in the quark density matrix does not contribute, because after the integration over the azimuthal angle it vanishes.
One can easily obtain the differential cross section for the subprocess $`\gamma u\stackrel{~}{w}^+\stackrel{~}{d}`$ as follows
$`{\displaystyle \frac{d\widehat{\sigma }}{d\widehat{t}}}={\displaystyle \frac{\pi \alpha ^2}{\widehat{s}^2sin^2\theta _w}}(1+\lambda _q)[{\displaystyle \frac{d\widehat{\sigma _0}}{d\widehat{t}}}+\lambda (y){\displaystyle \frac{d\widehat{\sigma }_1}{d\widehat{t}}}].`$ (10)
Performing the $`d\widehat{t}`$ integration from $`t_{min}`$ to $`t^{max}`$ which are given by
$`t_{min}^{max}={\displaystyle \frac{1}{2}}(m_{\stackrel{~}{w}}^2+m_{\stackrel{~}{q}}^2\widehat{s})[1\sqrt{14m_{\stackrel{~}{w}}^2m_{\stackrel{~}{q}}^2/(m_{\stackrel{~}{w}}^2+m_{\stackrel{~}{q}}^2\widehat{s})^2}]`$ (11)
we immediately get the total cross section as
$`\widehat{\sigma }(m_{\stackrel{~}{w}},m_{\stackrel{~}{q}},\widehat{s},\lambda (y))={\displaystyle \frac{\pi \alpha ^2}{\widehat{s}^2sin^2\theta _w}}(1+\lambda _q)[\widehat{\sigma _0}(m_{\stackrel{~}{w}},m_{\stackrel{~}{q}},\widehat{s})+\lambda (y)\widehat{\sigma _1}(m_{\stackrel{~}{w}},m_{\stackrel{~}{q}},\widehat{s})].`$ (12)
Note that the cross sections ( Eqs(6) and (8)) are zero for $`\lambda _q=1`$ because of the fact that we ignore the quark mass. Integrating the subprocess cross-section $`\widehat{\sigma }`$ over the quark and photon distributions we obtain the total cross-section for the physical process $`\gamma p\stackrel{~}{w}\stackrel{~}{q}X`$ (the new variables are defined by $`\widehat{s}s_{\gamma q}=xys`$, $`xy=\tau `$ and $`ss_{ep}`$):
$`\sigma ={\displaystyle _{(m_{\stackrel{~}{w}}+m_{\stackrel{~}{q}})^2/s}^{0.83}}𝑑\tau {\displaystyle _{\tau /0.83}^1}{\displaystyle \frac{dx}{x}}f_{\gamma /e}(\tau /x)f_q(x)\widehat{\sigma }(m_{\stackrel{~}{w}},m_{\stackrel{~}{q}},\widehat{s},\lambda (\tau /x))`$ (13)
where the photon distribution function, $`f_{\gamma /e}(y)`$, is actually the normalized differential cross-section of the Compton backscattering, Eq.(1) ; $`f_q(x)`$ is the distribution of quarks inside the proton. We set $`\lambda _q=+1`$ and $`f_q(x)u^+(x)=\frac{1}{2}(u_{unp}+\mathrm{\Delta }u_{pol})`$ for the $`u`$-type valence quark distribution. In our numerical calculations, we use the distribution functions given in Ref. and Ref. for the unpolarized and polarized up-quarks, respectively:
$`u_{unp}(x)=2.751x^{0.412}(1x)^{2.69}`$ (14)
$`\mathrm{\Delta }u_{pol}(x)=2.132x^{0.2}(1x)^{2.40}`$ (15)
Performing the integrations in Eq.(9) numerically we obtain the total cross-section for the associated wino-squark production. We plot the dependence of the total cross-sections on the masses of the SUSY particles for various proposed $`\gamma p`$ colliders in Figs. 1(a-c) for $`\lambda _0\lambda _e=+1`$ and $`1`$. By taking 100 events per running year as observation limit for a SUSY particle, one can easily find the upper discovery mass limits from these figures using the luminosities of the proposed $`\gamma p`$ colliders given in Table 1. These discovery limits are tabulated in the same table.
## IV asymmetry
It may be more interesting to use a polarization asymmetry in determining the masses of SUSY particles. Such an asymmetry can be defined with respect to the product of the polarizations of the laser photon and the electron as follows
$`A={\displaystyle \frac{\sigma _{}\sigma _+}{\sigma _{}+\sigma _+}}`$ (16)
where $`\sigma _+`$ and $`\sigma _{}`$ are the polarized total cross-sections given in Figs. 1(a-c). The results of the polarization asymmetry are shown in Figs. 2(a-c) for three colliders.
## V conclusion
If one compares the curves $`\sigma _+`$ and $`\sigma _{}`$ in Figs. 1(a-c) for each collider one sees that the polarized cross-sections for different polarization do not differ much from each other and also from the unpolarized ones. Therefore, the discovery mass limits for SUSY partners obtained with polarized beams are nearly equal to those obtained with unpolarized beams. But the polarization asymmetry is highly sensitive to the wino and the squark masses and as high as 0.4 for all cases. Especially in the case of $`m_{\stackrel{~}{q}}=250`$ $`GeV`$, the asymmetry parameter A is around 0.6 for the higher wino masses.
The signature of the associated $`\stackrel{~}{w}^+\stackrel{~}{d}`$ production will depend on the mass spectrum of SUSY particles. It is generally assumed that the photino and sneutrino are the lightest SUSY particles and that the hierarchy of the squark masses is similar to that of quarks. With these assumption we have the following decays for the case $`m_{\stackrel{~}{q}}=m_{\stackrel{~}{w}}`$:
$`\stackrel{~}{d}d\stackrel{~}{\gamma }`$, $`d\stackrel{~}{g}`$ and $`\stackrel{~}{w}l^+\stackrel{~}{\nu }`$, $`\nu \stackrel{~}{l}^+`$, $`W^+\stackrel{~}{\gamma }`$
By taking into account the further decays $`\stackrel{~}{l}^+l\stackrel{~}{\gamma }`$ and $`W^+l^+\nu ,q\overline{q}`$ we arrive at the ultimate final states as
$`l^++n`$ $`jets(n=1,3,5)+`$ large missing energy and missing $`P_T`$
The main background for the final state $`l^++jet+P_T^{miss}`$ will come from the process $`\gamma qWql^+\nu q`$; but this background may be reduced, in principle,by the cut $`P_T^{miss}>45`$ $`GeV`$ if $`m_{\stackrel{~}{w}}m_W`$. |
warning/0002/math0002203.html | ar5iv | text | # On the isotriviality of families of projective manifolds over curves
## 1. Hodge bundles and the Kodaira Spencer map
Let $`Y`$ be a complex manifold and let $`S`$ be a normal crossing divisor. A variation $`𝕍_0`$ of polarized Hodge structures of weight $`k`$ on $`Y_0=YS`$ gives rise to
$$E_0=\mathrm{gr}_F(𝕍𝒪_{Y_0})=\underset{p+q=k}{}E_0^{p,q},$$
together with a Higgs structure $`\theta _0=\theta _{p,q}:E_0E_0\mathrm{\Omega }_{Y_0}^1`$.
###### Lemma 1.1.
If $`𝒩E_0^{p,q}`$ is a sub-bundle with $`\theta _{p,q}(𝒩)=0,`$ then the curvature of the restricted Hodge metric on $`𝒩`$ is negative semidefinite.
In fact, the negativity of the curvature of the restricted Hodge metric on $`det(𝒩)`$, which will be the only case used in this note, as well as the next lemma 1.2, follow from and can also be found in , Lemma 1. For the convenience of the reader we sketch the proof.
###### Proof.
Let $`\mathrm{\Theta }(E_0,h)`$ denote the curvature form of the Hodge metric $`h`$ on $`E_0.`$ Then by , chapter II, we have
$$\mathrm{\Theta }(E_0,h)+\theta \overline{\theta }_h+\overline{\theta }_h\theta =0,$$
where $`\overline{\theta }_h`$ is the complex conjugation of $`\theta `$ with respect to $`h.`$ $`h`$ restricts to a metric $`h|_𝒩,`$ and induces a $`𝒞^{\mathrm{}}`$decomposition $`E_0=𝒩𝒩^{}.`$ One obtains
$$\mathrm{\Theta }(𝒩,h)=\mathrm{\Theta }(E_0,h)|_𝒩+\overline{A}_hA=\theta \overline{\theta }_h|_𝒩\overline{\theta }_h\theta |_𝒩+\overline{A}_hA,$$
where $`AA^{1,0}(\text{Hom}(𝒩,𝒩^{}))`$ is the second fundamental form of the sub-bundle $`𝒩E_0,`$ and $`\overline{A}_h`$ is its complex conjugate with respect to $`h.`$
Since $`\theta (𝒩)=0,`$ we have $`\overline{\theta }_h\theta |_𝒩=0,`$ hence
$$\mathrm{\Theta }(𝒩,h)=\theta \overline{\theta }_h|_𝒩+\overline{A}_hA.$$
$`\theta \overline{\theta }_h`$ is positive semidefinite and $`\overline{A}_hA`$ is negative semidefinite, so $`\mathrm{\Theta }(𝒩,h)`$ is negative semidefinite. ∎
Suppose the local monodromy of $`𝕍_0`$ around the components of $`S`$ are unipotent and let $`𝒱`$ be the Deligne extension of $`𝕍_0𝒪_{Y_0}`$. By the F-filtration extends to a filtration of $`𝒱`$ by subbundles, hence there exists a canonical extension $`E`$ of $`E_0`$ to $`Y,`$ and $`\theta _0`$ extends to
$$\theta =\underset{p+q=k}{}\theta _{p,q}:E=\underset{p+q=k}{}E^{p,q}E\mathrm{\Omega }_Y^1(\mathrm{log}S)=\underset{p+q=k}{}E^{p,q}\mathrm{\Omega }_Y^1(\mathrm{log}S).$$
###### Lemma 1.2.
Keeping the assumptions made above, suppose $`Y`$ is a smooth projective curve. If $`𝒩E^{p,q}`$ is a sub-bundle with $`\theta _{p,q}(𝒩)=0,`$ then $`\mathrm{deg}(𝒩)0.`$
###### Proof.
Let $`𝒩^{}`$ be the dual of $`𝒩`$. We have the projection $`E_{}^{p,q}{}_{}{}^{}𝒩^{}.`$ Note that $`E_{}^{p,q}{}_{}{}^{}=E_{}^{}{}_{}{}^{q,p}`$ as a Hodge bundle of the system of Hodge bundles corresponding to the dual variation of Hodge structures $`𝕍_0^{}`$. The monodromy of $`𝕍_0^{}`$ around $`S`$ is again unipotent. We have the projection
$$F^qE^{q,p}𝒩^{},$$
where $`F^q`$ is the q-th subbundle in the extended Hodge filtration of $`𝒱^{}.`$
This presentation of $`𝒩^{}`$ as a quotient of a subbundle of a variation of Hodge structures allows to apply , 5.20. So the Chern forms of the induced Hodge metric on $`(𝒩|_{Y_0})^{}`$ represent the corresponding Chern classes of $`𝒩^{}.`$ From 1.1 we get in particular $`\mathrm{deg}(𝒩^{})0`$, and hence $`\mathrm{deg}(𝒩)0.`$
Let $`g:ZY`$ be a surjective morphism between a projective $`n`$-dimensional manifold $`Z`$ and a non-singular curve $`Y`$, both defined over the complex numbers. Let $`SY`$ be a divisor such that $`g`$ is smooth outside of $`\mathrm{\Pi }=g^1(S)`$. We will assume $`\mathrm{\Pi }`$ to be a normal crossing divisor. The smooth projective morphism
$$g_0:Z_0=Z\mathrm{\Pi }YS$$
obtained by restricting $`g`$ gives rise to variations of Hodge structures $`𝕍_0=R^kg_{0}^{}{}_{}{}^{}_{Z_0}`$. As explained in , p. 423, the primitive decomposition of $`𝕍_0`$ allows to define a polarization on $`𝕍_0`$. If the fibres of $`g`$ are connected and if $`g`$ is semistable, i.e. if $`\mathrm{\Pi }`$ is reduced, the local monodromies around points in $`S`$ are unipotent.
Using the notations introduced above we find
$$E^{p,q}=R^qg_{}\mathrm{\Omega }_{Z/Y}^p(\mathrm{log}\mathrm{\Pi }).$$
$`\theta _{p,q}:E^{p,q}E^{p1,q+1}\mathrm{\Omega }_Y^1(\mathrm{log}S)`$, which we will call the Kodaira Spencer map, is the edge-morphism induced by the tautological exact sequence
(1.2.1)
$$0g^{}\mathrm{\Omega }_Y^1(\mathrm{log}S)\mathrm{\Omega }_{Z/Y}^{p1}(\mathrm{log}\mathrm{\Pi })\mathrm{\Omega }_Z^p(\mathrm{log}\mathrm{\Pi })\mathrm{\Omega }_{Z/Y}^p(\mathrm{log}\mathrm{\Pi })0.$$
It is given by the cup product with the Kodaira Spencer class, induced by $`g`$.
###### Proposition 1.3.
Let $`𝒩`$ be an invertible subsheaf of $`E^{p,q}`$ with $`\theta _{p,q}(𝒩)=0`$. Then $`\mathrm{deg}(𝒩)0.`$
###### Proof.
If the fibres of $`g`$ are connected and if $`g`$ is semistable, this is nothing but 1.2.
In general, let $`L`$ be a finite extension of the function field $`(Y)`$, containing the Galois hull of the algebraic closure of $`(Y)`$ in $`(Z)`$, and let $`Y^{}Y`$ be the normalization of $`Y`$ in $`L`$. Consider the normalization $`\stackrel{~}{Z}`$ of $`Z\times _YY^{}`$, a desingularization $`\phi ^{}:Z^{}\stackrel{~}{Z}`$ and the induced morphisms
| $`Z^{}`$ | $`\stackrel{\phi ^{}}{}`$ | $`\stackrel{~}{Z}`$ | $`\stackrel{\stackrel{~}{\phi }}{}`$ | $`Z\times _YY^{}`$ | $`\stackrel{p_1}{}`$ | $`Z`$ |
| --- | --- | --- | --- | --- | --- | --- |
| | | $`\stackrel{~}{g}`$ | | | | |
| | | $`Y^{}`$ | $`\stackrel{\psi }{}`$ | $`Y`$ |
$`\phi =\stackrel{~}{\phi }\phi ^{}`$ and $`\psi ^{}=p_1\phi `$. We will enlarge $`S`$ such that $`Y^{}Y`$ is étale over $`YS`$, hence for $`S^{}=\phi ^{}S`$
$$\psi ^{}\mathrm{\Omega }_Y^1(\mathrm{log}S)=\mathrm{\Omega }_Y^{}^1(\mathrm{log}S^{}).$$
If one chooses $`L`$ large enough, $`Z^{}`$ will be the disjoint union of semistable families over $`Y^{}`$, hence $`\mathrm{\Pi }^{}=g_{}^{}{}_{}{}^{}S^{}`$ is a reduced normal crossing divisor, and $`\phi |_{g_{}^{}{}_{}{}^{1}(Y^{}S^{})}`$ is an isomorphism. By the generalized Hurwitz formula , 3.20,
$$\psi _{}^{}{}_{}{}^{}\mathrm{\Omega }_Z^p(\mathrm{log}\mathrm{\Pi })\mathrm{\Omega }_Z^{}^1(\mathrm{log}\mathrm{\Pi }^{}),$$
and by , Lemme 1.2,
$$R^q\phi _{}^{}\mathrm{\Omega }_Z^{}^p(\mathrm{log}\mathrm{\Pi }^{})=\{\begin{array}{cc}\stackrel{~}{\phi }^{}p_1^{}\mathrm{\Omega }_Z^p(\mathrm{log}\mathrm{\Pi })\hfill & \text{for }q=0\hfill \\ 0\hfill & \text{for }q>0.\hfill \end{array}$$
The exact sequence (1.2.1), for $`Z^{}`$ instead of $`Z`$, and induction on $`p`$ allow to show the same for the relative differential forms, i.e.
$$R^q\phi _{}^{}\mathrm{\Omega }_{Z^{}/Y^{}}^p(\mathrm{log}\mathrm{\Pi }^{})=\{\begin{array}{cc}\stackrel{~}{\phi }^{}p_1^{}\mathrm{\Omega }_{Z/Y}^p(\mathrm{log}\mathrm{\Pi })\hfill & \text{for }q=0\hfill \\ 0\hfill & \text{for }q>0.\hfill \end{array}$$
Hence the pullback of the exact sequence (1.2.1) to $`\stackrel{~}{Z}`$ is isomorphic to
(1.3.1)
$$\begin{array}{c}0\phi _{}^{}{}_{}{}^{}(g_{}^{}{}_{}{}^{}\mathrm{\Omega }_Y^{}^1(\mathrm{log}S^{})\mathrm{\Omega }_{Z^{}/Y^{}}^{p1}(\mathrm{log}\mathrm{\Pi }^{}))\hfill \\ \hfill \phi _{}^{}{}_{}{}^{}\mathrm{\Omega }_Z^{}^p(\mathrm{log}\mathrm{\Pi }^{})\phi _{}^{}{}_{}{}^{}\mathrm{\Omega }_{Z^{}/Y^{}}^p(\mathrm{log}\mathrm{\Pi }^{})0.\end{array}$$
Writing
$$E_{}^{}{}_{}{}^{p,q}=R^qg_{}^{}\mathrm{\Omega }_{Z^{}/Y^{}}^p(\mathrm{log}\mathrm{\Pi }^{}),$$
and $`\theta _{p,q}^{}`$ for the edge-morphism, we find
$$E_{}^{}{}_{}{}^{p,q}=R^qp_{2}^{}{}_{}{}^{}(\stackrel{~}{\phi }_{}𝒪_{\stackrel{~}{Z}}p_1^{}\mathrm{\Omega }_{Z/Y}^p(\mathrm{log}\mathrm{\Pi })).$$
Moreover, the inclusion $`𝒪_{Z\times _YY^{}}\stackrel{~}{\phi }_{}𝒪_{\stackrel{~}{Z}}`$ and flat base change give an inclusion
$$\psi ^{}E^{p,q}=R^qp_{2}^{}{}_{}{}^{}p_1^{}\mathrm{\Omega }_Z^p(\mathrm{log}\mathrm{\Pi }))E^{}^{p,q}$$
and the diagram
| $`\psi ^{}E^{p,q}`$ | $`\stackrel{\psi ^{}\theta _{p,q}}{}`$ | $`\psi ^{}E^{p1,q+1}\mathrm{\Omega }_Y^1(\mathrm{log}S)`$ |
| --- | --- | --- |
| $``$ | | $``$ |
| $`E_{}^{}{}_{}{}^{p,q}`$ | $`\stackrel{\theta _{p,q}^{}}{}`$ | $`E_{}^{}{}_{}{}^{p1,q+1}\mathrm{\Omega }_Y^{}^1(\mathrm{log}S^{})`$ |
commutes. In particular, if $`𝒩`$ lies in the kernel of $`\theta _{p,q}`$, the sheaf $`\psi ^{}𝒩`$ lies in the kernel of $`\theta _{p,q}^{}`$. Since we already know 1.3 for semistable morphisms with connected fibres, we find $`\mathrm{deg}(𝒩)0`$. ∎
## 2. Positivity of direct image sheaves
As in or a second ingredient in the proof of 0.1 and 0.3 will be explicit bounds for the positivity of certain direct image sheaves.
###### Definition 2.1.
Let $``$ be a locally free sheaf and $`𝒜`$ an invertible sheaf on the curve $`Y`$.
1. $``$ is nef if for all finite morphisms $`\pi :ZY`$ and all invertible quotients $``$ of $`\pi ^{}`$ the degree of $``$ is non-negative.
2. For $`\alpha ,\beta \{0\}`$ we write
$$\frac{\alpha }{\beta }𝒜$$
if $`S^\beta ()𝒜^\alpha `$ is nef. This is welldefined, since obviously the latter holds true if and only if
$$S^{\beta \mu }()𝒜^{\alpha \mu }$$
is nef, for some $`\mu >0`$.
Nef locally free sheaves on curves, have already been used in to study the height of points of curves over function fields. In , §2, and in the higher dimensional birational classification theory, one needs positive coherent torsionfree sheaves over higher dimensional manifolds, and there one often considers weakly positive sheaves, instead of nef sheaves. In the one-dimensional case, both notions coincide, and all the properties of weakly positive sheaves, listed in or carry over to nef sheaves on curves. Let us recall one property:
###### Lemma 2.2.
Given $`d`$, assume that for all $`\mu `$, sufficiently large and divisible, there exists a covering $`\tau :Y^{}Y`$ of degree $`\mu `$ such that $`\tau ^{}`$ is nef, for one, hence for all invertible sheaves of degree $`d`$. Then $``$ is nef.
###### Proof.
Let $`\pi :ZY`$ and $``$ be as in 2.1, a), and let $`Z^{}`$ be a component of the normalization of $`Z\times _YY^{}`$. If
| $`Z^{}`$ | $`\stackrel{\tau ^{}}{}`$ | $`Z`$ |
| --- | --- | --- |
| $`\pi ^{}`$ | | $`\pi `$ |
| $`Y^{}`$ | $`\stackrel{\tau }{}`$ | $`Y`$ |
are the induced morphisms, then
$$\begin{array}{c}0\mathrm{deg}(\tau {}_{}{}^{}{}_{}{}^{}\pi {}_{}{}^{}{}_{}{}^{})=\mathrm{deg}(\tau ^{})\mathrm{deg}()+\mathrm{deg}(\pi ^{})d\hfill \\ \hfill \mu \mathrm{deg}()+\mathrm{deg}(\pi )d.\end{array}$$
This, for all $`\mu \{0\}`$, implies that $`\mathrm{deg}()0`$. ∎
Next recall the definition of the (algebraic) multiplier sheaves. We consider a surjective morphism $`f:XY`$, with connected general fibre $`F`$, where $`X`$ is an (n+1)-dimensional complex projective manifold, and $`Y`$ a non-singular projective curve. If $`\mathrm{\Gamma }`$ is an effective divisor on $`X`$,
$$\omega _{X/Y}\left\{\frac{\mathrm{\Gamma }}{N}\right\}=\tau _{}\left(\omega _{X^{}/Y}\left(\left[\frac{\mathrm{\Gamma }^{}}{N}\right]\right)\right)$$
where $`\tau :X^{}X`$ is any blowing up with $`\mathrm{\Gamma }^{}=\tau ^{}\mathrm{\Gamma }`$ a normal crossing divisor (see for example , 7.4, or , section 5.3).
Fujita’s positivity theorem (today an easy corollary of Kollár’s vanishing theorem) says that $`f_{}\omega _{X/Y}`$ is nef. A direct consequence is the following.
###### Lemma 2.3.
Let $`𝒩`$ be an invertible sheaf on $`X`$ and $`\mathrm{\Gamma }`$ an effective divisor. Assume that for some $`N>0`$ there exists a nef locally free sheaf $``$ on $`Y`$ and a surjection $`f^{}𝒩^N(\mathrm{\Gamma }).`$ Then
$$f_{}\left(𝒩\omega _{X/Y}\left\{\frac{\mathrm{\Gamma }}{N}\right\}\right)$$
is nef.
###### Proof.
Let $`pY`$ be a point. Then $`𝒪_Y(Np)`$ is ample, hence
$$𝒩^N(\mathrm{\Gamma })f^{}𝒪_Y(Np)$$
is semi-ample. By , 7.16, the sheaf
$$f_{}\left(𝒩\omega _{X/Y}\left\{\frac{\mathrm{\Gamma }}{N}\right\}\right)𝒪_Y(p)$$
is nef. Since the same holds true over all $`Y^{}`$, finite over $`Y`$ and unramified in $`S`$, one obtains 2.3 from 2.2
As an application of 2.3 one obtains, as explained in ,
###### Lemma 2.4.
1. $`f_{}\omega _{X/Y}^\nu `$ is nef, for all $`\nu 0`$.
2. If $`\lambda _\nu =det(f_{}\omega _{X/Y}^\nu )`$ is ample, for some $`\nu >1`$, then there exists a positive rational number $`\eta `$ with $`f_{}\omega _{X/Y}^\nu \eta \lambda _\nu `$.
As in , we will need an explicit bound for the rational number $`\eta `$ in 2.4, ii). To this aim recall the following definition, used in , , § 7 and , section 5.3.
###### Definition 2.5.
Let $``$ be an invertible sheaf on $`F`$ with $`H^0(F,)0`$, and let $`\mathrm{\Gamma }`$ be an effective divisor. Then
$$e(\mathrm{\Gamma })=\mathrm{Min}\left\{N\{0\};\omega _F\left\{\frac{\mathrm{\Gamma }}{N}\right\}=\omega _F\right\}\text{ and}$$
$$e()=\mathrm{Max}\left\{e(\mathrm{\Gamma });\mathrm{\Gamma }\text{the zero set of}\sigma H^0(F,)\{0\}\right\}.$$
###### Notations 2.6.
For $`f:XY`$, we choose $`\nu >1`$ with $`f_{}\omega _{X/Y}^\nu =0`$, and a blowing up $`\tau :X^{}X`$ such that the fibres of $`f^{}=f\tau `$ are normal crossing divisors, such that
$$=\mathrm{Im}(f{}_{}{}^{}{}_{}{}^{}f_{}^{}\omega _{X^{}/Y}^\nu =f{}_{}{}^{}{}_{}{}^{}\omega _{X^{}/Y}^\nu )$$
is invertible and $`\omega _{X^{}/Y}^\nu =(B)`$, for a normal crossing divisor $`B`$. Let $`F^{}`$ be a general fibre of $`f^{}`$. We define:
$`e=e(|_F^{})`$
$`r=\mathrm{rank}(f_{}^{}\omega _{X^{}/Y}^\nu )=\mathrm{rank}()`$
$`\lambda =det(f_{}^{}\omega _{X^{}/Y}^\nu )=det()`$.
###### Proposition 2.7.
If $`\lambda `$ is ample, $`f_{}\omega _{X/Y}^\nu \frac{1}{re}\lambda `$.
###### Proof.
If $`e=1`$, the sheaf $`|_F^{}`$ is trivial, hence $`f_{}\omega _{X/Y}^\nu =\lambda `$ and 2.7 obviously holds true. Hence we will assume $`e2`$. For some $`\mu 0`$ there exists an effective divisor $`\mathrm{\Sigma }_1`$, disjoint from $`S`$ with $`\lambda ^\mu =𝒪_Y(\mathrm{\Sigma }_1)`$. By 2.2 and by flat base change, we are free to replace $`Y`$ by any $`Y^{}`$, finite over $`Y`$ and unramified over a neighborhood of $`S`$. Hence we are allowed to assume that $`\mathrm{\Sigma }_1=(\nu 1)e\mu \mathrm{\Sigma }`$ or that
$$\lambda =𝒪_Y((\nu 1)e\mathrm{\Sigma }).$$
Consider the $`r`$-fold fibre product
$$f^r:X^r=X^{}\times _YX^{}\mathrm{}\times _YX^{}Y.$$
$`f^r`$ is flat and Gorenstein and smooth over some open subscheme. Let
$$\pi :X^{(r)}X^r$$
be a desingularization such that the general fibre $`F^{(r)}`$ of $`f^{(r)}=f^r\pi `$ is isomorphic to $`F\times \mathrm{}\times F`$. For
$$=\pi ^{}\underset{i=1}{\overset{r}{}}pr_i^{}\pi ^{}\omega _{X^r/Y}^\nu ,$$
using flat base change, and the natural maps
$$𝒪_{X^r}\pi _{}𝒪_{X^{(r)}}\text{ and }\pi _{}\omega _{X^{(r)}}\omega _{X^{(r)}},$$
one finds
(2.7.1)
$$\stackrel{r}{}f_{}^{}\omega _{X^{}/Y}^\nu =\stackrel{r}{}=\stackrel{r}{}f_{}f_{}^{(r)}\text{ and }$$
(2.7.2)
$$f_{}^{(r)}((\pi ^{}\omega _{X^r/Y}^{\nu 1})\omega _{X^{(r)}/Y})f_{}^r\omega _{X^r/Y}^\nu =f_{}^{}\omega _{X^{}/Y}^\nu ,$$
and both are isomorphism over some open dense subset of $`Y`$. (2.7.1) induces a surjection
$$f^{(r)}\stackrel{r}{}=\pi ^{}\underset{i=1}{\overset{r}{}}pr_i^{}f_{}^{}{}_{}{}^{}\pi ^{}\underset{i=1}{\overset{r}{}}pr_i^{}=.$$
In particular, since $`\lambda ^r`$, the sheaf $`f_{}^{(r)}{}_{}{}^{}\lambda `$ is a subsheaf of $``$. Let $`\mathrm{\Gamma }`$ denote the divisor with $`(\mathrm{\Gamma })=f^{(r)}\lambda `$. For some divisor $`C`$, supported in fibres of $`f^{(r)}`$ one has
$$\pi ^{}\omega _{X^r/Y}=\omega _{X^{(r)}/Y}(C).$$
Blowing up $`X^{(r)}`$ with centers in fibres of $`f^{(r)}`$ we find a normal crossing divisor $`D`$ with
$$D\pi ^{}\underset{i=1}{\overset{r}{}}pr_i^{}B$$
and such that $`(D)=\omega _{X^{(r)}/Y}(C)^\nu `$. For
$$=e(\nu 1)D+\nu \mathrm{\Gamma }+e\nu (\nu 1)f^{(r)}(\mathrm{\Sigma })$$
one obtains
$$\omega _{X^{(r)}/Y}(C)^{e\nu (\nu 1)}()=^{e(\nu 1)}(\nu \mathrm{\Gamma })f^{(r)}\lambda ^\nu =^{e(\nu 1)\nu }.$$
Since we assumed $`e,\nu 2`$, the exponent of $``$ is non-negative. By 2.4, i), the sheaf $`^r`$ is nef and 2.3 implies that
$$=f_{}^{(r)}\left(\omega _{X^{(r)}/Y}(C)^{\nu 1}\omega _{X^{(r)}/Y}\left\{\frac{}{e\nu }\right\}\right)$$
is nef. $``$ is contained in
$$^{}=f_{}^{(r)}(\pi ^{}\omega _{X^r/Y}^{\nu 1}\omega _{X^{(r)}/Y})𝒪_Y((\nu 1)\mathrm{\Sigma }),$$
and using (2.7.2) ones finds
$$(\stackrel{r}{}f_{}^{}\omega _{X^{}/Y}^\nu )𝒪_Y((\nu 1)\mathrm{\Sigma }).$$
On the other hand, $``$ contains
$$^{\prime \prime }=f_{}^{(r)}\left(\omega _{X^{(r)}/Y}(C)^{\nu 1}\omega _{X^{(r)}/Y}\left\{\frac{eD}{e\nu }\right\}\right).$$
Over some sufficiently small open dense subset $`UY`$
$$^{\prime \prime }|_U=f_{}^{(r)}\left(\omega _{X^{(r)}/Y}^{\nu 1}\omega _{X^{(r)}/Y}\left\{\frac{\nu \mathrm{\Gamma }}{\nu e}\right\}𝒪_{X^{(r)}}(D)\right)|_U.$$
By definition $`e=e(|_F)`$, and from or , 5.21, one has
$$e=e(|_F)=e(\underset{i=1}{\overset{r}{}}pr_i^{}|_F)=e(|_{F^{(r)}}).$$
The semicontinuity of $`e`$ in , 5.14, implies that for $`U`$ small enough,
$$^{\prime \prime }|_U=f_{}^{(r)}(\omega _{X^{(r)}/Y}^\nu 𝒪_{X^{(r)}}(D))=^{}|_U.$$
Hence $`,^{}`$ and $`^{\prime \prime }`$ have the same rank and $`^{}`$ is nef. We obtain
$$\stackrel{r}{}f_{}^{}\omega _{X^{}/Y}^\nu 𝒪_Y((\nu 1)\mathrm{\Sigma })=\frac{1}{e}\lambda $$
or $`f_{}^{}\omega _{X^{}/Y}^\nu \frac{1}{re}\lambda `$, as claimed. ∎
## 3. Differential forms and cyclic coverings
Let $`Y`$ be a non-singular projective curve of genus $`g`$, and let $`SY`$ be a reduced divisor of degree $`s`$. We consider again an $`(n+1)`$-dimensional manifold $`X`$ and a surjective morphism $`f:XY`$ with connected general fibre $`F`$. For $`\mathrm{\Delta }=f^1(S)`$ we will assume that $`f_0=f|_{X\mathrm{\Delta }}`$ is smooth.
For some $`\nu >1`$, with $`f_{}\omega _{X/Y}^\nu 0`$, we choose, as in 2.6, a blowing up $`\tau :X^{}X`$ and a normal crossing divisor $`B`$. Hence writing $`f^{}=f\tau `$ and $`=\omega _{X^{}/Y}^\nu (B)`$ one has $`f_{}^{}=f_{}^{}\omega _{X^{}/Y}^\nu `$ and $`f{}_{}{}^{}{}_{}{}^{}f_{}^{}`$. Choose $`S^{}S`$, such that $`f^{}`$ is smooth outside of $`\mathrm{\Delta }^{}=f{}_{}{}^{}{}_{}{}^{1}(S^{})`$, and such that $`B(B\mathrm{\Delta }^{})`$ is a relative normal crossing divisor over $`YS^{}`$. Blowing up with centers in $`\mathrm{\Delta }^{}`$, we may also assume $`B+\mathrm{\Delta }^{}`$ to be a normal crossing divisor.
Let $`𝒜`$ be an ample invertible sheaf with
(3.0.1)
$$f_{}^{}\omega _{X^{}/Y}^\nu 𝒜^\nu =f_{}\omega _{X/Y}^\nu 𝒜^\nu \text{ ample.}$$
For $`N=\nu \mu `$ and $`\mu `$ sufficiently large, and for $`=\omega _{X^{}/Y}f{}_{}{}^{}{}_{}{}^{}𝒜_{}^{1}`$ the sheaf
$$^N(\mu B)=^\mu f{}_{}{}^{}{}_{}{}^{}𝒜_{}^{N}$$
is generated by global sections. For a general section $`\sigma `$ with zero divisor $`V(\sigma )`$, the divisor $`M=\mu B+V(\sigma )`$ as well as $`M+\mathrm{\Delta }^{}`$ are normal crossing divisors. Enlarging $`S^{}`$, if necessary, and replacing $`X^{}`$ by a blowing up with centers in $`\mathrm{\Delta }^{}`$, we may assume that all the fibres of
$$f\tau :X^{}\mathrm{\Delta }^{}YS^{}$$
intersect $`M`$ transversely.
As explained in , § 2, $`\sigma `$ defines a cyclic covering $`Z`$ of $`X^{}`$. In explicit terms, writing $`^{(i)}=^i([\frac{iM}{N}])`$, the covering is given by
$$\gamma :Z=\mathrm{𝐒𝐩𝐞𝐜}\left(\underset{i=0}{\overset{N1}{}}^{(i)}\right)X^{}.$$
Let us write $`g=f^{}\gamma `$ for the induced morphism from $`Z`$ to $`Y`$. The discriminant of $`\gamma `$ is $`M^{}=(MN[\frac{M}{N}])_{\mathrm{red}}`$, and defining
$$\mathrm{\Omega }_Z^p(\mathrm{log}\gamma ^1(\mathrm{\Delta }^{}+M^{}))\text{ and }\mathrm{\Omega }_{Z/Y}^p(\mathrm{log}\gamma ^1(\mathrm{\Delta }^{}+M^{}))$$
to be the reflexive hulls of the corresponding sheaves on the smooth locus of $`Z`$, both sheaves are locally free and, as explained in , §2, one has
(3.0.2)
$$\begin{array}{c}\hfill \mathrm{\Omega }_Z^p(\mathrm{log}\gamma ^1(\mathrm{\Delta }^{}+M^{}))=\gamma ^{}\mathrm{\Omega }_X^{}^p(\mathrm{log}(\mathrm{\Delta }^{}+M^{}))\\ \hfill \text{and }\mathrm{\Omega }_{Z/Y}^p(\mathrm{log}\gamma ^1(\mathrm{\Delta }^{}+M^{}))=\gamma ^{}\mathrm{\Omega }_{X^{}/Y}^p(\mathrm{log}(\mathrm{\Delta }^{}+M^{})).\end{array}$$
Hence $`\mathrm{\Omega }_X^{}^p(\mathrm{log}(\mathrm{\Delta }^{}+M^{}))^{(1)}`$ is a direct factor of $`\gamma _{}\mathrm{\Omega }_Z^p(\mathrm{log}\gamma ^1(\mathrm{\Delta }^{}+M^{}))`$. The image of the induced map
$$\gamma ^{}\mathrm{\Omega }_X^{}^p(\mathrm{log}\mathrm{\Delta }^{}+M^{})^{(1)}\mathrm{\Omega }_Z^p(\mathrm{log}\gamma ^1(\mathrm{\Delta }^{}+M^{}))$$
lies in the subsheaf $`\mathrm{\Omega }_Z^p(\mathrm{log}\gamma ^1(\mathrm{\Delta }^{}))\stackrel{~}{}`$, where $`()\stackrel{~}{}`$ denotes the reflexive hull. Again, similar inclusions hold true for the relative differential forms over $`Y`$.
Kawamata’s covering construction (see , 2.2) allows to choose a finite covering $`Z^{\prime \prime }Z`$ with $`Z^{\prime \prime }`$ non singular. Blowing up centers in fibres over $`Y`$, we obtain a non-singular variety $`Z^{}`$ and a generically finite map $`\eta :Z^{}Z`$, such that all the fibres of
$$g^{}=g\eta :Z^{}Y$$
are normal crossing divisors. By abuse of notations we will add some points to $`S^{}`$, hence some fibres to $`\mathrm{\Delta }^{}`$, and assume that $`g^{}`$ is smooth outside of $`\mathrm{\Pi }^{}=g{}_{}{}^{}{}_{}{}^{1}(S^{})`$ and that $`\eta `$ is finite outside of $`\mathrm{\Pi }^{}`$. Let $`\gamma ^{}:Z^{}X^{}`$ be the induced map. Since $`\gamma ^{}`$ factors through a non-singular variety, finite over $`X^{}`$, the higher direct images $`R^q\gamma _{}^{}𝒪_Z^{}`$ are zero for $`q>0`$.
Let $`M^{\prime \prime }`$ denote the proper transform of $`M^{}`$ in $`Z^{}`$. Since $`\gamma ^{}`$ is finite outside of $`\mathrm{\Pi }^{}`$ one obtains natural inclusions (see , 3.20, for example)
(3.0.3)
$$\begin{array}{c}\hfill \eta ^{}\mathrm{\Omega }_Z^p(\mathrm{log}\gamma ^1(\mathrm{\Delta }^{}+M^{}))\mathrm{\Omega }_Z^{}^p(\mathrm{log}(\mathrm{\Pi }^{}+M^{\prime \prime }))\text{ and}\\ \hfill \eta ^{}\mathrm{\Omega }_{Z/Y}^p(\mathrm{log}\gamma ^1(\mathrm{\Delta }^{}+M^{}))\mathrm{\Omega }_{Z^{}/Y}^p(\mathrm{log}(\mathrm{\Pi }^{}+M^{\prime \prime })).\end{array}$$
(3.0.2) and (3.0.3) together induce inclusions
(3.0.4)
$$\begin{array}{c}\hfill \gamma _{}^{}{}_{}{}^{}(\tau ^{}\mathrm{\Omega }_{X/}^p(\mathrm{log}\mathrm{\Delta }))^{(1)}\gamma _{}^{}{}_{}{}^{}\mathrm{\Omega }_{X^{}/}^p(\mathrm{log}(\mathrm{\Delta }^{}+M^{}))^{(1)}\\ \hfill \eta ^{}\mathrm{\Omega }_{Z/}^p(\mathrm{log}\gamma ^1(\mathrm{\Delta }^{}+M^{}))\mathrm{\Omega }_{Z^{}/}^p(\mathrm{log}(\mathrm{\Pi }^{}+M^{\prime \prime })),\end{array}$$
where $``$ stands for $`Y`$ or for $`\mathrm{Spec}`$, respectively.
Again, the image of composite of the injections in 3.0.4 must have trivial residues along the components of $`M^{\prime \prime }`$, hence one finds a natural map
$$\iota _{}:\gamma _{}^{}{}_{}{}^{}(\tau ^{}\mathrm{\Omega }_{X/}^p(\mathrm{log}\mathrm{\Delta }))^{(1)}\stackrel{}{}\mathrm{\Omega }_{Z^{}/}^p(\mathrm{log}\mathrm{\Pi }^{}).$$
Consider the tautological sequence
(3.0.5)
$$\begin{array}{c}\hfill 0f^{}\mathrm{\Omega }_Y^1(\mathrm{log}S)\mathrm{\Omega }_{X/Y}^{p1}(\mathrm{log}\mathrm{\Delta })\mathrm{\Omega }_X^p(\mathrm{log}\mathrm{\Delta })\mathrm{\Omega }_{X/Y}^p(\mathrm{log}\mathrm{\Delta })0.\end{array}$$
Pulling it back to $`X^{}`$ and tensorizing with $`^{(1)}`$ one obtains
(3.0.6)
$$\begin{array}{c}\hfill 0f_{}^{}{}_{}{}^{}\mathrm{\Omega }_Y^1(\mathrm{log}S)\tau ^{}\mathrm{\Omega }_{X/Y}^{p1}(\mathrm{log}\mathrm{\Delta })^{(1)}\\ \hfill \tau ^{}\mathrm{\Omega }_X^p(\mathrm{log}\mathrm{\Delta })^{(1)}\tau ^{}\mathrm{\Omega }_{X/Y}^p(\mathrm{log}\mathrm{\Delta })^{(1)}0.\end{array}$$
The inclusions $`\iota _Y`$, $`\iota _{\mathrm{Spec}}`$ and
$$\iota :\mathrm{\Omega }_Y^1(\mathrm{log}S)\mathrm{\Omega }_Y^1(\mathrm{log}S^{}).$$
induce a morphism from the pullback of this exact sequence to
(3.0.7)
$$0g{}_{}{}^{}{}_{}{}^{}\mathrm{\Omega }_{Y}^{1}(\mathrm{log}S^{})\mathrm{\Omega }_{Z^{}/Y}^{p1}(\mathrm{log}\mathrm{\Pi }^{})\mathrm{\Omega }_Z^{}^p(\mathrm{log}\mathrm{\Pi }^{})\mathrm{\Omega }_{Z^{}/Y}^p(\mathrm{log}\mathrm{\Pi }^{})0$$
Let us define
$$E^{p,q}=R^qg_{}^{}\mathrm{\Omega }_{Z^{}/Y}^p(\mathrm{log}\mathrm{\Pi }^{})$$
$$\text{and }F^{p,q}=R^qf_{}^{}((\tau ^{}\mathrm{\Omega }_{X/Y}^p(\mathrm{log}\mathrm{\Delta }))^{(1)}).$$
The inclusion $`\iota _Y`$ gives a map
$$R^qg_{}^{}(\gamma _{}^{}{}_{}{}^{}(\tau ^{}\mathrm{\Omega }_{X/Y}^p(\mathrm{log}\mathrm{\Delta }))^{(1)})R^qg_{}^{}\mathrm{\Omega }_{Z^{}/Y}^p(\mathrm{log}\mathrm{\Pi }^{}).$$
Since the first sheaf is isomorphic to
$$R^qf_{}^{}((\gamma _{}^{}{}_{}{}^{}\gamma _{}^{}{}_{}{}^{}𝒪_Z^{})(\tau ^{}\mathrm{\Omega }_{X/Y}^p(\mathrm{log}\mathrm{\Delta }))^{(1)})$$
we obtain thereby a morphism $`\rho _{p,q}:F^{p,q}E^{p,q}`$. Obviously
$$\rho _{n,0}:f_{}^{}((\tau ^{}\mathrm{\Omega }_{X/Y}^n(\mathrm{log}\mathrm{\Delta }))^{(1)})g_{}^{}\mathrm{\Omega }_{Z^{}/Y}^n(\mathrm{log}\mathrm{\Pi }^{})$$
is injective, and
$$\rho _{0,n}:R^nf_{}^{}(^{(1)})R^nf_{}^{}(\gamma _{}𝒪_Z)R^nf_{}^{}(\gamma _{}^{}𝒪_Z^{})=R^ng_{}^{}𝒪_Z^{}$$
gives $`F^{0,n}`$ as a direct factor of $`E^{0,n}`$. The edge-morphism
$$E^{p,q}\stackrel{\theta _{p,q}}{}E^{p1,q+1}\mathrm{\Omega }_Y^1(\mathrm{log}S^{})$$
of the exact sequence (3.0.7) is the Kodaira Spencer map, studied in §1. Since the pullback of (3.0.6) to $`Z^{}`$ is a subsequence of (3.0.7) $`\theta _{p,q}`$ commutes with the edge-morphism
$$\begin{array}{c}R^qg_{}^{}(\gamma _{}^{}{}_{}{}^{}(\tau ^{}\mathrm{\Omega }_{X/Y}^p(\mathrm{log}\mathrm{\Delta }))^{(1)})\hfill \\ \hfill R^{q+1}g_{}^{}(\gamma _{}^{}{}_{}{}^{}(\tau ^{}\mathrm{\Omega }_{X/Y}^{p1}(\mathrm{log}\mathrm{\Delta }))^{(1)})\mathrm{\Omega }_Y^1(\mathrm{log}S).\end{array}$$
So the edge-morphism of the exact sequence (3.0.6), denoted by
$$F^{p,q}\stackrel{\tau _{p,q}}{}F^{p1,q+1}\mathrm{\Omega }_Y^1(\mathrm{log}S),$$
is compatible with $`\theta _{p,q}`$.
###### Lemma 3.1.
Assume for an ample invertible sheaf $`𝒜`$, and for $`\nu >1`$ (3.0.1) holds true. Then, using the notations introduced above,
1. the Kodaira Spencer map for $`g^{}:Z^{}Y`$ and the edge-morphism of the exact sequence (3.0.6) induce a commutative diagram
| $`E^{p,q}`$ | $`\stackrel{\theta _{p,q}}{}`$ | $`E^{p1,q+1}\mathrm{\Omega }_Y^1(\mathrm{log}S^{})`$ |
| --- | --- | --- |
| $`\rho _{p,q}`$ | | $`\rho _{p1,q+1}\iota `$ |
| $`F^{p,q}`$ | $`\stackrel{\tau _{p,q}}{}`$ | $`F^{p1,q+1}\mathrm{\Omega }_Y^1(\mathrm{log}S).`$ |
2. $`\rho _{n,0}`$ is injective.
3. $`F^{0,n}`$ is a direct factor of $`E^{0,n}`$.
4. the sheaf $`(F^{0,n})^{}`$ is ample.
5. the sheaf $`F^{n,0}`$ is invertible of degree
$$\mathrm{deg}(F^{n,0})\mathrm{deg}(𝒜)\delta ,$$
where $`\delta `$ denotes the number of non-semistable fibres.
###### Proof.
It remains to verify iv) and v). Comparing the first Chern classes of the sheaves in (3.0.5) one finds $`\mathrm{\Omega }_{X/Y}^n(\mathrm{log}\mathrm{\Delta })=\omega _{X/Y}(\mathrm{\Delta }_{\mathrm{red}}\mathrm{\Delta })`$. Hence for some invertible sheaf $``$ of degree $`\delta `$ one has an inclusion
$$\mathrm{\Omega }_{X/Y}^n(\mathrm{log}\mathrm{\Delta })\omega _{X/Y}f^{}^1.$$
Recall that $`=\omega _{X^{}/Y}f{}_{}{}^{}{}_{}{}^{}𝒜_{}^{1}`$, and that
$$^{(1)}=^1\left(\left[\frac{M}{N}\right]\right)=^1\left(\left[\frac{B}{\nu }\right]\right),$$
where $`B`$ is the relative fix locus of $`\omega _{X^{}/Y}^\nu `$. In particular, if $`E`$ denotes the effective divisor with $`\omega _{X^{}/Y}(E)=\tau ^{}\omega _{X/Y}`$, the divisor $`B`$ is larger than $`\nu E`$.
In order to prove that $`𝒜^1`$ is a subsheaf of $`F^{n,0}`$, and that the latter is invertible, we just have to show that
$$f_{}^{}(\tau ^{}\omega _{X/Y}^{(1)})=f_{}^{}\left((f{}_{}{}^{}{}_{}{}^{}𝒜)𝒪_X^{}\left(\left[\frac{B}{\nu }\right]E\right)\right)$$
is isomorphic to $`𝒜`$, or that $`𝒪_Yf_{}^{}𝒪_X^{}\left(\left[\frac{B}{\nu }\right]E\right)`$. This is obvious since $`0\left[\frac{B}{\nu }\right]EB`$ and since $`B`$ is the relative fix locus of an invertible sheaf.
1.3 says in particular, that $`E^{0,n}=\mathrm{Ker}(\tau _{0,n})`$ has no invertible subsheaf of positive degree. Using iii) one obtains the same for $`F^{0,n}=R^nf_{}^{}{}_{}{}^{}^{(1)}`$ and $`(F^{0,n})^{}`$ is nef. Serre duality and the projection formula imply
$$(F^{0,n})^{}=𝒜^1f_{}^{}{}_{}{}^{}(\omega _{X^{}/Y}^2𝒪_X^{}(\left[\frac{B}{\nu }\right])).$$
To prove the ampleness, claimed in iv), choose some $`\eta >0`$ such that
$$S^\eta (f_{}\omega _{X/Y}^\nu 𝒜^\nu )𝒜^1$$
is ample and consider a finite covering $`\phi :Y^{}Y`$ of degree $`\nu \eta `$, étale over a neighborhood of the discriminant divisor $`S^{}`$. Then $`\phi ^{}𝒜=𝒜_{}^{}{}_{}{}^{\nu \eta }`$, for some ample invertible sheaf $`𝒜^{}`$ on $`Y^{}`$, and both, $`f_{}\omega _{X/Y}^\nu `$ and $`F^{0,n}`$ are compatible with pullbacks.
Replacing $`Y`$ by $`Y^{}`$, we may assume thereby that $`𝒜=𝒜_{}^{}{}_{}{}^{\nu \eta }`$, for some $`𝒜^{}`$, and that $`f_{}\omega _{X/Y}^\nu 𝒜^\nu 𝒜_{}^{}{}_{}{}^{\nu }`$ is ample. Repeating the argument used above, for $`𝒜𝒜^{}`$ instead of $`𝒜`$, one finds
$$𝒜^1𝒜_{}^{}{}_{}{}^{1}f_{}^{}{}_{}{}^{}(\omega _{X^{}/Y}^2𝒪_X^{}(\left[\frac{B}{\nu }\right]))$$
to be nef, hence $`𝒜^1f_{}^{}{}_{}{}^{}(\omega _{X^{}/Y}^2𝒪_X^{}(\left[\frac{B}{\nu }\right]))`$ to be ample. ∎
## 4. The proof of 0.1, 0.2 and 0.3
Recall that starting from $`f:XY`$ we constructed a family $`g^{}:Z^{}Y`$, with discriminant locus $`S^{}S`$. As in §1, we consider the Hodge bundles
$$E^{p,q}=R^qg_{}^{}\mathrm{\Omega }_{Z^{}/Y}^p(\mathrm{log}\mathrm{\Pi }^{}).$$
In the last section we obtained a morphism
$$\rho _{p,q}:F^{p,q}=R^qf_{}^{}(\tau ^{}\mathrm{\Omega }_{X/Y}^p(\mathrm{log}\mathrm{\Delta })^{(1)})E^{p,q},$$
compatible with the Kodaira Spencer map
$$\theta _{p,q}:E^{p,q}E^{p1,q+1}\mathrm{\Omega }_Y^1(\mathrm{log}S^{}).$$
Moreover, by 3.1 the image of $`\rho _{p1,q1}\iota `$ is contained in $`E^{p1,q+1}\mathrm{\Omega }_Y^1(\mathrm{log}S)`$.
###### Proposition 4.1.
Let $`\delta `$ denote the number of those singular fibres of $`f`$ which are not reduced normal crossing divisors and let $`\nu >1`$ be an integer with $`f_{}\omega _{X/Y}^\nu 0`$. If $`𝒜`$ is an invertible sheaf, with $`\mathrm{deg}(𝒜)>\delta `$, and such that $`f_{}\omega _{X/Y}^\nu 𝒜^\nu `$ is ample, then
1. $`(2g2+s)0`$ implies that $`\mathrm{deg}(𝒜)n(2g2+s)+\delta `$.
2. $`(2g2+s)>0`$ implies that $`\mathrm{deg}(𝒜)<n(2g2+s)+\delta `$.
###### Proof.
To handle both cases at once, define $`ϵ=1`$, if $`2g2+s=0`$ and $`ϵ=0`$, otherwise. Assume that
$$\mathrm{deg}(𝒜)n(2g2+s)+ϵ+\delta .$$
Then by 3.1, v), $`F^{n,0}`$ is an invertible subsheaf of $`E^{n,0}`$ of degree
$$\mathrm{deg}(F^{n,0})n(2g2+s)+ϵ,$$
in particular it is ample. For $`0in`$, we will construct by induction an invertible subsheaf $`\stackrel{~}{F}^{ni,i}`$ of $`\rho _{ni,i}(F^{ni,i})E^{ni,i}`$ of degree
$$\mathrm{deg}(\stackrel{~}{F}^{ni,i})(ni)(2g2+s)+ϵ.$$
If $`i<n`$, the sheaf $`\stackrel{~}{F}^{ni,i}`$ is ample, and by 1.3 it can not lie in the kernel of $`\theta _{ni,i}`$. On the other hand, by 3.1, i),
$$\theta _{ni,i}(\stackrel{~}{F}^{ni,i})\theta _{ni,i}(\rho _{ni,i}F^{ni,i})\rho _{ni1,i+1}(F^{ni1,i+1})\mathrm{\Omega }_Y^1(\mathrm{log}S).$$
The invertible sheaf
$$\stackrel{~}{F}^{ni1,i+1}=\theta _{ni,i}(\stackrel{~}{F}^{ni,i})\mathrm{\Omega }_Y^1(\mathrm{log}S)^1$$
thereby is a subsheaf of
$$\rho _{ni1,i+1}(F^{ni1,i+1})E^{ni1,i+1}$$
of degree $`\mathrm{deg}(\stackrel{~}{F}^{ni1,i+1})(ni1)(2g2+s)+ϵ`$.
For $`i=n`$ we obtain a subsheaf $`\stackrel{~}{F}^{0,n}`$ of degree $`\mathrm{deg}(\stackrel{~}{F}^{0,n})ϵ0`$, contradicting 3.1, iv). ∎
Now everything is set to prove 0.1, 0.3, and 0.2. We will proceed in the following way: Adding one or two points to $`S`$, if necessary, hence declaring some of the smooth fibres to be “singular”, we are allowed to assume
(4.1.1)
$$(2g2+s)=\mathrm{deg}(\mathrm{\Omega }_Y^1(\mathrm{log}S))0.$$
If $`2g2+s=0`$ and if $`f`$ is semistable, one finds the degree of $`\lambda _\nu `$ to be zero. Then the assumptions a) and b) in 0.1 imply that $`f`$ is isotrivial. For $`Y=^1`$, independently of the additional assumptions on $`F`$, this will imply that $`\kappa (X)=\mathrm{}`$.
For $`2g2+s>0`$ the inequalities stated in 0.3 i), ii), will hold true, independently of the semi-ampleness of $`\omega _F`$, whenever $`\mathrm{deg}(\lambda _\nu )>0`$.
###### Proposition 4.2.
Let $`Y_0`$ be either an elliptic curve or $`^{}`$, and let $`\tau :Y^{}Y`$ be a finite morphism, étale over $`Y_0`$, such that $`X\times _YY^{}Y^{}`$ has a semistable model $`f^{}:X^{}Y^{}`$. Then, for all $`\nu >1`$ with $`H^0(F,\omega _F^\nu )0`$, the degree of $`det(f_{}^{}\omega _{X^{}/Y^{}}^\nu )`$ is zero.
###### Proof.
The sheaf $`\lambda ^{}:=det(f_{}^{}{}_{}{}^{}\omega _{X^{}/Y^{}}^\nu )`$ is nef, hence if 4.2 does not hold true, it is ample. Since $`f^{}:X^{}Y^{}`$ is semistable, $`\lambda ^{}`$ is compatible with further pullbacks. Replacing $`Y^{}`$ by a covering, we may assume that $`\mathrm{deg}(\lambda ^{})>\nu er`$. By Proposition 2.7
$$f_{}^{}\omega _{X^{}/Y^{}}^\nu \frac{1}{re}\lambda ^{},$$
hence $`f_{}^{}\omega _{X^{}/Y^{}}^\nu 𝒪_Y^{}(\nu y^{})`$ is ample, for a point $`y^{}Y^{}`$. 4.1, a) implies that $`\mathrm{deg}(𝒪_Y^{}(y^{}))0`$, obviously a contradiction. ∎
###### Proposition 4.3.
Assume that $`f_{}\omega _{X/Y}^\nu 0`$, and that $`\lambda =det(f_{}\omega _{X/Y}^\nu )`$ is ample, for some $`\nu >1`$. Let $`\delta `$ denote the number of non-semistable fibres, $`r=\mathrm{rank}(f_{}\omega _{X/Y}^\nu )`$, and let $`e`$ be the constant introduced in 2.6. If $`2g2+s>0`$, then $`\mathrm{deg}(\lambda )(n(2g2+s)+\delta )\nu er.`$
###### Proof.
Choose an invertible sheaf $`𝒜`$ on $`Y`$ of degree $`n(2g2+s)+\delta `$. If
$$\mathrm{deg}(\lambda )>(n(2g2+s)+\delta )\nu er,$$
one finds $`\mathrm{deg}(𝒜^{\nu er})<\mathrm{deg}(\lambda )`$. By Proposition 2.7
$$f_{}\omega _{X/Y}^\nu \frac{1}{re}\lambda ,$$
hence $`f_{}\omega _{X/Y}^\nu 𝒜^\nu `$ is ample. Proposition 4.1, b) implies that
$$\mathrm{deg}(𝒜)<n(2g2+s)+\delta ,$$
contradicting the choice of $`𝒜`$. ∎
Proof of 0.1. Assume $`f:XY`$ is a morphism, smooth over $`Y_0`$. If $`Y_0`$ is an elliptic curve, $`f`$ is smooth. For $`Y_0=^{}`$, there exists a finite covering $`\tau :Y^{}Y`$, étale over $`^{}`$, and a semistable family $`f^{}:X^{}Y^{}`$, birational to the pullback of $`f`$. Hence to show that $`f`$ is isotrivial, 4.2 allows to assume that $`\mathrm{deg}(det(f_{}\omega _{X/Y}^\nu ))=0`$, for all $`\nu >1`$.
The experts will have noticed that the assumption a) or b) in 0.1 are exactly those needed by Kawamata, Kollár and the first named author to prove the additivity of the Kodaira dimension, and even the stronger statement $`Q_{n,1}`$, saying that for non-isotrivial morphisms $`f`$
$$\kappa (\omega _{X/Y})=\kappa (F)+1.$$
Using 2.4, this is equivalent to the ampleness of $`det(f_{}\omega _{X/Y}^\nu )`$, for all positive multiples $`\nu `$ of some $`\nu _01`$ (see , 7.2 and 7.6, and the references given there). In particular, the morphism $`f:XY`$, considered above, is isotrivial. ∎
Proof of 0.2. Assume there exist a morphism $`f:X^1`$, smooth over $`^{}`$, and with $`\kappa (X)0`$. Let
$$X\stackrel{f^{}}{}Y^{}\stackrel{\tau }{}^1$$
be the Stein factorization. $`\tau `$ must be smooth over $`^{}`$, hence $`Y^{}=^1`$ and $`\tau ^1(S)=\{0,\mathrm{}\}`$.
Altogether we find a morphism, denoted again by $`f:X^1`$, which is smooth over $`^{}`$ and whose general fibre $`F`$ is connected.
For a finite morphism $`^1^1`$ of degree $`d`$, étale over $`^{}`$, let
$$\phi :ZX\times _^1^1$$
be a desingularization, and $`g=pr_2\phi :Z^1`$. Then $`\phi ^{}pr_1^{}\omega _X^\nu `$ is a subsheaf of $`\omega _Z^\nu `$, hence $`\kappa (Z)\kappa (X)0`$. In particular, for some $`\nu _0`$ and for all positive multiples $`\nu `$ of $`\nu _0`$ the sheaf
$$g_{}\omega _Z^\nu =𝒪_^1(2\nu )g_{}\omega _{Z/^1}^\nu $$
has a non trivial section. Therefore $`g_{}\omega _{Z/^1}^\nu `$ contains a non-trivial ample subsheaf. On the other hand, $`g_{}\omega _{Z/^1}^\nu `$ is nef, hence $`det(g_{}\omega _{Z/^1}^\nu )`$ must be ample. Since $`g`$ is smooth over $`^1\{0,\mathrm{}\}`$, and semistable for $`d`$ sufficiently large and divisible, this contradicts 4.2. ∎
Proof of 0.3. In fact, by 4.3 the inequalities i) and ii) hold true under the assumptions a) and b) stated in 0.1. However, the constant $`e`$, defined in 2.6 might depend on the general fibre of $`f`$ and not just on the Hilbert polynomial $`h(t)`$. For this reason we have to require in 0.3 the sheaf $`\omega _F`$ to be semi-ample.
Under this assumption, there exists a quasi-projective moduli scheme $`M_h`$, parameterizing pairs $`(F,)`$ where $`F`$ is a manifold with $`\omega _F`$ semi-ample and with $``$ a polarization with Hilbert polynomial $`h`$ (see ). Seshadri and Kollár constructed a finite covering $`Z`$ of $`M_h`$, which carries a universal family (see , 9.25), containing all $`(F,)`$ with Hilbert polynomial $`h(t)`$ as fibres. So we find some $`\nu >0`$, such that for all such $`(F,)`$, $`\omega _F^\nu `$ is generated by global sections, hence we may choose the sheaf $``$ in 2.6 in such a way that $`|_F=\omega _F^\nu )`$. Finally, by , 5.17, $`e(|_F)=e(\omega _F^\nu )`$ is upper semicontinous for the Zariski topology. Hence there exists some $`e`$, with $`ee(\omega _F^\nu )`$ for all $`F`$. Altogether, we can choose $`\nu `$ and $`e`$ in 0.3, to depend just on $`M_h`$, hence just on $`h(t)`$.∎ |
warning/0002/astro-ph0002248.html | ar5iv | text | # Big Bang Nucleosynthesis updated with the NACRE Compilation
## 1 Introduction
Besides the expansion of the Universe and the ubiquitous presence of the fossil cosmological radiation, the Big Bang Nucleosynthesis (BBN) is one pilar of modern cosmology. It allows, in principle, the derivation of the baryonic density of the universe (see for reviews Schramm and Turner 1998, Sarkar 1996, 1999 and Olive et al 2000). The determination of light element abundances has improved dramatically in the recent past and the planned observations of $`D`$, $`{}_{}{}^{4}He`$, $`{}_{}{}^{6}Li`$ and $`{}_{}{}^{7}Li`$ should allow a precise determination (10% is a reachable goal) of the universal baryonic density, provided the precision of the model is made compatible with this objective. Taking advantage of the release of a new compilation of thermonuclear reation rates, called NACRE (Nuclear Astrophysics Compilation of REaction rates) (Angulo et al 1999), we have updated the standard BBN model developed at the Institut d’Astrophysique de Paris, including the analysis of the $`Be`$ and $`B`$ production in the Big Bang.
On the other hand, observations of light isotopes have florished: i) $`D/H`$ has been observed in absorbing clouds on the sightline of remote quasars, ii) refined observations of $`{}_{}{}^{4}He`$ have been performed in extragalactic very metal poor HII regions and in the Small Magellanic Cloud (Peimbert and Peimbert 2000), iii) the $`{}_{}{}^{6}Li`$ abundance has been determined in two halo stars, iv) high quality $`{}_{}{}^{7}Li`$ observations in the halo stars have been accumulated. A review of the present data can be found in Olive et al (2000) and Tytler et al (2000). Thus, it is timely to reassess the determination of the baryonic density of the Universe in the light of advances in nuclear physics and astronomical observations.
In section II, we present the new compilation of the reaction rates and compare it with the classical Caughlan-Fowler (1988) ones; we evaluate the sensitivity of the different light element abundances to the change of each relevant reaction rate; BBN calculations are performed using i) recommended values of the reaction rates and ii) extreme values obtained from the low and high rate limits. In section III, we discuss the astrophysical status of each isotope, both observationally and theoretically in order to confront it to the BBN calculation; we deduce the baryonic density of the Universe. Finally, we draw conclusion in section IV and stress the importance of precise measurements of a few key nuclear cross sections and refined abundance determinations of $`D`$ and $`{}_{}{}^{6}Li`$.
## 2 Nuclear Physics and Big Bang nucleosynthesis
The new compilation of Angulo et al (1999) of thermonuclear reaction rates for astrophysics, includes 86 charged-particle induced reactions corresponding to the proton capture reactions involved in the cold pp-chain, CNO cycle, NeNa and MgAl chains. The BBN network is constituted by about 60 reactions which participate to primordial nucleosynthesis up to $`{}_{}{}^{11}B`$. 22 of these reactions are covered by the NACRE compilation. The main innovative features of NACRE with respect to the former compilation Caughlan and Fowler (1988) are the following: (1) detailed references are provided to the source of original data; (2) uncertainties are analyzed in detail, realistic lower and upper bounds of the rates are provided; (3) the rates are given in tabular form, available also electronically on the World–Wide–Web. For these reasons, we can adopt the NACRE recommended rates for the calculation of the yields and use the upper and lower limits of the rates to test the sensitivity of the abundances to the nuclear uncertainties. As the origins for these uncertainties are documented in Angulo et al (1999), we do not discuss them here unless they show a significant effect on yields. We calculated the isotopic abundances as a function of $`\eta _{10}`$ between 1 and 10, changing one single reaction at a time. For each reaction we made a calculation with the high and low NACRE limits while the remaining reaction rates were set to their recommended NACRE value. Then we calculated, the maximum of the quantity $`\mathrm{\Delta }N/NN_{high}/N_{low}1`$ within the range of $`\eta _{10}`$ variations for each of the 8 isotopes. Positive (resp. negative) values correspond to higher (resp. lower) isotope production when the high rate is used instead of the low one. Note however that following the $`\mathrm{\Delta }N/N`$ definition, the positive values are not bound while the negative values are bound by -1. Hence, for instance (see Table 1), $`\mathrm{\Delta }N/N=+10`$ (resp. -0.78) means that the isotope yield is 11 times higher (resp. 4.5 times lower) with the high rate than with the low one. The corresponding results are shown in Table 1.
For three of these reactions, the test has not been made because the NACRE compilation does not provide high and low limits for the reverse rate of endoenergic reactions ($`{}_{}{}^{7}Li(\alpha ,n)^{10}B`$ , $`{}_{}{}^{11}B(p,n)^{11}C`$ and $`{}_{}{}^{14}B(p,\alpha )^{11}C`$). The reverse recommended rate can be calculated from the formulas. However, the low and high rates are only tabulated and limited down to temperatures chosen in order that the reaction rate remains above the lower limit of $`N_\mathrm{A}\sigma v10^{25}`$ cm<sup>3</sup> mol<sup>-1</sup> s<sup>-1</sup>. For $`Q<0`$ reactions, the reverse rate is higher than the direct tabulated one but limited by $`N_\mathrm{A}\sigma v10^{25}`$ cm<sup>3</sup> mol<sup>-1</sup> s<sup>-1</sup> times the reverse ratio.
In the analysis of yield uncertainties, one should keep in mind that the guidelines for the NACRE compilation favoured conservative upper and lower limits for the rates in order that the actual rate be within these limits with a high degree of confidence. For instance, when incompatible data set were present, and if the differences could not be resolved by analysing the publications alone the high and low limits were set in order to incorporate all data sets. In some case (e.g. $`D(\alpha ,\gamma )^6Li`$) the incompatibility is between experimental data and theoretical results making very problematic the interpretation of the rate uncertainty in term of gaussian distributions. When only an upper limit is available experimentally, as in previous compilations, its contribution is weighted by a 0., 0.1 and 1. factor respectively. This again makes difficult the probabilistic interpretation of rate uncertainties. Nevertheless, few NACRE reaction rates are at the origin of a significant uncertainty (Table 1).
The $`D(p,\gamma )^3He`$ NACRE reaction rate is responsible for a 20–30% uncertainty on most isotopes. It comes from the dispersion of experimental results. Note that the incompatibility of the two data set at low energy reported in Angulo et al. (1999) has been removed after a correction factor (Schmidt et al. 1996) has been applied to Schmidt et al. (1995) data to account for an unsuspected experimental bias. To check the effect of this rate update, we reiterate the NACRE calculation by fitting the experimental data points up to 2 MeV but with the corrected Schmidt et al. (1995,1996) data. Since it affects only the lowest energies, it has a negligible effect in the domain of BBN (see Fig. 1).
The most dramatic effect comes from the $`D(\alpha ,\gamma )^6Li`$ reaction which induces uncertainties of a factor 22 and 11 on the $`{}_{}{}^{6}Li`$ and <sup>10</sup>B yields. This rate uncertainty originates from the discrepancy between theoretical low energy dependance of the S–factor and experimental data (Kiener et al. 1991) obtained with the coulomb break–up technique (see Kharbach and Descouvemont (1998) for a recent comparison between theories and experiment). This reaction clearly deserves further experimental effort. The reactions $`{}_{}{}^{3}He(\alpha ,\gamma )^7Be`$ and $`{}_{}{}^{7}Li(p,\alpha )^4He`$ induce a significant (25–40% each) uncertainty on $`{}_{}{}^{7}Li`$ production. For these reactions, the rate uncertainties comes from the dispersion (systematic errors) of non resonant experimental data at low energy.
Uncertainties on $`BeB`$ isotope yields remain negligible when compared with the gap between calculated values and observational limitations. At maximum, a factor of $`4`$ uncertainty on $`{}_{}{}^{11}B`$ at low $`\eta `$ arises from the influence of the $`{}_{}{}^{11}B(p,\alpha )^8Be`$ reaction. However, the NACRE compilation covers only 22 of the 60 reactions involved in Big Bang Nucleosynthesis. In particular, $`BeB`$ yield uncertainties are most likely dominated by uncertainties in reaction rates not included in the NACRE compilation.
In front of the large systematic uncertainties on the observational abundance data (section 3.1), it seems premature to elaborate complex statistical procedures to get very precise theoretical uncertainties on the primordial abundances.
Indeed, extensive studies have used Monte-Carlo techniques to estimate the theoretical uncertainties studies have used Monte-Carlo techniques (Krauss et al 1990, Smith et al 1993, Fiorentini et al 1998, Olive et al 2000, Nollet and Burles 2000). These powerfull methods have proven their efficiency in various domains (e.g. simulations of high energy physics experiment) provided that the probability distribution involved are known. Concerning our approach, since the values given by the NACRE compilation do not represent statistical confidence level, but upper and lower limits, they are not directly appropriate to Monte-Carlo calculations, but they include both statistical and systematic errors (see above).
We estimate the uncertainties performing to global calculations, one with all the reaction rates set to their lower limits, and the second one with all the reaction rates set to their higher limits (dashed lines in figures 4 and 5). This method could lead to compensation (according to the signs of individual uncertainties displayed in table 1) between production and destruction and therefore to underestimate the global uncertainties in some cases. The advantage of this technique is simplicity and transparency. But the disadvantage is that it does not allow to derive a confidence level, in the statistical sense.
The primordial abundances of the light elements $`D`$, $`{}_{}{}^{3,4}He`$ and $`{}_{}{}^{7}Li`$ are governed by the expansion rate of the Universe and the cooling it induces. Under the classical assumptions, (homogeneity and the isotropy of the Universe and standard particle physics: three light neutrino species, neutron lifetime equal to 887 seconds) the abundances depend only on the baryon to photon ratio $`\eta `$, related to the baryonic parameter by $`\eta _{10}=273\mathrm{\Omega }_B.h^2`$ with $`h=H/100kms^1Mpc^1`$ (see e.g. for details Olive et al 2000).
We do not include the small corrections on the $`{}_{}{}^{4}He`$ mass fraction due to Coulomb, radiative and finite temperature effects, finite nucleon mass effects and differential neutrino heating (Sarkar 1999, Lopez and Turner 1999) since these corrections lead to effects much less than the uncertainties on the observational data.
The network extends up to $`{}_{}{}^{11}C`$ (decaying into $`{}_{}{}^{11}B`$), the leakage is taken into account through the reaction $`{}_{}{}^{11}C(n,2\alpha )^4He`$.
In figures 2 and 3, we compare the results obtained with i) the NACRE recommended reaction rates (solid lines) and ii) the CF88 ones (dashed lines). There is no significant difference except for $`{}_{}{}^{7}Li`$ at high $`\eta `$. In this range, $`{}_{}{}^{7}Li`$ comes from $`D(p,\gamma ){}_{}{}^{3}He(\alpha ,\gamma ){}_{}{}^{7}Be`$ followed by electron capture. So changes in the first rate result directly in changes of the final $`{}_{}{}^{7}Li`$ yield. $`{}_{}{}^{10}B`$ presents the largest difference due to the $`{}_{}{}^{10}B(p,\alpha )^7Be`$ destruction reaction rate. The NACRE rate is several orders of magnitude higher due to the inclusion (Rauscher and Raimann 1997) of a 10 keV, 5/2<sup>+</sup> resonance. However, there is no astrophysical and cosmological consequences since $`{}_{}{}^{10}B`$ is essentially of spallative origin (Vangioni-Flam et al 2000).
Figures 4 and 5 present the updated theoretical primordial abundances from $`D`$ to $`{}_{}{}^{11}B`$ using the NACRE compilation (mean values and extreme ones). D and $`{}_{}{}^{3,4}He`$ are almost not affected. Due to the uncertainty of the $`D(p,\gamma )^3He`$, the $`{}_{}{}^{7}Li`$ abundance at $`\eta >3`$ is affected by a significant error (about 30% to be compared to the 42% one mentionned by Olive et al 2000 deduced from the Smith et al 1993 analysis). At $`\eta <3`$, the $`{}_{}{}^{7}Li`$ uncertainty is reduced due to improvements in the derivation of the S factor of the $`T(\alpha ,\gamma )^7Li`$ reaction (Angulo et al 1999). This is the result of the high precision data provided by Brune et al (1994) and spanning the entire energy range of interest to BBN nucleosynthesis. Considering only the precise Brune et al. (1994) data would even reduce the rate uncertainty. However, in this specific case, with our method, error compensation occurs between the $`{}_{}{}^{7}Li(p,\alpha )^4He`$ and $`T(\alpha ,\gamma )^7Li`$ reaction rates as shown in table 1. From the same table, one sees that the maximized uncertainties are $`\pm `$25%. Consequently on figure 4, the uncertainty is somewhat underestimated which does not affect significantly the general conclusions. On the contrary, for $`\eta >3`$ region, the errors on the reaction rates $`D(p,\gamma )^3He`$, $`{}_{}{}^{3}He(\alpha ,\gamma )^7Li`$ add up and therefore the uncertainties are not underestimated.
$`{}_{}{}^{6}Li`$ has a particularly large error bar due to the poor knowledge of $`D(\alpha ,\gamma )^6Li`$. The maximum value of the primordial $`{}_{}{}^{6}Li/H`$ (at low $`\eta `$) which is of the order of 5.$`\times 10^{13}`$ may not be out of reach of future measurements in halo stars; it represents a factor 10 below the 6/7 value measured at present in old stars (\[Fe/H\] = -2.3). The results presented here are in fair agreement with previous calculations (Thomas et al 1993, Schramm 1993, Delbourgo-Salvador and Vangioni-Flam 1993, Vangioni-Flam et al 1999). We confirm that primordial abundances of $`BeB`$ are negligible even in the most favorable case, and spallation remains the main mechanism to produce them in the course of the galactic evolution.
## 3 Astrophysical and cosmological discussion
### 3.1 Astrophysical aspects
In the following, we decline the astrophysical observational and theoretical status of each isotope of interest and their possible evolution since the Big Bang in order to define reasonable error boxes to prepare the confrontation to the theoretical calculations.
#### 3.1.1 Deuterium
$`D`$ is particularly sensitive to the baryon/photon ratio, $`\eta `$, and has been considered up to now as the best baryometer (e.g. Reeves 1994). However, due to a certain confusion on $`D/H`$ abundance evaluations, both at high redshifts and in the local Galaxy, some care has to be taken in the cosmological use of Deuterium. Let us present a brief overview of the observations.
$`D`$ is measured in three astrophysical and/or cosmological sites i) the local interstellar medium, ISM ($`D_{ISM}`$), ii) the protosolar nebula ($`D_{ps}`$), iii) the cosmological clouds ($`D_{cc}`$). These three values serve as signposts to follow the evolution of $`D`$ in the Universe and in the Galaxy. Deuterium, due to its fragility is completely burnt in stars. Thus, if no production mechanism is at work, we must have $`D_{cc}>D_{ps}>D_{ISM}`$.
The local $`D`$ abundance, inferred from UV observations of the nearby ISM, estimated to $`(1.6\pm 0.1)`$$`\times 10^5`$ (Linsky et al 1995), is probably not unique, ranging from 5.$`\times 10^6`$ to about 2.$`\times 10^5`$ (Vidal-Madjar et al 1998, Lemoine et al 1999, Vidal-Madjar 2000). These variations are lacking explanations. Thus, there is an ambiguity on the true local $`D/H`$ value which, by the way, serves as a normalisation for the chemical evolutionary models. These discrepancies weaken the predictive ability of the evolutionary models to derive the primordial $`D`$ abundance.
$`D/H`$ ratios are measured in the solar system (Jupiter, Saturn, Uranus, Neptun, comets). This allows to derive a precise protosolar value of $`(3\pm 0.3)`$$`\times 10^5`$ (Drouart et al 1999), somewhat higher than the estimate of Geiss and Gloeckler (1998) $`(2.1\pm 0.5)`$$`\times 10^5`$.
$`D/H`$ has also been determined in absorbing clouds on the sightlines of quasars. On one side, Tytler et al (2000) (and reference therein) have found three absorption systems in which $`D/H`$ are i) $`(3.24\pm 0.3)`$$`\times 10^5`$, ii) $`4_{0.6}^{+0.8}`$$`\times 10^5`$, iii) $`<6.7`$$`\times 10^5`$. As a fair representation of these data, we adopt the following range : 3. to 4.$`\times 10^5`$. This estimate, if identified with the primordial one, is unconfortably close to the protosolar one, since it implies a very small $`D`$ destruction all along the galactic evolution corresponding to a small variation of the star formation rate from the birth of the galaxy up to now, in contradiction with the general trend indicated by the strong increase of the cosmic star formation rate vs redshift ($`0<z<2`$) (Blain et al 2000, Madau 2000). On the other side, high values of $`D/H`$ have been reported concerning the quasar QSO 1718+4807 ($`z`$ = 0.701) namely, $`D/H=(2.5\pm 0.5)`$$`\times 10^4`$ according to Webb et al (1997) and 8.$`\times 10^5`$$`<D/H<`$5.7$`\times 10^4`$ according to Tytler et al (2000). Note that the analysis of Levshakov et al (1999) allowing non gaussian velocity distributions leads to lower values. Consequently, we adopt a second data box bounded by 8.$`\times 10^5`$$`<D/H<`$ 3.$`\times 10^4`$.
#### 3.1.2 Helium-3
This isotope is produced in comparable amount to that of deuterium, but at the opposite, its stellar and galactic story is not simple. Its production and destruction are model dependent (Vassiliadis and Woods 1993, Charbonnel 2000). In spite of great effort directed to its abundance determination in HII regions and planetary nebula (Balser et al 1999, Bania and Rood 2000) $`{}_{}{}^{3}He`$ cannot be used, at the moment, as a reliable cosmic baryometer (Olive et al 1995, Galli et al 1997) since it is very difficult to extrapolate its abundance back to its primordial value.
#### 3.1.3 Helium-4
The primordial abundance of $`{}_{}{}^{4}He`$ by mass, $`Y_p`$, is measured in low metallicity extragalactic HII regions (for a review see Kunth and Ostlin 2000). In addition to the primordial component, $`{}_{}{}^{4}He`$ is also produced in stars together with oxygen and nitrogen through global stellar nucleosynthesis. Therefore, in order to extract the primordial component from the observational data, it is necessary to extrapolate back the observed $`{}_{}{}^{4}He`$ value down to zero metallicity. Olive, Skillman and Steigman (1997) selected 62 blue compact galaxies and obtained $`Y_p`$ $``$ 0.234. Izotov and Thuan (1998) pointed out that the effect of the HeI stellar absorption has more importance that previously thought and they reported $`Y_p=0.245\pm 0.004`$. Recently, Fields and Olive (1998) reanalyzed the observational data and reported $`Y_p=0.238\pm (0.002)`$ stat, $`\pm (0.005)`$ syst where the errors are 1 $`\sigma `$ values. The new determination $`Y_p=0.2345\pm 0.0030`$ by Peimbert and Peimbert (2000) on the basis of observations of HII regions in the Small Magellanic Cloud points towards the lowest helium value proposed by Fields and Olive (1998). In this context, two data boxes emerge: i) $`Y_p=0.245\pm 0.004`$ and ii) $`Y_p=0.238\pm 0.005`$.
#### 3.1.4 Lithium-6,7
Recent advances on the determination of $`Li`$ in halo stars (Spite et al 1996, Bonifacio and Molaro 1997, Molaro 1999, Smith et al 1998, Ryan et al 1999a) indicate that the Spite plateau is exceptionally thin ($`<`$ 0.05 dex). This small dispersion, together with the presence of $`{}_{}{}^{6}Li`$ in two halo stars (Smith et al 1993, 1998, Hobbs and Thorburn 1994, 1997, Cayrel et al 1999) indicate that the stellar destruction of $`{}_{}{}^{7}Li`$ if any, is very limited (less than $``$ 0.1 dex, see Ryan et al 1999a). $`{}_{}{}^{6}Li`$ is however of cosmological interest since, being more fragile than $`{}_{}{}^{7}Li`$, its mere presence in the atmosphere of halo stars confirms that $`{}_{}{}^{7}Li`$ is essentially intact in these stars (Vangioni-Flam et al 1999, Fields and Olive 1999). This, combined with the very small dispersion around the average of the Spite plateau add confidence in interpreting it as indicative of the primordial $`Li`$ abundance (especially at the lowest metallicities, where the contamination by spallation is expected to be negligible (Ryan et al 1999b). Stellar modelisation should adapt to this constraint. Indeed, simple models of lithium evolution predict little or no depletion (Deliyannis et al 1990), thus conforting the primordial nature of lithium in metal poor halo stars. However, three different mechanisms of alteration of lithium in halo stars have been suggested i) diffusion/gravitational settling (Michaud 2000 and references therein), ii) rotational mixing (Chaboyer 1998, Pinsonneault et al 1999) and iii) stellar winds (Vauclair and Charbonnel 1995). Some combinations of these three mechanisms have to be envisioned. There is a paradox between the absence of dispersion and the number of the processes which could produce a potential dispersion. This implies either a curious statistical compensation or more radically, that these physical mechanisms are intrinsically irrelevant. However, metallicity independent depletion mechanisms for instance (mixing induced by gravity waves) cannot be totally excluded (Cayrel, private communication).
Consequently, we take the observed value of the Spite plateau including the observational dispersion ( $`A(Li)=2.2\pm 0.04`$) to which we add 0.1 dex to account the maximum $`Li`$ stellar depletion. This corresponds for the primordial lithium to the following range 1.4$`\times 10^{10}`$$`<^7Li/H<`$2.2$`\times 10^{10}`$. This evaluation is in fair agreement with that of Olive et al (2000) but it is narrower than that of Tytler et al (2000) who have enlarged the limit of depletion to allow agreement with their low $`D/H`$ measurement.
In the following, due to the high observational quality and the large sample analyzed (more than 70 objects), $`{}_{}{}^{7}Li`$ is used as a baryometer, keeping in mind that this isotope has the pecularity to allow to possible solutions depending on the side of the lithium valley.
#### 3.1.5 Beryllium and Boron
Abundance observations of elemental $`Be`$ and $`B`$ in very metal poor stars in the halo have made great progresses in the recent years (Gilmore et al 1992, Duncan et al 1992, Boesgaard and King 1993, Ryan et al 1994, Duncan et al 1997, Garcia-Lopez et al 1998, Primas et al 1999, Primas 2000). These observations concern indeed galactic evolution; the lowest observed values (of the order of $`10^{13}`$) are much higher than the BBN calculated abundances. We confirm that BBN calculated abundances of $`{}_{}{}^{9}Be`$, $`{}_{}{}^{10}B`$ and $`{}_{}{}^{11}B`$ are negligible with respect the measured ones in the more metal poor stars. The origin of these elements can be explained in term of spallation of fast carbon and oxygen in the early Galaxy (Vangioni-Flam et al 2000).
### 3.2 Baryonic density of the universe
Once the different error boxes corresponding to the observed isotopes abundances corrected for evolutionary effects are established, we can compare them to the predictions of the BBN calculations. As different $`D`$ and $`{}_{}{}^{4}He`$ measurements are dichotomic, contrary to $`{}_{}{}^{7}Li`$, we put emphasis on this last isotope to determine a possible range of baryonic densities. As shown in figure 4, considering $`{}_{}{}^{7}Li`$ alone, two possible ranges emerge: i) $`1.5<\eta _{10}<1.9`$, ii) $`3.3<\eta _{10}<5.1`$. For h = 0.65, we get i) $`0.013<\mathrm{\Omega }{}_{B}{}^{}<0.019`$ and ii) $`0.029<\mathrm{\Omega }{}_{B}{}^{}<0.045`$.
The first range is in good concordance with the error boxes related to high $`D`$ and low $`{}_{}{}^{4}He`$. However, the largest measured value of $`D/H`$ seems excluded ($`D/H<`$ 3.$`\times 10^4`$). The second one, on the right side of the diagram, is in fair agreement with a lower $`D/H`$ (except the lowest measured value, $`D/H`$ = 3.$`\times 10^5`$) and a higher $`{}_{}{}^{4}He`$ (except also the highest value 0.25). At this stage of the analysis, we have to admit that two ranges of baryonic density have to be into account, only future observations will help to remove the ambiguity.
It is worth comparing the baryonic density to that of the luminous matter ($`\mathrm{\Omega }_L`$) in the Universe to infer the amount of the baryonic dark matter. Recent estimates of $`\mathrm{\Omega }_L`$ ranges between 0.002 and 0.004 (Salucci and Persic 1999), which is lower than both $`\mathrm{\Omega }_B`$ obtained. The difference makes necessary baryonic dark matter. Focusing on spiral galaxies the amount of luminous matter is estimated to $`\mathrm{\Omega }_{LS}`$ = $`1.44_{0.2}^{+1.55}`$$`\times 10^3`$ ( Sallucci and Persic 1999). Considering that the dynamical mass of the halo of spiral galaxies is about ten times higher than that of the disk, the corresponding $`\mathrm{\Omega }_{HS}`$ is about 0.015. In both cases ( $`\mathrm{\Omega }_B`$0.015 or $`\mathrm{\Omega }_B`$0.04) all the dark matter in the halo of our Galaxy could in principle be baryonic. Note that since eight years of searches for microlensing events by the MACHO and EROS collaborations toward the Magellanic clouds have revealed only a few events. The fraction of the halo in the form of dark massive compact objects, using a typical model, is estimated to only 20% by Alcock et al (1999) in agreement with the limits given by Lasserre et al (2000). On the other hand, the observations of the Lyman alpha forest clouds between the redshifts 0 and 5, lead to a corresponding $`\mathrm{\Omega }_B`$ of about 0.03 ($`\pm 0.01`$), taking into account the uncertainty related to ionised hydrogen (Riediger, Petitjean and Mucket 1998). This value is thought to reflect the bulk of the baryons at large scale. It is more consistent with our high $`\mathrm{\Omega }_B`$ range.
## 4 Conclusion
Big bang nucleosynthesis has been studied since a long time, but it deserves permanent care since it gives access to the baryon density which is a key cosmological parameter. This work has been aimed at integrating the last development in both fields of nuclear physics and observational abundance determination of light elements.
1. The update of the reaction rates of the BBN using the NACRE compilation does not lead to crucial modifications of the general conclusions concerning the baryonic content of the Universe. The average values of the abundances of isotopes of cosmological interest are in general similar to that calculated on the basis of the Caughlan-Fowler (1988) compilation.
2. However, the modification of the $`D(p,\gamma )^3He`$ reaction rate leads to a $`{}_{}{}^{7}Li`$ abundance slightly lower at $`\eta >3`$. But the uncertainty on this reaction rate remains high ($`\pm `$30%). At $`\eta <3`$, the revision of the $`T(\alpha ,\gamma )^7Li`$ reaction rate leads to a very neat reduction of the uncertainty of the calculated $`{}_{}{}^{7}Li`$ abundance. However this reaction, together with the $`{}_{}{}^{7}Li(p,\alpha )^4He`$, remain the main sources of the $`{}_{}{}^{7}Li`$ uncertainty at low $`\eta `$.
3. The abundance of $`{}_{}{}^{10}B`$ is modified by the new $`{}_{}{}^{10}B(p,\alpha )^7Li`$ reaction rate, but there is no cosmological consequences.
4. $`{}_{}{}^{6}Li`$ is affected by the large uncertainty of the $`D(\alpha ,\gamma )^6Li`$ reaction. However, $`{}_{}{}^{6}Li`$ is essentially of spallative origin.
5. Owing to the high observational reliability of the $`{}_{}{}^{7}Li`$ abundance data with respect to the $`D`$ data avalaible both rare and debated, we choose it as the leading baryometer, since it appears that the primitive $`{}_{}{}^{7}Li`$ is almost intact in halo stars. Due to the competition between $`T(\alpha ,\gamma )^7Li`$ and $`D(p,\gamma )^3He(\alpha ,\gamma )^7Be`$ (e $`\nu `$) $`{}_{}{}^{7}Li`$, the curve of $`{}_{}{}^{7}Li`$ vs $`\eta `$ presents a valley shape. Consequently, the observational error box of $`{}_{}{}^{7}Li`$ leads to two ranges of $`\eta `$ : $`1.5<\eta <1.9`$ and $`3.3<\eta <5.1`$ (corresponding to $`0.013<\mathrm{\Omega }_B<0.019`$ and $`0.029<\mathrm{\Omega }_B<0.045`$ for h=0.65). In both cases, all the dark matter in the halo of our Galaxy could be baryonic. However, only a fraction of 20% is detected through microlensing events in the direction of the Magellanic clouds.
6. These two $`\eta `$ ranges are confronted to the other available cosmologically relevant isotopes, namely $`D`$ and $`{}_{}{}^{4}He`$. The first $`\eta `$ range agrees with a high $`D`$ and low $`{}_{}{}^{4}He`$ values, the second range is in concordance with a low $`D`$ and high $`{}_{}{}^{4}He`$ values. At present, none of these solutions can be excluded.
7. In the future, on the nuclear physics front, it would be important to (re-)measure the $`D(\alpha ,\gamma )^6Li`$ reaction to reduce the uncertainty on the calculated $`{}_{}{}^{6}Li`$ abundance. High precision measurements over the full energy range of interest to BBN of the $`D(p,\gamma )^3He`$, $`{}_{}{}^{3}He(\alpha ,\gamma )^7Be`$ and $`{}_{}{}^{7}Li(p,\alpha )^4He`$ reactions would also reduce the uncertainty on $`{}_{}{}^{7}Li`$ abundance calculations. On the astronomical front, more data on $`D`$ in absorbing clouds on the sightlines of quasars at different redshifts are mandatory to remove the ambiguity. However, if the large scale $`D`$ dichotomy remains, it will be time to invoke specific mechanisms of $`D`$ production and destruction like photodisintegration of $`D`$ and $`{}_{}{}^{4}He`$ by $`\gamma `$ ray quasars (blazars) (Cassé and Vangioni-Flam 1997) and/or nuclear spallation. The observations of $`{}_{}{}^{6}Li`$ in two halo stars has been a great progress and a determination of its abundance in very metal poor stars should be pursued. It will help to constrain even more stringently the possible lithium depletion in these stars and confort definitively the primordial status of the Spite plateau. Together with nuclear improvements, refined $`{}_{}{}^{6}Li`$ measurements in very metal poor stars (possibly via the VLT) could perhaps lead us towards the primordial $`{}_{}{}^{6}Li`$ abundance.
We warmly thank Roger Cayrel for illuminating discussions, Jurgen Kiener, Gilles Bogaert and Carmen Angulo for their comments on nuclear data. |
warning/0002/hep-ph0002241.html | ar5iv | text | # Nonperturbative Gluon Radiation and Energy Dependence of Elastic Scattering
## Abstract
The energy dependence of the total hadronic cross sections is caused by gluon bremsstrahlung which we treat nonperturbatively. It is located at small transverse distances about $`0.3fm`$ from the valence quarks. The cross section of gluon radiation is predicted to exponentiate and rise with energy as $`s^\mathrm{\Delta }`$ with $`\mathrm{\Delta }=0.17\pm 0.01`$. The total cross section also includes a large energy independent Born term which corresponds to no gluon radiation. The calculated total cross section and the slope of elastic scattering are in good agreement with the data.
The dynamics of energy dependence of the hadronic total cross sections is a long standing challenge since 1973 when this effect was first observed at the ISR. In DIS the source for the rising total cross section for interaction of highly virtual photons is well understood in QCD as caused by an intensive gluon bremsstrahlung . Indeed, radiation of each gluon supply an extra $`\mathrm{ln}s`$ and $`\mathrm{ln}Q^2`$. This is a specific regime of radiation when a $`\overline{q}q`$ fluctuation of the photon of a tiny size $`1/Q`$ radiates gluons at much larger transverse separations.
It is difficult to extend the perturbative results to soft hadronic collisions because it is quite a different regime where the approximations made in the perturbative case break down. Namely, gluon radiation giving rise to the energy dependence of the total cross section occurs at rather small transverse distances around the valence quarks, $`r_00.3fm`$ which are much smaller than the mean interquark spacing in light hadrons. This conclusion follows from the analysis of the data for diffractive gluon radiation based on the light-cone approach when the effective nonperturbative interaction of radiated gluons is included.
The smallness of the gluon clouds of the valence quarks is confirmed by the study of the gluon formfactor of the proton employing QCD sum rules . The $`Q^2`$ dependence of the formfactor turns out to be rather weak corresponding to a small radius of the gluon distribution which was estimated at the same value $`r_00.3fm`$.
Another evidence for a short gluon-gluon correlation length $`\lambda 0.3fm`$ arises in the stochastic vacuum model of Dosch and Simonov , as it was measured on the lattice . In the case of a poorly populated gluon cloud (only about one gluon is radiated by a valence quark at available energies, see below) this corresponds to the correlation radius between the quark and the gluon.
The same size $`0.3fm`$ emerges from the Shuryak’s instanton liquid model as the instanton size which controls the mean radius of the sea surrounding a valence quark and by many phenomenological analyses.
In the Gribov’s theory of confinement the same distance $`0.3fm`$ should correspond to the critical regime related to breaking of chiral symmetry. Namely, at smaller distances, a perturbative quark-gluon basis is appropriate, while at larger separations quasi-Goldstone pions emerge. The corresponding critical value of the QCD running constant $`\alpha _c=0.43`$ evaluated in turns out to be very close to our estimate (see below) of $`\alpha _s`$ corresponding to gluon radiation separated by $`0.3fm`$. The value of $`\alpha _s`$ is crucial for our evaluation of the energy dependence of the gluon bremsstrahlung.
It is quite plausible that all these observations are the manifestations of the same dynamics, however it is still unclear how to make a Lorentz boost in these approaches. This is the advantage of the light cone treatment of nonperturbative gluon radiation which seems to be best designed for calculating the energy dependence of the total cross section. We believe that the nonperturbative interaction of gluons introduced in as a light-cone potential is an effective manifestation of properties of the QCD vacuum. Similar scale $`0.3fm`$ found in all these approaches supports this conjecture.
An interesting attempt to implement the nonperturbative gluon interaction into the Pomeron ladder building was made recently by Kharzeev and Levin and Shuryak . They found that the radiation of colorless pairs of gluons is a part of the leading-log approximation since each extra power of the coupling $`\alpha _s`$ cancels due to the strong glue-glue interaction. The radiated glueballs are not clustering around the valence quarks, but spreading all over the hadron. The estimated $`\mathrm{\Delta }0.05`$ is about twice as small (and even more so if corrected for unitarity) as the data need. Although the scale for $`\alpha _P^{}1/M_0^2`$ seems to be correct, an extra factor $`\mathrm{\Delta }/4`$ makes it too small.
We start calculating the energy dependence of the total cross section summing up the contributions of different Fock components of the incident hadron,
$$\sigma _{tot}^{hN}=\underset{n=0}{}\sigma _n^{hN}.$$
(1)
To avoid double-counting, we sum over cross sections $`\sigma _n`$ of physical processes corresponding to the radiation of $`n`$ gluons.
The lowest Fock component of a hadron contains only valence quarks. The corresponding Born term in the total cross section has the form (for the sake of simplicity we assume that the incident hadron is a meson),
$$\sigma _0^{hN}=\underset{0}{\overset{1}{}}𝑑\alpha _qd^2R\left|\mathrm{\Psi }_{\overline{q}q}^h(\alpha _q,R)\right|^2\sigma _{\overline{q}q}^N(R).$$
(2)
Here the Fock state wave function $`\mathrm{\Psi }_{\overline{q}q}^h(\alpha _q,R)`$ depends on the transverse $`q\overline{q}`$ separation $`R`$ and on the fraction $`\alpha _q`$ of the light-cone momentum of the pair carried by the quark. The cross section $`\sigma _{\overline{q}q}^N(R)`$ of interaction of the valence $`\overline{q}q`$ dipole with a nucleon cannot be calculated perturbatively since the separation $`R`$ is large. According to this energy independent term has no relation to the smallness of the spots (gluon clouds) in the hadron.
The next contribution to $`\sigma _{tot}^{hN}`$ comes from the radiation of a single gluon. The radiation is possible only due to the difference between the cross sections for the $`\overline{q}q`$ and $`\overline{q}qG`$ Fock components, otherwise no new state can be produced . The cross section of radiation of a single gluon reads ,
$`\sigma _1^{hN}={\displaystyle \underset{0}{\overset{1}{}}}𝑑\alpha _q{\displaystyle d^2R\left|\mathrm{\Psi }_{\overline{q}q}^h(R,\alpha _q)\right|^2}`$ (3)
$`\times `$ $`{\displaystyle \frac{9}{4}}{\displaystyle \underset{\alpha _G1}{}}{\displaystyle \frac{d\alpha _G}{\alpha _G}}{\displaystyle }d^2r\{\left|\mathrm{\Psi }_{\overline{q}G}(\stackrel{}{R}+\stackrel{}{r},\alpha _G)\right|^2\sigma _{\overline{q}q}^N(\stackrel{}{R}+\stackrel{}{r})`$ (4)
$`+`$ $`\left|\mathrm{\Psi }_{qG}(\stackrel{}{r},\alpha _G)\right|^2\sigma _{\overline{q}q}^N(r)\mathrm{Re}\mathrm{\Psi }_{qG}^{}(\stackrel{}{r},\alpha _G)\mathrm{\Psi }_{\overline{q}G}(\stackrel{}{R}+\stackrel{}{r},\alpha _G)`$ (5)
$`\times `$ $`[\sigma _{\overline{q}q}^N(\stackrel{}{R}+\stackrel{}{r})+\sigma _{\overline{q}q}^N(r)\sigma _{\overline{q}q}^N(R)]\}`$ (6)
Here $`\alpha _G`$ is the fraction of the quark momentum carried by the gluon, and $`\stackrel{}{r}`$ is the quark-gluon transverse separation. The three terms in the curly brackets correspond to the radiation of the gluon by the quark, by the antiquark and to their interference respectively.
The nonperturbative wave function for a quark-gluon Fock component is derived in . Neglecting the quark mass, the wave function reads,
$$\mathrm{\Psi }_{qG}(\stackrel{}{r},\alpha _G1)=\frac{2i}{\pi }\sqrt{\frac{\alpha _s}{3}}\frac{\stackrel{}{e}^{}\stackrel{}{r}}{r^2}e^{r^2b_0^2/2},$$
(7)
where $`\stackrel{}{e}`$ is the polarization vector of the massless gluon. The parameter $`b_0=0.65GeV`$ characterizing the nonperturbative quark gluon interaction is fixed by the data on large mass diffractive dissociation corresponding to the triple-Pomeron limit. It leads to quite a short mean quark-gluon separation $`r_0=\sqrt{r^2}=1/b_00.3fm`$, which is small relative to the hadronic size. Therefore, only one or the other of the first two terms in (6) can be large, while the interference one can always be neglected. In this case, the integration in (6) is easily performed,
$$\sigma _1^{hN}=N\frac{4\alpha _s}{3\pi }\mathrm{ln}\left(\frac{s}{s_0}\right)\frac{9C}{4b_0^2}.$$
(8)
Here we assume that the approximation $`\sigma _{\overline{q}q}^N(r)=Cr^2`$ is valid for $`r1/b_0`$. $`N`$ is the number of valence quarks, $`\mathrm{ln}(s/s_0)=\mathrm{ln}[(\alpha _G)_{max}/(\alpha _G)_{min}]`$, where $`(\alpha _G)_{min}=2b_0^2/s`$, but $`(\alpha _G)_{max}`$ is ill defined. It should be sufficiently small to use the wave functions (6). This leads to the condition to $`s_03GeV^2`$. At high energy $`\sigma _1`$ has little sensitivity on $`s_0`$ which we fix at $`s_0=30GeV^2`$ for further applications.
The radiation of each new, $`n`$-th gluon can be treated as radiation by a color triplet which is an effective quark surrounded by $`n1`$ gluons. It should be resolved by the soft interaction with the target to be different from the radiation of $`n1`$ gluons i.e. the radiation cross section is proportional to the difference between the total cross sections of the two subsequent Fock states which is $`9C/4b_0^2`$. This can be also proved using a $`1/N_c`$ expansion and the dipole representation of Mueller . Since the radiation of a gluon with $`\alpha _G1`$ does not affect the impact parameter of the radiating quark, all the quark lines in the final state cancel with the same lines in the initial state (see the prescription for calculating the radiative cross section in ), except for the radiation of the $`n`$-th gluon. Thus, $`\sigma _n`$ for quark-proton interaction in the leading-log approximation reads,
$$\sigma _n^{qN}=\frac{1}{n!}\left[\frac{4\alpha _s}{3\pi }\mathrm{ln}\left(\frac{s}{s_0}\right)\right]^n\frac{9C}{4b_0^2}.$$
(9)
Summing up the powers of logarithms in (1) we arrive at the following expression for the total cross section,
$$\sigma _{tot}^{hp}=\stackrel{~}{\sigma }_0^{hp}+N\frac{9C}{4b_0^2}\left(\frac{s}{s_0}\right)^\mathrm{\Delta },$$
(10)
with
$$\mathrm{\Delta }=\frac{4\alpha _s}{3\pi },$$
(11)
and $`\stackrel{~}{\sigma }_0^{hp}=\sigma _0^{hp}9C/4b_0^2`$. The soft Pomeron intercept, $`\alpha _P(0)=1+\mathrm{\Delta }`$, and can be evaluated provided that the QCD coupling $`\alpha _s`$ is known.
In Gribov’s confinement scenario, chiral symmetry breaking occurs when the running coupling $`\alpha _s`$ exceeds the critical value $`\alpha _s=\alpha _c0.43`$ . This should happen at a distance of the order of the size of a constituent quark $`0.3fm`$. Therefore, this value can be used in (11).
One can also calculate the mean $`\alpha _s`$ for nonperturvative gluon radiation averaging over transverse momenta $`k_T`$ of the radiated gluons. The popular way to extend the running QCD coupling $`\alpha _s(k_T^2)`$ down to small $`k_T`$ is a shift of the variable $`k_T^2k_T^2+k_0^2`$, where $`k_0^20.25GeV^2`$ was evaluated in using the dispersive approach to calculating higher twist effects in hard reactions . The nonperturbative interaction of the radiated gluons drastically suppresses small transverse momenta, pushing $`k_T^2`$ to higher values which lowers $`\alpha _s`$. We use the transverse momentum gluon distribution calculated in in the light-cone approach in terms of the universal color dipole cross section . We calculated $`\alpha _s`$ with a simple parameterization $`\sigma (\rho )1\mathrm{exp}(\rho ^2/\rho _0^2)`$. For a reasonable variation of $`\rho _0=0.31fm`$ the mean coupling is in the range $`\alpha _s=0.380.43`$ which is very close the the critical value mentioned above. Taking the mid value $`\alpha _s=0.4`$ we get from (11),
$$\mathrm{\Delta }=0.17\pm 0.01.$$
(12)
This value is about twice as large as the one suggested by the data for the energy dependence of total hadronic cross sections . However, the radiative part is a rather small fraction of the total cross section (at medium high energies). A structure similar to (10) with a large $`\mathrm{\Delta }`$ was suggested in (with quite a different motivation) and proved to agree well with the data.
The factor $`C`$ in the second term in (10) can also be evaluated. We calculated the dipole cross section with the gluon effective mass $`0.15GeV`$ (to incorporate confinement) and $`\alpha _s=0.4`$ and found $`C=2.3`$ at $`\rho =1/b_0`$. Thus, the energy dependent term in (10) is fully determined.
The cross section (10) apparently violates the Froissart bound and one should perform unitarity corrections. Indeed, the partial elastic amplitude shows a precocious onset of unitarity restrictions at small impact parameters important even at medium high energies .
Following we assume that the $`t`$-dependence of the $`pp`$ elastic amplitude is given by the Dirac electromagnetic formfactor squared. Correspondingly, the mean square radius $`\stackrel{~}{r}_{ch}^2`$ evaluated in should be smaller than $`r_{ch}^2`$.
For the dipole parameterization of the formfactor the partial elastic amplitude which is related via unitarity to $`\sigma _n^{pp}`$, given by (2), (9), takes the form,
$$\mathrm{Im}\gamma _n^{pp}(b,s)=\frac{\sigma _n^{pp}(s)}{8\pi B_n}y^3K_3(y),$$
(13)
where $`K_3(y)`$ is the third order modified Bessel function and $`y=b\sqrt{8/B_n}`$. The slope parameter grows linearly with $`n`$ due to the random walk of radiated gluons with a step $`1/b_0^2`$ in the impact parameter plane, $`B_n=2\stackrel{~}{r}_{ch}^2/3+n/2b_0^2`$.
We unitarize the partial amplitude $`\mathrm{Im}\gamma _P(s,b)=\underset{n=0}{}\mathrm{Im}\gamma _n(s,b)`$ using the quasi-eikonal model ,
$$\mathrm{Im}\mathrm{\Gamma }_P(b,s)=\frac{1\mathrm{exp}\left[D(s)\mathrm{Im}\gamma _P(b,s)\right]}{D(s)}$$
(14)
where $`D(s)1=\sigma _{sd}(s)/\sigma _{el}(s)`$ is the ratio of the single diffractive to elastic cross sections. It is approximately equal to $`0.25`$ at the lowest ISR energy and slightly decreases with energy $`s^{0.04}`$ . Note that good results can be also achieved with a different unitarization scheme similar to one suggested in . The details will be presented elsewhere.
In order to calculate the total cross section, $`\sigma _{tot}=2d^2b\mathrm{Im}\mathrm{\Gamma }(b,s)`$, one needs to fix the energy independent term with $`n=0`$ in (13). This can be done comparing with the data for $`\sigma _{tot}`$ at any energy sufficiently high to neglect Reggeon contributions. We used the most precise data at $`\sqrt{s}=546GeV`$ and fixed $`\stackrel{~}{\sigma }_0=39.7mb`$.
The predicted energy dependence of $`\sigma _{tot}^{pp}`$ is shown by the dashed curve in Fig. 1 which is in good agreement with the data at high energy , but apparently needs Reggeon corrections towards low energies.
We added a Reggeon term $`\mathrm{Im}\mathrm{\Gamma }_R(s,b)[1\mathrm{Im}\mathrm{\Gamma }_P(s,b)]`$ screened by unitarity corrections, which was fitted independently for $`pp`$ and $`\overline{p}p`$, $`\sigma _R^{pp}=17.8mb/\sqrt{s/s_0}`$, $`\sigma _R^{\overline{p}p}=32.8mb/\sqrt{s/s_0}`$. The fitted Reggeon slope is $`B_R=R_R^2+2\alpha _R^{}\mathrm{ln}(s/s_0)`$, where $`\alpha _R^{}=0.9GeV^2`$ and $`R_R^2=3GeV^2`$.
The results are shown by the solid curves of Fig. 1 ($`pp`$ bottom curve and $`\overline{p}p`$ upper curve).
As soon as the partial amplitude (14) is known, we are in position to predict the slope of elastic scattering at $`t=0`$, $`B_{el}(s)=b^2/2`$, where averaging is weighted by the partial amplitude (14). The results exhibit good agreement when compared with the $`pp`$ and $`\overline{p}p`$ data in Fig. 2.
Although the value of the slope essentially depends on our choice of $`\stackrel{~}{r}_{ch}^2`$ in (13), the predicted energy dependence, i.e. the effective value $`\alpha _P^{}`$ is fully defined by the parameter $`b_0`$ fixed in . Indeed, each radiated gluon makes a “step” $`1/b_0^2=(0.3fm)^2`$ in the impact parameter plane leading to the rising energy dependence of the elastic slope. Eventually, at very high energies the approximation of small gluon clouds breaks down but the mean number of gluons in a quark $`n=\mathrm{\Delta }\mathrm{ln}(s/s_0)`$ remains quite small even in the energy range of colliders. It is only $`n=0.71`$ at the ISR and reaches about two gluons at the Tevatron. Correspondingly, the mean square of the quark radius grows from $`0.06fm^2`$ to $`0.18fm^2`$ which is still rather small compared to the mean square of the charge radius of the proton.
Summarizing, the strong nonperturbative interaction of radiated gluons substantially shrinks the gluon clouds around of valence quarks. This spots are small ($`0.3fm`$) compared to the hadronic radius, but the gluon radiation grows with energy as $`s^\mathrm{\Delta }`$ where $`\mathrm{\Delta }=0.17\pm 0.01`$. Such a steep rise does not contradict the data since this fraction of the total cross section is rather small (it contains a factor $`1/b_0^21mb`$). A large energy independent fraction comes from the Born term which corresponds to scattering of the valence quark skeleton without gluon radiation. A very soft interaction which cannot resolve and excite the small spots contributes to this term. It cannot be reliably predicted and is fixed by data , while the energy dependent term is fully calculated. The results are in good agreement with the data for total $`pp`$ and $`\overline{p}p`$ cross sections and elastic slopes.
Note that although we have some room for fine-tuning in the the parameters ($`C,s_0,\stackrel{~}{r}_{ch}^2`$), the results are rather insensitive and the agreement with data is always pretty good. We have also tried a different unitarization scheme suggested in arriving to similar results.
The details of calculations and further comparison with elastic scattering data will be published elsewhere.
Acknowledgments: We are thankful to Jörg Hüfner, Andreas Schäfer and Sasha Tarasov for illuminating and very helpful discussions. We are grateful to Jörg Raufeisen who has read the paper and made many improving comments. This work was partially supported by the grant INTAS-97-OPEN-31696, by the European Network: Hadronic Physics with Electromagnetic Probes, Contract No. FMRX-CT96-0008 and by the INFN and MURST of Italy. |
warning/0002/cond-mat0002143.html | ar5iv | text | # Extremal-point Densities of Interface Fluctuations
## I Introduction and Motivation
The aim of statistical mechanics is to relate the macroscopic observables to the microscopic properties of the system. Before attempting any such derivation one always has to specify the spectrum of length-scales the analysis will comprise: while ‘macroscopic’ is usually defined in a unique way by the every-day-life length scale, the ‘microscopic’ is never so obvious, and the choice of the best lower-end scale is highly non-universal, it is system dependent, usually left to our physical ‘intuition’, or it is set by the limitations of the experimental instrumentation. It is obvious that in order to derive the laws of the gaseous matter we do not need to employ the physics of elementary particles, it is enough to start from an effective microscopic model (or Hamiltonian) on the level of molecular interactions. Then starting from the equations of motion on the microscopic level and using a statistic and probabilistic approach, the macroscale physics is derived. In this ‘long wavelength’ approach most of the microscopic, or short wavelength information is usually redundant and it is scaled away.
Sometimes however, microscopic quantities are important and directly contribute to macroscopic observables, e.g., the nearest-neighbor correlations in driven systems determine the current, in model B the mobility, in kinetic Ising model the domain-wall velocity, in parallel computation the utilization (efficiency) of conservative parallel algorithms, etc. Once a lower length scale is set on which we can define an effective microscopic dynamics, it becomes meaningful to ask questions about local properties at this length scale, e.g. nearest neighbor correlations, contour distributions, extremal-point densities, etc. These quantities are obviously not universal, however their behavior against the variation of the length scales can present qualitative and universal features. Here we study the dynamics of macroscopically rough surfaces via investigating an intriguing miscroscopic quantity: the density of extrema (minima) and its finite size effects. We derive a number of analytical results about these quantities for a large class of non-equilibrium surface fluctuations described by linear Langevin equations, and solid-on-solid (SOS) lattice-growth models. Besides their obvious relevance to surface physics our technique can be used to show the asymptotic scalability of conservative massively parallel algorithms for discrete-event simulation, i.e., the fact that the efficiency of such computational schemes does not vanish with increasing the number of processing elements, but it has a lower non-trivial bound. The solution of this problem is not only of practical importance from the point of view of parallel computing, but it has important consequences for our understanding of systems with asynchronous parallel dynamics, in general. There are numerous dynamical systems both man-made, and found in the nature, that contain a “substantial amount” of parallelism. For example,
1) in wireless cellular communications the call arrivals and departures are happening in continuos time (Poisson arrivals), and the discrete events (call arrivals) are not synchronized by a global clock. Nevertheless, calls initiated in cells substantially far from each other can be processed simultaneously by the parallel simulator without changing the poissonian nature of the underlying process. The problem of designing efficient dynamic channel allocation schemes for wireless networks is a very difficult one and currently it is done by modelling the network as a system of interacting continuous time stochastic automata on parallel architectures .
2) in magnetic systems the discrete events are the spin flip attempts (e.g., Glauber dynamics for Ising systems). While traditional single spin flip dynamics may seem inherently serial, systems with short range interactions can be simulated in parallel: spins far from each other with different local simulated times can be updated simultaneously. Fast and efficient parallel Monte-Carlo algorithms are extremely welcome when studying metastable decay and hysteresis of kinetic Ising ferromagnets below their critical temperature, see and references therein.
3) financial markets, and especially the stock market is an extremely dynamic, high connectivity network of relations, thousands of trades are being made asynchronously every minute.
4) the brain. The human brain, in spite of its low weight of approx. 1kg, and volume of 1400 $`cm^3`$, it contains about 100 billion neurons, each neuron being connected through synapses to approximately 10,000 other neurons. The total number of synapses in a human brain is about 1000 trillion ($`10^{15}`$). The neurons of a single human brain placed end-to-end would make a “string” of an enourmous lenght: 250,000 miles . Assuming that each neuron of a single human cortex can be in two states only (resting or acting), the total number of different brain configurations would be $`2_{}^{10}{}_{}{}^{11}`$. According to Carl Sagan, this number is greater than the total number of protons and electrons of the known universe, . The brain does an incredible amount of parallel computation: it simultaneously manages all of our body functions, we can talk and walk at the same time, etc.
5) evolution of networks such as the internet, has parallel dynamics: the local network connectivity changes concurrently as many sites are attached or removed in different locations. As a matter of fact the physics of such dynamic networks is a currently heavily investigated and a rapidly emerging field .
In order to present the basic ideas and notions in the simplest way, in the following we will restrict ourselves to one dimensional interfaces that have no overhangs. The restriction on the overhangs may actually be lifted with a proper parametrization of the surface, a problem to which we will return briefly in the concluding section. The first visual impression when we look at a rough surface $`h(x,t)`$ is the extent of the fluctuations perpendicular to the substrate, in other words, the width of the fluctuations. The width (or the rms of the height $`h`$ of the fluctuations) is probably the most extensively studied quantity in interface physics, due to the fact that its definition is simply quantifiable and therefore measurable:
$$w(L,t)=\sqrt{\overline{\left[h(x,t)\right]^2}\left[\overline{h(x,t)}\right]^2},$$
(1)
where the overline denotes the average over the substrate. It is well-known that this quantity characterizes the long wavelength behavior of the fluctuations, the high frequency components being averaged out in (1). The short wavelength end of the spectrum has been ignored in the literature mainly because of its non-universal character, and also because it seemed to lack such a simple quantifiable definition as the width $`w`$.
In the following we will present a quantity that is almost as simple and intuitive as the width $`w`$ but it characterizes the high frequency components of the fluctuations and it is simply quantifiable. For illustrative purposes let us consider the classic Weierstrass function defined as the $`M\mathrm{}`$ limit of the smooth functions $`W_M(a,b;x)`$:
$$W(a,b;x)=\underset{M\mathrm{}}{lim}W_M(a,b;x)=\underset{M\mathrm{}}{lim}\underset{m=0}{\overset{M}{}}a^m\mathrm{cos}\left(b^mx\right),a,b>1$$
(2)
Figure 1a shows the graph of $`W_M`$ at $`a=2`$ and $`b=3`$ (arbitrary values) for $`M=0,1,2,3,4`$ in the interval $`x[0,4\pi ]`$. As one can see, by increasing $`M`$ we are adding more and more detail to the graph of the function on finer and finer length-scales. Thus $`M`$ plays the role of a regulator for the microscopic cut-off length which is $`b^M`$, and for $`M=\mathrm{}`$ and $`b>a`$, the function becomes nowhere differentiable as it was shown by Hardy .
Comparing the graphs of $`W_M`$ for lower $`M`$ values with those for higher $`M`$-s we observe that the width effectively does not change, however the curves look qualitatively very different. This is obvious from (2): adding an extra term will not change the long wavelength modes, but adds a higher frequency component to the Fourier spectrum of the graph. We need to operationally define a quantifiable parameter which makes a distinction between a much ‘fuzzier’ graph, such as $`M=4`$ and a smoother one, such as $`M=1`$. The natural choice based on Figure 1a is the number of local minima (or extrema) in the graph of function. In Figure 1b we present the number of local minima $`u_M`$ vs. $`M`$ for two different values of $`b`$, $`b=2.8`$ and $`b=1.8`$, while keeping $`a`$ at the same value of $`a=2`$. For all $`b`$ values (not only for these two) the leading behaviour is exponential: $`u_M\lambda ^M`$. The inset in Figure 1b shows the dependence of the rate $`\lambda `$ as a function of $`b`$ for fixed $`a`$. We observe that for $`b>a`$, $`\lambda =b`$, but below $`b=a`$ the dependence crosses over to another, seemingly linear function. For $`b>a`$ the amplitude of the extra added term in $`W_{M+1}`$ is large enough to prevent the cancellation of the newly appearing minima by the drop in the local slope of $`W_M`$. At $`ba`$ the number of cancelled extrema starts to increase drastically with an exponential trend, leading to the crossover seen in the inset of Figure 2b. It has been shown that the fractal dimension of the Weierstrass function for $`b>a`$ is given by $`D_0=\mathrm{ln}b/\mathrm{ln}a`$ , . For $`ba`$ the Weierstrass curve becomes non-fractal with a dimension of $`D_0=1`$. By varying $`b`$ with respect to $`a`$, we are crossing a fractal-smooth transition at $`b=a`$. The very intriguing observation we just come across is that even though we are in the smooth regime ($`b<a`$) the density of minima is still a diverging quantity (the $`b=1.8`$ curve in Figure 1b). It is thus possible to have an infinite number of ‘wrinkles’ in the Weierstrass function without having a diverging length, without having a fractal in the classical sense. The transition from fractal to smooth, as $`b`$ is lowered appears as a non-analyticity in the divergence rate of the curve’s wrinkledness. A rigorous analytic treatment of this problem seems to be highly non-trivial and we propose it as an open question to the interested reader.
The simple example shown above suggests that there is novel and non-trivial physics lying behind the analysis of extremal point densities, and it gives extra information on the morphology of interfaces. Given an interface $`h(x,t)`$, we propose a quantitative form that characterizes the density of minima via a ‘partition-function’ like expression, which is hardly more complex than Eq. (1) and gives an alternate description of the surface morphology:
$$u_q(L,[h])=\frac{1}{L}\underset{i}{}\left[K(x_i)\right]^q,q>0,x_i\text{ are non-degenerate minima of }h$$
(3)
with $`K(x_i)`$ denoting the curvature of $`h`$ at the local (non-degenerate) minimum point $`x_i`$. The variable $`q`$ can be conceived as ‘inverse temperature’. Obviously, for $`q=0`$ we obtain the number of local minima per unit interface length. The rigorous mathematical description and definitions lying at the basis of (3) is being presented in Section IV. The quantity in Eq. (3) is reminiscent to the partition function used in the definition of the thermodynamical formalism of one dimensional chaotic maps and also to the definition of the dynamical or Rényi entropies of these chaotic maps. In that case, however the curvatures at the minima are replaced by cylinder intervals and/or the visiting probabilities of these cylinders.
We present a detailed analysis for the above quantity in case of a large class of linear Langevin equations of type $`h/t=\nu (^2h)^{z/2}+\eta (x,t)`$, where $`\eta `$ is a Gaussian noise term, and $`z`$ a positive real number. These Langevin equations are found to describe faithfully the fluctuations of monoatomic steps on various substrates, see for a review Ref. . One of the interesting conclusions we came to by studying the extremal-point densities on such equations is that depending on the value of $`z`$ the typical surface morphology can be fractal, or locally smooth, and the two regimes are separated by a critical $`z`$ value, $`z_c`$. In the fractal case, the interface will have infinitely many minima and cusps just as in the case of the nondifferentiable Weierstrass function (2), and the extremal point densities become infinite, or if the problem is discretized onto a lattice with spacing $`a`$, a power-law diverging behavior is observed as $`a0^+`$ for these densities. This sudden change of the ‘intrinsic roughness’ may be conceived as a phase transition even in an experimental situation. Changing a parameter, such as the temperature, the law describing the fluctuations can change since the mechanism responsible for the fluctuations can change character as the temperature varies. For example, it has been recently shown using Scanning Tunneling Microscope (STM) measurements , that the fluctuations of single atom layer steps on $`Cu`$ (111) below $`T=300\text{ }^oC`$ correspond to the perifery diffusion mechanism ($`z=4`$), but above this temperature (such as $`T=500\text{ }^oC`$ in their measurements) the mechanism is attachement-detachment where $`z=2`$, see also Ref. .
The paper is organized as follows: In Sections II and III we define and investigate on several well known on-lattice models the minimum point density and derive exact results in the steady-state ($`t\mathrm{}`$) including finite size effects. As a practical application of these on-lattice results, we briefly present in Subsection III.B a lattice surface-growth model which exactly describes the evolution of the simulated time-horizon for conservative massively parallel schemes in parallel computation, and solve a long-standing asymptotic scalability question for these update schemes. In Section IV we lay down a more rigorous mathematical treatment for extremal point densities, and stochastic extremal point densities on the continuum, with a detailed derivation for a large class of linear Langevin equations (which are in fact the continuum counterparts of the discrete ones from Section II). The more rigorous treatment allows for an exact analytical evaluation not only in the steady-state, but for all times! We identify novel characteristic dimensions that separate regimes with divergent extremal point densities from convergent ones and which give a novel understanding of the short wavelength physics behind these kinetic roughening processes.
## II Linear surface growth models on the lattice
In the present Section we focus on discrete, one dimensional models from the linear theory of kinetic roughening . Let us consider a one dimensional substrate consisting of $`L`$ lattice sites, with periodic boundary conditions. For simplicity the lattice constant is taken as unity, which clearly, represents the lower cut-off length for the effective equation of motion. For the moment let us study the discretized counterpart of the general Langevin equation that describes the linear theory of Molecular Beam Epitaxy (MBE) :
$$_th_i(t)=\nu ^2h_i(t)\kappa ^4h_i(t)+\eta _i(t),$$
(4)
where $`\eta _i(t)`$ is Gaussian white noise with
$$\eta _i(t)\eta _i(t^{})=2D\delta _{i,j}\delta (tt^{}),$$
(5)
and $`^2`$ is the discrete Laplacian operator, i.e., $`^2f_j=f_{j+1}+f_{j1}2f_j`$, applied to an arbitrary lattice function $`f_j`$. This equation arrises in MBE with both surface diffusion mechanism (the 4th order, or curvature term) and desorption mechanism (the 2nd order, or diffusive term) present and it has been studied extensively by several authors . Stability requires $`\nu 0`$ and $`\kappa 0`$ (as a matter of fact, on the lattice is enough to have $`\nu >0`$ and $`\kappa \nu /2`$). Starting from a completely flat initial condition, the interface roughens until the correlation length $`\xi `$ reaches the size of the system $`\xi L`$, when the roughening saturates over into a steady-state regime. The process of kinetic roughening is controlled by the intrinsic length scale , $`\sqrt{\kappa /\nu }`$. Below this lengthscale the roughening is dominated by the surface diffusion or Mullins term (the 4th order operator) but above it is characterized by the evaporation piece (the diffusion) or Edwards-Wilkinson term. Since eq. (4) is translationaly invariant and linear in $`h`$ it can be solved via the discrete Fourier-transform:
$$\stackrel{~}{h}_k=\underset{j=0}{\overset{L1}{}}e^{ikj}h_i,k=\frac{2\pi n}{L},n=0,1,2,\mathrm{},L1.$$
(6)
Then eq. (4) translates into
$$_t\stackrel{~}{h}_k(t)=\left[2\nu (1\mathrm{cos}(k))+4\kappa (1\mathrm{cos}(k))^2\right]\stackrel{~}{h}_k(t)+\stackrel{~}{\eta }_k(t),$$
(7)
with
$$\stackrel{~}{\eta }_k(t)\stackrel{~}{\eta }_k^{}(t^{})=2DL\delta _{(k+k^{})\text{mod}\mathrm{\hspace{0.33em}2}\pi ,0}\delta (tt^{})$$
(8)
Following the definition of the equal-time structure factor for $`S(k,t)`$, namely
$$S^h(k,t)L\delta _{(k+k^{})\text{mod}\mathrm{\hspace{0.33em}2}\pi ,0}\stackrel{~}{h}_k(t)\stackrel{~}{h}_k^{}(t),$$
(9)
one obtains for an initially flat interface:
$$S^h(k,t)=S^h(k)\left(1e^{(4\nu (1\mathrm{cos}(k))+8\kappa (1\mathrm{cos}(k))^2)t}\right).$$
(10)
In the above equation
$$S^h(k)\underset{t\mathrm{}}{lim}S^h(k,t)=\frac{D}{2\nu (1\mathrm{cos}(k))+4\kappa (1\mathrm{cos}(k))^2}$$
(11)
is the steady-state structure factor.
For $`\nu 0`$, and in the asymptotic scaling limit where $`L\sqrt{\kappa /\nu }`$, the model belongs to the EW universality class and the roughening exponent is $`\alpha =1/2`$ (it is defined through the scaling $`L^{2\alpha }`$ of the interface width $`w_L^2(t)=(1/L)_{i=1}^L[h_i(t)\overline{h}]^2`$ in the steady state). The presence of the curvature term does not change the universal scaling properties for the surface width, and one finds the same exponents as for the pure EW ($`\kappa =0`$) case in eq. (4). For $`\nu =0`$ the surface is purely curvature driven ($`z=4`$) and the model belongs to a different universality class where the steady-state width scales with a roughness exponent of $`\alpha =3/2`$.
In the following we will be mostly interested in some local steady-state properties of the surface $`h_i`$. In particular, we want to find the density of local minima for the surface described by (4). The operator which measures this quantity is
$$u=\frac{1}{L}\underset{i=1}{\overset{L}{}}\mathrm{\Theta }(h_{i1}h_i)\mathrm{\Theta }(h_{i+1}h_i).$$
(12)
This expression motivates the introduction of the local slopes, $`\varphi _i=h_{i+1}h_i`$. In this representation the operator for the density of local minima (for the original surface) is
$$u=\frac{1}{L}\underset{i=1}{\overset{L}{}}\mathrm{\Theta }(\varphi _{i1})\mathrm{\Theta }(\varphi _i),$$
(13)
and its steady-state average is $`u=\mathrm{\Theta }(\varphi _{i1})\mathrm{\Theta }(\varphi _i)=\mathrm{\Theta }(\varphi _1)\mathrm{\Theta }(\varphi _2)`$, due to translational invariance. The average density of local minima is the same as the probability that at a randomly chosen site of the lattice the surface exhibits a local minimum. It is governed by the nearest-neighbor two-slope distribution, which is also Gaussian and fully determined by $`\varphi _1^2=\varphi _2^2`$ and $`\varphi _1\varphi _2`$:
$$P^{nn}(\varphi _1,\varphi _2)e^{\frac{1}{2}\varphi _jA_{jk}^{\mathrm{nn}}\varphi _k},j,k=1,2,$$
(14)
where
$$A^{\mathrm{nn}}=\left(\begin{array}{cc}\varphi _1^2& \varphi _1\varphi _2\\ \varphi _1\varphi _2& \varphi _1^2\end{array}\right)^1$$
(15)
As we derive in Appendix A, the density of local minima only depends on the ratio $`\varphi _1\varphi _2/\varphi _1^2`$:
$$u=\frac{1}{2\pi }\mathrm{arccos}\left(\frac{\varphi _1\varphi _2}{\varphi _1^2}\right).$$
(16)
Finite-size effects of $`u`$ are obviously carried through those of the correlations. First we find the steady-state structure factor for the slopes. Since $`\stackrel{~}{\varphi }_k=(1e^{ik})\stackrel{~}{h}_k`$, we have $`S^\varphi (k)=2(1\mathrm{cos}(k))S^h(k)`$. Then from (11) one obtains:
$$S^\varphi (k)=\frac{D}{\nu +2\kappa (1\mathrm{cos}(k))},\text{for}k0,\text{and}S^\varphi (k)=0,\text{for}k=0.$$
(17)
The former automatically follows from the $`_{i=1}^L\varphi _i=0`$ relation. Then we obtain the slope-slope correlations
$$C_L^\varphi (l)\varphi _i\varphi _{i+l}=\frac{1}{L}\underset{n=1}{\overset{L1}{}}e^{i\frac{2\pi n}{L}l}S^\varphi \left(\frac{2\pi n}{L}\right).$$
(18)
With the help of Poisson summation formulas, in Appendix B we show a derivation for the exact spatial correlation function, which yields
$$C_L^\varphi (l)=\frac{D}{\nu +2\kappa }\left\{\frac{b^{|l|}}{\sqrt{1a^2}}\frac{1}{1a}\frac{1}{L}+\frac{b^L}{1b^L}\frac{b^l+b^l}{\sqrt{1a^2}}\right\},|l|L,$$
(19)
with
$$a\frac{2\kappa }{\nu +2\kappa },\text{and}b\frac{1\sqrt{1a^2}}{a}.$$
(20)
We have $`|a|1`$ and $`b1`$. The second term in the bracket in (19) gives a uniform power law correction, while the third one gives an exponential correction to the correlation function in the thermodynamic limit. For $`\nu 0`$ and $`L\mathrm{}`$ one obtains
$$C_{\mathrm{}}^\varphi (l)=\frac{D}{\nu +2\kappa }\frac{b^{|l|}}{\sqrt{1a^2}}=\frac{D}{\nu +2\kappa }\frac{e^{|l|/\xi _{\mathrm{}}^\varphi }}{\sqrt{1a^2}}$$
(21)
where we define the correlation length of the slopes for an infinite system as:
$$\xi _{\mathrm{}}^\varphi \frac{1}{\mathrm{ln}(b)}.$$
(22)
In the $`\nu 0`$ limit it becomes the intrinsic correlation length which diverges as $`\nu ^{1/2}`$: $`\xi _{\mathrm{}}^\varphi \stackrel{\nu 0}{}\sqrt{\kappa /\nu }`$ and
$$C_{\mathrm{}}^\varphi (l)\stackrel{\nu 0}{}\frac{D}{2\kappa }\left(\sqrt{\frac{\kappa }{\nu }}|l|\right)\frac{D}{2\kappa }\left(\xi _{\mathrm{}}^\varphi |l|\right).$$
(23)
In this limit the slopes (separated by any finite distance) become highly correlated, and one may start to anticipate that the density of local minima will vanish for the original surface $`\{h_i\}`$. In the following two subsections we investigate the density of local minima and its finite-size effects for the Edwards-Wilkinson and the Mullins cases.
### A Density of local minima for Edwards-Wilkinson term dominated regime
To study the finite size effects for the local minimum density, we neglect the exponentially small correction in (19), so in the asymptotic limit, where $`L\xi _{\mathrm{}}^\varphi `$, $`C_L^\varphi (l)`$ decays exponentially with uniform finite-size corrections:
$$C_L^\varphi (l)\frac{D}{\nu +2\kappa }\left\{\frac{b^{|l|}}{\sqrt{1a^2}}\frac{1}{1a}\frac{1}{L}\right\}$$
(24)
This holds for the special case $`\kappa =0`$ as well, (in fact, there the exponential correction exatly vanishes) leaving
$$C_L^\varphi (l)=\frac{D}{\nu }\left(\delta _{l,0}\frac{1}{L}\right).$$
(25)
Now, emplying eq. (16), we can obtain the density of minima as:
$$u_L=\frac{1}{2\pi }\mathrm{arccos}\left(\frac{C_L^\varphi (1)}{C_L^\varphi (0)}\right)\frac{1}{2\pi }\mathrm{arccos}(b)+\frac{1}{2\pi }\sqrt{\frac{1b}{1+b}}\sqrt{\frac{1+a}{1a}}\frac{1}{L},$$
(26)
Again, for the $`\kappa =0`$ case one has a compact exact expression and the corresponding large $`L`$ behavior:
$$u_L=\frac{1}{2\pi }\mathrm{arccos}\left(\frac{1}{L1}\right)\frac{1}{4}+\frac{1}{2\pi }\frac{1}{L},$$
(27)
which can also be obtained by taking the $`\kappa 0`$ limit in (26). To summarize, as long as $`\nu 0`$, the model belongs to the EW universality class, and in the steady state, the density of local minima behaves as
$$u_Lu_{\mathrm{}}+\frac{\mathrm{const}.}{L},$$
(28)
where $`u_{\mathrm{}}`$ is the value of the density of local minima in the thermodynamic limit:
$$u_{\mathrm{}}=\frac{1}{2\pi }\mathrm{arccos}(b).$$
(29)
Note that this quantity can be small, but does not vanish if $`\nu `$ is close but not equal to $`0`$. Further, the system exhibits the scaling (28) for asymptotically large systems, where $`L\xi _{\mathrm{}}^\varphi `$. It is important to see in detail how $`u_{\mathrm{}}`$ behaves as $`\nu 0`$:
$$u_{\mathrm{}}\stackrel{\nu 0}{}\frac{1}{2\pi }\mathrm{arccos}\left(1\sqrt{2}\sqrt{1a}\right)\frac{1}{2\pi }\mathrm{arccos}\left(1\sqrt{\frac{\nu }{\kappa }}\right)\frac{1}{2\pi }\left(2\sqrt{\frac{\nu }{\kappa }}\right)^{1/2}\frac{\sqrt{2}}{2\pi }\frac{1}{\sqrt{\xi _{\mathrm{}}^\varphi }}.$$
(30)
Thus, the density of local minima for an infinite system vanishes as we approach the purely curvature driven ($`\nu 0`$) limit. Simply speaking, the local slopes become “infinitely” correlated, such that $`C_{\mathrm{}}^\varphi (l)`$ diverges \[according to eq. (23)\], and the ratio $`C_{\mathrm{}}^\varphi (l)/C_{\mathrm{}}^\varphi (0)`$ for any fixed $`l`$ tends to $`1`$. This is the physical picture behind the vanishing density of local minima.
### B Density of local minima for the Mullins term dominated regime
Here we take the $`\nu 0`$ limit first and then study the finite size effects in the purely curvature driven model. The slope correlations are finite for finite $`L`$ as can be seen from eq. (18), since the $`n=0`$ term is not included in the sum! Thus, in the exact closed formula (19) a careful limiting procedure has to be taken which indeed yields the internal cancellation of the apparently divergent terms. Then one obtains the exact slope correlations for the $`\nu =0`$ case:
$$C_L^\varphi (l)=\frac{D}{2\kappa }\left\{\frac{L}{6}\left(1\frac{1}{L^2}\right)\right|l|\left(1\frac{|l|}{L}\right)\}$$
(31)
and for the local minimum density:
$$u_L=\frac{1}{2\pi }\mathrm{arccos}\left(1\frac{6}{L+1}\right)\frac{\sqrt{3}}{\pi }\frac{1}{\sqrt{L}}$$
(32)
It vanishes in the thermodynamic limit, and hence, one observes that the limits $`\nu 0`$ and $`L\mathrm{}`$ are interchangable. For $`\nu =0`$, $`L`$ is directly associated with the correlation length and we can define $`\xi _L^\varphi L/6`$. Then the correlations and the density of local minima takes the same scaling form as eqs. (23) and (30):
$$C_L^\varphi (l)\frac{D}{2\kappa }\left(\xi _L^\varphi |l|\right),$$
(33)
and
$$u_L\frac{\sqrt{2}}{2\pi }\frac{1}{\sqrt{\xi _L^\varphi }}.$$
(34)
### C Scaling considerations for higher order equations
Let us now consider another equation but with a generalized relaxational term that includes the Edwards Wilkinson and the noisy Mullins equation as particluar cases:
$$_th_i(t)=\nu \left(^2\right)^{z/2}h+\eta _i(t).$$
(35)
where $`z`$ is a positive real number (not necessarily integer). Other $`z`$ values of experimental interest are $`z=1`$, relaxation through plastic flow, ), and $`z=3`$ terrace-diffusion mechanism . For early times, such that $`tL^z`$, the interface width $`w_L^2(t)`$ increases with time as
$$w_L^2(t)t^{2\beta },$$
(36)
where $`\beta =(z1)/2z`$ . In the $`t\mathrm{}`$ limit, where $`tL^z`$, the interface width saturates for a finite system, but diverges with $`L`$ according to $`w_L^2(\mathrm{})L^{2\alpha }`$ where $`\alpha =(z1)/2`$ is the roughness exponent .
For $`z=4`$ (curvature driven interface) we saw that the slope fluctuation behaves as $`C_L^\varphi (0)=\varphi _i^2L`$. For higher $`z`$ for the slope-slope correlation function one can deduce
$$C_L^\varphi (l)=\frac{D}{L}\underset{n=1}{\overset{L1}{}}\frac{e^{i\frac{2\pi n}{L}l}}{\nu \left[2\left(1\mathrm{cos}\left(\frac{2\pi n}{L}\right)\right)\right]^{\frac{z2}{2}}}.$$
(37)
It is divergent in the $`L\mathrm{}`$ limit, as a result of infinitely small wave-vectors $`1/L`$, and we can see that
$$C_L^\varphi (0)L^{z3}.$$
(38)
It is also useful to define the slope difference correlation function
$$G_L^\varphi (l)(\varphi _{i+l}\varphi _i)^2$$
(39)
for which one can write
$$G_L^\varphi (l)=\frac{D}{L}\underset{n=1}{\overset{L1}{}}\frac{2\left(1\mathrm{cos}\left(\frac{2\pi n}{L}l\right)\right)}{\nu \left[2\left(1\mathrm{cos}\left(\frac{2\pi n}{L}\right)\right)\right]^{\frac{z2}{2}}}.$$
(40)
For the small wave-vector behavior we can again deduce that for $`z>5`$
$$G_L^\varphi (l)L^{z5}l^2.$$
(41)
One may refer to this form as “anomalous” scaling for the slope difference correlation function in the following sense. For $`z<5`$ the scaling form for $`G_L^\varphi (l)`$ follows that of $`C_L^\varphi (0)`$ \[eq. (38)\], i.e., $`G_L^\varphi (l)l^{z3}`$. For $`z>5`$ \[eq. (41)\] it obviously features a different $`l`$ dependence and an additional power of $`L`$, and it diverges in the $`L\mathrm{}`$ limit.
Having these scaling functions for large $`L`$, we can easily obtain the scaling behavior for the average density of local minima. Exploiting the identity
$$C_L^\varphi (l)=C_L^\varphi (0)\frac{1}{2}G_L^\varphi (l)$$
(42)
we use the general form for the local minimum density:
$$u=\frac{1}{2\pi }\mathrm{arccos}\left(\frac{C_L^\varphi (1)}{C_L^\varphi (0)}\right)=\frac{1}{2\pi }\mathrm{arccos}\left(1\frac{1}{2}\frac{G_L^\varphi (1)}{C_L^\varphi (0)}\right)\frac{1}{2\pi }\mathrm{arccos}\left(1\frac{\mathrm{const}.}{L^2}\right)\frac{1}{L}$$
(43)
Note that this is the scaling behavior for all $`z>5`$. It simply shows the trivial lower bound for $`u`$: since there is always at least one minima (and one maxima) among the $`L`$ sites, it can never be smaller than $`1/L`$.
### D The average curvature at local minima
The next natural question to ask is how the average curvature, $`K`$ at the minimum points scales with the system size for the general system described by eq. (35). This can be evaluated as the conditional average of the local curvature at the local minima:
$$K_{\mathrm{min}}=(\varphi _i\varphi _{i1})_{\mathrm{min}}=\frac{(\varphi _i\varphi _{i1})\mathrm{\Theta }(\varphi _{i1})\mathrm{\Theta }(\varphi _i)}{\mathrm{\Theta }(\varphi _{i1})\mathrm{\Theta }(\varphi _i)}=\frac{(\varphi _2\varphi _1)\mathrm{\Theta }(\varphi _1)\mathrm{\Theta }(\varphi _2)}{u}$$
(44)
where translational invariance is exploited again. The numerator in (44) can be obtained after performing the same basis transformation (Appendix A) that was essential to find $`u`$. Then after elementary integrations we find
$$K_{\mathrm{min}}=\frac{1}{u}\frac{1}{\sqrt{2\pi }}\frac{C_L^\varphi (0)C_L^\varphi (1)}{\sqrt{C_L^\varphi (0)}}=\frac{\sqrt{2\pi }}{\sqrt{C_L^\varphi (0)}}\frac{C_L^\varphi (0)C_L^\varphi (1)}{\mathrm{arccos}(C_L^\varphi (1)/C_L^\varphi (0))}$$
(45)
Using the explicit results for the slope correlation function for $`z=2`$ and $`z=4`$, and the scaling forms for it for higher $`z`$ given in the previous subsections, one can easily deduce the following. For $`z<5`$ the average curvature at the local minimum points on a lattice tends to a constant in the thermodynamic limit. For $`z=2`$
$$K_{\mathrm{min}}\frac{2\sqrt{2}}{\sqrt{\pi }}\sqrt{\frac{D}{\nu }}+𝒪\left(\frac{1}{L}\right),$$
(46)
and for $`z=4`$
$$K_{\mathrm{min}}\sqrt{\pi }\sqrt{\frac{D}{2\nu }}+𝒪\left(\frac{1}{L}\right).$$
(47)
The behavior of this quantity drastically changes for $`z>5`$, where it diverges with the system size as:
$$K_{\mathrm{min}}L^{\frac{z5}{2}}$$
(48)
.
## III Other lattice models and an application to parallel computing
### A The single-step model
In the single-step model the height differences (i.e., the local slopes) are restricted to $`\pm 1`$, and the evolution consists of particles of height $`2`$ being deposited at the local minima. While the full dynamic behavior of the model belongs to the KPZ universality class, in one dimension the steady state is governed by the EW Hamiltonian . Thus, the roughness exponent is $`\alpha =1/2`$, and we expect the finite-size effects for $`u`$ to follow eq. (28). The advantage of this model is that it can be mapped onto a hard-core lattice gas for which the steady-state probability distribution of the configurations is known exactly . This enables us to find arbitrary moments of the local minimum density operator. Since $`\varphi _i=\pm 1`$, it can be simly written as
$$u=\frac{1}{L}\underset{i=1}{\overset{L}{}}\frac{1\varphi _{i1}}{2}\frac{1+\varphi _i}{2}=\frac{1}{L}\underset{i=1}{\overset{L}{}}(1n_{i1})n_i,$$
(49)
where $`n_i=(1+\varphi _i)/2`$, corresponds to the hard core lattice gas occupation number. The constraint $`_{i=1}^L\varphi _i=0`$ translates to $`_{i=1}^Ln_i=L/2`$. Note that here $`u=(1n_{i1})n_i`$ is proportional to the average current. Knowing the exact steady-state probability distribution , one can easily find that
$$n_i=\frac{1}{2},n_in_j_{ij}=\frac{1}{4}\frac{L2}{L1}$$
(50)
Thus the exact finite-size effects for the local minimum density:
$$u_L=\frac{1}{4}\frac{L}{L1}=\frac{1}{4}+\frac{1}{4L}+𝒪(L^2),$$
(51)
in qualitative agreement with (27).
### B The Massively Parallel Exponential Update model
One of the most challenging areas in parallel computing is the efficient implementation of dynamic Monte-Carlo algorithms for discrete-event simulations on massively parallel architectures. As already mentioned in the Introduction, it has numerous practical applications ranging from magnetic systems (the discrete events are spin flips) to queueing networks ( the discrete events are job arrivals). A parallel architecture by definition contains (usually) a large number of processors, or processing elements (PE-s). During the simulation each processor has to tackle only a fraction of the full computing task (e.g., a specific block of spins), and the algorithm has to ensure through synchronization that the underlying dynamics is not altered. In a wide range of models the discrete events are Poisson arrivals. Since this stochastic process is reproducible (the sum of two Poisson processes is a Poisson process again with a new arrival frequency), the Poisson streams can be simulated simultaneously on each subsystem carried by each PE. As a consequence, the simulated time is local and random, incremented by exponentially distributed random variables on each PE. However, the algorithm has to ensure that causality accross the boundaries of the neighboring blocks is not violated. This requires a comparison between the neighboring simulated times, and waiting, if necessary (conservative approach). In the simplest scenario (one site/PE), this means that only those PEs will be allowed to attempt the update the state of the underlying site and increment their local time, where the local simulated time is a local minimum regarding the full simulated time horizon of the system, $`\{\tau _i\}`$, $`i=1,..,L`$ (for simplicity we consider a chain-like connectivity among the PE-s but connectivities of higher degree can be treated as well). One can in fact think of the time horizon as a fluctuating surface with height variable $`\tau _i`$. Other examples where the update attempts are independent Poisson arrivals include arriving calls in the wireless cellular network of a large metropolitan area , or the spin flip attempts in an Ising ferromagnet. This extremely robust parallel scheme was introduced by Lubachevsky, and it is applicable to a wide range of stochastic cellular automata with local dynamics where the discrete events are Poisson arrivals. The local random time increments is, in the language of the associated surface, equivalent to depositing random amounts of ‘material’ (with an exponential distribution) at the local minima of the surface, see Figure 2. This in fact defines a simple surface growth model which we shall refer to as ‘the massively parallel exponential update model’ (MPEU). The main concern about a parallel implementation is its efficiency. Since in the next time step only a fraction of PE-s will get updated, i.e., those that are in the local minima of the time horizon, while the rest are in idle, the efficiency is nothing but the average number of non-idling PE-s divided by the total number of PE-s ($`L`$), i.e., the average number of minima per unit length, or the minimum-point density, $`u`$. The fundamental question of the so called scalability arises: will the efficiency of the algorithm go to zero as the number of PE-s is increased ($`L\mathrm{}`$) indefinitely, or not? If the efficiency has a non-zero lower bound for $`L\mathrm{}`$ the algorithm is called scalable, and certainly this is the preferred type of scheme. Can one design in principle such efficient algorithms?
As mentioned in the Introduction, we know of one example that nature provides with an efficient algorithm for a very large number of processing elements: the human brain with its $`10^{11}`$ PE-s is the largest parallel computer ever built. Although the intuition suggests that indeed there are scalable parallel schemes, it has only been proved recently, see for details Ref. , by using the aforementioned analogy with the simple MPEU surface growth model. While the MPEU model exactly mimics the evolution of the simulated time-horizon, it can also be considered as a primitive model for ion sputtering of surfaces (etching dynamics): to see this, define a new height variable via $`h_i\tau _i`$, i.e. flip Figure 2 upside down. This means that instead of depositing material we have to take, ‘etch’, and this has to be done at the local maxima of the $`\{h_i\}`$ surface. In sputtering of surfaces by ion bombardement an incoming ion-projectile will most likely ‘break off’ a piece from the top of a mound instead from a valley, very similar to our ‘reversed’ MPEU model. It was shown that the sputtering process is described by the KPZ equation, . This qualitative argument is in complete agrrement with the extensive MC simulations and a coarse-grained approximation of Ref. that MPEU, similar to the single-step model, it also belongs to the KPZ dynamic universality class; in one dimension the macroscopic landscape is governed by the EW Hamiltonian.
The slope varaibles $`\varphi _i`$ for MPEU are not independent in the $`L\mathrm{}`$ limit, but short-ranged. This already guarantees that the steady-state behavior is governed by the EW Hamiltonian, and the density of local minima does not vanish in the thermodynamic limit. Our results confirm that the finite-size effects for $`u`$ follow eq. (28):
$$u_Lu_{\mathrm{}}+\frac{\mathrm{const}.}{L}$$
(52)
with $`u_{\mathrm{}}=0.24641(7)`$, see Fig. 3.
We conclude that the basic algorithm (one site per PE) is scalable for one-dimensional arrays. The same correspondence can be applied to model the performance of the algorithm for higher-dimensional logical PE topologies. While this will involve the typical difficulties of surface-growth modeling, such as an absence of exact results and very long simulation times, it establishes potentially fruitful connections between two traditionally separate research areas.
### C The larger curvature model
In this subsection we briefly present a curvature driven SOS surface deposition model known in the literature as the larger curvature model, and show a numerical analysis of the density of minima on this model. This model was originally introduced by Kim and Das Sarma and Krug independently, as an atomistic deposition model which fully conforms to the behaviour of the continuum fourth order linear Mullins equation ($`\nu =0`$, $`\kappa >0`$ in Eq. 4). Note that the discrete analysis we presented in Section II is based on the discretization of the continuum equation using the simplest forward Euler differencing scheme. The larger curvature model, however, is a growth model where the freshly deposited particles diffuse on the surface according to the rules of the model until they are embedded. Since in all the quantities studied so far, the correspondence (on the level of scaling) between the larger curvature model and the Langevin equation is very good, we would expect that the dynamic scaling properties of the density of minima for both the model and the equation to be identical.
The large curvature model has rather simple rules: a freshly deposited atom (let us say at site $`i`$) will be incorporated at the nearest neighbor site which has largest curvature (i.e., $`K_i=h_{i+1}+h_{i1}2h_i`$ is maximum). If there are more neighbors with the same maximum curvature, then one is chosen randomly. If the original site ($`i`$) is among those with maximum curvature, then the atom is incorporated at $`i`$.
Figure 4 shows the scaling of the density of minima $`<u>_L`$ in the steady state, vs. $`1/\sqrt{L}`$. According to Eq. (32), for the fourth order equation on the lattice, the behavior of the density of minima in the steady state scales with system size as $`1/\sqrt{L}`$. And indeed, Figure 4 shows the same behaviour for the larger curvature model, as expected. Note that this behavior sets in at rather small system sizes already, at about $`L=100`$, meaning that the finite system size effects are rather small for the larger curvature model. This is a very fortunate property since increasing the system size means decreasing the density of minima, therefore relative statistical errors will increase.
This can only be improved by better statistics, i.e., with averages over larger number of runs. This becomes however quickly a daunting task, since the cross-over time toward the steady state scales with system size as $`L^4`$. As we shall see in Section V.A, a matematically rigorous approach to the continuum equation yields the same $`1/\sqrt{L}`$ behaviour. Since the density of minima does decay to zero, an algorithm corresponding to the larger curvature model (or the Mullins equation) would not be asymptotically scalable.
Finally, we would like to make a brief note about the observed morphologies in the steady state for the Mullins equation , or the related models. It has been shown previously that in the steady state the morphology tipically shows a single large mound (or macroscopic groove). At first sight this may appear as a surprise, since we have shown that the number of minima (or maxima) diverges as $`\sqrt{L}`$ (the density vanishes as $`1/\sqrt{L}`$. There is however no contradiction, because that refers to a a mound that expands throughout the system, i.e. it is a long wavelenght structure, whereas the number of minima measures all the minima, and thus it is a short wavelenght characteristic. In the steady state we indeed have a single large, macroscopic groove, however, there are numerous small dips and humps generated by the constant coupling to the noise.
## IV Extremal-point densities on the continuum
Let us consider a continuous and at least two times differentiable function $`f:[0,L]`$. We are interested in counting the total number of extrema of $`f`$ in the $`[0,L]`$ interval. The topology of continuous curves in one dimension allows for three possibilities on the nature of a point $`x_i`$ for which $`f^{}|_{x_i}=0`$. Namely, $`x_i`$ is a local minimum if $`f^{\prime \prime }|_{x_i}>0`$, a local maximum if $`f^{\prime \prime }|_{x_i}<0`$ and it is degenerate if $`f^{\prime \prime }|_{x_i}=0`$. We call the point $`x_i`$ a degenerate flat of order $`k`$, if $`f^{(j)}|_{x_i}=0`$ for $`j=1,2,..,k`$ and $`f^{(k+1)}|_{x_i}0`$, $`k2`$, assuming that the higher order derivatives $`f^{(j)}`$ implied exist. The counter-like quantity
$$c(L,[f])\frac{1}{L}\underset{0}{\overset{L}{}}𝑑x|f^{\prime \prime }|\delta (f^{})$$
(53)
where $`\delta `$ is the Dirac-delta, gives the number of extremum points per unit length in the interval $`[0,L]`$, which in the limit of $`L0`$ is the extremum point density of $`f`$ in the origin. For our purposes $`L`$ will always be a finite number, however, for the sake of briefness we shall refer to $`c`$ simply as the density of extrema. Note that counting the extrema of a function $`f`$ is equivalent to counting the zeros of its derivate $`f^{}`$. The divergence of $`c`$ for finite $`L`$ implies either the existence of completely flat regions (infinitely degenerate), or an “infinitely wrinkled” region, such as for the truncated Weierstrass function shown in Fig. 2 ( in this latter case the divergence is understood by taking the limit $`M\mathrm{}`$). As already explained in the Introduction this infinitely wrinkled region does not necessarily imply that the curve is fractal, but if the curve is fractal, then regions of infinite wrinkledness must exist. The divergence or non-divergence of $`c`$ can be used as an indicator of the existence of such regions (completely flat or infinitely wrinkled).
One can make the following precise statement related to the counter $`c`$: if $`x_i`$ is an extremum point of $`f`$ of at most finite degeneracy $`k`$, and if there exist a small enough $`ϵ`$, such that $`f`$ is analytic in the neighborhood $`[x_iϵ,x_i+ϵ]`$, and there are no other extrema in this neighborhood, then
$$I(x_i)\underset{x_iϵ}{\overset{x_i+ϵ}{}}𝑑x|f^{\prime \prime }|\delta (f^{})=1,\mathrm{\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}0}<ϵ1$$
(54)
In the following we give a proof to this statement.
Using Taylor-series expansions around $`x_i`$, one writes:
$`f^{}(x)={\displaystyle \frac{a_k}{k!}}(xx_i)^k+{\displaystyle \frac{a_{k+1}}{(k+1)!}}(xx_i)^{k+1}+\mathrm{}`$ (55)
$`f^{\prime \prime }(x)={\displaystyle \frac{a_k}{(k1)!}}(xx_i)^{k1}+{\displaystyle \frac{a_{k+1}}{k!}}(xx_i)^k+\mathrm{}`$ (56)
where we introduced the shorthand notation $`a_jf^{(j+1)}|_{x_i}`$. For the non-degenerate case of $`k=1`$, (54) follows from a classical property of the delta function, namely:
$$\delta (g(x))=\underset{i}{}|g^{}(x_i)|^1\delta (xx_i),\text{ }x_i\text{ are }\text{simple}\text{ zeros of }g.$$
(57)
Let us now assume that $`x_i`$ is degenerate of order $`k`$ ($`k2`$). Using the expansions (55), (56), the variable change $`u=xx_i`$, and the well-known property $`\delta (ax)=|a|^1\delta (x)`$, we obtain:
$$I(x_i)=k\underset{ϵ}{\overset{ϵ}{}}𝑑u|u|^{k1}\left|1+\underset{j=1}{\overset{\mathrm{}}{}}\frac{(k1)!}{(k1+j)!}\frac{a_{k+j}}{a_k}u^j\right|\delta \left(|u|^k\left[1+\underset{j=1}{\overset{\mathrm{}}{}}\frac{k!}{(k+j)!}\frac{a_{k+j}}{a_k}u^j\right]\right)$$
(58)
Next we split the integral (58) in two: $`_ϵ^ϵ\mathrm{}=_ϵ^0\mathrm{}+_0^ϵ`$, make the variable change $`uu`$ in the first one, and then $`z=u^k`$ in both integrals. The final expression can then be written in the form:
$$I(x_i)=\underset{ϵ^k}{\overset{ϵ^k}{}}𝑑z|A(z)|\delta (zB(z)),$$
(59)
where
$$A(z)=1+\underset{j=1}{\overset{\mathrm{}}{}}\frac{(k1)!}{(k1+j)!}\frac{a_{k+j}}{a_k}z^j|z|^{\frac{j}{k}j},\text{and}B(z)=1+\underset{j=1}{\overset{\mathrm{}}{}}\frac{k!}{(k+j)!}\frac{a_{k+j}}{a_k}z^j|z|^{\frac{j}{k}j}$$
(60)
We have $`A(0)=B(0)=1`$, and
$$[zB(z)]^{}=1+\underset{j=1}{\overset{\mathrm{}}{}}\frac{k!}{(k+j)!}\frac{a_{k+j}}{a_k}(\frac{j}{k}+1)z^j|z|^{\frac{j}{k}j},[zB(z)]^{}|_{z=0}=1.$$
(61)
(Take the derivatives separately to the right and to the left of $`z=0`$). Thus, since $`z=0`$ is a simple zero of $`zB(z)`$, property (57) can be applied for sufficiently small $`ϵ`$:
$$I(x_i)=|A(0)|=1$$
(62)
proving our assertion. Note that because of (54), $`c`$ counts all the non-degenerate and the finitely degenerate points as well, giving the equal weight of unity to each. Can we count the non-degenerate extrema separately? The answer is affirmative, if one considers instead of (53) the following quantity:
$$c_q(L,[f])\frac{1}{L}\underset{0}{\overset{L}{}}𝑑x|f^{\prime \prime }|^{q+1}\delta (f^{}),q>0$$
(63)
Performing the same steps as above we obtain for a degenerate point:
$$I_q(x_i)\underset{x_iϵ}{\overset{x_i+ϵ}{}}𝑑x|f^{\prime \prime }|^{q+1}\delta (f^{})=\left[\frac{|a_k|}{(k1)!}\right]^q\underset{ϵ^k}{\overset{ϵ^k}{}}𝑑z|z|^{q\left(1\frac{1}{k}\right)}|A(z)|^{q+1}\delta (zB(z)).$$
(64)
Since $`k2`$, $`q\left(1\frac{1}{k}\right)\frac{1}{2}q>0`$, i.e.,
$$I_q(x_i)=0,\text{for }x_i\text{ degenerate}.$$
(65)
This means, that $`q>0`$ eliminates the degenerate points from the count. To non-degenerate points ($`k=1`$) (63) gives the weight of
$$I_q(x_i)=|a_1|^q=|f^{\prime \prime }|_{x_i}|^q,\text{for }x_i\text{ non-degenerate}.$$
(66)
In other words,
$$c_q(L,[f])=\frac{1}{L}\underset{i}{}\left|K(x_i)\right|^q,q>0\text{}x_i\text{ non-degenerate extrema of }f$$
(67)
where $`K(x)=f^{\prime \prime }`$ is the curvature of $`f`$ at $`x`$. The limit $`q0^+`$ in (67) gives the extremum point density $`\overline{c}(L,[f])`$ of $`f`$ of non-degenerate extrema:
$$\overline{c}(L,[f])=\underset{q0^+}{lim}c_q(L,[f])=\underset{q0^+}{lim}\frac{1}{L}\underset{0}{\overset{L}{}}𝑑x|f^{\prime \prime }|^{q+1}\delta (f^{})$$
(68)
It is important to note, that taking the $`q0^+`$ limit in (67) is not equivalent to taking $`q=0`$ in (63), i.e., the limit and the integral on the rhs of (68) are not interchangeable! The difference is the set of degenerate points!
Until now, we did not make any distiction between maxima and minima. In a natural way, we expect that the quantity:
$$u(L,[f])\frac{1}{L}\underset{0}{\overset{L}{}}𝑑xf^{\prime \prime }\delta (f^{})\theta (f^{\prime \prime })$$
(69)
where $`\theta (x)`$ is the Heaviside step-function, will give the density of minima (due to the step function, here we can drop the absolute values). However, performing a similar derivation as above, one concludes that (69) is a little bit ill-defined, in the sense that the weight given to degenerate points depends on the definition of the step-function in the origin (however, $`u(L,[f])`$ is bounded). Introducing a $`qregulator`$ as above, the weight of degenerate points is pulled down to zero:
$$u_q(L,[f])\frac{1}{L}\underset{0}{\overset{L}{}}𝑑x[f^{\prime \prime }]^{q+1}\delta (f^{})\theta (f^{\prime \prime }),q>0.$$
(70)
and
$$u_q(L,[f])=\frac{1}{L}\underset{i}{}\left[K(x_i)\right]^q,q>0\text{}x_i\text{ non-degenerate minima of }f$$
(71)
Note that in the equation above the absolute values are not needed, since we are summing over the curvatures of all local minima. The density $`\overline{u}(L,[f])`$ of non-degenerate minima of $`f`$ is obtained after taking the limit $`q0^+`$:
$$\overline{u}(L,[f])=\underset{q0^+}{lim}u_q(L,[f])$$
(72)
and the limit is not interchangeable with the integral in (70). To obtain densities for maxima, one only has to replace the argument $`f^{\prime \prime }`$ of the Heaviside function with $`f^{\prime \prime }`$.
### A Stochastic extremal-point densities
We are interested to explore the previously introduced quantities for a stochastic function, subject to time evolution, $`h(x,t)`$. This function may be for example the solution to a Langevin equation. We define the two basic quantities in the same way as before, except that now one performs a stochastic average over the noise, as well:
$`C_q(L,t)={\displaystyle \frac{1}{L}}{\displaystyle \underset{0}{\overset{L}{}}}𝑑x\left|{\displaystyle \frac{^2h}{x^2}}\right|^{q+1}\delta \left({\displaystyle \frac{h}{x}}\right),\text{and}U_q(L,t)={\displaystyle \frac{1}{L}}{\displaystyle \underset{0}{\overset{L}{}}}𝑑x\left[{\displaystyle \frac{^2h}{x^2}}\right]^{q+1}\delta \left({\displaystyle \frac{h}{x}}\right)\theta \left({\displaystyle \frac{^2h}{x^2}}\right)`$ (73)
For systems preserving translational invariance, the stochastic average of the integrand becomes $`x`$-independent, and the integrals can be dropped:
$`C_q(L,t)=\left|{\displaystyle \frac{^2h}{x^2}}\right|^{q+1}\delta \left({\displaystyle \frac{h}{x}}\right)`$ (74)
$`U_q(L,t)=\left[{\displaystyle \frac{^2h}{x^2}}\right]^{q+1}\delta \left({\displaystyle \frac{h}{x}}\right)\theta \left({\displaystyle \frac{^2h}{x^2}}\right)`$ (75)
According to (67) and (71), $`C_q(L,t)`$ and $`U_q(L,t)`$ can be thought of as time dependent “partition functions” for the non-degenerate extremal-point densities of the underlying stochastic process, with $`q`$ playing the role of “inverse temperature”:
$`C_q(L,t)={\displaystyle \frac{1}{L}}{\displaystyle \underset{i}{}}\left|K(x_i)\right|^q,q>0\text{}x_i\text{ non-degenerate extrema}`$ (76)
$`U_q(L,t)={\displaystyle \frac{1}{L}}{\displaystyle \underset{i}{}}\left[K(x_i)\right]^q,q>0\text{}x_i\text{ non-degenerate minima}`$ (77)
It is important to mention that in the above equations the average $`<..>`$ and the summation are not interchangeable: particular realizations of $`h`$ have particular sets of minima.
Two values for $`q`$ are of special interest: when $`q0^+`$ and $`q=1`$. In the first case we obtain the stochastic average of the density of non-degenerate extrema and minima:
$$\overline{C}(L,t)=\underset{q0^+}{lim}C_q(L,t),\text{and}\overline{U}(L,t)=\underset{q0^+}{lim}U_q(L,t),$$
(78)
and in the second case we obtain the stochastic average of the mean curvature at extrema and minima:
$$\overline{K}_{ext}(L,t)=\frac{C_1(L,t)}{\overline{C}(L,t)},\text{and}\overline{K}_{min}(L,t)=\frac{U_1(L,t)}{\overline{U}(L,t)},$$
(79)
(we need to normalize with the number of extrema/minima per unit length to get the curvature per extremum/minimum).
In the following we explore the quantities (74)-(79) for a large class of linear Langevin equations. To simplify the calculations, we will assume that $`q`$ is a positive integer. Then we will attempt analytic continuation on the final result as a function of $`q`$. In the calculations we will make extensive use of the standard integral representations of the delta and step functions:
$`\delta (y)={\displaystyle \underset{\mathrm{}}{\overset{\mathrm{}}{}}}{\displaystyle \frac{dz}{2\pi }}e^{izy}={\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{\mathrm{}}{\overset{\mathrm{}}{}}}{\displaystyle \frac{dz}{2\pi }}{\displaystyle \frac{(iz)^n}{n!}}y^n,`$ (80)
$`\theta (y)=\underset{ϵ0^+}{lim}{\displaystyle \underset{\mathrm{}}{\overset{\mathrm{}}{}}}{\displaystyle \frac{dz}{2\pi }}{\displaystyle \frac{e^{izy}}{ϵ+iz}}=\underset{ϵ0^+}{lim}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{\mathrm{}}{\overset{\mathrm{}}{}}}{\displaystyle \frac{dz}{2\pi }}{\displaystyle \frac{1}{ϵ+iz}}{\displaystyle \frac{(iz)^n}{n!}}y^n`$ (81)
If $`q`$ is a positive integer, we may drop the absolute value signs in (75). In (74) we can only do that for odd $`q`$. The absolute values make the calculation of stochastic averages very difficult. We can get around this problem by employing the following identity:
$$|y|^n=y^n\left\{(1)^n+\theta (y)\left[1(1)^n\right]\right\}$$
(82)
This brings (74) to
$$C_q(L,t)=\left[1(1)^q\right]U_q(L,t)+(1)^{q+1}B_q(L,t)$$
(83)
where
$$B_q(L,t)=\frac{1}{L}\underset{0}{\overset{L}{}}𝑑x\left(\frac{^2h}{x^2}\right)^{q+1}\delta \left(\frac{h}{x}\right)=\left(\frac{^2h}{x^2}\right)^{q+1}\delta \left(\frac{h}{x}\right)$$
(84)
Obviously for $`q`$ odd integer, $`B_q=C_q`$. For $`q`$ even, $`B_q`$ is an interesting quantity by itself. In this case the weight of an extremum $`x_i`$ is $`sgn(K(x_i))|K(x_i)|^q`$. If the analytic continuation can be performed, then the $`q0^+`$ limit will tell us if there are more non-degenerate maxima than minima (or otherwise) in average. Using the integral representations (80) and (81):
$`B_q(L,t)={\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{\mathrm{}}{\overset{\mathrm{}}{}}}{\displaystyle \frac{dz}{2\pi }}{\displaystyle \frac{(iz)^n}{n!}}\left({\displaystyle \frac{^2h}{x^2}}\right)^{q+1}\left({\displaystyle \frac{h}{x}}\right)^n`$ (85)
$`U_q(L,t)=\underset{ϵ0^+}{lim}{\displaystyle \underset{n_1=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{n_2=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{\mathrm{}}{\overset{\mathrm{}}{}}}{\displaystyle \frac{dz_1}{2\pi }}{\displaystyle \frac{(iz_1)^{n_1}}{n_1!}}{\displaystyle \underset{\mathrm{}}{\overset{\mathrm{}}{}}}{\displaystyle \frac{dz_2}{2\pi }}{\displaystyle \frac{(iz_2)^{n_2}}{n_2!}}{\displaystyle \frac{1}{ϵ+iz_2}}\left({\displaystyle \frac{^2h}{x^2}}\right)^{n_2+q+1}\left({\displaystyle \frac{h}{x}}\right)^{n_1}`$ (86)
## V Extremal-point densities of linear stochastic evolution equations
Next we calculate the densities (85), (86) for the following type of linear stochastic equations:
$$\frac{h}{t}=\nu \left(^2\right)^{z/2}h+\eta (x,t),\nu ,D,z>0,x[0,L]$$
(87)
with initial condition $`h(x,0)=0`$, for all $`x[0,L]`$. $`\eta `$ is a white noise term drawn from a Gaussian distribution with zero mean $`\eta (x,t)=0`$, and covariance :
$$\eta (x,t)\eta (x^{},t^{})=2D\delta (xx^{})\delta (tt^{}),$$
(88)
We also performed our calculations with other noise types, such as volume conserving and long-range correlated, however the details are to lengthy to be included in the present paper, it will be the subject of a future publication. As boundary condition we choose periodic boundaries:
$`h(x+nL,t)=h(x,t),\eta (x+nL,t)=\eta (x,t),\text{for all }n`$ (89)
The general solution to (87) is obtained simply with the help of Fourier series . The Fourier series and its coefficients for a function $`f`$ defined on $`[0,L]`$ is
$`f(x)={\displaystyle \underset{k}{}}\stackrel{~}{f}(k)e^{ikx},\stackrel{~}{f}(k)={\displaystyle \frac{1}{L}}{\displaystyle \underset{L}{\overset{L}{}}}𝑑xf(x)e^{ikx}`$ (90)
where $`k=\frac{2\pi }{L}n`$, $`n=..,2,1,0,1,2,..`$. The Fourier coefficients of the general solution to (87) are:
$$\stackrel{~}{h}(k,t)=\underset{0}{\overset{t}{}}𝑑t^{}e^{\nu |k|^z(tt^{})}\stackrel{~}{\eta }(k,t^{})$$
(91)
The correlations of the noise in momentum space are:
$$\stackrel{~}{\eta }(k,t)\stackrel{~}{\eta }(k^{},t^{})=\frac{2D}{L}\delta _{k,k^{}}\delta (tt^{}).$$
(92)
Due to the Gaussian character of the noise, the two-point correlation of the solution (91) is also delta-correlated and it completely characterizes the statistical properties of the stochastic dynamics (87). It is given by:
$$\stackrel{~}{h}(k,t)\stackrel{~}{h}(k^{},t^{})=S(k,t)\delta _{k,k^{}}$$
(93)
where $`S(k,t)`$ is the structure factor given by:
$$S(k,t)=\frac{D}{\nu L|k|^z}\left[1e^{2\nu |k|^zt}\right].$$
(94)
Equation (87) has been analyzed in great detail by a number of authors, see Ref. for a review. It was shown that there exist un upper critical dimension $`d_c=z`$ for the noisy case of Eq. (87) which separates the rough regime with $`d<z`$ from the non-roughening regime $`d>z`$. In one dimension, the rough regime corresponds to the condition $`z>1`$, which we shall assume from now on, since this is where the interesting physics lies.
Next, we evaluate the quantities (74)-(79) via directly calculating the expressions in (85) and (86). This amounts to computing averages of type:
$$Q_{N,M}=\left(\frac{^2h}{x^2}\right)^N\left(\frac{h}{x}\right)^M$$
(95)
Expressing $`h`$ with its Fourier series according to (90), we write:
$`\left({\displaystyle \frac{h}{x}}\right)^M=i^M{\displaystyle \underset{k_1}{}}\mathrm{}{\displaystyle \underset{k_M}{}}k_1\mathrm{}k_M\stackrel{~}{h}(k_1,t)\mathrm{}\stackrel{~}{h}(k_M,t)e^{i(k_1+\mathrm{}+k_M)x}`$ (96)
$`\left({\displaystyle \frac{^2h}{x^2}}\right)^N=(1)^N{\displaystyle \underset{k_1^{}}{}}\mathrm{}{\displaystyle \underset{k_N^{}}{}}k_{1}^{}{}_{}{}^{2}\mathrm{}k_{N}^{}{}_{}{}^{2}\stackrel{~}{h}(k_1^{},t)\mathrm{}\stackrel{~}{h}(k_N^{},t)e^{i(k_1^{}+\mathrm{}+k_N^{})x}`$ (97)
which then is inserted in (95). Thus in Fourier space one needs to calculate averages of type $`\stackrel{~}{h}(k_1,t)\mathrm{}\stackrel{~}{h}(k_M,t)\stackrel{~}{h}(k_1^{},t)\mathrm{}\stackrel{~}{h}(k_N^{},t)`$. According to (93) $`\stackrel{~}{h}`$ is anti-delta-correlated, therefore these averages can be performed in the standard way which is by taking all the possible pairings of indices and employing (93). In our case there are three types of pairings: $`\{k_j,k_l\}`$, $`\{k_j,k_l^{}\}`$, and $`\{k_j^{},k_l^{}\}`$. Let us pick a ‘mixed’ pair $`\{k_j,k_l^{}\}`$ containing a primed and a non-primed index. The corresponding contribution in the $`Q_{N,M}`$ will be:
$$\underset{k_j}{}\underset{k_l^{}}{}k_jk_{l}^{}{}_{}{}^{2}S(k_l^{},t)e^{i(k_j+k_l^{})x}\delta _{k_j,k_l^{}}$$
(98)
Since the structure factor $`S(k,t)`$ is an even function in $`k`$, (98) becomes $`_{k_j}k_{j}^{}{}_{}{}^{3}S(k_j,t)=0`$, because the summand is an odd function of $`k_j`$ and the summation is symmetric around zero. Thus, it is enough to consider non-mixed index-pairs, only. This means, that $`Q_{N,M}`$ decouples into:
$$Q_{N,M}=\left(\frac{^2h}{x^2}\right)^N\left(\frac{h}{x}\right)^M$$
(99)
The averages are calculated easily, and we find:
$`\left({\displaystyle \frac{h}{x}}\right)^M=\{\begin{array}{c}(M1)!!\left[\varphi _2(L,t)\right]^{M/2},\text{for }M\text{ even },\hfill \\ \\ 0,\text{for }M\text{ odd }\hfill \end{array}`$ (103)
and
$`\left({\displaystyle \frac{^2h}{x^2}}\right)^N=\{\begin{array}{c}(N1)!!\left[\varphi _4(L,t)\right]^{N/2},\text{for }N\text{ even },\hfill \\ \\ 0,\text{for }M\text{ odd }\hfill \end{array}`$ (107)
where
$$\varphi _m(L,t)\underset{k}{}|k|^mS(k,t)$$
(108)
Employing (103), and (107) in (85), it follows that if $`q`$ is an even integer, $`q=2s`$, $`s=1,2,..`$:
$$B_{2s}(t)=0,s=1,2,\mathrm{}$$
(109)
whereas for $`q`$ odd integer, $`q=2s1`$, $`s=1,2,..`$:
$$B_{2s1}(t)=C_{2s1}(t)=\frac{2^{s\frac{1}{2}}}{\pi }\mathrm{\Gamma }\left(s+\frac{1}{2}\right)\frac{\left[\varphi _4(L,t)\right]^s}{\sqrt{\varphi _2(L,t)}},s=1,2,\mathrm{}$$
(110)
where we used the identity $`2^p(2p1)!!/(2p)!=1/p!`$, and performed the Gaussian integral.
The calculation of $`U_q`$ is a bit trickier. The sum over $`n_1`$ in (86) is easy and leads to the Gaussian $`e^{\varphi _2(L,t)z_1^2/2}`$. However, the sum over $`n_2`$ is more involved. Let us make the temporary notation for the sum over $`n_2`$:
$$R_q=\underset{n_2=0}{\overset{\mathrm{}}{}}\frac{(iz_2)^{n_2}}{n_2!}(2r1)!!\left[\varphi _4(L,t)\right]^r,n_2+q+1=2r$$
(111)
We have to distinguish two cases according to the parity of $`q`$:
1) $`q`$ is odd, $`q=2s1`$, $`s=1,2,..`$. In this case $`R_q`$ becomes
$$R_{2s1}=\underset{r=s}{\overset{\mathrm{}}{}}\frac{(iz_2)^{2(rs)}}{[2(rs)]!}\frac{(2r)!}{r!}\left[\frac{1}{2}\varphi _4(L,t)\right]^r=(z_2)^{2s}(1)^s\left\{\frac{^{2s}}{x^{2s}}\left[e^{\varphi _4(L,t)z_2^2x^2/2}\right]\right\}_{x=1}$$
(112)
The Hermite polynomials are defined via the Rodrigues formula as:
$$H_n(x)=(1)^ne^{x^2}\frac{d^n}{dx^n}\left(e^{x^2}\right)$$
(113)
Using this, we can express $`R_{2s1}`$ with the help of Hermite polynomials:
$$R_{2s1}=(1)^s\left[\frac{1}{2}\varphi _4(L,t)\right]^sH_{2s}\left(\sqrt{\frac{1}{2}\varphi _4(L,t)}z_2\right)e^{\varphi _4(L,t)z_2^2x^2/2}$$
(114)
2) $`q`$ is even, $`q=2s`$, $`s=1,2,..`$. The calculations are analogous to the odd case:
$$R_{2s}=\underset{r=s+1}{\overset{\mathrm{}}{}}\frac{(iz_2)^{2(rs)1}}{[2(rs)1]!}\frac{(2r)!}{r!}\left[\frac{1}{2}\varphi _4(L,t)\right]^r=(iz_2)^{2s1}\left\{\frac{^{2s+1}}{x^{2s+1}}\left[e^{\varphi _4(L,t)z_2^2x^2/2}\right]\right\}_{x=1}$$
(115)
or via Hermite polynomials:
$$R_{2s}=i(1)^s\left[\frac{1}{2}\varphi _4(L,t)\right]^{s+\frac{1}{2}}H_{2s+1}\left(\sqrt{\frac{1}{2}\varphi _4(L,t)}z_2\right)e^{\varphi _4(L,t)z_2^2x^2/2}$$
(116)
In order to obtain $`U_q`$ we have to do the integral over $`z_2`$ in (86). This can be obtained after using the formula:
$$\underset{\mathrm{}}{\overset{\mathrm{}}{}}𝑑x(x\pm ic)^\nu H_n(x)e^{x^2}=2^{n1\nu }\sqrt{\pi }\frac{\mathrm{\Gamma }\left(\frac{n\nu }{2}\right)}{\mathrm{\Gamma }(\nu )}e^{\pm \frac{i\pi }{2}(\nu +n)},c0^+.$$
(117)
Finally, the densities for the minima read as:
$`U_{2s}(t)={\displaystyle \frac{2^{s1}}{\pi }}\mathrm{\Gamma }(s+1){\displaystyle \frac{\left[\varphi _4(L,t)\right]^{s+\frac{1}{2}}}{\sqrt{\varphi _2(L,t)}}}`$ (118)
$`U_{2s1}(t)={\displaystyle \frac{2^{s\frac{3}{2}}}{\pi }}\mathrm{\Gamma }\left(s+{\displaystyle \frac{1}{2}}\right){\displaystyle \frac{\left[\varphi _4(L,t)\right]^s}{\sqrt{\varphi _2(L,t)}}}`$ (119)
Formulas (109), (110), (118), and (119) combined with (83) can be condensed very simply, and we obtain the general result as:
$`U_q(L,t)={\displaystyle \frac{2^{\frac{q}{2}1}}{\pi }}\mathrm{\Gamma }\left({\displaystyle \frac{q}{2}}+1\right){\displaystyle \frac{\left[\varphi _4(L,t)\right]^{\frac{q+1}{2}}}{\sqrt{\varphi _2(L,t)}}}`$ (120)
$`C_q(L,t)=2U_q(L,t)`$ (121)
Equations (120), (121) with together with (119) fully solve the problem for the density of non-degenerate extrema. Eq. (121) is an expected result in one dimension, because Eq. (87) preserves the up-down symmetry. The density of non-degenerate minima is:
$$\overline{U}(L,t)=\underset{q0^+}{lim}U_q(L,t)=\frac{1}{2\pi }\sqrt{\frac{\varphi _4(L,t)}{\varphi _2(L,t)}}$$
(122)
and the stochastic average of the mean curvature at a minimum point is:
$$\overline{K}(L,t)=\frac{U_1(L,t)}{\overline{U}(L,t)}=\sqrt{\frac{\pi }{2}}\sqrt{\varphi _4(L,t)}$$
(123)
i.e., the average curvature at a minimum is proportional to the square root of the fourth moment of the structure factor. In the following section we exploit the physical information behind the above expressions for the stochastic process (87). At some parameter values a few, or all the quantities above may diverge. In this case we introduce a microscopic lattice cut-off $`0<a1`$, and analyze the limit $`a0^+`$ in the final formulas. This in fact corresponds to placing the whole problem on a lattice with lattice constant $`a`$. It has been shown in Ref. that for the class of equations (87) there are three important length-scales that govern the statistical behavior of the interface $`h`$: the lattice constant $`a`$, the system size $`L`$, and the dynamical correlation length $`\xi `$ defined by:
$$\xi (t)(2\nu t)^{1/z}$$
(124)
According to (108) and (94) the function $`\varphi _m(L,t)`$ becomes:
$$\varphi _m(L,t)=\frac{2D}{\nu L}\underset{n=0}{\overset{\mathrm{}}{}}\left(\frac{2\pi n}{L}\right)^{mz}\left[1e^{\left(\xi {\scriptscriptstyle \frac{2\pi n}{L}}\right)^z}\right],m=2,4$$
(125)
The $`n=0`$ term can be dropped from the sum above, because it is zero even for $`m<z`$ (expand the exponential and then take $`n=0`$). However, the whole sum may diverge depending on $`m`$ and $`z`$. In order to handle all the cases, including the divergent ones we introduce a microscopic lattice cut-off $`a`$, $`0<a1`$, and then analyze the limit $`a0^+`$ in the final expressions. This is in fact equivalent to putting the whole problem on a lattice of lattice spacing $`a`$. Appropiately, (125) becomes:
$$\varphi _m(L,t)=\frac{2D}{\nu L}\underset{n=1}{\overset{\frac{L}{2a}}{}}\left(\frac{2\pi n}{L}\right)^{mz}\left[1e^{\left(\xi {\scriptscriptstyle \frac{2\pi n}{L}}\right)^z}\right],m=2,4$$
(126)
### A Steady-state regime.
Putting $`\xi =\mathrm{}`$ in (126) $`\varphi _m`$ takes a simpler form:
$$\varphi _m(L,\mathrm{})=\frac{2D}{\nu L}\left(\frac{2\pi }{L}\right)^{mz}\underset{n=1}{\overset{\frac{L}{2a}}{}}n^{mz},m=2,4$$
(127)
As $`a0^+`$, $`\varphi _m`$ becomes proportional to $`\zeta (zm)`$. For $`zm>1`$ $`\varphi _m`$ is convergent, otherwise it is divergent. In the divergent case we quote the following results:
$$\underset{n=1}{\overset{N}{}}n^s=\mathrm{ln}N+𝒞+𝒪(1/N),\text{if}s=1$$
(128)
and
$$\underset{n=1}{\overset{N}{}}n^s=\frac{N^{s+1}}{s+1}\left[1+𝒪(1/N)\right],\text{if}s>1$$
(129)
which we will use to derive the leading behaviour of the extremal point densities when $`L/a\mathrm{}`$. From equations (127), (120), (122) and (123) follows:
$$U_q(L,\mathrm{})=\mathrm{\Gamma }\left(\frac{q}{2}+1\right)\left(\frac{2D}{\pi \nu }\right)^{\frac{q}{2}}\left(2\pi \right)^{\frac{q}{2}(5z)}L^{1\frac{q}{2}(5z)}\left[\underset{n=1}{\overset{L/2a}{}}n^{4z}\right]^{\frac{q+1}{2}}\left[\underset{n=1}{\overset{L/2a}{}}n^{2z}\right]^{\frac{1}{2}},$$
(130)
$$\overline{U}(L,\mathrm{})=\frac{1}{L}\sqrt{\underset{n=1}{\overset{L/2a}{}}n^{4z}\left(\underset{n=1}{\overset{L/2a}{}}n^{2z}\right)^1},$$
(131)
and
$$\overline{K}(L,\mathrm{})=\sqrt{\frac{D}{2}\left(2\pi \right)^{5z}L^{z5}\underset{n=1}{\overset{L/2a}{}}n^{4z}}.$$
(132)
The convergency (divergency) properties of the sums in Eqs. (130-132) for $`a0^+`$ generate two critical values for $`z`$, namely $`z=3`$ and $`z=5`$. In the three regions separated by these values we obtain qualitatively different behaviors for the extremal-point densities.
i) $`z>5`$. All quantities are convergent as $`a0^+`$. We have:
$$U_q(L,\mathrm{})=\mathrm{\Gamma }\left(\frac{q}{2}+1\right)\left(\frac{2D}{\pi \nu }\right)^{\frac{q}{2}}(2\pi )^{\frac{q}{2}(5z)}\frac{[\zeta (z4)]^{\frac{q+1}{2}}}{[\zeta (z2)]^{\frac{1}{2}}}L^{1+\frac{q}{2}(z5)},z>5$$
(133)
$$\overline{U}(L,\mathrm{})=\frac{1}{L}\sqrt{\frac{\zeta (z4)}{\zeta (z2)}},$$
(134)
$$\overline{K}(L,\mathrm{})=(2\pi )^{\frac{5z}{2}}\sqrt{\frac{D}{2\nu }\zeta (z4)}L^{\frac{z5}{2}},$$
(135)
Eq. (134) shows that there are a finite number of minima ($`\sqrt{\zeta (z4)/\zeta (z2)}`$) in the steady state, independently of the system size $`L`$. ($`\overline{U}(L,\mathrm{})`$ is the number of minima per unit length, and $`L\overline{U}(L,\mathrm{})`$ is the number of minima on the substrate of size $`L`$). The mean curvature $`\overline{K}(L,\mathrm{})`$ diverges with system size as $`L^{(z5)/2}`$. This is consistent with the fact that the system size grows as $`L`$, the width grows as $`L^{(z1)/2}`$, i.e., faster than $`L`$, and thus the peaks and minima should become sleeker and sharper as $`L\mathrm{}`$, expecting diverging curvatures in minima and maxima. However, this is not always true, since the sleekness of the humps and mounds does not necessarily imply large curvatures in minima and maxima if the shape of the humps also changes as $`L`$ changes, i.e., there is lack of self-affinity. The existence of $`z=5`$ as a critical value is a non-trivial results coming from the presented analysis.
ii) $`z=5`$. According to (128), $`\varphi _4(L,\mathrm{})`$ diverges logarithmically as $`a0^+`$. One obtains:
$$U_q(\mathrm{})\mathrm{\Gamma }\left(\frac{q}{2}+1\right)\left(\frac{2D}{\nu \pi }\right)^{\frac{q}{2}}\frac{1}{\sqrt{\zeta (3)}}\frac{1}{L}\left(\mathrm{ln}\frac{L}{2a}+𝒞\right)^{\frac{q+1}{2}},$$
(136)
$$\overline{U}(L,\mathrm{})=\frac{1}{\sqrt{\zeta (3)}}\frac{1}{L}\sqrt{\mathrm{ln}\frac{L}{2a}+𝒞},$$
(137)
$$\overline{K}(L,\mathrm{})=\sqrt{\frac{D}{2\nu }\left(\mathrm{ln}\frac{L}{2a}+𝒞\right)}.$$
(138)
Eq. (137) shows that the although the density of minima vanishes, the number of minima is no longer a constant but diverges logarithmically with system size $`L`$. The mean curvature still diverges, but logarithmically, when compared to the power law divergence of (135).
For the mean curvature $`\overline{K}(\mathrm{})`$ in (132) $`z=5`$ is the only critical value, since it only depends on $`\varphi _4`$. For $`z<5`$, using Eq. (129) we arrive to the result that the mean curvature in a minimum point approaches to an $`L`$-independent constant for $`L/a\mathrm{}`$ with corrections on the order of $`a/L`$:
$$\overline{K}(L,\mathrm{})\left(\frac{\pi }{a}\right)^{\frac{5z}{2}}\sqrt{\frac{D}{2\nu (5z)}},z<5$$
(139)
We arrived to the same conclusion in Section II.D when we studied the steady state of the discretized version of the continuum equation. Coincidentally, for $`z=4`$ the two constant values from (139) and (47) are identical ($`a=1`$ by definition in (47).
iii) $`3<z<5`$. In this case $`\varphi _4(L,\mathrm{})\mathrm{}`$ and $`\varphi _2(L,\mathrm{})<\mathrm{}`$ as $`a0^+`$, and:
$$U_q(L,\mathrm{})\mathrm{\Gamma }\left(\frac{q}{2}+1\right)\left(\frac{2D}{\nu \pi }\right)^{\frac{q}{2}}\left(\frac{\pi }{a}\right)^{\frac{q}{2}(5z)}\left(\frac{1}{2a}\right)^{\frac{5z}{2}}\frac{L^{\frac{z3}{2}}}{(5z)^{\frac{q+1}{2}}\sqrt{\zeta (z2)}},$$
(140)
and
$$\overline{U}(L,\mathrm{})\left(\frac{1}{2a}\right)^{\frac{5z}{2}}\frac{L^{\frac{z3}{2}}}{\sqrt{(5z)\zeta (z2)}},$$
(141)
and the mean curvature is just given by (139).
Comparing Eqs. (133), (136), and (140) we can make an interesting observation: while for $`z5`$ the dependence on the system size $`L`$ is coupled to the ‘inverse temperature’ $`q`$, for $`3<z<5`$ the dependence on $`L`$ decouples from $`q`$, i.e., it becomes independent of the inverse temperature! Eq. (141) shows that the density of minima vanishes with system size as a power law with an exponent $`(z3)/2`$ but the number of minima of the substrate diverges as a power law with an exponent of $`(5z)/2`$.
iv) $`z=3`$. In this case $`\varphi _4(L,\mathrm{})\mathrm{}`$ and $`\varphi _2(L,\mathrm{})\mathrm{}`$ logarithmically as $`a0^+`$. One obtains:
$$U_q(L,\mathrm{})\frac{1}{2\sqrt{2}a}\mathrm{\Gamma }\left(\frac{q}{2}+1\right)\left(\frac{\pi D}{\nu a^2}\right)^{\frac{q}{2}}\frac{1}{\sqrt{\mathrm{ln}\frac{L}{2a}+𝒞}},$$
(142)
and
$$\overline{U}(L,\mathrm{})\frac{1}{2\sqrt{2}a}\frac{1}{\sqrt{\mathrm{ln}\frac{L}{2a}+𝒞}},$$
(143)
with a logarithmically vanishing density of minima, and the dependence on the system size in (142) is not coupled to $`q`$.
v) $`1<z<3`$. Now both $`\varphi _4`$ and $`\varphi _2`$ diverge as $`a0^+`$. Employing (129), yields:
$$U_q(L,\mathrm{})\frac{1}{2a}\mathrm{\Gamma }\left(\frac{q}{2}+1\right)\left(\frac{2D}{\pi \nu }\right)^{\frac{q}{2}}\left(\frac{\pi }{a}\right)^{\frac{q}{2}(5z)}\frac{\sqrt{3z}}{(5z)^{\frac{q+1}{2}}},$$
(144)
and
$$\overline{U}(L,\mathrm{})\frac{1}{2a}\sqrt{\frac{3z}{5z}}.$$
(145)
Note, that in leading order, both $`U_q(L,\mathrm{})`$ and the density of minima $`\overline{U}(L,\mathrm{})`$ become system size independent! The system size dependence comes in as corrections on the order of $`a/L`$ and higher. The fact that the efficiency of the massively parallel algorithm presented in Section III.B is not vanishing is due precisely to the above phenomenon: the fluctuations of the time horizon in the steady state belong to the $`z=2`$ class (Edwards-Wilkinson universality), and according to the results under iv), the density of minima (or the efficiency of the parallel algorithm) converges to a non-zero constant, as $`L\mathrm{}`$, ensuring the scalability of the algorithm. An algorithm that would map into a $`z3`$ class would have a vanishing efficiency with increasing the number of processing elements. In particular, for $`z=2`$, one obtains from (145) $`\overline{U}(L,\mathrm{})(a2\sqrt{3})^1=0.2886\mathrm{}/a`$. Note that the utilization we obtained is somewhat different from the discrete case which was $`0.25`$. This is due to the fact that this number is non-universal and it may show differences depending on the discretization scheme used, however it cannot be zero.
Another important conclusion can be drawn from the final results enlisted above: at and below $`z=5`$, all the quantities diverge when $`a0^+`$, and keep $`L`$ fixed. This means that the higher the resolution the more details we find in the morphology, just as for an infinitely wrinkled, or a fractal-like surface. We call this transition accross $`z=5`$ a ‘wrinkle’ transition. As shown in the Introduction, wrinkledness can assume two phases depending on whether the curve is a fractal or not and the transition between these two pases may be conceived as a phase transition. However, one may be able to scale the system size $`L`$ with $`a`$ such that the quantities calculated will not diverge in this limit. This is possible only in the regime $`3<z<5`$, when we impose:
$$L^{z3}a^{(q+1)(5z)}=const.$$
(146)
This shows that the rescaling cannot be done for all inverse temperatures $`q`$ at the same time. In particular, for the density of minima and $`z=4`$, $`La=const`$.
### B Scaling regime
In order to obtain the temporal behavior of the extremal-point densities we will use the Poisson summation formula (B4) from Appendix B on (126). After simple changes of variables in the integrals this leads to:
$`\varphi _m(L,t)={\displaystyle \frac{D}{\nu L}}\left({\displaystyle \frac{\pi }{a}}\right)^{mz}\left[1e^{\left(\xi \frac{\pi }{a}\right)^z}\right]+{\displaystyle \frac{D}{\pi \nu }}\xi ^{(mz+1)}{\displaystyle \underset{0}{\overset{\pi \xi /a}{}}}𝑑xx^{mz}\left(1e^{x^z}\right)+`$ (147)
$`{\displaystyle \frac{2D}{\pi \nu }}\xi ^{(mz+1)}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \underset{0}{\overset{\pi \xi /a}{}}}𝑑xx^{mz}\mathrm{cos}\left({\displaystyle \frac{L}{\xi }}nx\right)\left(1e^{x^z}\right)`$ (148)
This expression shows, that the scaling properties of the dynamics are determined by the dimensionless ratios $`L/\xi `$ and $`\xi /a`$. The scaling regime is defined by $`a\xi L`$.
As we have seen in the previous section, $`\varphi _m`$ is convergent for $`z>m+1`$ but diverges when $`zm+1`$, as $`a0^+`$. In the convergent case, the lattice spacing $`a`$ can be taken as zero, and thus the first term on the rhs. of (148) vanishes and the time dependence of the infinite system-size piece of $`\varphi _m`$ (the first integral term in (148)) assumes the clean power-law behaviour of $`t^{(zm1)/z}`$ (with a positive exponent). In the divergent case, however, the non-integral term of (148) does not vanish, and the time-dependence will not be a clean power-law. Even the integral terms will present corrections to the power-law $`t^{(m+1z)/z}`$ (which has now a negative exponent), since the limits for integration contain $`\xi `$. The first intergal on the rhs of (148) for $`zm+1`$ can be calculated exactly:
$$\underset{0}{\overset{\pi \xi /a}{}}𝑑xx^{mz}\left(1e^{x^z}\right)=\frac{1}{mz+1}\left(\frac{\pi \xi }{a}\right)^{mz+1}\frac{1}{z}\mathrm{\Gamma }\left(\frac{mz+1}{z}\right)+\frac{1}{z}\mathrm{\Gamma }(\frac{mz+1}{z},\left(\frac{\pi \xi }{a}\right)^z),zm+1$$
(149)
where $`\mathrm{\Gamma }(\alpha ,x)`$ is the incomplete Gamma function. In our case $`(\pi \xi /a)^z`$ is a large number, and therefore we can use the asymptotic representation of $`\mathrm{\Gamma }(\alpha ,x)`$ for large $`x`$, see , pp. 951, equation 8.357. According to this, for large $`x`$, $`\mathrm{\Gamma }(\alpha ,x)x^{\alpha 1}e^x`$, i.e., it can become arbitrarily small, with an exponential decay. This term can therefore be neglected from (149), compared to the other two terms, even in the divergent case. Inserting (149) into (148), we will see that also the non-integral piece of (148) can be neglected compared to the term generated by the first on the rhs of (149), since in the scaling regime $`a\xi L`$, and thus the ratio $`a/L`$ can be neglected compared to $`(mz+1)^1`$. (This is needed only in the divergent regime, $`z<m+1`$.) Thus, one obtains:
$$\varphi _m(L,t)\frac{D}{\pi \nu (mz+1)}\left(\frac{\pi }{a}\right)^{mz+1}\frac{D}{\pi \nu z}\mathrm{\Gamma }\left(\frac{mz+1}{z}\right)\xi ^{zm1}\left[1E_m(\frac{L}{\xi },\frac{\xi }{a})\right],zm+1$$
(150)
where
$$E_m(\lambda ,\rho )=\frac{2z}{\mathrm{\Gamma }\left(\frac{mz+1}{z}\right)}\underset{n=1}{\overset{\mathrm{}}{}}\underset{0}{\overset{\pi \rho }{}}𝑑xx^{mz}\mathrm{cos}(\lambda nx)\left(1e^{x^z}\right),zm+1$$
(151)
The oscillating terms condensed in $`E_m`$ will give the finite-size corrections, as long as $`L/\xi 1`$.
The $`z=m+1`$ case (divergent) can also be calculated, however, instead of (149) now we have:
$$\underset{0}{\overset{\pi \xi /a}{}}\frac{dx}{x}\left(1e^{x^z}\right)=\mathrm{ln}\left(\frac{\pi \xi }{a}\right)+\frac{𝒞}{z}\frac{1}{z}\text{Ei}\left(\left(\frac{\pi \xi }{a}\right)^z\right),z=m+1.$$
(152)
where $`\text{Ei}(x)`$ is the exponential integral function. According to the large-$`x`$ expansion of the exponential integral function, see , pp. 935, equation 8.215, $`\text{Ei}(x)x^1e^x`$, it is vanishing exponentially fast, thus it can be neglected in the expression of $`\varphi _m`$ in the scaling limit:
$$\varphi _m(L,t)\frac{D}{\pi \nu }\mathrm{ln}\left(\frac{\pi \xi }{a}\right)+\frac{D𝒞}{\pi \nu z}+\frac{D}{\pi \nu }F_m(\frac{L}{\xi },\frac{\xi }{a}),z=m+1$$
(153)
where
$$F_m(\lambda ,\rho )=2\underset{n=1}{\overset{\mathrm{}}{}}\underset{0}{\overset{\pi \rho }{}}\frac{dx}{x}\mathrm{cos}(\lambda nx)\left(1e^{x^z}\right),z=m+1$$
(154)
Thus in the scaling limit, the temporal behavior of $`\varphi _m`$ becomes a logarithmic time dependence plus a constant, as long as $`L/\xi 1`$.
Observe that for $`z<m+1`$ the first term on the rhs. of (150), reproduces exactly the diverging term (as $`a0^+`$) of the steady-state expression (127) which can be seen after employing (128) in (127). This means that for $`\xi \mathrm{}`$, $`E_m(L/\xi ,\xi /a)`$ diverges slower than $`\xi ^{m+1z}`$ (this is how the saturation occurs). Similarly, for $`z=m+1`$ the first term on the rhs of (153) (after replacing $`\xi `$ with $`L`$ ) reproduces the diverging term (as $`a0^+`$) of the steady-state expression (127) which can be seen after employing (129) in (127). This means, that in the saturation (or steady-state) regime the remaining terms from (153) must behave as $`const.+𝒪(a/L)+\mathrm{ln}(L/\xi )`$, as $`\xi \mathrm{}`$ while keeping $`L`$ and $`a`$ fixed.
Just as in the case of steady-state one has to distinguish 5 cases depending on the values of $`z`$, with respect to the critical values 3 and 5. For the sake of simplicity of writing, we will omit the arguments of $`E_m(\lambda ,\rho )`$ and $`F_m(\lambda ,\rho )`$.
i) $`z>5`$. We have:
$$\varphi _m(L,t)=\frac{D}{\pi \nu (zm1)}\mathrm{\Gamma }\left(\frac{m+1}{z}\right)\xi ^{zm1}(1E_m),m=2,4$$
(155)
From Eqs. (120), (122) and (123), it follows:
$$U_q(L,t)=\frac{1}{2\pi }\mathrm{\Gamma }\left(\frac{q}{2}+1\right)\left(\frac{2D}{\pi \nu }\right)^{\frac{q}{2}}\left[\frac{\mathrm{\Gamma }\left(\frac{5}{z}\right)}{z5}\right]^{\frac{q+1}{2}}\left[\frac{z3}{\mathrm{\Gamma }\left(\frac{3}{z}\right)}\right]^{\frac{1}{2}}[\xi (t)]^{1\frac{q}{2}(z5)}\frac{(1E_4)^{\frac{q+1}{2}}}{(1E_2)^{\frac{1}{2}}},$$
(156)
$$\overline{U}(L,t)=\frac{1}{2\pi }\sqrt{\frac{(z3)\mathrm{\Gamma }\left(\frac{5}{z}\right)}{(z5)\mathrm{\Gamma }\left(\frac{3}{z}\right)}}\left[\xi (t)\right]^1\sqrt{\frac{1E_4}{1E_2}},$$
(157)
and
$$\overline{K}(L,t)=\sqrt{\frac{D\mathrm{\Gamma }\left(\frac{5}{z}\right)}{2\nu (z5)}}\left[\xi (t)\right]^{\frac{z5}{2}}\sqrt{1E_4}$$
(158)
and therefore the time-behaviour is a clean power-law: $`U_q(L,t)`$ decays as $`t^{[2+q(z5)]/2z}`$, $`\overline{U}(L,t)t^{1/z}`$, and $`\overline{K}(L,t)`$ diverges as $`t^{(z5)/2z}`$, for $`L/\xi 1`$.
ii) $`z=5`$. In this case $`\varphi _4`$ takes the form (153) but $`\varphi _2`$ is still given by (150). The quantities of interest become:
$$U_q(L,t)\frac{\mathrm{\Gamma }\left(\frac{q}{2}+1\right)}{2\pi }\left(\frac{2D}{\pi \nu }\right)^{\frac{q}{2}}\sqrt{\frac{2}{\mathrm{\Gamma }\left(\frac{3}{5}\right)}}\xi ^1\left[\mathrm{ln}\left(\frac{\pi \xi }{a}\right)\right]^{\frac{q+1}{2}}\frac{\left\{1+\left[\mathrm{ln}\left(\frac{\pi \xi }{a}\right)\right]^1\left(\frac{𝒞}{5}+F_4\right)\right\}^{\frac{q+1}{2}}}{\sqrt{1E_2}},$$
(159)
$$\overline{U}(L,t)\frac{1}{2\pi }\sqrt{\frac{2}{\mathrm{\Gamma }\left(\frac{3}{5}\right)}}\xi ^1\sqrt{\mathrm{ln}\left(\frac{\pi \xi }{a}\right)}\sqrt{\frac{1+\left[\mathrm{ln}\left(\frac{\pi \xi }{a}\right)\right]^1\left(\frac{𝒞}{5}+F_4\right)}{\sqrt{1E_2}}},$$
(160)
and
$$\overline{K}(L,t)=\sqrt{\frac{D}{2\nu }}\sqrt{\mathrm{ln}\left(\frac{\pi \xi }{a}\right)}\sqrt{1+\left[\mathrm{ln}\left(\frac{\pi \xi }{a}\right)\right]^1\left(\frac{𝒞}{5}+F_4\right)}$$
(161)
One can observe that the leading temporal behaviour has logarithmic component due to the borderline situation: $`U_q(L,t)`$ decays as $`t^{1/5}(\mathrm{ln}t)^{(q+1)/2}`$, $`\overline{U}(L,t)t^{1/5}(\mathrm{ln}t)^{1/2}`$, and $`\overline{K}(L,t)`$ diverges as $`(\mathrm{ln}t)^{1/2}`$.
iii) $`3<z<5`$.
$$U_q(L,t)\frac{\mathrm{\Gamma }\left(\frac{q}{2}+1\right)}{2\pi }\left(\frac{2D}{\pi \nu }\right)^{\frac{q}{2}}\sqrt{\frac{z3}{\mathrm{\Gamma }\left(\frac{3}{z}\right)}}(5z)^{\frac{q+1}{2}}\left(\frac{\pi }{a}\right)^{\frac{q+1}{2}(5z)}\xi ^{\frac{z3}{2}}\frac{\left[1\mathrm{\Gamma }\left(\frac{5}{z}\right)\left(\frac{a}{\pi \xi }\right)^{5z}(1E_4)\right]^{\frac{q+1}{2}}}{\sqrt{1E_2}},$$
(162)
$$\overline{U}(L,t)\frac{1}{2\pi }\sqrt{\frac{z3}{(5z)\mathrm{\Gamma }\left(\frac{3}{z}\right)}}\left(\frac{\pi }{a}\right)^{\frac{5z}{2}}\xi ^{\frac{z3}{2}}\sqrt{\frac{1\mathrm{\Gamma }\left(\frac{5}{z}\right)\left(\frac{a}{\pi \xi }\right)^{5z}(1E_4)}{1E_2}},$$
(163)
$$\overline{K}(L,t)\sqrt{\frac{D}{2\nu (5z)}}\left(\frac{\pi }{a}\right)^{\frac{5z}{2}}\sqrt{1\mathrm{\Gamma }\left(\frac{5}{z}\right)\left(\frac{a}{\pi \xi }\right)^{5z}(1E_4)}$$
(164)
An important conclusion that can be drawn from these expressions is that below $`z=5`$, the leading time-dependence of the partition function $`U_q(L,t)`$ becomes independent of the inverse temperature $`q`$ and it presents a clean power-law decay $`t^{(z3)/2z}`$ which is the same also for $`\overline{U}(L,t)`$. In particular, for $`z=4`$ this means a $`t^{1/8}`$ decay which is very well verified by the larger curvature model from Section III.C, see Figure 5. Also notice from Eq. (163) that the leading term is system size independent. And indeed, this property is also in a very good agreement with the numerics on the larger curvature model from Figure 5, where the two data sets for $`L=100`$ and $`L=120`$ practically coincide.
Since the mean curvature depends on $`\varphi _4`$, only, for all cases below $`z=5`$ the dependence is given by the same formula (164) (just need to replace the corresponding value for $`z`$).
iv) $`z=3`$. This is another borderline situation, the corresponding expressions are found easily:
$$U_q(L,t)\frac{\mathrm{\Gamma }\left(\frac{q}{2}+1\right)}{2\sqrt{2}\pi }\left(\frac{2D}{\pi \nu }\right)^{\frac{q}{2}}\left(\frac{\pi }{a}\right)^{q+1}\left[\mathrm{ln}\left(\frac{\pi \xi }{a}\right)\right]^{\frac{1}{2}}\frac{\left[1\mathrm{\Gamma }\left(\frac{5}{3}\right)\left(\frac{a}{\pi \xi }\right)^2(1E_4)\right]^{\frac{q+1}{2}}}{\sqrt{1+\left[\mathrm{ln}\left(\frac{\pi \xi }{a}\right)\right]^1\left(\frac{𝒞}{3}+F_2\right)}},$$
(165)
$$\overline{U}(L,t)\frac{1}{2\sqrt{2}a}\left[\mathrm{ln}\left(\frac{\pi \xi }{a}\right)\right]^{\frac{1}{2}}\sqrt{\frac{1\mathrm{\Gamma }\left(\frac{5}{3}\right)\left(\frac{a}{\pi \xi }\right)^2(1E_4)}{1+\left[\mathrm{ln}\left(\frac{\pi \xi }{a}\right)\right]^1\left(\frac{𝒞}{3}+F_2\right)}}$$
(166)
and the leading time dependences are: $`U_q(L,t)(\mathrm{ln}t)^{1/2}`$, $`\overline{U}(L,t)(\mathrm{ln}t)^{1/2}`$.
v) $`1<z<3`$.
$$U_q(L,t)\frac{\mathrm{\Gamma }\left(\frac{q}{2}+1\right)}{2\sqrt{2}\pi }\left(\frac{2D}{\pi \nu }\right)^{\frac{q}{2}}\sqrt{\frac{3z}{(5z)^{q+1}}}\left(\frac{\pi }{a}\right)^{1+\frac{q}{2}(5z)}\frac{\left[1\mathrm{\Gamma }\left(\frac{5}{z}\right)\left(\frac{a}{\pi \xi }\right)^{5z}(1E_4)\right]^{\frac{q+1}{2}}}{\sqrt{1\mathrm{\Gamma }\left(\frac{3}{z}\right)\left(\frac{a}{\pi \xi }\right)^{3z}(1E_2)}},$$
(167)
$$\overline{U}(L,t)\frac{1}{2a}\sqrt{\frac{3z}{(5z)}}\sqrt{\frac{1\mathrm{\Gamma }\left(\frac{5}{z}\right)\left(\frac{a}{\pi \xi }\right)^{5z}(1E_4)}{1\mathrm{\Gamma }\left(\frac{3}{z}\right)\left(\frac{a}{\pi \xi }\right)^{3z}(1E_2)}},$$
(168)
In this case the partition function and the density of minima all converge to a constant which in leading order is independent of the system size. The density of minima was shown in Section II to have this property in the steady-state. Here we see not only that but also the fact that all $`q`$-moments show the same behavior, and even more, the time behavior before reaching the steady-state constant is not a clean power-law, but rather a decaying correction in the approach to this constant. The leading term in the temporal correction is of $`t^{(3z)/z}`$ and the next-to-leading has $`t^{(5z)/z}`$.
## VI Conclusions and outlook
In summary, based on the analytical results presented, a short wavelength based analysis of interface fluctuations can provide us with novel type of information and give an alternative description of surface morphologies. This analysis gives a more detailed characterization and can be used to distinguish interfaces that are ‘fuzzy’ from those that locally appear to be smooth, and the central quantities, the extremal-point densities are numerically and analytically accessible. The partition function-like formalism enables us to access a wide range of $`q`$-momenta of the local curvatures distribution. In the case of the stochastic evolution equations studied we could exactly relate these $`q`$-momenta to the structure function of the process via the simple quantities $`\varphi _2`$ and $`\varphi _4`$. The wide spectrum of results accessed through this technique shows the richness of short wavelength physics. This physics is there, and the long wavelength approach just simply cannot reproduce it, but instead may suggest an oversimplified picture of the reality. For example, the MPEU model has been shown to belong in the steady state to the EW universality class, however, it cannot be described exactly by the EW equation in all respects, not even in the steady-state! For example, the utilization (or density of minima) of the MPEU model is 0.24641 which for the EW model on a lattice is 0.25. Also, if one just simply looks at the steady-state configuration, one observes high skewness for the MPEU model , whereas the EW is completely up-down symmetric. This can also be shown by comparing the calculated two-slope correlators. For a number of models that belong to the KPZ equation universality class, this broken-symmetry property vs. the EW case has been extensively investigated by Neergard and den Nijs . The difference on the short wavelength scale between two models that otherwise belong to the same universality class lies in the existence of irrelevant operators (in the RG sense). Although these operators do not change universal properties, the quantities associated with them may be of very practical interest. The parallel computing example shows that the fundamental question of algorithmic scalability is answered based on the fact that the simulated time horizon in the steady state belongs to the EW universality class, thus it has a finite density of local minima. The actual value of the density of local minima in the thermodynamic limit, however, strongly depends on the details of the microscopics, which in principle can be described in terms of irrelevant operators .
The extremal-point densities introduced in the present paper may actually have a broader application than stochastic surface fluctuations. The main geometrical characterization of fractal curves is based on the construction of their Haussdorff-Besikovich dimension, or the ‘box-counting’ dimension: one covers the set with small boxes of linear size $`ϵ`$ and then track the divergence of the number of boxes needed to cover in a minimal way the whole set as $`ϵ`$ is lowered to zero. For example, a smooth line in the plane has a dimension of unity, but the Weierstrass curve of (2) has a dimension of $`\mathrm{ln}b/\mathrm{ln}a`$ (for $`b>a`$). The actual length of a fractal curve whose dimension is larger than unity will diverge when $`ϵ0^+`$. The total length at a given resolution $`ϵ`$ is a global property of the fractal, it does not tell us about the way ‘it curves’. The novel measure we propose in (3) is meant to characterize the distribution of a local property of the curve, its bending which in turn is a measure of the curve’s wrinkledness. For simplicity we formulated it for functions, i.e., for curves which are single-valued in a certain direction. This can be remedied and generalized by introducing a parametrization $`\gamma [0,1]`$ of the curve, and then plotting the local curvature vs. this parameter $`K(\gamma )`$. The plot will be a single valued function on which now (3) is easily defined.
Other desirable extensions of the present technique are: 1) to include a statistical description of the degeneracies of higher order, and 2) to repeat the analysis for higher (such as $`d=2`$) substrate dimensions. The latter is promising an even richer spectrum of novelties, since in higher dimensions there is a plethora of singular points ($`f=\mathrm{𝟎}`$) which are classified by the eigenvalues of the Hessian matrix of the function in the singular point. Deciphering the statistical behaviour of these various singularities for randomly evolving surfaces is an interesting challenge. The studies performed by Kondev and Henley on the distribution of contours on random Gaussian surfaces should come to a good aid in achieving this goal. In particular we may find the method developed here useful in studying the spin-glass ground state, and the spin-glass transition problem. And at last but not the least, we invite the reader to consider instead of the Langevin equations studied here, noisy wave equations, with a second derivative of the time component, or other stochastic evolution equations.
## Acknowledgements
We thank S. Benczik, M.A. Novotny, P.A. Rikvold, Z. Rácz, B. Schmittmann, T. Tél, E.D. Williams, and I. Žutić for stimulating discussions. This work was supported by NSF-MRSEC at University of Maryland, by DOE through SCRI-FSU, and by NSF-DMR-9871455.
## A $`\mathrm{\Theta }(x_1)\mathrm{\Theta }(x_2)`$ for general coupled Gaussian variables
The expression we derive in this appendix, despite its simplicity, is probably the most important one concerning the extremal-point densities of one-dimensional Gaussian interfaces on a lattice. If the correlation matrix for two possibly coupled Gaussian variables is given by
$`x_1^2=x_1^2`$ $`=`$ $`d>0`$ (A1)
$`x_1x_2`$ $`=`$ $`c`$ (A2)
then the distribution follows as
$$P(x_1,x_2)=\frac{1}{2\pi \sqrt{𝒟}}\mathrm{exp}\left\{\frac{1}{2𝒟}\left(dx_1^2+dx_2^22cx_1x_2\right)\right\}=\frac{1}{2\pi \sqrt{𝒟}}\mathrm{exp}\left\{\frac{d}{2𝒟}\left(x_1^2+x_2^22\frac{c}{d}x_1x_2\right)\right\},$$
(A3)
where $`𝒟d^2c^2>0`$. We aim to find the average of the stochastic variable $`u=\mathrm{\Theta }(x_1)\mathrm{\Theta }(x_2)`$:
$$u=\mathrm{\Theta }(x_1)\mathrm{\Theta }(x_2)=_{\mathrm{}}^{\mathrm{}}_{\mathrm{}}^{\mathrm{}}𝑑x_1𝑑x_2\mathrm{\Theta }(x_1)\mathrm{\Theta }(x_2)P(x_1,x_2)$$
(A4)
which is simply the total weight of the density $`P(x_1,x_2)`$ in the $`x_1<0`$, $`x_2>0`$ quadrant. If $`c=0`$, the density is isotropic, and $`u=1/4`$. In the general case it is convenient to find a new set of basis vectors, where the probability density is isotropic (of course the shape of the original quadrant will transform accordingly). Introducing the following linear transformation
$`x_1`$ $`=`$ $`\sqrt{{\displaystyle \frac{𝒟}{2}}}\left({\displaystyle \frac{y_1}{\sqrt{d+c}}}+{\displaystyle \frac{y_2}{\sqrt{dc}}}\right)`$ (A5)
$`x_2`$ $`=`$ $`\sqrt{{\displaystyle \frac{𝒟}{2}}}\left({\displaystyle \frac{y_1}{\sqrt{d+c}}}+{\displaystyle \frac{y_2}{\sqrt{dc}}}\right),`$ (A6)
and exploiting that $`\mathrm{\Theta }(\lambda x)=\mathrm{\Theta }(x)`$ for $`\lambda >0`$ we have
$$u=_{\mathrm{}}^{\mathrm{}}_{\mathrm{}}^{\mathrm{}}𝑑y_1𝑑y_2\mathrm{\Theta }\left(\frac{y_1}{\sqrt{d+c}}\frac{y_2}{\sqrt{dc}}\right)\mathrm{\Theta }\left(\frac{y_1}{\sqrt{d+c}}+\frac{y_2}{\sqrt{dc}}\right)\frac{1}{2\pi }\mathrm{exp}\left\{\frac{1}{2}(y_1^2+y_2^2)\right\}$$
(A7)
Now the probability density for the new variables, $`y_1,y_2`$, is isotropic, and $`u=\theta /(2\pi )`$, where $`\theta `$ is the angle enclosed by the following two unit vectors:
$$𝐯_1=\frac{1}{\sqrt{2d}}\left(\begin{array}{c}\hfill \sqrt{d+c}\\ \hfill \sqrt{dc}\end{array}\right),𝐯_2=\frac{1}{\sqrt{2d}}\left(\begin{array}{c}\hfill \sqrt{d+c}\\ \hfill \sqrt{dc}\end{array}\right).$$
(A8)
From their dot product one obtains
$$\mathrm{cos}(\theta )=\frac{𝐯_1𝐯_2}{|𝐯_1||𝐯_2|}=\frac{c}{d}.$$
(A9)
and, thus, for $`u`$:
$$u=\frac{1}{2\pi }\mathrm{arccos}\left(\frac{c}{d}\right).$$
(A10)
## B Poisson summation formulas
In this Appendix we recall the well-known Poisson summation formula and adapt it for functions with finite support in $``$. In the theory of generalized functions the following identity is proven:
$$\underset{m=\mathrm{}}{\overset{\mathrm{}}{}}\delta (xm)=\underset{m=\mathrm{}}{\overset{\mathrm{}}{}}e^{2\pi imx}$$
(B1)
Let $`f:[\alpha ,\beta ]`$ be a continuous function with continuous derivative on the interval $`[\alpha ,\beta ]`$. Multiply Eq. (B1) on both sides with $`f(x)`$, then integrate both sides from $`\alpha `$ to $`\beta `$. In the evaluation of the left hand side we have to pay attention to whether any, or both the numbers $`\alpha `$ and $`\beta `$ are integers, or non-integers. In the integer case the contribution of the end-point is calculated via the identity:
$$\underset{n}{\overset{n+r}{}}𝑑x\delta (xn)f(x)=\frac{1}{2}f(n),r>0$$
(B2)
Assuming that $`f`$ is absolutely integrable if $`\beta =\mathrm{}`$, and choosing $`\alpha =0`$, the classical Poisson summation formula is obtained:
$$\underset{n=0}{\overset{\mathrm{}}{}}f(n)=\frac{1}{2}f(0)+\underset{0}{\overset{\mathrm{}}{}}𝑑xf(x)+2\underset{m=1}{\overset{\mathrm{}}{}}\underset{0}{\overset{\mathrm{}}{}}𝑑xf(x)\mathrm{cos}(2\pi mx)$$
(B3)
Let us write also explicitely out the case when both $`\alpha `$ and $`\beta `$ are integers:
$$\underset{n=\alpha }{\overset{\beta }{}}f(n)=\frac{1}{2}[f(\alpha )+f(\beta )]+\underset{\alpha }{\overset{\beta }{}}𝑑xf(x)+2\underset{m=1}{\overset{\mathrm{}}{}}\underset{\alpha }{\overset{\beta }{}}𝑑xf(x)\mathrm{cos}(2\pi mx),\text{when}\alpha ,\beta $$
(B4)
Next we apply these equations to give an exact closed expression for the slope correlation function for finite $`L`$ \[eq. (18)\]:
$$C_L^\varphi (l)=\frac{D}{L}\underset{n=1}{\overset{L1}{}}\frac{e^{i\left(\frac{2\pi n}{L}\right)l}}{\nu +2\kappa \left[1\mathrm{cos}\left(\frac{2\pi n}{L}\right)\right]}$$
(B5)
where $`l\{0,1,2,..,L1\}`$, $`\nu ,\kappa ^+`$. Let us denote
$$a=\frac{2\kappa }{\nu +2\kappa }.$$
(B6)
We have $`|a|<1`$, and
$$C_L^\varphi (l)=\frac{Da}{2\kappa L}\underset{n=1}{\overset{L1}{}}\frac{e^{i\left(\frac{2\pi n}{L}\right)l}}{1a\mathrm{cos}\left(\frac{2\pi n}{L}\right)}$$
(B7)
In order to apply the Poisson summation formula (B4), we introduce the function:
$$f(x)=\frac{a}{2\kappa L}\underset{n=1}{\overset{L1}{}}\frac{e^{i\left(\frac{2\pi x}{L}\right)l}}{1a\mathrm{cos}\left(\frac{2\pi x}{L}\right)},\mathrm{\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}1}xL1$$
(B8)
and identify in (B4) $`\alpha 1`$ and $`\beta L1`$. The non-integral terms of (B4) give:
$$\frac{1}{2}[f(1)+f(L1)]=\frac{a}{2\kappa L}\frac{\mathrm{cos}\left(\frac{2\pi }{L}l\right)}{1a\mathrm{cos}\left(\frac{2\pi }{L}l\right)}$$
(B9)
The next term becomes:
$$\underset{1}{\overset{L1}{}}𝑑xf(x)=\frac{a}{2\kappa \sqrt{1a^2}}\left(\frac{1\sqrt{1a^2}}{a}\right)^l\frac{a}{2\pi \kappa }\underset{0}{\overset{2\pi /L}{}}𝑑x\frac{\mathrm{cos}xl}{1a\mathrm{cos}x}$$
(B10)
where during the evaluation of the integral we made a simple change of variables and used a well-known integral from random walk theory , :
$$\underset{\pi }{\overset{\pi }{}}𝑑x\frac{e^{ixl}}{1a\mathrm{cos}x}=\frac{2\pi }{\sqrt{1a^2}}\left(\frac{1\sqrt{1a^2}}{a}\right)^l,l0$$
(B11)
The sum over the integrals in (B4) can also be evaluated, and one obtains:
$$2\underset{n=1}{\overset{\mathrm{}}{}}\underset{1}{\overset{L1}{}}𝑑xf(x)\mathrm{cos}(2\pi nx)=\frac{a\left(b^l+b^l\right)}{2\kappa \sqrt{1a^2}}\frac{b^L}{1b^L}\frac{a}{2\pi \kappa }\underset{n=1}{\overset{\mathrm{}}{}}\underset{2\pi /L}{\overset{2\pi /L}{}}𝑑x\mathrm{cos}(nLx)\frac{e^{ilx}}{1a\mathrm{cos}x}$$
(B12)
where
$$b=\frac{1\sqrt{1a^2}}{a},\text{and}|b|<1$$
(B13)
To compute the sum on the rhs of (B12) we recall another identity from the theory of generalized functions (see Ref. , page 155):
$$\underset{n=1}{\overset{\mathrm{}}{}}e^{inx}=\pi \underset{m=\mathrm{}}{\overset{\mathrm{}}{}}\delta (x2m\pi )+\frac{i}{2}ctg\left(\frac{x}{2}\right)\frac{1}{2}$$
(B14)
Combining (B14) and identity (B1), one obtains:
$$\underset{n=1}{\overset{\mathrm{}}{}}\mathrm{cos}(nx)=\pi \underset{m=\mathrm{}}{\overset{\mathrm{}}{}}\delta (x2m\pi )+\frac{1}{2}$$
(B15)
Peforming the sum over $`n`$ directly in the rhs of (B12) via (B15), yields:
$$\frac{a}{2\pi \kappa }\underset{n=1}{\overset{\mathrm{}}{}}\underset{2\pi /L}{\overset{2\pi /L}{}}𝑑x\frac{\mathrm{cos}(nLx)e^{ilx}}{1a\mathrm{cos}x}=\frac{a}{2\kappa L}\underset{m=\mathrm{}}{\overset{\mathrm{}}{}}\underset{2\pi }{\overset{2\pi }{}}𝑑y\frac{e^{ily}\delta (y2m\pi )}{1a\mathrm{cos}y}+\frac{a}{2\pi \kappa }\underset{0}{\overset{2\pi /L}{}}𝑑x\frac{\mathrm{cos}lx}{1a\mathrm{cos}x}$$
(B16)
Only $`m=\pm 1,0`$ contribute in (B16). With the help of (B2):
$$\frac{a}{2\pi \kappa }\underset{n=1}{\overset{\mathrm{}}{}}\underset{2\pi /L}{\overset{2\pi /L}{}}𝑑x\frac{\mathrm{cos}(nLx)e^{ilx}}{1a\mathrm{cos}x}=\frac{a}{2\kappa L}\left\{\frac{1}{1a}+\frac{\mathrm{cos}\left(\frac{2\pi }{L}l\right)}{1a\mathrm{cos}\left(\frac{2\pi }{L}\right)}\right\}+\frac{a}{2\pi \kappa }\underset{0}{\overset{2\pi /L}{}}𝑑x\frac{\mathrm{cos}lx}{1a\mathrm{cos}x}$$
(B17)
Using (B17) in (B12), we can add the result to the rest of the contributions (B9) and (B10) to obtain the final expression \[eq. (19)\] after the cancellations. |
warning/0002/quant-ph0002013.html | ar5iv | text | # 1 Families of Nonlinear Schrödinger Equations
## 1 Families of Nonlinear Schrödinger Equations
About nine years ago, H.-D. Doebner and I introduced a certain family of nonlinear Schrödinger equations. We were led to these equations not by any prior inclination to study nonlinear quantum mechanics, but by our desire to interpret quantum-mechanically a class of representations of an infinite-dimensional, nonrelativistic current algebra, and the corresponding group . We proposed these equations as candidates for describing quantum systems with dissipation.
To review the development briefly, we sought self-adjoint representations of the infinite-dimensional Lie algebra of densities and currents, given at arbitrary time $`t`$ by
$$[\rho _{op}(f_1),\rho _{op}(f_2)]=\mathrm{\hspace{0.17em}0},[\rho _{op}(f),J_{op}(𝐠)]=i\mathrm{}\rho _{op}(𝐠f),$$
$$[J_{op}(𝐠_1),J_{op}(𝐠_2)]=i\mathrm{}J_{op}([𝐠_1,𝐠_2]),$$
(1)
where the $`f`$’s are real-valued $`C^{\mathrm{}}`$ functions on the physical space $`𝐑^n`$, the $`𝐠`$’s are $`C^{\mathrm{}}`$ vector fields on $`𝐑^n`$, and $`[𝐠_1,𝐠_2]=𝐠_1𝐠_2𝐠_2𝐠_1`$ is the usual Lie bracket . The $`N`$-particle Bose or Fermi representations of (1) may be written
$`\rho _{op}^N(f)\psi ^{(s,a)}(𝐱_1,\mathrm{}𝐱_N)=m{\displaystyle \underset{j=1}{\overset{N}{}}}f(𝐱_j)\psi ^{(s,a)}(𝐱_1,\mathrm{}𝐱_N),`$
$`J_{op}^N(𝐠)\psi ^{(s,a)}(𝐱_1,\mathrm{}𝐱_N)={\displaystyle \frac{\mathrm{}}{2i}}{\displaystyle \underset{j=1}{\overset{N}{}}}\{𝐠(𝐱_j)_j\psi ^{(s,a)}(𝐱_1,\mathrm{}𝐱_N)`$
$`+_j[𝐠(𝐱_j)\psi ^{(s,a)}(𝐱_1,\mathrm{}𝐱_N)]\},`$ (2)
where the $`\psi ^{(s,a)}`$ are (respectively) symmetric or antisymmetric square-integrable functions of the $`N`$ particle coordinate variables. There exists a family of related but unitarily inequivalent representations of (1), parameterized by the real number $`D`$, leading to physically distinct quantizations :
$$J_{op}^{N,D}(𝐠)=J_{op}^N(𝐠)+D\rho _{op}^N(𝐠).$$
(3)
Here $`D`$ is a constant with the dimensions of a diffusion coefficient. Even in the case of one-particle quantum mechanics, interpreting these representations posed a challenge.
In the usual notation for operator-valued distributions, write (suppressing the superscripts) $`\rho _{op}(f)=_X\rho _{op}(𝐱)f(𝐱)𝑑𝐱`$ and $`J_{op}(𝐠)=_X𝐉_{op}(𝐱)𝐠(𝐱)𝑑𝐱`$. Then, for a single particle at time $`t`$, take the expectation values $`m\rho (𝐱,t)=\psi _t|\rho _{op}(𝐱)|\psi _t`$ and $`m𝐣(𝐱,t)=\psi _t|𝐉_{op}(𝐱)|\psi _t`$. When $`D=0`$ the usual expressions are recovered for the probability density and flux in the Schrödinger representation:
$$\rho =\overline{\psi }\psi ,𝐣=\frac{\mathrm{}}{2mi}[\overline{\psi }\psi (\overline{\psi })\psi ].$$
(4)
For arbitrary $`D`$, one obtains instead $`𝐣^D=𝐣D(\overline{\psi }\psi )`$. Imposing the equation of continuity $`_t\rho =𝐣^D`$ then gives, as a kinematical constraint on the time-evolution of $`\psi `$, a Fokker-Planck type of equation: $`_t\rho =𝐣+D^2\rho `$.
No linear time-evolution equation for $`\psi `$ obeys this constraint. Rather we derived an interesting family of nonlinear Schrödinger equations, with the purely imaginary functional $`i\mathrm{}(D/2)^2\rho /\rho `$ multiplying $`\psi `$ on the right-hand side. That is, this particular form of nonlinearity was forced on us by the current algebra representation. And without linearity as an axiom, we also could not eliminate a priori the possibility of additional, real nonlinear functionals multiplying $`\psi `$. Doebner and I restricted these to homogeneous rational expressions with no more than two derivatives in the numerator. Defining (for convenience) $`\widehat{𝐣}=(m/\mathrm{})𝐣=(1/2i)[\overline{\psi }\psi (\overline{\psi })\psi ]`$, we introduced the real, homogeneous functionals $`R_1[\psi ],\mathrm{},R_5[\psi ]`$ given by
$$R_1=\frac{\widehat{𝐣}}{\rho },R_2=\frac{^{\mathrm{\hspace{0.17em}2}}\rho }{\rho },R_3=\frac{\widehat{𝐣}^{\mathrm{\hspace{0.17em}2}}}{\rho ^2},R_4=\frac{\widehat{𝐣}\rho }{\rho ^2},R_5=\frac{(\rho )^2}{\rho ^2}.$$
(5)
The family of nonlinear Schrödinger equations became then:
$$i\mathrm{}\frac{\psi }{t}=H_0\psi +\frac{i}{2}\mathrm{}DR_2[\psi ]\psi +\mathrm{}D^{}\underset{j=1}{\overset{5}{}}c_jR_j[\psi ]\psi ,$$
(6)
where $`D^{}`$ is another diffusion coefficient, the $`c_j`$ are real and dimensionless, and
$$H_0\psi =\frac{1}{2m}[i\mathrm{}(e/c)𝐀(𝐱,t)]^2\psi +[V+e\mathrm{\Phi }(𝐱,t)]\psi .$$
(7)
Below we shall see how an important subclass of (6), and certain more general nonlinear Schrödinger equations, can be obtained from the linear Schro‘ödinger equation via nonlinear gauge transformations. Eq. (6) contains as special cases a remarkable variety of nonlinear modifications of quantum mechanics proposed independently by other researchers , though without our fundamental motivation for the nonlinearity and typically without the above local, pure imaginary nonlinear functional multiplying $`\psi `$.
Using the expansion $`^2\psi /\psi =iR_1[\psi ]+(1/2)R_2[\psi ]R_3[\psi ](1/4)R_5[\psi ]`$, let us rewrite this family of equations as in Ref. , with some additional terms:
$$i\frac{\dot{\psi }}{\psi }=i\left[\underset{j=1}{\overset{2}{}}\nu _jR_j[\psi ]+\frac{(𝒜(𝐱,t)\rho )}{\rho }\right]+$$
$$\left[\underset{j=1}{\overset{5}{}}\mu _jR_j[\psi ]+U(𝐱,t)+\frac{(𝒜_1(𝐱,t)\rho )}{\rho }+\frac{𝒜_2(𝐱,t)\widehat{𝐣}}{\rho }+\alpha _1\mathrm{ln}\rho +\alpha _2S\right].$$
(8)
Here $`S`$ is the phase of $`\psi `$, $`U`$ is a (sufficiently smooth) external, real-valued, time-dependent scalar function; and $`𝒜,𝒜_1,`$ and $`𝒜_2`$ are distinct (sufficiently smooth) external, real-valued, time-dependent vector fields. Eq. (6) is obtained from Eq. (8) with the following substitutions:
$$\nu _1=\frac{\mathrm{}}{2m},\nu _2=\frac{1}{2}D,𝒜=\frac{e}{2mc}𝐀,$$
$$\mu _1=D^{}c_1,\mu _2=\frac{\mathrm{}}{4m}+D^{}c_2,\mu _3=\frac{\mathrm{}}{2m}+D^{}c_3,\mu _4=D^{}c_4,\mu _5=\frac{\mathrm{}}{8m}+D^{}c_5,$$
$$U(𝐱,t)=\frac{1}{\mathrm{}}[V(𝐱,t)+e\mathrm{\Phi }]+\frac{e^2}{2m\mathrm{}c^2}𝐀^2,𝒜_1=0,𝒜_2=\frac{e}{mc}𝐀,$$
$$\alpha _1=\alpha _2=\mathrm{\hspace{0.17em}0}.$$
(9)
The coefficients $`\nu _j(j=1,2),`$ $`\mu _j(j=1,\mathrm{},5),`$ and $`\alpha _j(j=1,2)`$ are taken to be continuously differentiable, real-valued functions of $`t`$. The motivation for this expansion, the reason behind the introduction of terms with $`\alpha _1`$, $`\alpha _2`$, and $`𝒜_1\mathrm{\hspace{0.17em}0}`$, and the reason for permitting the coefficients to be time-dependent, all stem from the discussion of nonlinear gauge transformations in the next section.
Finally, let us introduce here a further, natural generalization of Eq. (8). Let us insert into the imaginary part of the right-hand side the terms $`\nu _3R_3,\nu _4R_4`$, and $`\nu _5R_5`$, as well as new external scalar and vector fields, to achieve full symmetry between the real and imaginary parts . Thus we have, in effect, allowed for complexification of all the coefficients and external fields. The equation becomes:
$$i\frac{\dot{\psi }}{\psi }=i\left[\underset{j=1}{\overset{5}{}}\nu _jR_j[\psi ]+𝒯(𝐱,t)+\frac{(𝒜(𝐱,t)\rho )}{\rho }+\frac{𝒟(𝐱,t)\widehat{𝐣}}{\rho }+\delta _1\mathrm{ln}\rho +\delta _2S\right]+$$
$$\left[\underset{j=1}{\overset{5}{}}\mu _jR_j[\psi ]+U(𝐱,t)+\frac{(𝒜_1(𝐱,t)\rho )}{\rho }+\frac{𝒜_2(𝐱,t)\widehat{𝐣}}{\rho }+\alpha _1\mathrm{ln}\rho +\alpha _2S\right],$$
(10)
where $`𝒯`$ is a new external scalar field, and $`𝒟`$ a new external vector field. Note that the heat equation and other interesting equations of mathematical physics fall within this family. Some equations with soliton-like solutions are also included . But the equation of continuity relating $`\rho `$ and $`𝐣^D`$ no longer holds. Evidently when $`\nu _3=\nu _4=\nu _5=\delta _1=\delta _2=0`$, $`𝒯=0`$, and $`𝒟=0`$, we recover Eq. (8). When the remaining values are as in Eq. (9) with $`D=D^{}=0`$, we are back with the linear Schrödinger equation.
We shall see that the generalization of Eq. (8) to Eq. (10) follows from a further, natural extension of the notion of nonlinear gauge transformation.
## 2 Time-Dependent Nonlinear Gauge <br>Transformations
Let us write $`\psi =R\mathrm{exp}[iS]`$, where the amplitude $`R`$ and the phase $`S`$ are real. Then $`\rho =R^{\mathrm{\hspace{0.17em}2}}`$ and $`𝐣=(\mathrm{}/m)R^2S`$. While $`R`$ is gauge invariant, $`S`$ is not: under the usual, unitary gauge transformations of quantum mechanics, $`R^{}=R`$ but $`S^{}=S+\theta (𝐱,t)`$. Then $`\rho ^{}=\rho `$, while $`𝐣^{}=𝐣+(\mathrm{}/m)R^{\mathrm{\hspace{0.17em}2}}\theta `$.
If we begin with the linear Schrödinger equation in the absence of a vector potential, i.e., $`i\mathrm{}_t\psi =(\mathrm{}^2/2m)^{\mathrm{\hspace{0.17em}2}}\psi +V\psi `$, then the transformed wave function $`\psi ^{}=R^{}\mathrm{exp}[iS^{}]`$ satisfies $`i\mathrm{}_t\psi ^{}=(\mathrm{}^2/2m)[i\mathrm{g}rad\theta ]^2\psi ^{}+[V\mathrm{}\dot{\theta }]\psi ^{}`$. This observation can actually motivate introduction of the external electromagnetic gauge potentials $`𝐀`$ and $`\mathrm{\Phi }`$, and the “minimally coupled” Schrödinger equation whose Hamiltonian is given by Eq. (7). When we begin with (7), we have that $`\psi ^{}`$ satisfies the transformed equation obtained by substituting the gauge-transformed potentials: $`𝐀^{}=𝐀+(\mathrm{}c/e)\mathrm{g}rad\theta `$ and $`\mathrm{\Phi }^{}=\mathrm{\Phi }(\mathrm{}/e)\dot{\theta }.`$ A gauge-invariant current can now be written $`𝐉^{\mathrm{gi}}=𝐣(e/mc)\rho 𝐀`$, with $`_t\rho =𝐉^{\mathrm{gi}}`$. The physical fields $`𝐁=\times 𝐀`$ and $`𝐄=\mathrm{\Phi }(1/c)_t𝐀`$ are likewise gauge invariant. All this is elementary, and standard. It sets the pattern for consideration of nonlinear gauge transformations for nonlinear Schrödinger equations.
In the latter context we (necessarily) abandon the usual, tacit assumption that gauge transformations act linearly and unitarily. Doebner and I introduced a group of nonlinear transformations leaving our class of equations invariant as a family ,
$$R^{}=R,S^{}=\mathrm{\Lambda }S+\gamma \mathrm{ln}R+\theta ,$$
(11)
where in general $`\gamma `$ and $`\mathrm{\Lambda }`$ are continuously differentiable, real-valued functions of $`t`$, $`\mathrm{\Lambda }\mathrm{\hspace{0.17em}0}`$, and $`\theta `$ is a continuously differentiable, real-valued function of $`𝐱`$ and $`t`$. Then $`(\mathrm{\Lambda }_1,\gamma _1,\theta _1)(\mathrm{\Lambda }_2,\gamma _2,\theta _2)=(\mathrm{\Lambda }_1\mathrm{\Lambda }_2,\gamma _1+\mathrm{\Lambda }_1\gamma _2,\theta _1+\mathrm{\Lambda }_1\theta _2)`$. The original justification for taking these to be gauge transformations was the argument, put forth by many theorists, that any physical quantum-mechanical measurement could be reduced to a sequence of positional measurements at different times; with the system subjected to external force fields between measurements . Under Eq. (11),
$$\rho ^{}=\overline{\psi ^{}}\psi ^{}=\rho ,$$
$$\widehat{𝐣}^{}=\frac{1}{2i}[\overline{\psi ^{}}\psi ^{}(\overline{\psi ^{}})\psi ^{}]=\mathrm{\Lambda }\widehat{𝐣}+\frac{\gamma }{2}\rho +\rho \theta .$$
(12)
Keeping the interpretation of $`\rho =|\psi |^2`$ as the positional probability density, and writing invariant force fields in terms of the external potentials, the outcomes of all measurements do remain invariant. Eq. (11) also has other nice properties: it is strictly local, and it respects a certain separation condition for (many-particle) product wave functions . If $`\psi `$ obeys a Schrödinger equation of the type in Eq. (8), then $`\psi ^{}`$ transformed by (11) obeys another equation in the family, with transformed coefficients and external fields. The coefficients are given by:
$$\nu _1^{}=\frac{\nu _1}{\mathrm{\Lambda }},\nu _2^{}=\frac{\gamma }{2\mathrm{\Lambda }}\nu _1+\nu _2,$$
$$\mu _1^{}=\frac{\gamma }{\mathrm{\Lambda }}\nu _1+\mu _1,\mu _2^{}=\frac{\gamma ^2}{2\mathrm{\Lambda }}\nu _1\gamma \nu _2\frac{\gamma }{2}\mu _1+\mathrm{\Lambda }\mu _2,$$
$$\mu _3^{}=\frac{\mu _3}{\mathrm{\Lambda }},\mu _4^{}=\frac{\gamma }{\mathrm{\Lambda }}\mu _3+\mu _4,\mu _5^{}=\frac{\gamma ^2}{4\mathrm{\Lambda }}\mu _3\frac{\gamma }{2}\mu _4+\mathrm{\Lambda }\mu _5,$$
$$\alpha _1^{}=\mathrm{\Lambda }\alpha _1\frac{\gamma }{2}\alpha _2+\frac{1}{2}\left(\frac{\dot{\mathrm{\Lambda }}}{\mathrm{\Lambda }}\gamma \dot{\gamma }\right),\alpha _2^{}=\alpha _2\frac{\dot{\mathrm{\Lambda }}}{\mathrm{\Lambda }},$$
(13)
while the transformed vector and scalar fields are
$$𝒜^{}=𝒜\frac{\nu _1}{\mathrm{\Lambda }}\theta ,$$
$$𝒜_1^{}=\mathrm{\Lambda }𝒜_1\gamma 𝒜\frac{\gamma }{2}𝒜_2+\left(\frac{\gamma }{\mathrm{\Lambda }}\nu _1\mu _1+\frac{\gamma }{\mathrm{\Lambda }}\mu _3\mu _4\right)\theta ,$$
$$𝒜_2^{}=𝒜_2\frac{2\mu _3}{\mathrm{\Lambda }}\theta ,$$
$$U^{}=\mathrm{\Lambda }U\dot{\theta }+\left(\frac{\dot{\mathrm{\Lambda }}}{\mathrm{\Lambda }}\alpha _2\right)\theta +\frac{\mu _3}{\mathrm{\Lambda }}[\theta ]^{\mathrm{\hspace{0.17em}2}}+$$
$$\left(\mu _4\mu _3\frac{\gamma }{\mathrm{\Lambda }}\right)^2\theta +\frac{\gamma }{2}𝒜_2𝒜_2\theta .$$
(14)
Regarding Eqs. (13), note how the time-dependence of $`\gamma `$ and $`\mathrm{\Lambda }`$ in Eq. (11) requires that the $`\nu _j,\mu _j`$, and $`\alpha _j`$ in Eq. (8) be time-dependent, and that the $`\alpha _j`$ be allowed nonzero values. The terms with $`\alpha _1`$ and $`\alpha _2`$ were, respectively, first introduced by Bialynicki-Birula and Micielski and by Kostin . Likewise, we see in (14) how the $`𝒜_1`$ and $`𝒜_2`$ terms in Eq. (8) are needed. Nonlinear Schrödinger equations with arbitrary values of $`𝒜_2`$ were considered by Haag and Bannier , while as far as I know the field $`𝒜_1`$ was first considered in Ref. . An important subclass of Eq. (8) is linearizable by means of nonlinear gauge transformations; for this subclass, the physics is unchanged from ordinary quantum mechanics.
The coefficients, the external fields, and many of the nonlinear functionals in Eq. (8) are not gauge invariant. But we do have a current $`𝐉^{\mathrm{g}i},`$ invariant under nonlinear gauge transformations, that enters the continuity equation $`\dot{\rho }=𝐉^{\mathrm{g}i}`$, given by
$$𝐉^{\mathrm{g}i}=\mathrm{\hspace{0.17em}2}\nu _1\widehat{𝐣}\mathrm{\hspace{0.17em}2}\nu _2\rho \mathrm{\hspace{0.17em}2}\rho 𝒜.$$
(15)
This reduces, of course, to the usual gauge-invariant current in the linear case . Now, the existence of $`𝐉^{\mathrm{g}i}`$ means that our earlier assumption about all measurements being reducible to a succession of positional measurements is unnecessarily restrictive. It is sufficient that all measurements be expressible in terms of gauge-invariant quantities; and we have available for this the density $`\rho `$, the current $`𝐉^{\mathrm{g}i}`$, and gauge-invariant force fields (see below).
Doebner and I also introduced gauge-invariant parameters:
$$\tau _1=\nu _2\frac{1}{2}\mu _1,\tau _2=\nu _1\mu _2\nu _2\mu _1,\tau _3=\frac{\mu _3}{\nu _1},\tau _4=\mu _4\mu _1\frac{\mu _3}{\nu _1},$$
$$\tau _5=\nu _1\mu _5\nu _2\mu _4+\nu _2^{\mathrm{\hspace{0.17em}2}}\frac{\mu _3}{\nu _1},$$
$$\beta _1=\nu _1\alpha _1\nu _2\alpha _2+\nu _2\frac{\dot{\nu }_1}{\nu _1}\dot{\nu }_2,\beta _2=\alpha _2\frac{\dot{\nu }_1}{\nu _1}.$$
(16)
Some discussion of the physics behind these parameters may found in Ref. ; in particular, $`\tau _10`$, $`\tau _40`$, or $`\beta _20`$ violates time-reversal invariance; $`\tau _31`$ or $`\tau _40`$ breaks Galileian invariance; and in all these cases $`\tau _2`$ corresponds to the observed value of $`\mathrm{}^2/8m^2`$ (no longer can we identify the gauge-dependent quantity $`\nu _1`$ with the gauge-independent, observable constant $`\mathrm{}/2m`$). Thus the classical limit can be taken in a gauge-invariant manner by letting $`\tau _20`$.
Let me also remark here that the gauge-invariant parameter $`\beta _2`$ is naturally interpreted as a coefficient of friction, as it contributes (see below) a term $`\beta _2(𝐉^{\mathrm{g}i}/\rho )`$ to the expression for $`_t(𝐉^{\mathrm{g}i}/\rho )`$.
Continuing the discussion in Ref. we have also gauge-invariant fields. Set
$$\widehat{U}=\nu _1U\tau _3𝒜^{\mathrm{\hspace{0.17em}2}}(\tau _42\tau _1\tau _3)𝒜+𝒜𝒜_2\nu _2𝒜_2,$$
(17)
so that under nonlinear gauge transformation,
$$\widehat{U}^{}=\widehat{U}+\frac{\nu _1}{\mathrm{\Lambda }}\dot{\theta }+\frac{\nu _1}{\mathrm{\Lambda }}\alpha _2\theta \nu _1\frac{\dot{\mathrm{\Lambda }}}{\mathrm{\Lambda }^2}\theta .$$
(18)
Eq. (17) corrects algebraic errors in Ref. . The field $`\widehat{U}`$ is easily reduced to $`(1/2m)(V+e\mathrm{\Phi })`$ for the linear Schrödinger equation. We have the new gauge-invariant vector fields,
$$𝒜_1^{gi}=\nu _1𝒜_1+\left(\frac{2\nu _2\mu _3}{\nu _1}\mu _1\mu _4\right)𝒜\nu _2𝒜_2,$$
$$𝒜_2^{gi}=\frac{\nu _1}{2\mu _3}𝒜_2𝒜,$$
(19)
as well as magnetic and (generalized) electric plus other potential force fields,
$$=\times 𝒜=\frac{e}{2mc}𝐁,$$
$$=\widehat{U}\frac{𝒜}{t}\beta _2𝒜=\frac{1}{2m}V+\frac{e}{2m}𝐄.$$
(20)
Thus $`\widehat{U}=(1/2m)(V+e\mathrm{\Phi })`$ in general, and $`𝐄=\mathrm{\Phi }(1/c)_t𝐀(\beta _2/c)𝐀`$. Notice the extra term associated with Kostin’s nonlinearity; without it, $``$ is not gauge invariant. This leads in turn to an interesting modification of one of Maxwell’s equations:
$$\times 𝐄=\frac{1}{c}\frac{𝐁}{t}\frac{\beta _2}{c}𝐁.$$
(21)
## 3 Gauge-Invariant Equations of Motion
Using the (hydrodynamical) variables $`\rho `$ and $`𝐕=𝐉^{\mathrm{g}i}/\rho `$, it is straightforward to write down in manifestly gauge-invariant form the equations of motion corresponding to Eq. (8). We have in all cases the useful relation $`\times 𝐕=2=(e/mc)𝐁`$, and the continuity equation $`_t\rho =𝐉^{\mathrm{g}i}`$. In addition,
$$\frac{}{t}\left(\frac{𝐉^{\mathrm{g}i}}{\rho }\right)=\left[\mathrm{\hspace{0.17em}2}\tau _1\left(\frac{𝐉^{\mathrm{g}i}}{\rho }\right)+\mathrm{\hspace{0.17em}2}\tau _2\frac{^2\rho }{\rho }+\frac{1}{2}\tau _3\left(\frac{𝐉^{\mathrm{g}i}}{\rho }\right)^2\right]$$
$$+\left[(\mathrm{\hspace{0.17em}2}\tau _1[1+\tau _3]\tau _4)\left(\frac{𝐉^{\mathrm{g}i}}{\rho }\right)\frac{\rho }{\rho }+\mathrm{\hspace{0.17em}2}\tau _5\frac{(\rho )^2}{\rho ^{\mathrm{\hspace{0.17em}2}}}\right]$$
$$+\left[\mathrm{\hspace{0.17em}\hspace{0.17em}2}\frac{(𝒜_1^{gi}\rho )}{\rho }\mathrm{\hspace{0.17em}2}\tau _3𝒜_2^{gi}\left(\frac{𝐉^{\mathrm{g}i}}{\rho }\right)+\mathrm{\hspace{0.17em}2}\beta _1\mathrm{ln}\rho \right]$$
$$\beta _2\left(\frac{𝐉^{\mathrm{g}i}}{\rho }\right)\frac{1}{m}V+\frac{e}{m}𝐄.$$
(22)
Now we have the expected values of position, velocity, and acceleration:
$$<𝐱>=𝐱\rho (𝐱)𝑑𝐱,$$
$$<𝐯>=\frac{<𝐱>}{t}=\rho \left(\frac{𝐉^{\mathrm{g}i}}{\rho }\right)𝑑𝐱,$$
(23)
$$<𝐚>=\frac{<𝐯>}{t}=\rho [\frac{1}{2}\left(\frac{𝐉^{\mathrm{g}i}}{\rho }\right)^2+\left(\frac{𝐉^{\mathrm{g}i}}{\rho }\right)\times \frac{e}{mc}𝐁+\frac{}{t}\left(\frac{𝐉^{\mathrm{g}i}}{\rho }\right)]d𝐱.$$
Note that in Eqs. (22)-(23), the force laws governing interaction with the external electric and magnetic fields are unchanged from linear quantum mechanics.
## 4 The Enlarged Gauge Group
To this point, the amplitude $`R`$ and the phase $`S`$ have a fundamentally different status, both in linear quantum mechanics and in our nonlinear variations: $`R`$ is gauge invariant, and physically observable; while $`S`$ is not. This asymmetry seems more and more puzzling as one comes to appreciate the flexibility of description offered by nonlinear quantum time-evolutions, allowing for instance linear quantum mechanics to be written in a nonlinear gauge. Why should we be required to combine the gauge field $`S`$ with the physical field $`R`$ into a single complex-valued function $`\psi `$, and then through the Schrödinger equation couple both $`R`$ and $`S`$ to the gauge potentials? Why not instead try to couple gauge-dependent quantitites to each other, and correspondingly, physical fields to each other?
In addition, we remark that just as the formula (15) for the gauge-invariant current $`𝐉^{\mathrm{g}i}`$ depended on two coefficients and one external potential in the nonlinear time-evolution equation (8), there is no a priori principle that forbids the formula for the gauge-invariant probability density from likewise depending on coefficients and external potentials in the time-evolution equation. This is important as we consider enlarging the nonlinear gauge group further.
To achieve the desired generalization, define $`T=\mathrm{ln}R`$, so that $`\mathrm{ln}\psi =T+iS`$, and consider the transformations
$$\left(\begin{array}{c}S^{}\\ T^{}\end{array}\right)=\left(\begin{array}{cc}\mathrm{\Lambda }& \gamma \\ \lambda & \kappa \end{array}\right)\left(\begin{array}{c}S\\ T\end{array}\right)+\left(\begin{array}{c}\theta \\ \varphi \end{array}\right),$$
(24)
where $`\mathrm{\Lambda },\gamma ,\lambda ,`$ and $`\kappa `$ depend on $`t`$, and where $`\theta `$ and $`\varphi `$ depend on $`𝐱`$ and $`t`$. In place of the condition $`\mathrm{\Lambda }\mathrm{\hspace{0.17em}0}`$, we impose that $`\mathrm{\Delta }=\kappa \mathrm{\Lambda }\lambda \gamma 0`$, so that (24) is invertible. This is the transformation group $`𝒢`$, modeled on $`GL(2,𝐑)`$, with which we shall now work; the earlier gauge group is the subgroup with $`\lambda \mathrm{\hspace{0.17em}0}`$, $`\kappa \mathrm{\hspace{0.17em}1}`$, and $`\varphi \mathrm{\hspace{0.17em}0}`$. We thus treat the phase and the logarithm of the amplitude on an equal footing. The logarithmic variables $`T`$ and $`S`$ are, of course, familiar from earlier hydrodynamical and stochastic versions of quantum mechanics ; but they normally are treated quite asymmetrically.
We immediately see that Eq. (8) must be generalized further for it to be invariant under $`𝒢`$. This is accomplished by complexifying the coefficients and external potentials, to obtain Eq. (10)—a procedure that is natural, as Eq. (24) can be obtained by complexifying $`\mathrm{\Lambda }`$, $`\gamma `$, and $`\theta `$ in the transformation from $`\psi `$ to $`\psi ^{}`$.
Since so many terms in our equations involve logarithmic derivatives, let us continue with the variables $`S`$ and $`T`$. The operation of multiplying $`\psi `$ by a complex scalar is then to add real constants to $`S`$ and to $`T`$. The homogeneous terms in Eq. (5) become, $`R_1=^2S+\mathrm{\hspace{0.17em}2}ST`$, $`R_2=\mathrm{\hspace{0.17em}2}^2T+\mathrm{\hspace{0.17em}4}(T)^2`$, $`R_3=(S)^2`$, $`R_4=\mathrm{\hspace{0.17em}2}ST`$, and $`R_5=\mathrm{\hspace{0.17em}4}(T)^2`$. We now write the new, general nonlinear Schrödinger equation (10) as a pair of coupled partial differential equations for the extended real-valued functions $`S`$ and $`T`$, which are first order in time but have general second-order and quadratic terms:
$`\dot{S}`$ $`=`$ $`a_1^2S+a_2^2T+a_3(S)^2+a_4ST+a_5(T)^2`$
$`+a_6S+a_7T+u_0+𝐮_1S+𝐮_2T,`$
$`\dot{T}`$ $`=`$ $`b_1^2S+b_2^2T+b_3(S)^2+b_4ST+b_5(T)^2`$ (25)
$`+b_6S+b_7T+v_0+𝐯_1S+𝐯_2T.`$
The relation between Eq. (25) and and Eq. (10) is straightforward:
$$\begin{array}{cc}a_1=\mu _1,& b_1=\nu _1,\\ a_2=2\mu _2,& b_2=2\nu _2,\\ a_3=\mu _3,& b_3=\nu _3,\\ a_4=2\mu _12\mu _4,& b_4=2\nu _1+2\nu _4,\\ a_5=4\mu _24\mu _5,& b_5=4\nu _2+4\nu _5,\\ a_6=\alpha _2,& b_6=\delta _2,\\ a_7=2\alpha _1,& b_7=2\delta _1,\\ u_0=U𝒜_1,& v_0=𝒯+𝒜,\\ 𝐮_1=𝒜_2,& 𝐯_1=𝒟,\\ 𝐮_2=2𝒜_1,& 𝐯_2=2𝒜\end{array}.$$
(26)
Of course Eq. (8) is embedded in (25), as are many other interesting equations of mathematical physics. For reference, the usual, linear Schrödinger equation (7) corresponds to
$$a_1=\mathrm{\hspace{0.17em}0},a_2=\frac{\mathrm{}}{2m},a_3=,\frac{\mathrm{}}{2m},a_4=\mathrm{\hspace{0.17em}0},a_5=\frac{\mathrm{}}{2m},a_6=a_7=\mathrm{\hspace{0.17em}0},$$
$$u_0=\frac{1}{\mathrm{}}(V+e\mathrm{\Phi })\frac{e^2}{2m\mathrm{}c^2}𝐀^2,𝐮_1=\frac{e}{mc}𝐀,𝐮_2=\mathrm{\hspace{0.17em}0},$$
$$b_1=\frac{\mathrm{}}{2m},b_2=\mathrm{\hspace{0.17em}0}b_3=\mathrm{\hspace{0.17em}0}b_4=\frac{\mathrm{}}{m},b_5=\mathrm{\hspace{0.17em}0},b_6=b_7=\mathrm{\hspace{0.17em}0},$$
$$v_0=\frac{e}{2mc}𝐀,𝐯_1=\mathrm{\hspace{0.17em}0},𝐯_2=\frac{e}{mc}𝐀.$$
(27)
Now the coefficients $`a_j`$, $`b_j`$ obey the following transformation laws under (24), with the determinant $`\mathrm{\Delta }=\kappa \mathrm{\Lambda }\lambda \gamma `$:
$$\left[\begin{array}{c}a_1^{}\\ a_2^{}\\ b_1^{}\\ b_2^{}\end{array}\right]=\mathrm{\Delta }^1\left[\begin{array}{cccc}\kappa \mathrm{\Lambda }& \lambda \mathrm{\Lambda }& \kappa \gamma & \lambda \gamma \\ \gamma \mathrm{\Lambda }& \mathrm{\Lambda }^2& \gamma ^2& \gamma \mathrm{\Lambda }\\ \kappa \lambda & \lambda ^2& \kappa ^2& \kappa \lambda \\ \lambda \gamma & \lambda \mathrm{\Lambda }& \kappa \gamma & \kappa \mathrm{\Lambda }\end{array}\right]\left[\begin{array}{c}a_1\\ a_2\\ b_1\\ b_2\end{array}\right];$$
(28)
$$\left[\begin{array}{c}a_3^{}\\ a_4^{}\\ a_5^{}\\ b_3^{}\\ b_4^{}\\ b_5^{}\end{array}\right]=\mathrm{\Delta }^2\left[\begin{array}{c}a_3\\ a_4\\ a_5\\ b_3\\ b_4\\ b_5\end{array}\right],\mathrm{where}$$
(29)
$$=\left[\begin{array}{cccccc}\kappa ^2\mathrm{\Lambda }& \kappa \lambda \mathrm{\Lambda }& \lambda ^2\mathrm{\Lambda }& \kappa ^2\gamma & \kappa \lambda \gamma & \lambda ^2\gamma \\ 2\kappa \gamma \mathrm{\Lambda }& \mathrm{\Lambda }(\kappa \mathrm{\Lambda }+\lambda \gamma )& 2\lambda \mathrm{\Lambda }^2& 2\kappa \gamma ^2& \gamma (\kappa \mathrm{\Lambda }+\lambda \gamma )& 2\lambda \gamma \mathrm{\Lambda }\\ \gamma ^2\mathrm{\Lambda }& \gamma \mathrm{\Lambda }^2& \mathrm{\Lambda }^3& \gamma ^3& \gamma ^2\mathrm{\Lambda }& \gamma \mathrm{\Lambda }^2\\ \kappa ^2\lambda & \kappa \lambda ^2& \lambda ^3& \kappa ^3& \kappa ^2\lambda & \kappa \lambda ^2\\ 2\kappa \lambda \gamma & \lambda (\kappa \mathrm{\Lambda }+\lambda \gamma )& 2\lambda ^2\mathrm{\Lambda }& 2\kappa ^2\gamma & \kappa (\kappa \mathrm{\Lambda }+\lambda \gamma )& 2\kappa \lambda \mathrm{\Lambda }\\ \lambda \gamma ^2& \lambda \gamma \mathrm{\Lambda }& \lambda \mathrm{\Lambda }^2& \kappa \gamma ^2& \kappa \gamma \mathrm{\Lambda }& \kappa \mathrm{\Lambda }^2\end{array}\right];$$
and
$$\left[\begin{array}{c}a_6^{}\\ a_7^{}\\ b_6^{}\\ b_7^{}\end{array}\right]=\mathrm{\Delta }^1\left[\begin{array}{cccc}\kappa \mathrm{\Lambda }& \lambda \mathrm{\Lambda }& \kappa \gamma & \lambda \gamma \\ \gamma \mathrm{\Lambda }& \mathrm{\Lambda }^2& \gamma ^2& \gamma \mathrm{\Lambda }\\ \kappa \lambda & \lambda ^2& \kappa ^2& \kappa \lambda \\ \lambda \gamma & \lambda \mathrm{\Lambda }& \kappa \gamma & \kappa \mathrm{\Lambda }\end{array}\right]\left[\begin{array}{c}a_6\\ a_7\\ b_6\\ b_7\end{array}\right]+\mathrm{\Delta }^1\left[\begin{array}{c}\kappa \dot{\mathrm{\Lambda }}\lambda \dot{\gamma }\\ \mathrm{\Lambda }\dot{\gamma }\gamma \dot{\mathrm{\Lambda }}\\ \kappa \dot{\lambda }\lambda \dot{\kappa }\\ \mathrm{\Lambda }\dot{\kappa }\gamma \dot{\lambda }\end{array}\right].$$
(30)
The behavior of the external fields under generalized gauge transformation is more complicated. The transformed vector fields $`𝐮_1^{}`$, $`𝐮_2^{}`$, $`𝐯_1^{}`$, and $`𝐯_2^{}`$ are linear combinations of the six coefficients $`a_3`$, $`a_4`$, $`a_5`$, $`b_3`$, $`b_4`$, $`b_5`$ and the four vector fields $`𝐮_1`$, $`𝐮_2`$, $`𝐯_1`$, and $`𝐯_2`$; for example, the matrix element of $`𝐮_1^{}`$ by $`a_3`$ is $`\mathrm{\Delta }^2(2\kappa ^2\mathrm{\Lambda }\theta +\mathrm{\hspace{0.17em}2}\kappa \gamma \mathrm{\Lambda }\varphi )`$, and its matrix element by $`𝐯_2`$ is $`\mathrm{\Delta }^1(\lambda \gamma )`$. The transformed scalar fields $`u_0^{}`$ and $`v_0^{}`$ are linear combinations of all fourteen coefficients $`a_1\mathrm{}a_7`$ and $`b_1\mathrm{}b_7`$, the scalar fields $`u_0`$ and $`v_0`$, and the four vector fields, plus affine terms that depend on the time-derivatives of $`\mathrm{\Lambda }`$ $`\gamma `$, $`\lambda `$, $`\kappa `$, $`\theta `$, and $`\varphi `$. Probably little insight would be added by reproducing all the equations here.
Now we come to the main point. The generalization that is proposed will work (i.e., allow a gauge-invariant theory of measurement) only if it is possible to write combinations formed from $`S`$ and $`T`$ that are invariant under Eq. (24)—just as the earlier combinations $`\rho =\mathrm{exp}[2T]`$ and $`𝐉^{\mathrm{g}i}/\rho =2\nu _1S4\nu _2T2𝒜`$ are invariant under the smaller group. Consider for simplicity only the matrix part of (24); that is, set $`\theta =\varphi =0`$; call the gauge transformation matrix $`A`$. Suppose that $`d_1`$, $`d_2`$ are some coefficients depending on the $`a_j`$ and the $`b_j`$. Then $`d_1S+d_2T`$ is invariant under $`A`$ if and only if $`[d_1d_2]A^1=[d_1^{}d_2^{}]`$. From (29), we observe that the choice $`d_1=2a_3+b_4`$ and $`d_2=a_4+2b_5`$ obeys this condition. Hence $`d_1S+d_2T`$ can serve as one of the desired invariant combinations. Next let $`L_1=a_1S+a_2T`$ and $`L_2=b_1S+b_2T`$. Then the pair $`(L_1,L_2)`$ transforms under $`A`$ exactly as does the pair $`(S,T)`$, whence $`d_1L_1+d_2L_2`$ is also an invariant. In fact, any combination $`d_1(\sigma L_1+\tau S)+d_2(\sigma L_2+\tau T)`$, where $`\sigma `$ and $`\tau `$ are fully invariant combination of the coefficients, will be invariant; and, of course, any function of invariants is invariant. It is straightforward to verify that $`a_1+b_2=2\tau _1`$ and $`a_1b_2a_2b_1=2\tau _2`$, which were earlier identified as gauge invariants for (11), are also invariants under (24). We shall interpret $`\tau _2>0`$ as characterizing the class of Eqs. (25) that pertain to quantum mechanics, with $`\tau _20`$ defining the classical limit in a gauge-independent way.
To conclude, the desired invariant combinations of $`S`$ and $`T`$ exist. There is enough flexibility to permit a choice that reduces to the usual formulas in the case of the linear Schrödinger equation. In this way we can construct a positive definite, gauge-invariant probability density $`𝒫^{\mathrm{g}i}`$ and gauge-invariant current $`𝒥^{\mathrm{g}i}`$. A large subfamily of Eqs. (10) have solutions for which $`𝒫^{\mathrm{g}i}`$ and $`𝒥^{\mathrm{g}i}`$ obey the desired continuity equation, so that the total probability is conserved. And it is important to stress that a (smaller) subclass of Eqs. (10) is equivalent to ordinary quantum mechanics by way of generalized nonlinear gauge transformations, so that we are assured the new formalism is consistent. We can even exchange $`S`$ and $`\mathrm{ln}R`$ in ordinary quantum mechanics, by taking $`\gamma =\lambda =1`$, $`\kappa =\mathrm{\Lambda }=0`$.
It is clear that in this wider framework, many of the tacit assumptions of quantum mechanics no longer hold. For instance, integrability of the probability density function is only equivalent to square integrability of the wave function in certain gauges, so that we are often outside the usual Hilbert space of quantum mechanics.
Further details of these results will be presented elsewhere.
## Acknowledgments
I wish to thank the Alexander von Humboldt Foundation for generous support of this work during my 1998-99 sabbatical year in Germany, and the Arnold Sommerfeld Institute for Mathematical Physics, Technical University of Clausthal, for hospitality. |
warning/0002/hep-ph0002228.html | ar5iv | text | # Magnetic Moments of Δ Baryons in Light Cone QCD Sum Rules
## 1 Introduction
The extraction of the fundamental parameters of hadrons from experimental data requires some information about physics at large distances and they can not be calculated directly from fundamental QCD Lagrangian because at large distance strong coupling constant, $`\alpha _s`$, becomes large and perturbation theory is invalid. For this reason for determination of hadron parameters, a reliable non-perturbative approach is needed. Among other non-perturbative approaches, QCD sum rules is an especially powerful method in studying the properties of low-lying hadrons. In this method, deep connection between the hadron parameters and the QCD vacuum structure is established via a few condensate parameters. This method is adopted and extended in many works (see for example Refs. and references therein). One of characteristic parameters of the hadrons is their magnetic moments. Calculation of the nucleon magnetic moments in the framework of QCD sum rules method using external fields technique, first suggested in , was carried out in . They were later refined and extended to the entire baryon octet in .
Magnetic moments of the decuplet baryons are calculated in within the framework of QCD sum rules using external field. Note that in , from the decuplet baryons, only the magnetic moments of $`\mathrm{\Delta }^{++}`$ and $`\mathrm{\Omega }^{}`$ were calculated. At present, the magnetic moments of $`\mathrm{\Delta }^{++}`$ , $`\mathrm{\Delta }^0`$ and $`\mathrm{\Omega }^{}`$ are known from experiments. The experimental information provides new incentives for theoretical scrutiny of these physical quantities.
In this letter, we present an independent calculation of the magnetic moments of $`\mathrm{\Delta }^{++}`$, $`\mathrm{\Delta }^+`$, $`\mathrm{\Delta }^0`$, and $`\mathrm{\Delta }^{}`$ within the framework of an alternative approach to the traditional sum rules, i.e. the light cone QCD sum rules (LCQSR). Comparison of the predictions of this approach on magnetic moments with the results of other methods existing in the literature, and the experimental results is also presented.
The LCQSR is based on the operator product expansion on the light cone, which is an expansion over the twists of the operators rather than dimensions as in the traditional QCD sum rules. The main contribution comes from the lower twist operator. The matrix elements of the nonlocal operators between the vacuum and hadronic state defines the hadronic wave functions. (More about this method and its applications can be found in and references therein). Note that magnetic moments of the nucleon using LCQSR approach was studied in .
The paper is organized as follows. In Sect. II, the light cone QCD sum rules for the magnetic moments of $`\mathrm{\Delta }^{++}`$, $`\mathrm{\Delta }^+`$, $`\mathrm{\Delta }^0`$, and $`\mathrm{\Delta }^{}`$ are derived. In Sect. III, we present our numerical analysis and conclusion.
## 2 Sum Rules for the Magnetic Moments of $`\mathrm{\Delta }`$ baryons
According to the QCD sum rules philosophy, a quantitative estimate for the $`\mathrm{\Delta }`$ magnetic moment can be obtained by equating two different representations of the corresponding correlator, written in terms of hadrons and quark-gluons. For this aim, we consider the following correlation function
$`\mathrm{\Pi }_{\mu \nu }=i{\displaystyle 𝑑xe^{ipx}0|𝒯\eta _\mu ^B(x)\overline{\eta }_\nu ^B(0)|0_\gamma },`$ (1)
where $`𝒯`$ is the time ordering operator, $`\gamma `$ means external electromagnetic field. In this expression the $`\eta _\mu ^B`$’s are the interpolating currents for the baryon B. This correlation function can be calculated on one side phenomenologically, in terms of the hadron properties and on the other side by the operator product expansion (OPE) in the deep Eucledian region of the correlator momentum $`p^2\mathrm{}`$ using QCD degrees of freedom. By equating both expressions, we construct the corresponding sum rules.
On the phenomenological side, by inserting a complete set of one hadron states into the correlation function, Eq. (1), one obtains:
$`\mathrm{\Pi }_{\mu \nu }(p_1^2,p_2^2)={\displaystyle \underset{B_1,B_2}{}}{\displaystyle \frac{0|\eta _\mu ^B|B_1(p_1)}{p_1^2M_1^2}}B_1(p_1)|B_2(p_2)_\gamma {\displaystyle \frac{B_2(p_2)|\eta _\nu ^B|0}{p_2^2M_2^2}},`$ (2)
where $`p_2=p_1+q`$, $`q`$ is the photon momentum, $`B_i`$ form a complete set of baryons having the same quantum numbers as $`B`$, with masses $`M_i`$.
The matrix elements of the interpolating currents between the ground state and the state containing a single baryon, $`B`$, with momentum $`p`$ and having spin $`s`$ is defined as:
$`0|\eta _\mu |B(p,s)=\lambda _Bu_\mu (p,s),`$ (3)
where $`\lambda _B`$ is a phenomenological constant parametrizing the coupling strength of the baryon to the current, and $`u_\mu `$ is the Rarita-Schwinger spin-vector satisfying $`(\overline{)}pM_B)u_\mu =0`$, $`\gamma _\mu u_\mu =p_\mu u_\mu =0`$. (For a discussion of the properties of the Rarita-Schwinger spin-vector see e.g. ). In order to write down the phenomenological part of the sum rules, one also needs an expression for the matrix element $`B(p_1)|B(p_2)_\gamma `$. In the general case, the electromagnetic vertex of spin $`3/2`$ baryons can be written as
$`B(p_1)|B(p_2)_\gamma =ϵ_\rho \overline{u}_\mu (p_1)𝒪^{\mu \rho \nu }(p_1,p_2)u_\nu (p_2),`$ (4)
where $`ϵ_\rho `$ is the polarization vector of the photon and the Lorentz tensor $`𝒪^{\mu \rho \nu }`$ is given by:
$`𝒪^{\mu \rho \nu }(p_1,p_2)`$ $`=`$ $`g^{\mu \nu }\left[\gamma _\rho (f_1+f_2)+{\displaystyle \frac{(p_1+p_2)_\rho }{2M_B}}f_2+q_\rho f_3\right]`$ (5)
$``$ $`{\displaystyle \frac{q_\mu q_\nu }{(2M_B)^2}}\left[\gamma _\rho (g_1+g_2)+{\displaystyle \frac{(p_1+p_2)_\rho }{2M_B}}g_2+q_\rho g_3\right]`$
where the form factors $`f_i`$ and $`g_i`$ are (in the general case) functions of $`q^2=(p_1p_2)^2`$. In our problem, the values of the formfactors only at one point, $`q^2=0`$, are needed.
In our calculation, we also performed summation over spins of the Rarita-Schwinger spin vector,
$`{\displaystyle \underset{s}{}}u_\sigma (p,s)\overline{u}_\tau (p,s)={\displaystyle \frac{(\overline{)}p+M_B)}{2M_B}}\left\{g_{\sigma \tau }{\displaystyle \frac{1}{3}}\gamma _\sigma \gamma _\tau {\displaystyle \frac{2p_\sigma p_\tau }{3M_B^2}}+{\displaystyle \frac{p_\sigma \gamma _\tau p_\tau \gamma _\sigma }{3M_B}}\right\}`$ (6)
Using Eqs. (3-6), the correlation function can be expressed as the sum of various structures, not all of them independent. To remove the dependencies, an ordering of the gamma matrices should be chosen. For this purpose the structure $`\gamma _\mu \overline{)}p_1\overline{)}ϵ\overline{)}p_2\gamma _\nu `$ is chosen. With this ordering, the correlation function becomes:
$`\mathrm{\Pi }_{\mu \nu }`$ $`=`$ $`\lambda _B^2{\displaystyle \frac{1}{(p_1^2M_B^2)(p_2^2M_B^2)}}[g_{\mu \nu }\overline{)}p_1\overline{)}ϵ\overline{)}p_2{\displaystyle \frac{g_M}{3}}+`$ (7)
$`+`$ other structures with $`\gamma _\mu `$ at the beginning and $`\gamma _\nu `$ at the end $`]`$
where $`g_M`$ is the magnetic form factor, $`g_M/3=f_1+f_2`$. $`g_M`$ evaluated at $`q^2=0`$ gives the magnetic moment of the baryon in units of its natural magneton, $`e\mathrm{}/2m_Bc`$. The appearance of the factor $`3`$ can be understood from the fact that in the nonrelativistic limit, the maximum energy of the baryon in the presence of a uniform magnetic field with magnitude $`H`$ is $`3(f_1+f_2)Hg_MH`$ .
The reason for choosing this structure can be explained as follows. In general the interpolation current might also have a non-zero overlap with spin $`\frac{1}{2}`$ baryons, but spin $`\frac{1}{2}`$ baryons do not contribute to the structure $`g_{\mu \nu }\overline{)}p_1\overline{)}ϵ\overline{)}p_2`$ since their overlap is given by:
$`0|\eta _\mu |J=1/2=(Ap_\mu +B\gamma _\mu )u(p)`$ (8)
where $`(\overline{)}pm)u(p)=0`$ and $`(Am+4B)=0`$ .
In order to calculate the correlator (1) from the QCD side, first, appropriate currents should be chosen. For the case of the $`\mathrm{\Delta }`$ baryons, they can be chosen as (see for example ):
$`\eta _\mu ^{\mathrm{\Delta }^{++}}`$ $`=`$ $`ϵ^{abc}(u^{aT}C\gamma _\alpha u^b)u^c,`$
$`\eta _\mu ^{\mathrm{\Delta }^+}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{3}}}ϵ^{abc}[2(u^{aT}C\gamma _\alpha d^b)u^c+(u^{aT}C\gamma _\alpha u^b)d^c],`$
$`\eta _\mu ^{\mathrm{\Delta }^0}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{3}}}ϵ^{abc}[2(d^{aT}C\gamma _\alpha u^b)d^c+(d^{aT}C\gamma _\alpha d^b)u^c],`$
$`\eta _\mu ^\mathrm{\Delta }^{}`$ $`=`$ $`ϵ^{abc}(u^{aT}C\gamma _\alpha u^b)u^c,`$ (9)
where $`C`$ is the charge conjugation operator, $`a,b,c`$ are color indices. It should be noted that these baryon currents are not unique, one can choose an infinite number of currents with the $`\mathrm{\Delta }`$ baryon quantum numbers .
On the QCD side, for the same correlation functions we obtain:
$`\mathrm{\Pi }_{\mu \nu }^{\mathrm{\Delta }^{++}}=`$ $`\mathrm{\Pi }_{\mu \nu }^{\mathrm{\Delta }^{++}}+{\displaystyle \frac{1}{2}}ϵ^{abc}ϵ^{def}{\displaystyle d^4xe^{ipx}\gamma (q)|\overline{u}^fA_iu^a}`$ (10)
$`\{2S_u^{cd}\gamma _\nu S_u^{be}\gamma _\mu A_i+2S_u^{cd}\gamma _\nu A_i^{}\gamma _\mu S_u^{be}+`$
$`+2A_i\gamma _\nu S_u^{cd}\gamma _\mu S_u^{be}+S_u^{cd}\text{Tr}(\gamma _\nu S_{u}^{}{}_{}{}^{be}\gamma _\mu A_i)+`$
$`+S_u^{cd}\text{Tr}(\gamma _\nu A_i^{}\gamma _\mu S_u^{be})+A_i\text{Tr}(\gamma _\nu S_{u}^{}{}_{}{}^{cd}\gamma _\mu S_u^{be})\}|0`$
$`\mathrm{\Pi }_{\mu \nu }^{\mathrm{\Delta }^+}=`$ $`\mathrm{\Pi }_{\mu \nu }^{\mathrm{\Delta }^+}{\displaystyle \frac{1}{6}}ϵ^{abc}ϵ^{def}{\displaystyle d^4xe^{ipx}\gamma (q)|\overline{u}^dA_iu^a}`$ (11)
$`\{2A_i\gamma _\nu S_{d}^{}{}_{}{}^{be}\gamma _\mu S_u^{cf}+2A_i\gamma _\nu S_{u}^{}{}_{}{}^{cf}\gamma _\mu S_d^{be}+`$
$`+2S_d^{be}\gamma _\nu A_i^{}\gamma _\mu S_u^{cf}+2A_i\text{Tr}(\gamma _\nu S_{u}^{}{}_{}{}^{cf}\gamma _\mu S_d^{be})+`$
$`+S_d^{be}\text{Tr}(\gamma _\nu A_i^{}\gamma _\mu S_u^{cf})+`$
$`+2S_u^{cf}\gamma _\nu S_{d}^{}{}_{}{}^{be}\gamma _\mu A_i+2S_u^{cf}\gamma _\nu A_i^{}\gamma _\mu S_d^{be}+`$
$`+2S_d^{be}\gamma _\nu S_{u}^{}{}_{}{}^{cf}\gamma _\mu A_i+2S_u^{cf}\text{Tr}(\gamma _\nu A_i^{}\gamma _\mu S_d^{be})+`$
$`+S_d^{be}\text{Tr}(\gamma _\nu S_{u}^{}{}_{}{}^{cf}\gamma _\mu A_i)\}+\overline{d}^eA_id^b`$
$`\{2S_u^{ad}\gamma _\nu A_i^{}\gamma _\mu S_u^{cf}+2S_u^{ad}\gamma _\nu S_{u}^{}{}_{}{}^{cf}\gamma _\mu A_i+`$
$`+2A_i\gamma _\nu S_{u}^{}{}_{}{}^{ad}\gamma _\mu S_u^{cf}+2S_u^{ad}\text{Tr}(\gamma _\nu S_{u}^{}{}_{}{}^{cf}\gamma _\mu A_i)+`$
$`+A_i\text{Tr}(\gamma _\nu S_{u}^{}{}_{}{}^{ad}\gamma _\mu S_u^{ad})\}|0`$
where $`A_i=1,\gamma _\alpha ,\sigma _{\alpha \beta }/\sqrt{2},i\gamma _\alpha \gamma _5,\gamma _5`$, a sum over $`A_i`$ implied, $`S^{}CS^TC`$, $`A_i^{}=CA_i^TC`$, with $`T`$ denoting the transpose of the matrix, and $`S_q`$ is the full light quark propagator with both perturbative and non-perturbative contributions:
$`S_q`$ $`=`$ $`0|𝒯\overline{q}(x)q(0)|0`$ (12)
$`=`$ $`{\displaystyle \frac{i\overline{)}x}{2\pi ^2x^4}}{\displaystyle \frac{\overline{q}q}{12}}{\displaystyle \frac{x^2}{192}}m_0^2\overline{q}q`$
$``$ $`ig_s{\displaystyle _0^1}𝑑v\left[{\displaystyle \frac{\overline{)}x}{16\pi ^2x^2}}G_{\mu \nu }(vx)\sigma _{\mu \nu }vx_\mu G_{\mu \nu }(vx)\gamma _\nu {\displaystyle \frac{i}{4\pi ^2x^2}}\right]`$
The $`\mathrm{\Pi }_{\mu \nu }^\mathrm{\Delta }`$s in Eqs. (10) and (11) describe diagrams in which the photon interact with the quark lines perturbatively. Their explicit expressions can be obtained from the remaining terms by substituting all occurances of
$`\overline{q}^a(x)A_iq^bA_{i}^{}{}_{\alpha \beta }{}^{}2\left({\displaystyle d^4yF_{\mu \nu }y_\nu S_q^{pert}(xy)\gamma _\mu S_q^{pert}(y)}\right)_{\alpha \beta }^{ba}`$ (13)
where the Fock-Schwinger gauge is used, and $`S_q^{pert}`$ is the perturbative quark propogator, i.e. the first term in Eq. (12).
The corresponding expressions for the correlation functions for the $`\mathrm{\Delta }^{}`$ and $`\mathrm{\Delta }^0`$ baryons can be obtained from Eqs. (10) and (11), if one exchanges $`u`$-quarks by $`d`$-quarks and vice versa, respectively.
In order to be able to calculate the QCD part of the sum rules, one needs to know the matrix elements $`\gamma (q)|\overline{q}A_iq|0`$. Upto twist-4, the non-zero matrix elements given in terms of the photon wave functions are \[22-24\]:
$`\gamma (q)|\overline{q}\gamma _\alpha \gamma _5q|0`$ $`=`$ $`{\displaystyle \frac{f}{4}}e_qϵ_{\alpha \beta \rho \sigma }ϵ^\beta q^\rho x^\sigma {\displaystyle _0^1}𝑑ue^{iuqx}\psi (u)`$
$`\gamma (q)|\overline{q}\sigma _{\alpha \beta }q|0`$ $`=`$ $`ie_q\overline{q}q{\displaystyle _0^1}𝑑ue^{iuqx}`$ (14)
$`\times `$ $`\{(ϵ_\alpha q_\beta ϵ_\beta q_\alpha )[\chi \varphi (u)+x^2[g_1(u)g_2(u)]]`$
$`+`$ $`[qx(ϵ_\alpha x_\beta ϵ_\beta x_\alpha )+ϵx(x_\alpha q_\beta x_\beta q_\alpha )]g_2(u)\}`$
where $`\chi `$ is the magnetic susceptibility of the quark condensate and $`e_q`$ is the quark charge. The functions $`\varphi (u)`$ and $`\psi (u)`$ are the leading twist-2 photon wave functions, while $`g_1(u)`$ and $`g_2(u)`$ are the twist-4 functions.
Note that, since we have assumed massless quarks, $`m_u=m_d=0`$, and exact SU(2) flavor symmetry, which implies $`\overline{u}u=\overline{d}d`$, the $`u`$ and $`d`$ quark propagators are identical, $`S_u=S_d`$, whereas for the wave functions, the only difference is due to the different charges of the two quarks. The general expressions given by Eqs. (10) and (11), under exact $`SU(2)`$ flavor symmetry lead to the following results:
$`\mathrm{\Pi }_{\mu \nu }^{\mathrm{\Delta }^{++}}`$ $`=`$ $`{\displaystyle \frac{1}{2}}e_u𝒞`$
$`\mathrm{\Pi }_{\mu \nu }^{\mathrm{\Delta }^+}`$ $`=`$ $`{\displaystyle \frac{1}{6}}(2e_u+e_d)𝒞`$
$`\mathrm{\Pi }_{\mu \nu }^{\mathrm{\Delta }^0}`$ $`=`$ $`{\displaystyle \frac{1}{6}}(2e_d+e_u)𝒞`$
$`\mathrm{\Pi }_{\mu \nu }^\mathrm{\Delta }^{}`$ $`=`$ $`{\displaystyle \frac{1}{2}}e_d𝒞`$ (15)
where $`𝒞`$ is a factor independent of the quark charges. From Eq. (15), the following exact relations between theoretical parts of the correlator functions immediately follows:
$`\mathrm{\Pi }_{\mu \nu }^{\mathrm{\Delta }^+}`$ $`=`$ $`\mathrm{\Pi }_{\mu \nu }^\mathrm{\Delta }^{}={\displaystyle \frac{1}{2}}\mathrm{\Pi }_{\mu \nu }^{\mathrm{\Delta }^{++}}`$
$`\mathrm{\Pi }_{\mu \nu }^{\mathrm{\Delta }^0}`$ $`=`$ $`0`$ (16)
Hence, from now on, only the results for $`\mathrm{\Delta }^{++}`$ will be reported and for the other $`\mathrm{\Delta }`$’s, the corresponding results can be obtained from the Eqs. (16). Using Eqs. (12) and (14), from Eq. (10) and after some algebra and after performing Fourier transformation, for the coefficient of the structure
$`g_{\mu \nu }\overline{)}p_1\overline{)}ϵ\overline{)}p_2`$, we get:
$`\mathrm{\Pi }`$ $`=`$ $`e_u{\displaystyle _0^1}du\{{\displaystyle \frac{f\psi (u)}{48\pi ^2}}[4\mathrm{ln}(P^2)+{\displaystyle \frac{g^2G^2}{12P^4}}]+`$ (17)
$`+{\displaystyle \frac{8}{3P^4}}\overline{u}u^2[g_1(u)g_2(u)]+{\displaystyle \frac{\chi \varphi (u)\overline{u}u^2}{6P^2}}\left({\displaystyle \frac{m_0^2}{P^2}}+4\right)+`$
$`+{\displaystyle \frac{2\overline{u}u^2}{3P^4}}{\displaystyle \frac{g^2G^2}{768\pi ^4P^2}}{\displaystyle \frac{3P^2\mathrm{ln}(P^2)}{64\pi ^4}}\}`$
where $`P=p+qu`$. In Eq. (17), polynomials in $`P^2`$ are omitted since they do not contribute after the Borel transformation.
As stated earlier, in order to obtain the sum rules, one equates the phenomenological and theoretical expressions obtained within QCD. After performing the Borel transformation on the variables $`p^2`$ and $`(p+q)^2`$ in order to suppress the contributions of the higher resonances and the continuum, the following sum rules is obtained for the magnetic moment of $`\mathrm{\Delta }^{++}`$:
$`g_M`$ $`=`$ $`{\displaystyle \frac{3e_u}{\lambda _\mathrm{\Delta }^2}}e^{\frac{M_\mathrm{\Delta }^2}{M^2}}\{{\displaystyle \frac{f\psi (u_0)}{12\pi ^2}}[{\displaystyle \frac{g^2G^2}{48}}M^4f_1({\displaystyle \frac{s_0}{M^2}})]+`$ (18)
$`+`$ $`{\displaystyle \frac{8}{3}}\overline{u}u^2[g_1(u_0)g_2(u_0)]+`$
$`+`$ $`{\displaystyle \frac{\chi \varphi (u_0)\overline{u}u^2}{6}}\left[m_0^24M^2f_0({\displaystyle \frac{s_0}{M^2}})\right]`$
$`+{\displaystyle \frac{2\overline{u}u^2}{3}}+{\displaystyle \frac{g^2G^2M^2}{768\pi ^4}}f_0({\displaystyle \frac{s_0}{M^2}})+{\displaystyle \frac{3M^6}{64\pi ^4}}f_2({\displaystyle \frac{s_0}{M^2}})\}`$
where the functions
$`f_n(x)=1e^x{\displaystyle \underset{k=0}{\overset{n}{}}}{\displaystyle \frac{x^k}{k!}}`$ (19)
are used to subtract the contributions of the continuum. In Eq. (18), $`s_0`$ is the continuum threshold,
$`u_0`$ $`=`$ $`{\displaystyle \frac{M_1^2}{M_1^2+M_2^2}}`$
$`{\displaystyle \frac{1}{M^2}}`$ $`=`$ $`{\displaystyle \frac{1}{M_1^2}}+{\displaystyle \frac{1}{M_2^2}}`$
As we are dealing with just a single baryon, the Borel parameters $`M_1^2`$ and $`M_2^2`$ can be taken to be equal, i.e. $`M_1^2=M_2^2`$, from which it follows that $`u_0=1/2`$.
## 3 Numerical Analysis
From Eq. (18), one sees that, in order to calculate the numerical value of the magnetic moment of the $`\mathrm{\Delta }^{++}`$, besides several numerical constants, one requires expressions for the photon wave functions. It was shown in that the leading photon wave functions receive only small corrections from the higher conformal spin, so they do not deviate much from the asymptotic form. Following , we shall use the following photon wave functions:
$`\varphi (u)`$ $`=`$ $`6u\overline{u}`$
$`\psi (u)`$ $`=`$ $`1`$
$`g_1(u)`$ $`=`$ $`{\displaystyle \frac{1}{8}}\overline{u}(3u)`$
$`g_2(u)`$ $`=`$ $`{\displaystyle \frac{1}{4}}\overline{u}^2`$
where $`\overline{u}=1u`$. The values of the other constants that are used in the calculation are: $`f=0.028GeV^2`$, $`\chi =4.4GeV^2`$ (in , $`\chi `$ is estimated to be $`\chi =3.3GeV^2`$), $`g^2G^2=0.474GeV^4`$, $`\overline{u}u=(0.243)^3GeV^3`$, $`m_0^2=(0.8\pm 0.2)GeV^2`$ , $`\lambda _\mathrm{\Delta }=0.038`$ .
Having fixed the input parameters, our next task is to find a region of Borel parameter, $`M^2`$, where dependence of the magnetic moments on $`M^2`$ and the continuum threshold $`s_0`$ is rather weak and at the same time the higher dimension operators, higher states and continuum contributions remain under control. We demand that these contributions are less then $`35\%`$. Under this requirement, the working region for the Borel parameter, $`M^2`$, is found to be $`1.1GeV^2M^21.4GeV^2`$. Finally, in this range of the Borel parameter, the magnetic moment of $`\mathrm{\Delta }^{++}`$ is found to be $`(4.55\pm 0.03)\mu _N`$. This prediction on the magnetic moment is obtained at $`s_0=4.4GeV^2`$. Choosing $`s_0=3.8GeV^2`$ or $`s_0=4.2GeV^2`$ changes the result at most by $`6\%`$(see Fig. 1). The calculated magnetic moment is in good agreement with the experimental result $`(4.52\pm 0.50\pm 0.45)\mu _N`$ . Our results on the magnetic moments for $`\mathrm{\Delta }^+`$, $`\mathrm{\Delta }^0`$ and $`\mathrm{\Delta }^{}`$ are presented in Table 1. For completeness, in this table, we have also presented the predictions of other approaches. Comparing the values presented in Table 1, it is seen that our predictions on magnetic moments is larger than the QCDSR predictions.
Finally, for the calculation of the magnetic moments of other members of the decuplet (which contain at least one $`s`$-quark), the correction due to the strange quark mass should be taken into account. Their calculations would be presented in a future work. |
warning/0002/math-ph0002019.html | ar5iv | text | # Contents
## 1 Introduction
One of the most interesting and important multidimensional integrable equations is the self-dual Yang-Mills equation (SDYME) . This four - dimensional equation arises in relatiivity and in field theory . The SDYM equations describe a connection for a bundle over the Grassmannian of two-dimensional subspaces of the twistor space. Integrability for a SDYM connection means that its curvature vanishes on certain two-planes in the tangent space of the Grassmannian. As shown in , This allows one to characterize SDYM connections in terms of the spliting problem for a transition function in a holomorphic bundle over the Riemann sphere, i.e. the trivialiization of the bundle .
Resently it has been shown that practically all known soliton equations in 1+1 and 2+1 dimensions may be obtained by reductions of the SDYME \[2, 4-5, 18-19, 34-37\] (see also the book and references therein). On the other hand, it is well known that almost all known integrable equations may be obtained from the some equations of soliton geometry by reductions. These equations are the integrability conditions of the system describing the moving orthogonal trihedral of a curve or surface \[8-24,38-40,43-44,46\]. For example, in 2+1 dimensions, the role of the such geometrical equations play the mM-LXII or M-LXII equations. In and in our previous notes of this series \[18-21\] we have considered some aspects of the multidimensional soliton geometry (see, also the refs. \[38-40, 43-44\]). Also we have studied the relation between the multidimensional soliton equations and the Self-Dual Yang-Mills equation. In this note we continue this work.
## 2 Soliton geometry in $`d=4`$ dimensions
It is well known that exist several equations describing the 4 - dimensional curves/”surfaces” in n - dimensional space. Here we present some of them.
### 2.1 The M-LXVIII equation
Consider the M-LXVIII equation
$$\left(\begin{array}{c}𝐞_1\\ 𝐞_2\\ 𝐞_3\\ \mathrm{}\\ 𝐞_n\end{array}\right)_{\xi _1}=A\left(\begin{array}{c}𝐞_1\\ 𝐞_2\\ 𝐞_3\\ \mathrm{}\\ 𝐞_n\end{array}\right)_{\xi _3}+B\left(\begin{array}{c}𝐞_1\\ 𝐞_2\\ 𝐞_3\\ \mathrm{}\\ 𝐞_n\end{array}\right)$$
$`(1a)`$
$$\left(\begin{array}{c}𝐞_1\\ 𝐞_2\\ 𝐞_3\\ \mathrm{}\\ 𝐞_n\end{array}\right)_{\xi _2}=C\left(\begin{array}{c}𝐞_1\\ 𝐞_2\\ 𝐞_3\\ \mathrm{}\\ 𝐞_n\end{array}\right)_{\xi _4}+D\left(\begin{array}{c}𝐞_1\\ 𝐞_2\\ 𝐞_3\\ \mathrm{}\\ 𝐞_n\end{array}\right)$$
$`(1b)`$
where $`𝐞_j^2=1,(𝐞_i𝐞_j)=\delta _{ij}`$ and $`A(\lambda ),B(\lambda ),C(\lambda ),D(\lambda )`$ are (n$`\times `$n)-matrices, $`\lambda `$ is some parameter, $`\xi _i`$ are coordinates. This equation describes some four dimensional curves and/or ”surfaces” in n-dimensional space. It is one of main equations of the multidimensional soliton geometry and admits several integrable reductions . The compatibility condition of these equations gives the M-LXX equation, which contents several interesting integrable nonlinear evolution equations (NEE). In this note we will present some of these integrable reductions.
### 2.2 The M-LXXI equation
Consider the M-LXXI equation
$$\left(\begin{array}{c}𝐞_1\\ 𝐞_2\\ 𝐞_3\\ \mathrm{}\\ 𝐞_n\end{array}\right)_{\xi _1}=A\left(\begin{array}{c}𝐞_1\\ 𝐞_2\\ 𝐞_3\\ \mathrm{}\\ 𝐞_n\end{array}\right)$$
$`(2a)`$
$$\left(\begin{array}{c}𝐞_1\\ 𝐞_2\\ 𝐞_3\\ \mathrm{}\\ 𝐞_n\end{array}\right)_{\xi _2}=B\left(\begin{array}{c}𝐞_1\\ 𝐞_2\\ 𝐞_3\\ \mathrm{}\\ 𝐞_n\end{array}\right)$$
$`(2b)`$
$$\left(\begin{array}{c}𝐞_1\\ 𝐞_2\\ 𝐞_3\\ \mathrm{}\\ 𝐞_n\end{array}\right)_{\xi _4}=C\left(\begin{array}{c}𝐞_1\\ 𝐞_2\\ 𝐞_3\\ \mathrm{}\\ 𝐞_n\end{array}\right)_{\xi _3}+D\left(\begin{array}{c}𝐞_1\\ 𝐞_2\\ 𝐞_3\\ \mathrm{}\\ 𝐞_n\end{array}\right).$$
$`(2c)`$
The compatibility condition of these equations gives some NEEs (see, e.g. the refs. \[8,18-19\]).
### 2.3 The M-LXI equation
Consider the 4-dimensional M-LXI equation
$$\left(\begin{array}{c}𝐞_1\\ 𝐞_2\\ 𝐞_3\\ \mathrm{}\\ 𝐞_n\end{array}\right)_{\xi _1}=A\left(\begin{array}{c}𝐞_1\\ 𝐞_2\\ 𝐞_3\\ \mathrm{}\\ 𝐞_n\end{array}\right)$$
$`(3a)`$
$$\left(\begin{array}{c}𝐞_1\\ 𝐞_2\\ 𝐞_3\\ \mathrm{}\\ 𝐞_n\end{array}\right)_{\xi _2}=B\left(\begin{array}{c}𝐞_1\\ 𝐞_2\\ 𝐞_3\\ \mathrm{}\\ 𝐞_n\end{array}\right)$$
$`(3b)`$
$$\left(\begin{array}{c}𝐞_1\\ 𝐞_2\\ 𝐞_3\\ \mathrm{}\\ 𝐞_n\end{array}\right)_{\xi _3}=C\left(\begin{array}{c}𝐞_1\\ 𝐞_2\\ 𝐞_3\\ \mathrm{}\\ 𝐞_n\end{array}\right)$$
$`(3c)`$
$$\left(\begin{array}{c}𝐞_1\\ 𝐞_2\\ 𝐞_3\\ \mathrm{}\\ 𝐞_n\end{array}\right)_{\xi _4}=D\left(\begin{array}{c}𝐞_1\\ 𝐞_2\\ 𝐞_3\\ \mathrm{}\\ 𝐞_n\end{array}\right).$$
$`(3d)`$
The compatibility condition of these equations gives the 4-dimensional M-LXII equation
$$A_{\xi _2}B_{\xi _1}+[A,B]=0$$
$`(4a)`$
$$A_{\xi _3}C_{\xi _1}+[A,C]=0$$
$`(4b)`$
$$A_{\xi _4}D_{\xi _1}+[A,D]=0$$
$`(4c)`$
$$C_{\xi _2}B_{\xi _3}+[C,B]=0$$
$`(4d)`$
$$D_{\xi _2}B_{\xi _4}+[D,B]=0$$
$`(4e)`$
$$C_{\xi _4}D_{\xi _3}+[C,D]=0.$$
$`(4f)`$
This equation contents many interesting NEEs (see, e.g. the refs. \[8,18-19\]).
## 3 The M-LXX equation
From (1) we get the following M-LXX equation
$$AD_{\xi _3}CB_{\xi _4}+B_{\xi _2}D_{\xi _1}+[B,D]=0$$
$`(5a)`$
$$A_{\xi _2}CA_{\xi _4}+[A,D]=0$$
$`(5b)`$
$$[A,C]=0$$
$`(5c)`$
$$C_{\xi _1}AC_{\xi _3}+[C,B]=0.$$
$`(5d)`$
If we choose
$$A=aI,C=bI,a,b=consts,$$
$`(6)`$
then the M-LXX equation (5) takes the form
$$aD_{\xi _3}bB_{\xi _4}+B_{\xi _2}D_{\xi _1}+[B,D]=0.$$
$`(7)`$
## 4 The SDYME as the particular case of the M-LXX equation
Now we assume that
$$B=A_1\lambda A_3,D=A_2\lambda A_4,a=b=\lambda .$$
$`(8)`$
So that the M-LXVIII equation (1) takes the form
$$\left(\begin{array}{c}𝐞_1\\ 𝐞_2\\ 𝐞_3\\ \mathrm{}\\ 𝐞_n\end{array}\right)_{\xi _1}=\lambda \left(\begin{array}{c}𝐞_1\\ 𝐞_2\\ 𝐞_3\\ \mathrm{}\\ 𝐞_n\end{array}\right)_{\xi _3}+(A_1\lambda A_3)\left(\begin{array}{c}𝐞_1\\ 𝐞_2\\ 𝐞_3\\ \mathrm{}\\ 𝐞_n\end{array}\right)$$
$`(9a)`$
$$\left(\begin{array}{c}𝐞_1\\ 𝐞_2\\ 𝐞_3\\ \mathrm{}\\ 𝐞_n\end{array}\right)_{\xi _2}=\lambda \left(\begin{array}{c}𝐞_1\\ 𝐞_2\\ 𝐞_3\\ \mathrm{}\\ 𝐞_n\end{array}\right)_{\xi _4}+(A_2\lambda A_4)\left(\begin{array}{c}𝐞_1\\ 𝐞_2\\ 𝐞_3\\ \mathrm{}\\ 𝐞_n\end{array}\right).$$
$`(9b)`$
From (7) we obtain the SDYME
$$A_{2\xi _1}A_{1\xi _2}+[A_2,A_1]=0$$
$`(10a)`$
$$A_{4\xi _3}A_{3\xi _4}+[A_4,A_3]=0$$
$`(10b)`$
$$A_{1\xi _4}A_{4\xi _1}+[A_1,A_4]=A_{2\xi _3}A_{3\xi _2}+[A_2,A_3]$$
$`(10c)`$
or
$$F_{\xi _1\xi _2}=0$$
$`(11a)`$
$$F_{\xi _3\xi _4}=0$$
$`(11b)`$
$$F_{\xi _4\xi _1}F_{\xi _3\xi _2}=0.$$
$`(11c)`$
Here
$$F_{\xi _i\xi _k}=A_{k\xi _i}A_{i\xi _k}+[A_k,A_i].$$
$`(12)`$
The SDYME (10) on a connection $`A`$ are the self-duality conditions on the curvature under the Hodge star operation
$$F=F$$
$`(13)`$
or in index notation
$$F_{\mu \nu }=ϵ\mu \nu \rho \delta F^{\rho \delta }$$
$`(14)`$
where $``$ is the Hodge operator, $`ϵ_{\mu \nu \rho \delta }`$ stands for the completely antisymmetric tensor in four dimensions with the convention: $`ϵ_{1234}=1`$. The SDYME is integrable by the Inverse Scattering Transform (IST) method (see, e.g. ). The Lax representation (LR) of the SDYME has the form
$$\mathrm{\Phi }_{\xi _1}\lambda \mathrm{\Phi }_{\xi _3}=(A_1\lambda A_3)\mathrm{\Phi }$$
$`(15a)`$
$$\mathrm{\Phi }_{\xi _2}\lambda \mathrm{\Phi }_{\xi _4}=(A_2\lambda A_4)\mathrm{\Phi }.$$
$`(15b)`$
## 5 The M-LXVII equation as the particular case of the M-LXVIII equation
In this section we consider the 3- dimensional curves/”surfaces”. Let variables in the M-LXVIII equation (1) are independent of $`\xi _3`$. Then we obtain the following M-LXVII equation
$$\left(\begin{array}{c}𝐞_1\\ 𝐞_2\\ 𝐞_3\\ \mathrm{}\\ 𝐞_n\end{array}\right)_{\xi _1}=B\left(\begin{array}{c}𝐞_1\\ 𝐞_2\\ 𝐞_3\\ \mathrm{}\\ 𝐞_n\end{array}\right)$$
$`(16a)`$
$$\left(\begin{array}{c}𝐞_1\\ 𝐞_2\\ 𝐞_3\\ \mathrm{}\\ 𝐞_n\end{array}\right)_{\xi _2}=C\left(\begin{array}{c}𝐞_1\\ 𝐞_2\\ 𝐞_3\\ \mathrm{}\\ 𝐞_n\end{array}\right)_{\xi _4}+D\left(\begin{array}{c}𝐞_1\\ 𝐞_2\\ 𝐞_3\\ \mathrm{}\\ 𝐞_n\end{array}\right).$$
$`(16b)`$
In this case we get the 3-dimensional M-LXX equation
$$bB_{\xi _4}+B_{\xi _2}D_{\xi _1}+[B,D]=0$$
$`(17)`$
and the 3-dimensional SDYME
$$A_{2\xi _1}A_{1\xi _2}+[A_2,A_1]=0$$
$`(18a)`$
$$A_{3\xi _4}+[A_4,A_3]=0$$
$`(18b)`$
$$A_{1\xi _4}A_{4\xi _1}+[A_1,A_4]=A_{3\xi _2}+[A_2,A_3].$$
$`(18c)`$
We note that for the 3-dimensional SDYME (18) the corresponding M-LXVII equation has the form
$$\left(\begin{array}{c}𝐞_1\\ 𝐞_2\\ 𝐞_3\\ \mathrm{}\\ 𝐞_n\end{array}\right)_{\xi _1}=(A_1\lambda A_3)\left(\begin{array}{c}𝐞_1\\ 𝐞_2\\ 𝐞_3\\ \mathrm{}\\ 𝐞_n\end{array}\right)$$
$`(19a)`$
$$\left(\begin{array}{c}𝐞_1\\ 𝐞_2\\ 𝐞_3\\ \mathrm{}\\ 𝐞_n\end{array}\right)_{\xi _2}=\lambda \left(\begin{array}{c}𝐞_1\\ 𝐞_2\\ 𝐞_3\\ \mathrm{}\\ 𝐞_n\end{array}\right)_{\xi _4}+(A_2\lambda A_4)\left(\begin{array}{c}𝐞_1\\ 𝐞_2\\ 𝐞_3\\ \mathrm{}\\ 𝐞_n\end{array}\right).$$
$`(19b)`$
So that the curve or ”surface” corresponding to the equation (19) is the integrable.
## 6 The Zakharov equation as the exact reduction of the SDYME
Now let us we consider the gauge condition
$$A_4=0.$$
$`(20)`$
Then the 3-dimensional SDYME takes the form
$$A_{2\xi _1}A_{1\xi _2}+[A_2,A_1]=0$$
$`(21a)`$
$$A_{3\xi _4}=0$$
$`(21b)`$
$$A_{1\xi _4}=A_{3\xi _2}+[A_2,A_3].$$
$`(21c)`$
If we take
$$A_3=const$$
$`(22)`$
then the 3-dimensional SDYME has the form
$$A_{2\xi _1}A_{1\xi _2}+[A_2,A_1]=0$$
$`(23a)`$
$$A_{1\xi _4}=[A_2,A_3].$$
$`(23b)`$
Now we consider the case $`n=3`$. And we take the following representations of the connections $`A_1,A_2,A_3`$
$$A_1=\left(\begin{array}{ccc}0& i(\varphi r^2\overline{\varphi })& (\varphi +r^2\overline{\varphi })\\ i(\varphi r^2\overline{\varphi })& 0& 0\\ (\varphi +r^2\overline{\varphi })& 0& 0\end{array}\right)$$
$`(24a)`$
$$A_2=\left(\begin{array}{ccc}0& (\varphi +r^2\overline{\varphi })_y& i(\varphi r^2\overline{\varphi })_y\\ (\varphi +r^2\overline{\varphi })_y& 0& v\\ i(\varphi +r^2\overline{\varphi })_y& v& 0\end{array}\right)$$
$`(24b)`$
$$A_3=\left(\begin{array}{ccc}0& 0& 0\\ 0& 0& 1\\ 0& 1& 0\end{array}\right)$$
$`(24c)`$
and let
$$\xi _1=x,\xi _2=t,\xi _4=y.$$
$`(25)`$
So in our case the equation (23) takes the form
$$A_{2x}A_{1t}+[A_2,A_1]=0$$
$`(26a)`$
$$A_{1y}=[A_2,A_3]$$
$`(26b)`$
or in elements
$$i\varphi _t=\varphi _{xy}+v\varphi $$
$`(27a)`$
$$v_x=2r^2_y|\varphi |^2.$$
$`(27b)`$
It is the Zakharov equation (ZE) . We note that in this case the corresponding M-LXVII equation (19) looks like
$$\left(\begin{array}{c}𝐞_1\\ 𝐞_2\\ 𝐞_3\end{array}\right)_x=(A_1\lambda A_3)\left(\begin{array}{c}𝐞_1\\ 𝐞_2\\ 𝐞_3\end{array}\right)$$
$`(28a)`$
$$\left(\begin{array}{c}𝐞_1\\ 𝐞_2\\ 𝐞_3\end{array}\right)_t=\lambda \left(\begin{array}{c}𝐞_1\\ 𝐞_2\\ 𝐞_3\end{array}\right)_y+A_2\left(\begin{array}{c}𝐞_1\\ 𝐞_2\\ 𝐞_3\end{array}\right).$$
$`(28b)`$
From (20), (22) and (28) we get the following LR of the ZE
$$\mathrm{\Phi }_x=(A_1\lambda A_3)\mathrm{\Phi }$$
$`(29a)`$
$$\mathrm{\Phi }_t\lambda \mathrm{\Phi }_y=A_2\mathrm{\Phi }$$
$`(29b)`$
where $`A_i`$ are $`\mathrm{𝐬𝐨}(\mathrm{𝟑})`$ or $`\mathrm{𝐬𝐨}(\mathrm{𝟐},\mathrm{𝟏})`$ matrices. It is convenient to use the well known isomorphism $`\mathrm{𝐬𝐨}(\mathrm{𝟑})\mathrm{𝐬𝐮}(\mathrm{𝟐})`$ or $`\mathrm{𝐬𝐨}(\mathrm{𝟐},\mathrm{𝟏})\mathrm{𝐬𝐮}(\mathrm{𝟏},\mathrm{𝟏})`$ and to rewrite the equations (29) in terms of 2$`\times `$2 matrices. As result we have the following standard LR of the ZE (27)
$$\mathrm{\Psi }_x=U\mathrm{\Psi }$$
$`(30a)`$
$$\mathrm{\Psi }_t=\lambda \mathrm{\Psi }_y+V\mathrm{\Psi }$$
$`(30b)`$
where
$$U=\frac{i\lambda }{2}\sigma _3+G,G=\left(\begin{array}{cc}0& \varphi \\ r^2\overline{\varphi }& 0\end{array}\right)$$
$`(31a)`$
$$V=i\sigma _3(\frac{1}{2}vI+G_y).$$
$`(31b)`$
We note that in the (1+1)-dimensional case, i.e. when $`y=x`$ instead of the equations (28) and (27) we obtain the (1+1)-dimensional M-LXVII equation
$$\left(\begin{array}{c}𝐞_1\\ 𝐞_2\\ 𝐞_3\end{array}\right)_x=(A_1\lambda A_3)\left(\begin{array}{c}𝐞_1\\ 𝐞_2\\ 𝐞_3\end{array}\right)$$
$`(32a)`$
$$\left(\begin{array}{c}𝐞_1\\ 𝐞_2\\ 𝐞_3\end{array}\right)_t=[\lambda (A_1\lambda A_3)+A_2]\left(\begin{array}{c}𝐞_1\\ 𝐞_2\\ 𝐞_3\end{array}\right)$$
$`(32b)`$
and the nonlinear Schrodinger equation
$$i\varphi _t=\varphi _{xx}+2r^|\varphi |^2\varphi =0$$
$`(33)`$
At last we note that if $`n>3`$ then we get the n-component generalisation of the Zakharov equation
$$i\varphi _{jt}=\varphi _{jxy}+v\varphi _j=0$$
$`(34a)`$
$$v_x=2r^2(\underset{k=1}{\overset{n}{}}|\varphi _k|^2)_y$$
$`(34b)`$
## 7 Integrable spin systems and the SDYME
Integrable spin systems in 2+1 dimensions also can be considered as exact reductions of the SDYME. As example we show that Myrzakulov I (M-I) equation is the reduction of the SDYME. Consider the following gauge condition
$$A_1=A_2=0.$$
$`(35)`$
In this case the 3-dimensional SDYME takes the form
$$A_{3\xi _4}+[A_4,A_3]=0$$
$`(36a)`$
$$A_{4\xi _1}=A_{3\xi _2}$$
$`(36b)`$
or in terms of $`x,y,t`$
$$A_y+[A_4,A_3]=0$$
$`(37a)`$
$$A_{4x}=A_{3t}.$$
$`(37b)`$
Let the connections $`A_4,A_3`$ have the forms
$$A_3=\left(\begin{array}{ccc}0& rS_1& irS_2\\ rS_1& 0& S_3\\ irS_2& S_3& 0\end{array}\right)$$
$`(38a)`$
$$A_4=$$
$$\left(\begin{array}{ccc}0& ir[2iS_3S_{2y}2iS_2S_{3y}+iuS_1]& r[2S_3S_{1y}2S_1S_{3y}uS_2]\\ ir[2iS_3S_{2y}2iS_2S_{3y}+iuS_1]& 0& [ir^2(S^+S_y^{}S^{}S_y^+)uS_3]\\ r[2S_3S_{1y}2S_1S_{3y}uS_2& ir^2(S^+S_y^{}S^{}S_y^+)uS_3& 0\end{array}\right)$$
$`(38b)`$
Substituting (38) into (37) we get the M-I equation
$$iS_t=([S,S_y]+2iuS)_x$$
$`(39a)`$
$$u_x=\frac{1}{2i}tr(S[S_x,S_y])$$
$`(39b)`$
where
$$S=\left(\begin{array}{cc}S_3& S^{}\\ S^+& S_3\end{array}\right),S^\pm =S_1\pm iS_2.$$
$`(40)`$
For the M-I equation the corresponding geomerical equation (19) looks like
$$\left(\begin{array}{c}𝐞_1^{}\\ 𝐞_2^{}\\ 𝐞_3^{}\end{array}\right)_x=\lambda A_3\left(\begin{array}{c}𝐞_1^{}\\ 𝐞_2^{}\\ 𝐞_3^{}\end{array}\right)$$
$`(41a)`$
$$\left(\begin{array}{c}𝐞_1^{}\\ 𝐞_2^{}\\ 𝐞_3^{}\end{array}\right)_t=\lambda \left(\begin{array}{c}𝐞_1^{}\\ 𝐞_2^{}\\ 𝐞_3^{}\end{array}\right)_y\lambda A_4\left(\begin{array}{c}𝐞_1^{}\\ 𝐞_2^{}\\ 𝐞_3^{}\end{array}\right)$$
$`(41b)`$
and called the M-LXVI equation. As known , the M-LXVII and M-LXVI equations (41), (28) and (19) are some integrable (2+1)-dimensional extensions of the Serret-Frenet equation (SFE)
$$\left(\begin{array}{c}𝐞_1^{}\\ 𝐞_2^{}\\ 𝐞_3^{}\end{array}\right)_x=\left(\begin{array}{ccc}0& k& 0\\ \beta k& 0& \tau \\ 0& \tau & 0\end{array}\right)\left(\begin{array}{c}𝐞_1^{}\\ 𝐞_2^{}\\ 𝐞_3^{}\end{array}\right)$$
$`(42)`$
or the Codazzi-Mainardi equation (CME) for the surfaces.
The LR of the M-I equation follows from the LR of the SDYME and has the form
$$\mathrm{\Phi }_x^{}=\lambda A_3\mathrm{\Phi }^{}$$
$`(43a)`$
$$\mathrm{\Phi }_t^{}\lambda \mathrm{\Phi }_y^{}=\lambda A_4\mathrm{\Phi }^{}.$$
$`(43b)`$
As for ZE, we can rewrite the LR of the M-I equation in the standart form, in terms of 2$`\times `$2 - matrices
$$\mathrm{\Psi }_x^{}=U^{}\mathrm{\Psi }^{}$$
$`(44a)`$
$$\mathrm{\Psi }_t^{}=\lambda \mathrm{\Psi }_y^{}+V^{}\mathrm{\Psi }^{}$$
$`(44b)`$
where
$$U^{}=\frac{i\lambda }{2}S$$
$`(45a)`$
$$V^{}=\frac{\lambda }{4}([S,S_y]+2iuS).$$
$`(45b)`$
As above for the ZE (27), in the 1+1 dimensions $`(y=x)`$ instead of the equations (41) and (39), we get the (1+1)-dimensional M-LXVI equation
$$\left(\begin{array}{c}𝐞_1^{}\\ 𝐞_2^{}\\ 𝐞_3^{}\end{array}\right)_x=\lambda A_3\left(\begin{array}{c}𝐞_1^{}\\ 𝐞_2^{}\\ 𝐞_3^{}\end{array}\right)$$
$`(46a)`$
$$\left(\begin{array}{c}𝐞_1^{}\\ 𝐞_2^{}\\ 𝐞_3^{}\end{array}\right)_t=(\lambda ^2A_3+\lambda A_4)\left(\begin{array}{c}𝐞_1^{}\\ 𝐞_2^{}\\ 𝐞_3^{}\end{array}\right)$$
$`(46b)`$
and the Landau-Lifshitz equation
$$iS_t=\frac{1}{2}[S,S_{xx}].$$
$`(47a)`$
## 8 Gauge equivalence
It is well known that the ZE (33) and the M-I equation (39) are gauge and Lakshmanan equivalent (G-equivalent and L-equivalent) to each other. In our case this fact realize by the transformation
$$\left(\begin{array}{c}𝐞_1^{}\\ 𝐞_2^{}\\ 𝐞_3^{}\end{array}\right)=G\left(\begin{array}{c}𝐞_1\\ 𝐞_2\\ 𝐞_3\end{array}\right)$$
$`(48)`$
or
$$\mathrm{\Phi }^{}=h^1\mathrm{\Phi }$$
$`(49a)`$
or
$$\mathrm{\Psi }^{}=h^1\mathrm{\Psi }$$
$`(50)`$
where $`h`$ is the solution of the equations (29) or (30) as $`\lambda =0`$.
## 9 A nonisospectral case and the breaking solutions of the SDYME
Usually for the SDYME (10) the spectral parameter $`\lambda `$=constant. But in general it satisfies the following set of nonlinear equations
$$\lambda _{\xi _1}=\lambda \lambda _{\xi _3}$$
$`(51a)`$
$$\lambda _{\xi _2}=\lambda \lambda _{\xi _4}.$$
$`(51b)`$
These equations have the following solutions
$$\lambda =\frac{n_1\xi _3+n_3}{n_4n_1\xi _1}$$
$`(52a)`$
$$\lambda =\frac{m_1\xi _4+m_3}{m_4m_1\xi _2}.$$
$`(52b)`$
So that the general solution of the set (51) has the form
$$\lambda =\frac{n_1\xi _3+n_3+m_1\xi _4}{n_4n_1\xi _1m_1\xi _2}$$
$`(53)`$
where $`m_i,n_i=constants`$. The corresponding solution of the SDYME (10) is called the breaking (overlapping) solutions. In the case (18), i.e. for the ZE and the M-I equation the set of equations (21) takes the form
$$\lambda _x=0$$
$`(54a)`$
$$\lambda _t=\lambda \lambda _y$$
$`(54b)`$
and the solution (9) has the form
$$\lambda =\frac{n_1y+n_3}{n_4n_1t}.$$
$`(55)`$
## 10 Hierarchy of the M-LXVIII and Self-Dual Yang-Mills equations
The higher hierarchy of the M-LXVIII equation for the SDYME case we write in the form
$$\left(\begin{array}{c}𝐞_1\\ 𝐞_2\\ 𝐞_3\\ \mathrm{}\\ 𝐞_n\end{array}\right)_{\xi _1}=\underset{i=0}{\overset{k_1}{}}A_i\lambda ^i\left(\begin{array}{c}𝐞_1\\ 𝐞_2\\ 𝐞_3\\ \mathrm{}\\ 𝐞_n\end{array}\right)_{\xi _3}+\underset{i=0}{\overset{k_2}{}}B_i\lambda ^i\left(\begin{array}{c}𝐞_1\\ 𝐞_2\\ 𝐞_3\\ \mathrm{}\\ 𝐞_n\end{array}\right)$$
$`(56a)`$
$$\left(\begin{array}{c}𝐞_1\\ 𝐞_2\\ 𝐞_3\\ \mathrm{}\\ 𝐞_n\end{array}\right)_{\xi _2}=\underset{i=0}{\overset{k_3}{}}C_i\lambda ^i\left(\begin{array}{c}𝐞_1\\ 𝐞_2\\ 𝐞_3\\ \mathrm{}\\ 𝐞_n\end{array}\right)_{\xi _4}+\underset{i=0}{\overset{k_4}{}}D_i\lambda ^i\left(\begin{array}{c}𝐞_1\\ 𝐞_2\\ 𝐞_3\\ \mathrm{}\\ 𝐞_n\end{array}\right).$$
$`(56b)`$
The compatibility condition of these equations yields the higher hierarchy SDYME. As example, we consider the 3-dimensional case, work the notation (25) and $`n=3`$. Then insteod of (56) we obtain the M-LXVII equation in the form
$$\left(\begin{array}{c}𝐞_1\\ 𝐞_2\\ 𝐞_3\end{array}\right)_x=\underset{i=0}{\overset{k_2}{}}B_i\lambda ^i\left(\begin{array}{c}𝐞_1\\ 𝐞_2\\ 𝐞_3\end{array}\right)$$
$`(57a)`$
$$\left(\begin{array}{c}𝐞_1\\ 𝐞_2\\ 𝐞_3\end{array}\right)_t=\underset{i=0}{\overset{k_3}{}}C_i\lambda ^i\left(\begin{array}{c}𝐞_1\\ 𝐞_2\\ 𝐞_3\end{array}\right)_y+\underset{i=0}{\overset{k_4}{}}D_i\lambda ^i\left(\begin{array}{c}𝐞_1\\ 𝐞_2\\ 𝐞_3\end{array}\right).$$
$`(57b)`$
It is remarkable that using the equation (57) we can show that some known (2+1)-dimensional soliton equations are exact reductions of the higher hierarchy of the SDYME. Here we present some of them.
### 10.1 The (2+1)-dimensional mKdV equation
.
In the equation (57) we assume that
$$k_2=1,k_3=2,k_4=2,C_2=1,C_1=C_0=0.$$
$`(58)`$
Thus in this case the M-LXVII equation looks like
$$\left(\begin{array}{c}𝐞_1\\ 𝐞_2\\ 𝐞_3\end{array}\right)_x=(A_1\lambda A_3)\left(\begin{array}{c}𝐞_1\\ 𝐞_2\\ 𝐞_3\end{array}\right)$$
$`(59a)`$
$$\left(\begin{array}{c}𝐞_1\\ 𝐞_2\\ 𝐞_3\end{array}\right)_t=\lambda ^2\left(\begin{array}{c}𝐞_1\\ 𝐞_2\\ 𝐞_3\end{array}\right)_y+(D_2\lambda ^2+D_1\lambda +D_0)\left(\begin{array}{c}𝐞_1\\ 𝐞_2\\ 𝐞_3\end{array}\right)$$
$`(59b)`$
where
$$A_1=\left(\begin{array}{ccc}0& (q+p)& i(qp)\\ q+p& 0& 0\\ i(qp)& 0& 0\end{array}\right)$$
$`(60)`$
and $`A_3,D_k`$ are some matrices . Then the complex functions $`q,p`$ satisfy the (2+1)-dimensional complex mKdV equation
$$q_t+q_{xxy}(qv_1)_xqv_2=0,$$
$`(61a)`$
$$v_{1x}=2E(\overline{q}q)_y$$
$`(61b)`$
$$v_{2x}=2E(\overline{q}q_{xy}\overline{q}_{xy}q)$$
$`(61c)`$
If $`p=q`$ is real, we get the following (2+1)- dimensional mKdV equation
$$q_t+q_{xxy}(qv_1)_x=0,$$
$`(62a)`$
$$v_{1x}=2E(q^2)_y=4Eqq_y.$$
$`(62b)`$
### 10.2 The (2+1)-dimensional derivative NLSE
.
In (57) now we put
$$k_2=k_3=k_4=2;C_2=2c,C_1=C_0=0,B_2=cA_3,B_2=2cA_1.$$
$`(63)`$
So that the M-LXVII equation takes the form
$$\left(\begin{array}{c}𝐞_1\\ 𝐞_2\\ 𝐞_3\end{array}\right)_x=(A_3\lambda ^2+A_1\lambda )\left(\begin{array}{c}𝐞_1\\ 𝐞_2\\ 𝐞_3\end{array}\right)$$
$`(64a)`$
$$\left(\begin{array}{c}𝐞_1\\ 𝐞_2\\ 𝐞_3\end{array}\right)_t=\lambda ^2\left(\begin{array}{c}𝐞_1\\ 𝐞_2\\ 𝐞_3\end{array}\right)_y+(D_2\lambda ^2+D_1\lambda +D_0)\left(\begin{array}{c}𝐞_1\\ 𝐞_2\\ 𝐞_3\end{array}\right)$$
$`(64b)`$
where $`A_1`$ is given by (60). Then the complex functions $`q,p`$ satisfy the (2+1)-dimensional derivative NLSE
$$iq_t=q_{xy}2ic(vq)_x$$
$`(65a)`$
$$ip_t=p_{xy}+2ic(vq)_x$$
$`(65b)`$
$$v_x=2(pq)_y$$
$`(65c)`$
which is the Strachan equation .
### 10.3 The M-III<sub>q</sub> equation
.
Now we consider the case
$$k_2=k_3=k_4=2,C_2=2cI,C_1=2dI,C_0=0,B_2=cA_3,B_1=dA_3+2cA_1,B_0=dA_1$$
$`(66)`$
for which the M-LXVII equation has the form
$$\left(\begin{array}{c}𝐞_1\\ 𝐞_2\\ 𝐞_3\end{array}\right)_x=[A_3(c\lambda ^2+d\lambda )+A_1(2c\lambda +d)]\left(\begin{array}{c}𝐞_1\\ 𝐞_2\\ 𝐞_3\end{array}\right)$$
$`(67a)`$
$$\left(\begin{array}{c}𝐞_1\\ 𝐞_2\\ 𝐞_3\end{array}\right)_t=2(c\lambda ^2+d\lambda )\left(\begin{array}{c}𝐞_1\\ 𝐞_2\\ 𝐞_3\end{array}\right)_y+(D_2\lambda ^2+D_1\lambda +D_0)\left(\begin{array}{c}𝐞_1\\ 𝐞_2\\ 𝐞_3\end{array}\right)$$
$`(67b)`$
where $`A_1`$ is given by (60). Then the complex functions $`q,p`$ satisfy the (2+1)-dimensional M-III<sub>q</sub> equation
$$iq_t=q_{xy}2ic(vq)_x+d^2vq$$
$`(68a)`$
$$ip_t=p_{xy}+2ic(vq)_x+d^2vp$$
$`(68b)`$
$$v_x=2(pq)_y.$$
$`(68c)`$
The M-III<sub>q</sub> equation (68) admits two integrable reductions: the Strachan equation (65) as $`d=0`$ and the ZE (27) as $`c=0`$. If we rewrite the equations (67) in terms of $`2\times 2`$ \- matrices, we get the following LR of the M-III<sub>q</sub> equation
$$\mathrm{\Psi }_x=U\mathrm{\Psi }$$
$`(69a)`$
$$\mathrm{\Psi }_t=2(c\lambda ^2+d\lambda )\mathrm{\Psi }_y+V\mathrm{\Psi }$$
$`(69b)`$
where
$$U=i[(c\lambda ^2+d\lambda )\sigma _3+(2c\lambda +d)G],G=\left(\begin{array}{cc}0& q\\ p& 0\end{array}\right)$$
$`(70a)`$
$$V=B_2\lambda ^2+B_1\lambda +B_0.$$
$`(70b)`$
Here
$$B_0=\frac{d}{2c}B_1\frac{d^2}{4c^2}B_2,B_2=2ic^2v\sigma _3,$$
$$B_1=2icdv\sigma _3+2cG_y\sigma _34ic^2vG.$$
$`(71)`$
### 10.4 The M-XXII<sub>q</sub> equation
Now let
$$k_2=k_3=k_4=2.B_2=A_3,B_1=A_1,B_0=\frac{pq}{4}A_3,C_2=2I,C_1=C_0=0.$$
$`(72)`$
Then the functions $`q,p`$ satisfy the following $`MXXII_q`$ equations
$$iq_t+q_{yx}+\frac{i}{2}[(v_1q)_xv_2qqpq_y]=0$$
$`(73a)`$
$$ip_tp_{yx}+\frac{i}{2}[(v_1p)_x+v_2pqpp_y]=0$$
$`(73b)`$
$$v_{1x}=(pq)_y$$
$`(73c)`$
$$v_{2x}=p_{yx}qpq_{yx}.$$
$`(73d)`$
This set of equations is the G- and L-equivalent counterpart of the M-XXII<sub>s</sub> equation (spin system). The LR of this equation has the form
$$\mathrm{\Psi }_{2x}=\{i(\lambda ^2\frac{pq}{4})\sigma _3+\lambda Q\}\mathrm{\Psi }_2$$
$`(74a)`$
$$\mathrm{\Psi }_{2t}=2\lambda ^2\mathrm{\Psi }_{2y}+(\lambda ^2B_2+\lambda B_1+B_0)\mathrm{\Psi }_2$$
$`(74b)`$
with
$$Q=\left(\begin{array}{cc}0& q\\ p& 0\end{array}\right),B_2=\frac{i}{2}v_1\sigma _3,B_1=i\sigma _3Q_y\frac{1}{2}v_1Q,B_0=\frac{1}{4}v_2\frac{i}{8}pqv_1.$$
$`(75)`$
Now let us consider the following transformation
$$q^{}=q\mathrm{exp}[\frac{i}{2}_x^1(pq)],p^{}=p\mathrm{exp}[\frac{i}{2}_x^1(pq)].$$
$`(76)`$
Then the new variables $`p^{},q^{}`$ satisfy the Strachan equation
$$iq_t^{}+q_{xy}^{}+i(v^{}q^{})_x=0,$$
$`(77a)`$
$$ip^{}p_{xy}^{}+i(v^{}p^{})=0,$$
$`(77b)`$
$$v_x^{}=E(p^{}q^{})_y.$$
$`(77c)`$
We see that the M-XXII<sub>q</sub> equation (73) and the Strachan equation (70) is gauge eqivalent to each other. The tranformation (76) induces the following tranformation of the Jost function$`\mathrm{\Psi }_1`$
$$\mathrm{\Psi }_1=f^1\mathrm{\Psi }_2$$
$`(78)`$
where $`\mathrm{\Psi }_1`$ is the solution of the equations (30) as $`d=0`$ and
$$f=\mathrm{exp}(\frac{i}{4}_x^1q^2\sigma _3)=\mathrm{\Psi }_1^1_{\lambda =0}.$$
$`(79)`$
Then the new Jost function $`\mathrm{\Psi }_2`$ satisfies the following set of equations
$$\mathrm{\Psi }_{2x}=\{i\lambda ^2\sigma _3+\lambda Q^{}\}\mathrm{\Psi }_2$$
$`(80a)`$
$$\mathrm{\Psi }_{2t}=2\lambda ^2\mathrm{\Psi }_{2y}+\{\lambda ^2B_2^{}+\lambda B_1^{}+B_0^{})\}\mathrm{\Psi }_2$$
$`(80b)`$
with
$$Q=\left(\begin{array}{cc}0& q^{}\\ p^{}& 0\end{array}\right)$$
$`(81)`$
and $`B_j^{}`$ are given in . Note that the Strachan (77), (65) and MXXII<sub>q</sub> (73) equations are the simplest (2+1)-dimensional extensions of the following known NLSE
$$iq_t+q_{xx}+i(pq^2)_x=0$$
$`(82a)`$
$$ip_tp_{xx}+i(qp^2)_x=0$$
$`(82b)`$
and
$$iq_t+q_{xx}+ipqq_x=0$$
$`(83a)`$
$$ip_tp_{xx}+ipqp_x=0$$
$`(83b)`$
respectively. It is well known that these equations are gauge equivalent to each other .
## 11 The M-LXI and M-LXII equations and soliton equations in $`d=3`$ dimensions
### 11.1 The M-LXI equation
The M-LXI equation in $`d=3`$ dimensions has the form
$$\left(\begin{array}{c}𝐞_1\\ 𝐞_2\\ 𝐞_3\\ \mathrm{}\\ 𝐞_n\end{array}\right)_{\xi _1}=A\left(\begin{array}{c}𝐞_1\\ 𝐞_2\\ 𝐞_3\\ \mathrm{}\\ 𝐞_n\end{array}\right)$$
$`(84a)`$
$$\left(\begin{array}{c}𝐞_1\\ 𝐞_2\\ 𝐞_3\\ \mathrm{}\\ 𝐞_n\end{array}\right)_{\xi _2}=B\left(\begin{array}{c}𝐞_1\\ 𝐞_2\\ 𝐞_3\\ \mathrm{}\\ 𝐞_n\end{array}\right)$$
$`(84b)`$
$$\left(\begin{array}{c}𝐞_1\\ 𝐞_2\\ 𝐞_3\\ \mathrm{}\\ 𝐞_n\end{array}\right)_{\xi _3}=C\left(\begin{array}{c}𝐞_1\\ 𝐞_2\\ 𝐞_3\\ \mathrm{}\\ 𝐞_n\end{array}\right).$$
$`(84c)`$
### 11.2 The M-LXII equation
The compatibility condition of the equations (84) gives the 3-dimensional M-LXII equation
$$A_{\xi _2}B_{\xi _1}+[A,B]=0$$
$`(85a)`$
$$A_{\xi _3}C_{\xi _1}+[A,C]=0$$
$`(85b)`$
$$C_{\xi _2}B_{\xi _3}+[C,B]=0.$$
$`(85c)`$
This equation is the particular case of the Bogomolny equation (BE)
$$\mathrm{\Psi }_{\xi _3}+[\mathrm{\Psi },C]+A_{\xi _2}B_{\xi _1}+[A,B]=0$$
$`(86a)`$
$$\mathrm{\Psi }_{\xi _2}+[\mathrm{\Psi },B]+A_{\xi _3}C_{\xi _1}+[A,C]=0$$
$`(86b)`$
$$\mathrm{\Psi }_{\xi _1}+[\mathrm{\Psi },A]+C_{\xi _2}B_{\xi _3}+[C,B]=0.$$
$`(86c)`$
In fact, frome hence as $`\mathrm{\Psi }=0`$ we obtain the M-LXII equation (85). As well known that the BE (86) is integrable (see, e.g. the book ). As the particular case of the integrable BE (86), the M-LXII equation is also integrable. The corresponding LR has the form
$$\mathrm{\Phi }_{\xi _1}\lambda \mathrm{\Phi }_{\xi _3}=[iC\lambda (A+iB)]\mathrm{\Phi }$$
$`(87a)`$
$$\mathrm{\Phi }_{\xi _2}\lambda \mathrm{\Phi }_{\xi _4}=[AiB\lambda iC]\mathrm{\Phi }.$$
$`(87b)`$
Let us consider the case $`n=3,\xi _1=x,\xi _2=y,\xi _3=t`$. Then the mM-LXI amd mM-LXII equations take the form
$$\left(\begin{array}{c}𝐞_1\\ 𝐞_2\\ 𝐞_3\end{array}\right)_x=A\left(\begin{array}{c}𝐞_1\\ 𝐞_2\\ 𝐞_3\end{array}\right)$$
$`(88a)`$
$$\left(\begin{array}{c}𝐞_1\\ 𝐞_2\\ 𝐞_3\end{array}\right)_y=B\left(\begin{array}{c}𝐞_1\\ 𝐞_2\\ 𝐞_3\end{array}\right)$$
$`(88b)`$
$$\left(\begin{array}{c}𝐞_1\\ 𝐞_2\\ 𝐞_3\end{array}\right)_t=C\left(\begin{array}{c}𝐞_1\\ 𝐞_2\\ 𝐞_3\end{array}\right)$$
$`(88c)`$
and
$$A_yB_x+[A,B]=0$$
$`(89a)`$
$$A_tC_x+[A,C]=0$$
$`(89b)`$
$$C_yB_t+[C,B]=0$$
$`(89c)`$
where
$$A=\left(\begin{array}{ccc}0& k& \sigma \\ \beta k& 0& \tau \\ \beta \sigma & \tau & 0\end{array}\right),B=\left(\begin{array}{ccc}0& m_3& m_2\\ \beta m_3& 0& m_1\\ \beta m_2& m_1& 0\end{array}\right)$$
$$D=\left(\begin{array}{ccc}0& \omega _3& \omega _2\\ \beta \omega _3& 0& \omega _1\\ \beta \omega _2& \omega _1& 0\end{array}\right).$$
$`(90)`$
The mM-LXII equation contents many (and perhaps all?) known soliton equations (see, e.g. the refs.\[8,18-24,38-40,43-44\]). We note that in the case $`\sigma =0`$ (89) is called the M-LXII equation. Some examples as follows.
#### 11.2.1 The Ishimori and DS equations
In this section, we obtain the Ishimori (IE) and DS equations from the M-LXI and M-LXII $`(\sigma =0)`$ equations as some exact reductions. The IE reads as
$$𝐒_t=𝐒(𝐒_{xx}+\alpha ^2𝐒_{yy})+u_x𝐒_y+_y𝐒_x$$
$`(91a)`$
$$u_{xx}\alpha ^2u_{yy}=2\alpha ^2𝐒(𝐒_x𝐒_y).$$
$`(91b)`$
We take the following identification
$$𝐒=𝐞_1.$$
$`(92)`$
In this case we have
$$m_1=_x^1[\tau _y\frac{\beta }{2\alpha ^2}M_2^{Ish}u]$$
$`(93a)`$
$$m_2=\frac{1}{2\alpha ^2k}M_2^{Ish}u$$
$`(93b)`$
$$m_3=_x^1[k_y+\frac{\tau }{2\alpha ^2k}M_2^{Ish}u]$$
$`(93c)`$
$$M_2^{IE}u=u_{xx}\alpha ^2u_{yy}$$
$`(94)`$
and
$$\omega _1=\frac{1}{k}[\omega _{2x}+\tau \omega _3]$$
$`(95a)`$
$$\omega _2=k_x\alpha ^2(m_{3y}+m_2m_1)+im_2u_x$$
$`(95b)`$
$$\omega _3=k\tau +\alpha ^2(m_{2y}m_3m_1)+iku_y+im_3u_x.$$
$`(95c)`$
Now let us introduce two complex functions $`q,p`$, according to the formulae
$$q=a_1e^{ib_1},p=a_2e^{ib_2}.$$
$`(96)`$
Let $`a_j,b_j`$ have the forms
$$a_1^2=\frac{1}{4}k^2+\frac{|\alpha |^2}{4}(m_3^2+m_2^2)\frac{1}{2}\alpha _Rkm_3\frac{1}{2}\alpha _Ikm_2$$
$`(97a)`$
$$b_1=_x^1\{\frac{\gamma _1}{2ia_1^2}(\overline{A}A+D\overline{D})\}$$
$`(97b)`$
$$a_2^2=\frac{1}{4}k^2+\frac{|\alpha |^2}{4}(m_3^2+m_2^2)+\frac{1}{2}\alpha _Rkm_3\frac{1}{2}\alpha _Ikm_2$$
$`(97c)`$
$$b_2=_x^1\{\frac{\gamma _2}{2ia_2^2}(A\overline{A}+\overline{D}D)\}$$
$`(97d)`$
where
$$\gamma _1=i\{\frac{1}{2}k^2\tau +\frac{|\alpha |^2}{2}(m_3km_1+m_2k_y)$$
$$\frac{1}{2}\alpha _R(k^2m_1+m_3k\tau +m_2k_x)+\frac{1}{2}\alpha _I[k(2k_ym_{3x})k_xm_3]\}.$$
$`(98a)`$
$$\gamma _2=i\{\frac{1}{2}k^2\tau +\frac{|\alpha |^2}{2}(m_3km_1+m_2k_y)+$$
$$\frac{1}{2}\alpha _R(k^2m_1+m_3k\tau +m_2k_x)+\frac{1}{2}\alpha _I[k(2k_ym_{3x})k_xm_3]\}.$$
$`(98b)`$
Here $`\alpha =\alpha _R+i\alpha _I`$. In this case, $`q,p`$ satisfy the following DS equation
$$iq_t+q_{xx}+\alpha ^2q_{yy}+vq=0$$
$`(99a)`$
$$ip_t+p_{xx}+\alpha ^2p_{yy}+vp=0$$
$`(99b)`$
$$v_{xx}\alpha ^2v_{yy}+2[(pq)_{xx}+\alpha ^2(pq)_{yy}]=0.$$
$`(99c)`$
It is means that the IE (91) and the DS equation (99) are L-equivalent to each other. As well known that these equations are G-equivalent to each other . A few comments are in order.
i) From these results, we get the Ishimori I and DS-I equations as $`\alpha _R=1,\alpha _I=0`$
ii) the Ishimori II and DS-II equations as $`\alpha _R=0,\alpha _I=1`$.
iii) For DS-II equation we have
$$pq=q^2=p^2$$
$`(100)`$
iv) at the same time, for the DS-I equation we obtain
$$pqq^2p^2$$
$`(101)`$
$$q^2=p^2km_3$$
$`(102)`$
$$pq=(pq)_R+i(pq)_I$$
$`(103)`$
so that $`pq`$ is the complex quantity.
#### 11.2.2 The KP and M-X equations
The (2+1)-dimensional mM-LXI equation in plane has the form
$$\left(\begin{array}{cc}𝐞_1& \\ 𝐞_2& \end{array}\right)_x=A_p\left(\begin{array}{cc}𝐞_1& \\ 𝐞_2& \end{array}\right),\left(\begin{array}{cc}𝐞_1& \\ 𝐞_2& \end{array}\right)_y=B_p\left(\begin{array}{cc}𝐞_1& \\ 𝐞_2& \end{array}\right),\left(\begin{array}{cc}𝐞_1& \\ 𝐞_2& \end{array}\right)_t=D_p\left(\begin{array}{cc}𝐞_1& \\ 𝐞_2& \end{array}\right)$$
$`(104)`$
where
$$A_p=\left(\begin{array}{cc}0& k\\ \beta k& 0\end{array}\right),B_p=\left(\begin{array}{cc}0& m_3\\ \beta m_3& 0\end{array}\right)$$
$$D_p=\left(\begin{array}{cc}0& \omega _3\\ \beta \omega _3& 0\end{array}\right).$$
$`(105)`$
In the plane case the mM-LXII equation takes the following simple form
$$k_y=m_{3x}$$
$`(106a)`$
$$k_t=\omega _{3x}$$
$`(106b)`$
$$m_{3t}=\omega _{3y}.$$
$`(106c)`$
Hence we get
$$m_3=_x^1k_y.$$
$`(107)`$
The NEE has the form (106b). Many (2+1)-dimensional integrable equations such as the Kadomtsev-Petviashvili, Novikov-Veselov (NV), mNV, KNV, (2+1)-KdV, mKdV equations are the integrable reductions of the M-LXII equation (106). For example, let us show that the KP and mKP equations are exact reductions of the mM-LXII equation (106). Consider the M-X equation
$$𝐒_t=\frac{\omega _3}{k}𝐒_x$$
$`(108)`$
where
$$\omega _3=k_{xx}3k^23\alpha ^2_x^1m_{3y}.$$
$`(109)`$
If we put $`𝐒=𝐞_1`$ then from (106) we obtain the L-equivalent counterpart of the M-X equation which is the KP equation
$$k_t+6kk_x+k_{xxx}+3\alpha ^2m_{3y}=0$$
$`(110a)`$
$$m_{3x}=k_y.$$
$`(110b)`$
As known the LR of this equation is given by
$$\alpha \psi _y+\psi _{xx}+k\psi =0$$
$`(111a)`$
$$\psi _t+4\psi _{xxx}+6k\psi _x+3(k_x\alpha m_3)\psi =0.$$
$`(111b)`$
#### 11.2.3 The Zakharov and M-IX equations
Now we find the connection between the Myrzakulov IX (M-IX) equation and the curves (the M-LXI equation). The M-IX equation reads as
$$𝐒_t=𝐒M_1𝐒+A_2𝐒_x+A_1𝐒_y$$
$`(112a)`$
$$M_2u=2\alpha ^2𝐒(𝐒_x𝐒_y)$$
$`(112b)`$
where $`\alpha ,b,a`$= consts and
$$M_1=\alpha ^2\frac{^2}{y^2}+4\alpha (ba)\frac{^2}{xy}+4(a^22abb)\frac{^2}{x^2},$$
$$M_2=\alpha ^2\frac{^2}{y^2}2\alpha (2a+1)\frac{^2}{xy}+4a(a+1)\frac{^2}{x^2},$$
$$A_1=i\{\alpha (2b+1)u_y2(2ab+a+b)u_x\},$$
$$A_2=i\{4\alpha ^1(2a^2b+a^2+2ab+b)u_x2(2ab+a+b)u_y\}.$$
$`(113)`$
The M-IX equation was introduced in and is integrable. It admits several integrable reductions: 1) the Ishimori equation as $`a=b=\frac{1}{2}`$; 2) the M-VIII equation as $`a=b=1`$ and so on . In this case we have
$$m_1=_x^1[\tau _y\frac{\beta }{2\alpha ^2}M_2u]$$
$`(114a)`$
$$m_2=\frac{1}{2\alpha ^2k}M_2u$$
$`(114b)`$
$$m_3=_x^1[k_y+\frac{\tau }{2\alpha ^2k}M_2u]$$
$`(114c)`$
and
$$\omega _1=\frac{1}{k}[\omega _{2x}+\tau \omega _3],$$
$`(115a)`$
$$\omega _2=4(a^22abb)k_x4\alpha (ba)k_y\alpha ^2(m_{3y}+m_2m_1)+m_2A_1$$
$`(115b)`$
$$\omega _3=4(a^22abb)k\tau 4\alpha (ba)km_1+\alpha ^2(m_{2y}m_3m_1)+kA_2+m_3A_1.$$
$`(115c)`$
Functions $`q,p`$ are given by (96) with
$$a_1^2=\frac{|a|^2}{|b|^2}a_1^^2=\frac{|a|^2}{|b|^2}\{(l+1)^2k^2+\frac{|\alpha |^2}{4}(m_3^2+m_2^2)(l+1)\alpha _Rkm_3(l+1)\alpha _Ikm_2\}$$
$`(116a)`$
$$b_1=_x^1\{\frac{\gamma _1}{2ia_1^^2}(\overline{A}A+D\overline{D})\}$$
$`(116b)`$
$$a_2^2=\frac{|b|^2}{|a|^2}a_2^^2=\frac{|b|^2}{|a|^2}\{l^2k^2+\frac{|\alpha |^2}{4}(m_3^2+m_2^2)l\alpha _Rkm_3+l\alpha _Ikm_2\}$$
$`(116c)`$
$$b_2=_x^1\{\frac{\gamma _2}{2ia_2^^2}(A\overline{A}+\overline{D}D)$$
$`(116d)`$
where
$$\gamma _1=i\{2(l+1)^2k^2\tau +\frac{|\alpha |^2}{2}(m_3km_1+m_2k_y)$$
$$(l+1)\alpha _R[k^2m_1+m_3k\tau +m_2k_x]+(l+1)\alpha _I[k(2k_ym_{3x})k_xm_3]\}$$
$`(117a)`$
$$\gamma _2=i\{2l^2k^2\tau +\frac{|\alpha |^2}{2}(m_3km_1+m_2k_y)$$
$$l\alpha _R(k^2m_1+m_3k\tau +m_2k_x)l\alpha _I[k(2k_ym_{3x})k_xm_3]\}.$$
$`(117b)`$
Here $`\alpha =\alpha _R+i\alpha _I`$. In this case, $`q,p`$ satisfy the following other Zakharov equation
$$iq_t+M_1q+vq=0$$
$`(118a)`$
$$ip_tM_1pvp=0$$
$`(118b)`$
$$M_2v=2M_1(pq)$$
$`(118c)`$
which is integrable and admits several particular cases.
As well known the M-IX equation admits several reductions: 1) the M-IXA equation as $`\alpha _R=1,\alpha _I=0`$; 2) the M-IXB equation as $`\alpha _R=0,\alpha _I=1`$; 3) the M-VIII equation as $`a=b=1`$ 4) the IE $`a=b=\frac{1}{2}`$ and so on. The corresponding versions of the ZE (118), we obtain as the corresponding values of the parameter $`\alpha `$.
## 12 Conclusion
In this paper we analyzed the M-LXVIII equation. We have found the some integrable reductions of this equation. Also we have shown that the Zakharov and its spin counterpart the M-I equation are exact reductions of the SDYME. The higher hierarchy of the SDYME was introduced. Using this hierarchy it was shown that several simplest soliton equations in 2+1 dimensions such as mKdV, derivative NLS and M-III<sub>q</sub> equations and so on are also its exact reductions.
Finally we would like ask you, dear colleaque, if you know, have or will have any results on multidimensional soliton equations, soliton geometry and the Yang-Mills equations, please, inform me (R.M.) or send me a hard copy of your papers. Also any comments, proposals and questions are welcome.
## 13 Acknowledgments
R.M. is grateful to M.Lakshmanan for hospitality during visits, many useful discussions and for financial support. Special thanks to Radha Balakrishnan and M.Daniel for discussions. We thank G.N.Nugmanova for helpful comments on the manuscript.
## 14 Tasks
Task-1: Please find the breaking solutions (instantons, monopols, dions and so on ) of the SDYME.
Task-2: Please find the breaking solutions of the Bogomolny equation.
Task-3: Please find the breaking solutions of the Ishimori and DS equations.
Task-4: Please find the breaking solutions of the M-IX and Zakharov equations.
Task-5: Please consider the above presented results from twistor point of view. |
warning/0002/gr-qc0002043.html | ar5iv | text | # Computing radiation from Kerr black holes: Generalization of the Sasaki-Nakamura equation
## I Introduction
In 1973, Teukolsky derived a single partial differential equation describing the evolution of perturbations to rotating (Kerr) black holes. This master equation gives the linearized evolution of fields that arise from a perturbing source of stress energy — the charge and current densities associated with the perturbation — to the (vacuum) black hole background. The solutions of the homogeneous version of this equation describe the propagation of radiation in black hole spacetimes. Thus, a common use of this formalism is to study the radiation emitted by matter in the environment of a black hole. In some cases, one can use such an analysis to study back reaction, determining how the perturbing source evolves as radiation carries away energy and angular momentum. A beautiful feature of the master equation is that it describes radiation fields of arbitrary spin weight $`s`$. It has been used extensively to study scalar ($`s=0`$), electromagnetic ($`s=\pm 1`$), and gravitational ($`s=\pm 2`$) radiation in Kerr spacetimes.
The master equation is often<sup>*</sup><sup>*</sup>*Separation is not always used. There is also a body of work that uses the master equation to evolve initial data. This approach has been extensively used to study the endpoint of binary black hole collisions; see Ref. and references therein. solved by introducing a multipolar decomposition of the radiation field. The solution separates into functions of the Boyer-Lindquist coordinates:
$${}_{s}{}^{}\mathrm{\Psi }=\underset{l,m,\omega }{}R_{lm\omega }(r){}_{s}{}^{}S_{lm}^{a\omega }(\theta )e^{im\varphi }e^{i\omega t};$$
(1)
each function is governed by an ordinary differential equation. (The precise meaning of $`{}_{s}{}^{}\mathrm{\Psi }`$ is described in Sec. II.) The $`t`$ and $`\varphi `$ dependences are trivial. The $`\theta `$ dependence is more involved, but can be evaluated in a straightforward matter. The functions $`{}_{s}{}^{}S_{lm}^{a\omega }`$ are spin-weighted spheroidal harmonics, which are generalizations of spin-weighted spherical harmonics to a spheroidal geometry. The spin-weighted spherical harmonics in turn are generalizations of spherical harmonics that encode the rotation properties of spin $`s`$ fields; see Refs. for further discussion. A detailed algorithm for computing $`{}_{s}{}^{}S_{lm}^{a\omega }(\theta )`$ is given in Ref. . For the purposes of this paper, the $`\theta `$ dependence is considered known.
The radial dependence, $`R_{lm\omega }(r)R(r)`$, on the other hand, can be rather difficult to calculate in practice, particularly in a numerical computation. The fundamental reason for this difficulty is the nature of the equation that governs $`R(r)`$: this equation (the Teukolsky equation) has a long-ranged potential. In source-free form, it can be written
$$\frac{d^2R}{dr^2}+F_T(r)\frac{dR}{dr^{}}+\left[\omega ^2U_T(r)\right]R=0,$$
(2)
where $`\omega `$ is the frequency of the radiation mode and
$$r^{}=r+\frac{2Mr_+}{r_+r_{}}\mathrm{ln}\frac{rr_+}{2M}\frac{2Mr_{}}{r_+r_{}}\mathrm{ln}\frac{rr_{}}{2M}$$
(3)
is the Kerr “tortoise coordinate”. The potentials $`U_T(r)`$ and $`F_T(r)`$ are rather complicated; they encode the most interesting features of wave propagation in black hole spacetimes, such as scatter from spacetime curvatureAs shown by Leonard and Poisson , the phenomenon of tails (delayed propagation due to scatter from spacetime curvature) is to leading order independent of $`s`$, and is encoded in the logarithmic behavior of $`r^{}`$, not the potentials. and superradiant scattering (radiation whose scattered amplitude exceeds the ingoing amplitude due to extraction of energy from the black hole’s spin). For large $`r`$,
$`F_T(r)`$ $`=`$ $`{\displaystyle \frac{2(1+s)}{r}}{\displaystyle \frac{2M(2+s)}{r^2}}+O(1/r^3),`$ (4)
$`U_T(r)`$ $`=`$ $`{\displaystyle \frac{4is\omega }{r}}+{\displaystyle \frac{\lambda +2am\omega +8IMs\omega }{r^2}}+O(1/r^3).`$ (5)
\[The quantity $`\lambda `$ is related to the eigenvalues of the $`\theta `$ dependence; see Ref. for details.\] For large $`r`$, $`F_T(r)`$ and $`U_T(r)`$ fall off only as $`1/r`$ — they are long-ranged, like the Coulomb potential. The solution of Eq. (2) for large $`r`$ is
$$R=C_1\frac{e^{i\omega r^{}}}{r}+C_2\frac{e^{i\omega r^{}}}{r^{2s+1}}.$$
(6)
The complex constants $`C_1`$ and $`C_2`$ are determined by boundary conditions. This asymptotic solution illustrates the difficulty in solving the Teukolsky equation: the coefficient of $`e^{i\omega r^{}}`$ differs from the coefficient of $`e^{i\omega r^{}}`$ by $`r^{2s}`$. For negative $`s`$, this becomes extremely large — large enough that the ingoing $`e^{i\omega r^{}}`$ piece will eventually be entirely lost in any numerical computation due to round-off error. Hence, for negative $`s`$, it is nearly impossible to set proper boundary conditions on the solution’s phase at large $`r`$. Similarly, for positive $`s`$ it is difficult to set boundary conditions near the event horizonOne can expand the potentials near the horizon and see that they die away slowly as $`r^{}\mathrm{}`$ (which corresponds to $`rr_+`$). The actual form is somewhat messy, and is not given here explicitly.. For $`r`$ very close to $`r_+=M+\sqrt{M^2a^2}`$ (the location of the event horizon in Boyer-Lindquist coordinates), the solution is
$$R=C_3\mathrm{\Delta }^se^{ipr^{}}+C_4e^{ipr^{}}.$$
(7)
I have introduced $`p=\omega m\omega _+`$, where $`\omega _+=a/2Mr_+`$ is the angular velocity at which observers at the horizon are seen to rotate. Because the Boyer-Lindquist coordinates $`t`$ and $`\varphi `$ become twisted and entangled near the horizon, $`p`$ describes the frequency of wave modes in that region. The factor $`\mathrm{\Delta }=r^22Mr+a^2`$ goes to zero at the event horizon. Hence, for positive $`s`$, the ingoing solution swamps the outgoing solution as one approaches the horizon. Whether one chooses positive or negative $`s`$, there exists a domain in which one cannot accurately compute numerical solutions by directly integrating the homogeneous Teukolsky equation.
Various approaches have been discussed to circumvent this difficulty. One of the first was introduced by Teukolsky and Press . Their approach used the fact that, for a given $`|s|`$, the solutions to Eq. (2) for $`s=+|s|`$ and $`s=|s|`$ are physically equivalent: there exist rules to take the positive $`s`$ solution to the negative $`s`$ solution, and vice versa. Thus, one can initially choose $`s`$ so that the solution is well-behaved in the initial $`r`$ domain, and then “switch horses” and integrate with the other sign of $`s`$ as the integration approaches the other asymptotic domain.
A somewhat more elegant way to compute $`R`$ was developed by Chandrasekhar . As already noted, the poor behavior of the solutions (6) and (7) is due to the long-ranged nature of the Teukolsky equation’s potentials. Rather than try to work with an equation that is simply not well-behaved, one should find transformations which relate the Teukolsky solution $`R`$ to the solution $`X`$ of some equation whose potentials are short-ranged. For example, when the black hole spin is zero, black hole perturbations can be described using the generalized Regge-Wheeler equation :
$$\left[\frac{d^2}{dr^2}+\omega ^2V_{RW}(r,s)\right]X=0,$$
(8)
where
$$V_{RW}(r,s)=f\left[\frac{l(l+1)}{r^2}\frac{2(s^21)M}{r^3}\right].$$
(9)
(Here, $`f=12M/r`$.) This potential dies away faster than $`1/r`$, and so is short ranged. Chandrasekhar showed (for specific choices of $`s`$) that solutions to Eq. (8) and solutions to Eq. (2) (with $`a=0`$) are related by simple rules. (Below I generalize these rules to any value of $`s`$.) This is extremely useful for numerical work: one can integrate Eq. (8) to accurately compute $`X`$ and then transform to $`R`$. Note that in Ref. Press and Teukolsky introduced a transformation that, in essence, transformed to a function governed by an equation with a better behaved potential. They did not, however, discuss the nature of the transformation in terms of the rangedness of the potentials; Chandrasekhar appears to have been the first to systematically approach this problem with the viewpoint that the long-ranged potential was the key issue. Note also that Chandrasekhar’s notation is rather different from that used here; I use a notation similar to that used in . The transformation given in Refs. and is for $`s=2`$; a rule for $`s=1`$ is given in .
For spinning black holes, perhaps the most elegant generalization of Chandrasekhar’s approach was given by Sasaki and Nakamura . They derive a transformation rule which relates $`R`$ for any physical spin $`a`$ to the solution $`X`$ of an equation whose potentials are short ranged. The transformation and short-ranged potentials are designed such that if $`a=0`$, the potentials reduce to the Regge-Wheeler potentials. This approach is very natural in the sense that its solutions are monotonic with respect to spin, ranging from the Schwarzschild value $`a=0`$ to the extreme Kerr limit $`a=M`$. Chandrasekhar and Detweiler also investigated several transformations relating $`R`$ to a short-ranged solution $`X`$ . In some (but not all) cases, these transformations are governed by equations which reduce to the Regge-Wheeler equation in the $`a=0`$ limit; however, the equations themselves are often not as “nice” to work with. For example, the equation for $`X`$ is sometimes given in terms of a frequency dependent variable $`\widehat{r}^{}(\omega )`$ (see Ref. ) which is different from the “usual” tortoise coordinate $`r^{}`$ \[cf. Eq. (3)\]. This coordinate can be doubly valued and mask features such as superradiant scattering. Also, the potentials of each of their equations are pathological for some set of frequencies . By using a set of multiple perturbation equations and transformation laws, one can always find a non-pathological tool for any given frequency. But, there is no single rule that works for all frequencies. (These difficulties do not mean that Chandrasekhar and Detweiler’s approaches are not useful. Campanelli and Lousto used rules very similar to Chandrasekhar and Detweiler’s in order to show that solutions to the Teukolsky equation are well-behaved even for sources that extend to infinity.)
Sasaki and Nakamura’s work is restricted to the choice $`s=2`$. This is an appropriate choice for studies of gravitational perturbations, and so the Sasaki-Nakamura equation has been extensively used in studies of gravitational-wave generation and gravitational radiation reaction . Other values of $`s`$ are interesting as well. For example, $`s=\pm 1`$ corresponds to electromagnetic radiation. The propagation and scatter of electromagnetic waves in black hole spacetimes is of great astrophysical interest. Also, much work is currently being directed toward understanding how one calculates self forces and radiation reaction forces in curved spacetimes . Implementations of the general formalism (Refs. ) to date have been restricted to scalar ($`s=0`$) or electromagnetic fields. They have also been restricted to spherically symmetric spacetimes. Tools for effective calculation of radiation fields in Kerr spacetimes will make it possible to extend these calculations to more realistic spinning black holes.
In this paper, I generalize the Sasaki-Nakamura equation to arbitrary integer spin weight $`s`$. The generalized Sasaki-Nakamura (GSN) potentials are given for any $`s`$, but in terms of two unknown functions, $`\alpha (r)`$ and $`\beta (r)`$. These functions are fixed by requiring that the transformation which relates the Teukolsky solution to the GSN solution reduces, in the Schwarzschild limit, to the transformation between the Teukolsky solution and the Regge-Wheeler solution. They also must be chosen so that the potentials they generate are of short range.
There is a great deal of freedom in how one chooses $`\alpha `$ and $`\beta `$: for each value of $`s`$, there are an infinite number of functions which lead to a transformation with the correct limiting value and that produces short-ranged potentials. It is thus most practical to develop $`\alpha `$ and $`\beta `$ on a case by case basis, rather than trying to develop generic formulas. Given the interest in scalar and electromagnetic radiation, I provide examples of $`\alpha `$ and $`\beta `$ for $`s=0`$ and $`s=1`$.
Throughout this paper, a prime denotes $`d/dr`$, where $`r`$ is either the Boyer-Lindquist or the Schwarzschild coordinate (which coordinate should be clear from context). An overbar denotes complex conjugation. The function $`\mathrm{\Delta }=r^22Mr+a^2`$, and $`f=12M/r`$. Thus, for Schwarzschild holes, $`\mathrm{\Delta }=r^2f`$. The function $`R`$ will always denote the solution to the homogeneous Teukolsky equation; $`X`$ will always refer to the solution of the equation with short-ranged potential. Section II reviews important aspects of the Teukolsky equation and its solutions, particularly how one can compute the solution given an appropriate source term and solutions to the homogeneous equation, and how those solutions are related to physical radiation fields. In Sec. III, I review the transformation rules for Schwarzschild black holes. These rules, which are particularly simple, will be used as guidelines for constructing transformations appropriate to radiation in Kerr spacetimes. Finally, in Sec. IV I construct the short-ranged equation for Kerr black holes and provide a recipe for specifying the transformation rule for a radiation field of arbitrary spin weight. I apply this recipe in Secs. V and VI to scalar and electromagnetic radiation fields, respectively. The resultant transformation rules and equations should form a useful basis for further studies of radiation in Kerr spacetimes. Some concluding discussion is given in Sec. VII.
## II Some properties of the Teukolsky equation and its solutions
As background for the calculations in this paper, I review in this section the most important properties of the Teukolsky equation and its solutions.
As discussed in the Introduction, Teukolsky showed that one can separate the wave equation for a field $`{}_{s}{}^{}\mathrm{\Psi }`$ of spin weight $`s`$ radiation propagating on a Kerr black-hole background using the multipolar decomposition given in Eq. (1). The $`\varphi `$ and $`t`$ dependence is trivial, and the $`\theta `$ dependence is straightforwardly dealt with. The $`r`$ dependence, on the other hand, can cause problems.
The radial function $`R(r)`$ is a solution to the Teukolsky equation, Eq. (2). Here I write the Teukolsky equation with its source term and in terms of derivatives with respect to $`r`$ rather than $`r^{}`$:
$$\mathrm{\Delta }^s\left(\mathrm{\Delta }^{s+1}R^{}\right)^{}V_T(r)R=𝒯(r).$$
(10)
This is the way the Teukolsky equation usually appears in the literature. The potential $`V_T(r)`$ is
$$V_T(r)=\lambda 4is\omega r\frac{K(r)^22is(rM)K(r)}{\mathrm{\Delta }};$$
(11)
the quantity $`\lambda =_{lm}2am\omega +a^2\omega ^2s(s+1)`$, where $`_{lm}`$ is the eigenvalue of the spheroidal harmonic \[see ; in the Schwarzschild limit, $`_{lm}=l(l+1)`$\]. The function $`K(r)=(r^2+a^2)\omega ma`$.
The source term $`𝒯(r)`$ depends upon the spin weight of the radiation. It is constructed by projecting the radiation source onto legs of the Newman-Penrose null tetrad, $`𝐥`$, $`𝐧`$, $`𝐦`$, and $`\overline{𝐦}`$. A useful representation of the tetrad in Boyer-Lindquist coordinates is
$`l_\alpha `$ $`=`$ $`[1,{\displaystyle \frac{\mathrm{\Sigma }}{\mathrm{\Delta }}},0,a\mathrm{sin}^2\theta ],`$ (12)
$`n_\alpha `$ $`=`$ $`{\displaystyle \frac{1}{2}}[{\displaystyle \frac{\mathrm{\Delta }}{\mathrm{\Sigma }}},1,0,{\displaystyle \frac{a\mathrm{\Delta }\mathrm{sin}^2\theta }{\mathrm{\Sigma }}}],`$ (13)
$`m_\alpha `$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}(r+ia\mathrm{cos}\theta )}}[ia\mathrm{sin}\theta ,0,\mathrm{\Sigma },i(r^2+a^2)\mathrm{sin}\theta ].`$ (14)
The tetrad legs $`𝐥`$ and $`𝐧`$ represent ingoing and outgoing null vectors, respectively. Quantities constructed by projecting onto $`𝐥`$ correspond to ingoing radiation and their sources; they map to positive $`s`$. Likewise, projection onto $`𝐧`$ corresponds to outgoing radiation and their sources<sup>§</sup><sup>§</sup>§Because one can transform between positive and negative $`s`$ solutions, one can actually develop both ingoing and and outgoing radiation with a single source term., and map to negative $`s`$. See Ref. for details.
Because Eq. (10) is in self-adjoint form, one can construct its solution by the method of Green’s functions . This means that one needs to know only the solutions to the homogeneous equation,
$$\mathrm{\Delta }^s\left(\mathrm{\Delta }^{s+1}R^{}\right)^{}V_T(r)R=0,$$
(15)
in addition to the source. One does this by adapting the generic solution, given in Eqs. (6) and (7), to the appropriate boundary conditions: no radiation may come in from infinity and none may come out from the event horizon. In other words, there exist two solutions, $`R^H(r)`$ and $`R^{\mathrm{}}(r)`$, whose asymptotic forms are
$`R^H(r)`$ $`=`$ $`B^{\mathrm{hole}}\mathrm{\Delta }^se^{ipr^{}},rr_+`$ (16)
$`=`$ $`B^{\mathrm{out}}{\displaystyle \frac{e^{i\omega r^{}}}{r^{2s+1}}}+B^{\mathrm{in}}{\displaystyle \frac{e^{i\omega r^{}}}{r}},r\mathrm{},`$ (17)
$`R^{\mathrm{}}(r)`$ $`=`$ $`D^{\mathrm{out}}e^{ipr^{}}+D^{\mathrm{in}}\mathrm{\Delta }^se^{ipr^{}},rr_+`$ (18)
$`=`$ $`D^{\mathrm{}}{\displaystyle \frac{e^{i\omega r^{}}}{r^{2s+1}}},r\mathrm{}.`$ (19)
The solution to the inhomogeneous equation (10) which one constructs from Eqs. (17), (19), and the source $`𝒯(r)`$ is conveniently written
$$R(r)=Z^H(r)R^{\mathrm{}}(r)+Z^{\mathrm{}}(r)R^H(r),$$
(20)
where
$`Z^H(r)`$ $`=`$ $`{\displaystyle \frac{1}{2i\omega B^{\mathrm{in}}D^{\mathrm{}}}}{\displaystyle _{r_+}^r}𝑑r^{}\mathrm{\Delta }(r^{})^sR^H(r^{})𝒯(r^{}),`$ (21)
$`Z^{\mathrm{}}(r)`$ $`=`$ $`{\displaystyle \frac{1}{2i\omega B^{\mathrm{in}}D^{\mathrm{}}}}{\displaystyle _r^{\mathrm{}}}𝑑r^{}\mathrm{\Delta }(r^{})^sR^{\mathrm{}}(r^{})𝒯(r^{}).`$ (22)
Using Eq. (20) one can construct $`{}_{s}{}^{}\mathrm{\Psi }`$. This quantity is related to a radiation field of spin weight $`s`$; the details of that relation depend upon the value of $`s`$. Typically, $`{}_{s}{}^{}\mathrm{\Psi }`$ is constructed by projecting a tensor describing the radiation onto legs of the Newman-Penrose null tetrad. For example, $`{}_{0}{}^{}\mathrm{\Psi }=\mathrm{\Phi }`$, a massless scalar field. No projections are needed in this case. For $`s=\pm 1`$, we have
$`{}_{1}{}^{}\mathrm{\Psi }`$ $`=`$ $`\varphi _0=F_{\mu \nu }l^\mu m^\nu ,`$ (23)
$`{}_{1}{}^{}\mathrm{\Psi }`$ $`=`$ $`(ria\mathrm{cos}\theta )^2\varphi _2=(ria\mathrm{cos}\theta )^2F_{\mu \nu }n^\mu \overline{m}^\nu ,`$ (24)
where $`F_{\mu \nu }`$ is the electromagnetic field tensor. \[There is a third projection, $`\varphi _1=(1/2)F_{\mu \nu }(l^\mu n^\nu +\overline{m}^\mu m^\nu )`$. It does not describe the radiative degrees of freedom of the electromagnetic field, and so is of less interest here.\] For $`s=\pm 2`$, the radiative quantities are
$`{}_{2}{}^{}\mathrm{\Psi }`$ $`=`$ $`\psi _0=C_{\alpha \beta \gamma \delta }l^\alpha m^\beta l^\gamma m^\delta ,`$ (25)
$`{}_{2}{}^{}\mathrm{\Psi }`$ $`=`$ $`(ria\mathrm{cos}\theta )^4\psi _4=(ria\mathrm{cos}\theta )^4C_{\alpha \beta \gamma \delta }n^\alpha \overline{m}^\beta n^\gamma \overline{m}^\delta .`$ (26)
The tensor $`C_{\alpha \beta \gamma \delta }`$ is the Weyl component of the spacetime’s curvature. The quantities $`\psi _i`$, with $`i`$ an integer from 0 to 4, are the Newman-Penrose projections of the Weyl curvature (see Ref. ). For unperturbed black hole spacetimes, all components except $`\psi _2=C_{\alpha \beta \gamma \delta }l^\alpha m^\beta \overline{m}^\gamma n^\delta =M/(ria\mathrm{cos}\theta )^3`$ can be set to zero with an appropriate choice of gauge. This is the non-radiative “background” component of the curvature; the perturbations $`\psi _0`$ and $`\psi _4`$ represent radiation on the background.
The solution for the (linear) evolution of radiation of spin weight $`s`$ in a Kerr black hole spacetime is thus completely described by construction of the source $`𝒯(r)`$ appropriate to that spin weight and construction of the homogeneous solutions $`R^H(r)`$ and $`R^{\mathrm{}}(r)`$. As discussed in the Introduction — and, as should be clear from the asymptotic solutions (17) and (19) — it is very difficult to build these solutions in a numerical integration. The remainder of this paper is devoted to methods for constructing $`R^H(r)`$ and $`R^{\mathrm{}}(r)`$ by finding transformations that relate the Teukolsky solution $`R(r)`$ to solutions of equations with short-ranged potentials.
## III Results for Schwarzschild holes
The Teukolsky equation for Schwarzschild black holes is
$$\mathrm{\Delta }^s\left(\mathrm{\Delta }^{s+1}R^{}\right)^{}V_{TS}(r)R=0,$$
(27)
where
$$V_{TS}(r)=\lambda 4isr\omega +[2is(rM)\omega (r\omega )^2]/f,$$
(28)
and $`\lambda =\lambda (a=0)=l(l+1)s(s+1)`$.
We would like to find rules that allow us to obtain $`R`$ given a solution $`X`$ of the Regge-Wheeler equation, (8). To do so, first define the quantity
$$\chi \frac{X}{r\sqrt{\mathrm{\Delta }^s}}.$$
(29)
If $`X`$ satisfies the Regge-Wheeler equation, it is straightforward to show that $`\chi `$ satisfies
$$\mathrm{\Delta }^s\left(\mathrm{\Delta }^{s+1}\chi ^{}\right)^{}U_{\chi S}(r)\chi =0$$
(30)
where
$$U_{\chi S}(r)=\lambda +\frac{1}{f}\left[s^2\left(\frac{3M^2}{r^2}\frac{2M}{r}\right)(r\omega )^2\right].$$
(31)
By direct substitution, one can show that $`\chi `$ can be transformed to $`R`$, and vice versa, via
$`s<0:\chi `$ $`=`$ $`\left(r\sqrt{\mathrm{\Delta }}\right)^{|s|}𝒟_{}^{|s|}\left[{\displaystyle \frac{R}{r^{|s|}}}\right],`$ (33)
$`R`$ $`=`$ $`\left({\displaystyle \frac{\mathrm{\Delta }}{r}}\right)^{|s|}𝒟_+^{|s|}\left[\left({\displaystyle \frac{r}{\sqrt{\mathrm{\Delta }}}}\right)^{|s|}\chi \right],`$ (34)
$`s=0:\chi `$ $`=`$ $`R,`$ (35)
$`s>0:\chi `$ $`=`$ $`\left({\displaystyle \frac{r}{\sqrt{\mathrm{\Delta }}}}\right)^s𝒟_+^s\left[\left({\displaystyle \frac{\mathrm{\Delta }}{r}}\right)^sR\right],`$ (36)
$`R`$ $`=`$ $`\left({\displaystyle \frac{1}{r}}\right)^s𝒟_{}^s\left[\left(r\sqrt{\mathrm{\Delta }}\right)^s\chi \right],`$ (37)
where
$$𝒟_\pm =d/dr\pm i\omega /f.$$
(38)
For $`s=1`$ and $`s=2`$, Eq. (34) reduces to the Chandrasekhar transformation (see Ref. for $`s=2`$, Ref. for $`s=1`$). Equations (34) – (37) serve as guidelines that will be used to fix the form of the transformation rules for Kerr black holes. Note that the transformations from $`R`$ to $`\chi `$ can be written
$$\chi =\alpha R+\beta \mathrm{\Delta }^{s+1}R^{}$$
(39)
by repeatedly using Eq. (27) to eliminate derivatives of second order and higher. The resulting functions $`\alpha `$ and $`\beta `$ may become rather complicated, particularly for large values of $`|s|`$, but the general operation is straightforward. (The factor $`\mathrm{\Delta }^{s+1}`$ is inserted for later convenience.)
## IV Perturbation equation for Kerr holes
Guided by Eq. (39), let us assume that functions $`\alpha `$ and $`\beta `$ can be found that transform the Kerr solution $`R`$ to solutions $`\chi `$ of some other equation. By generalizing the relation (29) to a form appropriate for the Kerr metric and rewriting all derivatives in terms of $`r^{}`$ we will come to an equation with short-ranged potentials governing the behavior of a function $`X(r)`$. This is the generalized Sasaki-Nakamura (GSN) equation. It will depend explicitly on the (currently unspecified) functions $`\alpha `$ and $`\beta `$. These functions will be specified by requiring that the transformation rule satisfy a form which reduces to Eqs. (34) – (37) when $`a=0`$. This guarantees that solutions to the GSN equation are equivalent to solutions of the Regge-Wheeler equation in the Schwarzschild limit.
To begin, differentiate $`\chi `$ and use Eq. (15) to eliminate the second derivative of $`R`$. The resulting equations for $`\chi `$ and $`\chi ^{}`$ can be gathered neatly into matrix form:
$$\left(\begin{array}{c}\chi \\ \chi ^{}\end{array}\right)=\left(\begin{array}{cc}\alpha & \beta \mathrm{\Delta }^{s+1}\\ \alpha ^{}+\beta V_T\mathrm{\Delta }^s& \alpha +\beta ^{}\mathrm{\Delta }^{s+1}\end{array}\right)\left(\begin{array}{c}R\\ R^{}\end{array}\right).$$
(40)
A nice feature of Eq. (40) is that the inverse solution is rather obvious:
$$\left(\begin{array}{c}R\\ R^{}\end{array}\right)=\frac{1}{\eta }\left(\begin{array}{cc}\alpha +\beta ^{}\mathrm{\Delta }^{s+1}& \beta \mathrm{\Delta }^{s+1}\\ (\alpha ^{}+\beta V_T\mathrm{\Delta }^s)& \alpha \end{array}\right)\left(\begin{array}{c}\chi \\ \chi ^{}\end{array}\right),$$
(41)
where
$$\eta =\alpha \left(\alpha +\beta ^{}\mathrm{\Delta }^{s+1}\right)\beta \mathrm{\Delta }^{s+1}\left(\alpha ^{}+\beta V_T\mathrm{\Delta }^s\right)$$
(42)
is the determinant of the matrix in Eq. (40).
Differentiating again and massaging the resultant expression gives us a second-order differential equation for $`\chi `$:
$$\mathrm{\Delta }^s\left(\mathrm{\Delta }^{s+1}\chi ^{}\right)^{}\mathrm{\Delta }{}_{s}{}^{}F_{1}^{}(r)\chi ^{}{}_{s}{}^{}U_{1}^{}(r)\chi =0.$$
(43)
The potentials $`{}_{s}{}^{}F_{1}^{}(r)`$ and $`{}_{s}{}^{}U_{1}^{}(r)`$ are given by
$`{}_{s}{}^{}F_{1}^{}(r)`$ $`=`$ $`\eta ^{}/\eta ,`$ (44)
$`{}_{s}{}^{}U_{1}^{}(r)`$ $`=`$ $`V_T+{\displaystyle \frac{1}{\beta \mathrm{\Delta }^2}}\left[\left(2\alpha +\beta ^{}\mathrm{\Delta }^{s+1}\right)^{}{\displaystyle \frac{\eta ^{}}{\eta }}\left(\alpha +\beta ^{}\mathrm{\Delta }^{s+1}\right)\right].`$ (45)
Next, generalize Eq. (29) to the Kerr form
$$\chi \frac{X}{\sqrt{(r^2+a^2)\mathrm{\Delta }^s}}.$$
(46)
Using this to replace $`\chi `$ for $`X`$ in Eq. (43) and then replacing derivatives in $`r`$ with derivatives in $`r^{}`$ with the rule
$$\frac{d}{dr}=\frac{(r^2+a^2)}{\mathrm{\Delta }}\frac{d}{dr^{}}$$
(47)
yields the GSN equation:
$$\frac{d^2X}{dr^2}{}_{s}{}^{}F(r)\frac{dX}{dr^{}}{}_{s}{}^{}U(r)X=0.$$
(48)
The potentials are
$`{}_{s}{}^{}F(r)`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Delta }{}_{s}{}^{}F_{1}^{}(r)}{r^2+a^2}},`$ (49)
$`{}_{s}{}^{}U(r)`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Delta }{}_{s}{}^{}U_{1}^{}(r)}{(r^2+a^2)^2}}+{}_{s}{}^{}G(r)^2+{\displaystyle \frac{\mathrm{\Delta }d{}_{s}{}^{}G/dr}{r^2+a^2}}{\displaystyle \frac{\mathrm{\Delta }{}_{s}{}^{}G(r){}_{s}{}^{}F_{1}^{}(r)}{r^2+a^2}}.`$ (50)
The function $`{}_{s}{}^{}G(r)`$ is
$${}_{s}{}^{}G(r)=\frac{r\mathrm{\Delta }}{(r^2+a^2)^2}+\frac{s(rM)}{r^2+a^2};$$
(51)
the functions $`{}_{s}{}^{}F_{1}^{}(r)`$ and $`{}_{s}{}^{}U_{1}^{}(r)`$ are from Eq. (45). When $`s=2`$, all functions reduce to those given by Sasaki and Nakamura (see Ref. ).
All of the quantities which have been derived to this point depend upon the as-yet-undetermined functions $`\alpha `$ and $`\beta `$. We fix these functions by requiring that they affect a transformation between $`R`$ and $`\chi `$ which, as $`a0`$, reduces to Eqs. (34) – (37). A useful generalization of these transformations is
$`s<0:\chi `$ $`=`$ $`\left(\sqrt{(r^2+a^2)\mathrm{\Delta }}\right)^{|s|}g_0(r)J_{}\left[g_1(r)J_{}\left[g_2(r)\mathrm{}J_{}\left[{\displaystyle \frac{g_{|s|}(r)R}{\left(\sqrt{r^2+a^2}\right)^{|s|}}}\right]\right]\right],`$ (53)
$`s=0:\chi `$ $`=`$ $`g_0(r)R,`$ (54)
$`s>0:\chi `$ $`=`$ $`\left(\sqrt{{\displaystyle \frac{r^2+a^2}{\mathrm{\Delta }}}}\right)^sg_0(r)J_+\left[g_1(r)J_+\left[g_2(r)\mathrm{}J_+\left[g_s(r)\left({\displaystyle \frac{\mathrm{\Delta }}{\sqrt{r^2+a^2}}}\right)^sR\right]\right]\right],`$ (55)
where the operator
$$J_\pm =d/dr\pm iK(r)/\mathrm{\Delta }$$
(56)
generalizes $`𝒟_\pm `$ to Kerr. The $`s=+2`$ transformation rule, for example, is
$$\chi =g_0(r)\frac{(r^2+a^2)}{\mathrm{\Delta }}J_+\left[g_1(r)J_+\left[g_2(r)\frac{\mathrm{\Delta }^2}{r^2+a^2}R\right]\right];$$
(57)
an example for $`s=1`$ is given in Sec. VI.
To now specify $`\alpha `$ and $`\beta `$, one must pick functions $`g_i(r)`$ and then repeatedly use Eq. (15) to eliminate derivatives of second order and higher in Eqs. (53) – (55). The resultant expressions for $`\alpha `$ and $`\beta `$ will be, in general, quite complicated; examples are discussed in Secs. V and VI. The functions $`g_i(r)`$ must be chosen so that they become constant in the Schwarzschild limit, and lead to potentials $`F(r)`$ and $`U(r)`$ which are short-ranged \[i.e., fall off at a rate $`O(1/r^2)`$ or faster as $`r\mathrm{}`$\]. In practice, choosing $`g_i(r)=1`$ or $`g_i(r)=(r^2+a^2)/r^2`$ appears to lead to well-behaved potentials; some experimentation may be needed to make useful choices.
## V Scalar radiation
For scalar radiation, $`s=0`$, the functions $`R`$, $`X`$, and $`\chi `$ have the following relationship:
$$g_0(r)R=\chi =\frac{X}{\sqrt{r^2+a^2}}.$$
(58)
A good choice is $`g_0(r)=1`$. The functions $`\alpha `$ and $`\beta `$ \[cf. Eq. (39)\] are then given by
$$\alpha =1,\beta =0.$$
(59)
From this, it follows that
$`\eta `$ $`=`$ $`1,`$ (61)
$`{}_{0}{}^{}F_{1}^{}`$ $`=`$ $`0,`$ (62)
$`{}_{0}{}^{}U_{1}^{}`$ $`=`$ $`V_T,`$ (63)
$`{}_{0}{}^{}G`$ $`=`$ $`{\displaystyle \frac{r\mathrm{\Delta }}{r^2+a^2}}.`$ (64)
The potentials $`{}_{0}{}^{}F(r)`$ and $`{}_{0}{}^{}U(r)`$ are given by substituting Eqs. (61) – (64) into Eq. (50). For large $`r`$,
$${}_{0}{}^{}U(r)=\omega ^2+\frac{\lambda +2am\omega }{r^2}+\frac{2M\left(1\lambda \right)}{r^3}+O(1/r^4);$$
(65)
clearly, $`{}_{0}{}^{}F(r)=0`$ for all $`r`$. Hence, the potentials are short-ranged. When $`a=0`$, $`{}_{0}{}^{}U(r)`$ reduces to $`\omega ^2+V_{\mathrm{RW}}(r,s=0)`$. Thus, it reduces to the Regge-Wheeler equation in the Schwarzschild limit, as it was supposed to.
The asymptotic solutions to the $`s=0`$ GSN equation are simple plane waves:
$`X^H(r)`$ $`=`$ $`e^{ipr^{}},rr_+,`$ (66)
$`=`$ $`A^{\mathrm{out}}\overline{P}_0(r)e^{i\omega r^{}}+A^{\mathrm{in}}P_0(r)e^{i\omega r^{}},r\mathrm{};`$ (67)
$`X^{\mathrm{}}(r)`$ $`=`$ $`C^{\mathrm{out}}e^{ipr^{}}+C^{\mathrm{in}}e^{ipr^{}},rr_+,`$ (68)
$`=`$ $`\overline{P}_0(r)e^{i\omega r^{}},r\mathrm{}.`$ (69)
The function
$$P_0(r)=1+\frac{𝒜_0}{\omega r}+\frac{_0}{(\omega r)^2}+\frac{𝒞_0}{(\omega r)^3}+\mathrm{}$$
(70)
allows us to more accurately describe the behavior of $`X^{H,\mathrm{}}`$ near infinity. This is useful both to improve numerical computations and to derive certain relations between the amplitudes of the Teukolsky solution and the GSN solution. The first three coefficients have the values
$`𝒜_0`$ $`=`$ $`{\displaystyle \frac{i}{2}}\left(\lambda +2am\omega \right),`$ (71)
$`_0`$ $`=`$ $`{\displaystyle \frac{1}{8}}\left[\lambda ^2\lambda \left(24am\omega \right)4\left[am\omega iM\omega am\omega \left(am\omega +2iM\omega \right)\right]\right],`$ (72)
$`𝒞_0`$ $`=`$ $`{\displaystyle \frac{i}{6}}[_0(\lambda 6+2am\omega +8iM\omega )4\left(M\omega \right)^22𝒜_0M\omega (\lambda 6)`$ (74)
$`\left(a\omega \right)^2(\lambda 1+m^2+2am\omega )].`$
## VI Electromagnetic radiation
For electromagnetic radiation, $`s=1`$, the functions $`R`$, $`X`$, and $`\chi `$ exhibit the following relationships:
$`\chi `$ $`=`$ $`\sqrt{{\displaystyle \frac{\mathrm{\Delta }}{r^2+a^2}}}X,`$ (76)
$`\chi `$ $`=`$ $`\alpha R+\beta R^{}`$ (77)
$`=`$ $`g_0(r)\sqrt{(r^2+a^2)\mathrm{\Delta }}J_{}\left[{\displaystyle \frac{g_1(r)R}{\sqrt{r^2+a^2}}}\right].`$ (78)
From this, we can read off
$`\alpha `$ $`=`$ $`g_0\sqrt{\mathrm{\Delta }}\left[g_1^{}{\displaystyle \frac{rg_1}{r^2+a^2}}{\displaystyle \frac{ig_1K}{\mathrm{\Delta }}}\right],`$ (79)
$`\beta `$ $`=`$ $`g_0g_1\sqrt{\mathrm{\Delta }}.`$ (80)
A useful choice for $`g_0`$ and $`g_1`$ is
$$g_0(r)=\frac{(r^2+a^2)}{r^2},g_1(r)=1.$$
(81)
The function $`\eta `$ that follows from these choices is
$$\eta =c_0+c_1/r+c_2/r^2+c_3/r^3+c_4/r^4,$$
(82)
where
$`c_0`$ $`=`$ $`\lambda ,`$ (83)
$`c_1`$ $`=`$ $`2iam,`$ (84)
$`c_2`$ $`=`$ $`a^2(12\lambda ),`$ (85)
$`c_3`$ $`=`$ $`2a^2(M+iam),`$ (86)
$`c_4`$ $`=`$ $`a^4(1\lambda ).`$ (87)
Using this $`\eta `$ and
$${}_{1}{}^{}G(r)=\frac{r\mathrm{\Delta }}{(r^2+a^2)^2}\frac{(rM)}{r^2+a^2},$$
(88)
it is straightforward to construct the functions $`{}_{1}{}^{}F_{1}^{}(r)`$, $`{}_{1}{}^{}U_{1}^{}(r)`$, $`{}_{1}{}^{}F(r)`$, and $`{}_{1}{}^{}U(r)`$. The results are rather complicated and are not given here. When $`r`$ is large,
$`{}_{1}{}^{}F(r)`$ $`=`$ $`{\displaystyle \frac{2iam}{\lambda r^2}}+{\displaystyle \frac{2a\left[2imM\lambda a\left(2m^2+2\lambda ^2\lambda \right)\right]}{\lambda ^2r^3}}+O(1/r^3),`$ (89)
$`{}_{1}{}^{}U(r)`$ $`=`$ $`\omega ^2+{\displaystyle \frac{\lambda ^2+2am\omega (\lambda +1)}{\lambda r^2}}{\displaystyle \frac{2\left[ia^2(2m^2\lambda )+M(\lambda ^3+2am\omega \lambda )\right]}{\lambda ^2r^3}}+O(1/r^4).`$ (90)
Both $`{}_{1}{}^{}F(r)`$ and $`{}_{1}{}^{}U(r)`$ are short ranged. When $`a=0`$, $`{}_{1}{}^{}F(r)=0`$ and $`{}_{1}{}^{}U(r)=\omega ^2+V_{RW}(r,s=1)`$.
The solutions $`X^{H,\mathrm{}}`$ are, in the limits $`rr_+`$ and $`r\mathrm{}`$, essentially identical to those given in Eqs. (67) and (69); one need only change the subscript on the $`P`$ function to $`1`$. The corresponding coefficients in $`P_1(r)`$ are
$`𝒜_1`$ $`=`$ $`{\displaystyle \frac{i}{2}}\left(\lambda +2am\omega \right),`$ (92)
$`_1`$ $`=`$ $`{\displaystyle \frac{1}{8}}\left[\lambda ^2\lambda \left(24am\omega \right)4a\omega \left[m2imM\omega +a\omega \left(2m^2\right)\right]\right],`$ (93)
$`𝒞_1`$ $`=`$ $`{\displaystyle \frac{i}{6}}[_1(\lambda 6+2am\omega +8iM\omega )+2𝒜_1[M\omega (5\lambda )+ia\omega (2a\omega +m/\lambda )]`$ (95)
$`(a\omega )^2[\lambda 5+8iM\omega +2am\omega +m^2(\lambda +2)/\lambda ]].`$
## VII Conclusion
In this paper, I have shown how the Sasaki-Nakamura short-ranged equation for gravitational perturbations to a rotating black hole may be generalized to arbitrary spin weight radiation. Of course, there is no particular physical motivation for choosing radiation spin of magnitude greater than $`2`$; this approach is taken simply so that one can write down a single rule which encompasses all physically interesting radiation fields, much as the Teukolsky equation itself encompasses all spin weights. Efficient numerical computation of Teukolsky equation solutions now reduces to a simple recipe. First, following the analysis in Sec. IV, develop the potentials needed in the GSN equation, Eq. (48). Examples are given for $`s=0`$ and $`s=1`$. Integrate Eq. (48) for the GSN solution $`X`$. Transform to the variable $`\chi `$ using Eq. (46). Then construct the Teukolsky solution $`R`$ using Eq. (40).
One application of these results may be to extend the mode sum regularization scheme described in Ref. to self forces computed in Kerr spacetimes. Calculations that employ scalar or electromagnetic charges and fields are generally simpler than the gravitational self force calculations, which are of great interest for researchers studying gravitational-wave sources. The electromagnetic perturbation equation given in Sec. VI may be of astrophysical interest, particularly when coupled to an appropriate source.
###### Acknowledgements.
I thank Lior Burko for encouraging me to find transformation rules and perturbation equations for generic $`s`$, Eric Poisson for valuable comments, and Manuela Campanelli for pointing me to some useful references. I am also very grateful to Steven Detweiler for valuable comments on the history of black hole perturbation studies which I have used to correct some of the discussion in the Introduction. The package Mathematica was used to aid some of the calculations. This research was supported by NSF Grant AST-9731698 and NASA Grants NAG5-7034 and NAGW-4268. |
warning/0002/nucl-th0002060.html | ar5iv | text | # Pairing correlation involving the continuum states
## Introduction
In neutron-rich nuclei, the pairing correlation significantly involves the continuum single-particle states. This makes the HF+BCS approximation inadequate due to the unlocalization of the neutron density distribution and demands one to solve the Hartree-Fock-Bogoliubov (HFB) equation without approximations. The solution in the coordinate space was first formulated in Ref. $`^{\text{1})}`$ using the quasi-particle states and performed for spherically symmetric states. However, its application to deformed states is difficult because there are quite a large number of quasiparticle states even for a moderate size of the normalization box (i.e., the cavity to confine the nucleons to discretize the single-particle states). Every HFB solution has an equivalent expression of BCS variational form. The single-particle states to construct the BCS type wavefunctions are called the HFB canonical basis or sometimes the natural orbitals. This expression was used to solve the HFB for spherical states originally in Ref. $`^{\text{2})}`$.
Although spherical solutions can be obtained easily with present computers (for zero-range forces), deformed solutions are still difficult to obtain. The two-basis method$`^{\text{3},\text{4},\text{5})}`$ is the only one implemented so far for neutron-rich nuclei. Some recent developments like Ref.$`^{\text{6})}`$ are also in progress.
We have applied the canonical-basis method to deformed states in a three dimensional cubic mesh representation with density dependent delta interactions.$`^{\text{7})}`$ It has turned out to be a very efficient alternative approach to obtain the solutions. The origin of its effectiveness is that the number of necessary single-particle basis states to describe the ground state of a nucleus is proportional to the volume of the nucleus in the canonical-basis method while it commensurates with the volume of the normalization box in conventional methods. The difference of the number of the basis states amounts to a factor of $`10^1`$ \- $`10^3`$.
In this paper we discuss the canonical-basis formulation of the HFB, the method to obtain the canonical-basis solutions, faster gradient-method paths than a naive imaginary-time evolution, the necessity of the cut-off of the pairing interaction, and an implementation of the cut-off in terms of an interaction dependent on the pairing density.
## HFB in the canonical representation
To begin with, let us formulate the HF and the HFB in the coordinate-space representation in order to elucidate a difficulty of the HFB and to suggest its possible solution in terms of the canonical-basis representation. For the sake of simplicity, we consider only one kind of nucleons and designate the number of nucleons by $`N`$ in Eqs. (1)-(14), which are in this section. The $`z`$-component of the spin of a nucleon is represented by $`s`$ (= $`\pm \frac{1}{2}`$).
In the HF, one should minimize $`\mathrm{\Psi }|H|\mathrm{\Psi }`$ for single Slater-determinant states,
$`|\mathrm{\Psi }`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}a_i^{}|0,`$ (1)
$`a_i^{}`$ $`=`$ $`{\displaystyle \underset{s}{}}{\displaystyle 𝑑\stackrel{}{r}\psi _i(\stackrel{}{r},s)a^{}(\stackrel{}{r},s)},`$ (2)
by varying $`\{\psi _i\}_{i=1,\mathrm{},N}`$ under orthonormality conditions $`\psi _i|\psi _j`$ = $`\delta _{ij}`$. The operator $`a_i^{}`$ creates a nucleon with a wavefunction $`\psi _i(\stackrel{}{r},s)`$. The distribution function of the density of the nucleons is related to the wavefunctions as
$$\rho (\stackrel{}{r})=\underset{s}{}\mathrm{\Psi }|a(\stackrel{}{r},s)a^{}(\stackrel{}{r},s)|\mathrm{\Psi }=\underset{i=1}{\overset{N}{}}\underset{s}{}|\psi _i(\stackrel{}{r},s)|^2.$$
(3)
In the HFB, the solution takes the following form,
$`|\mathrm{\Psi }`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{I}{}}}b_i|0,`$ (4)
$`b_i`$ $`=`$ $`{\displaystyle \underset{s}{}}{\displaystyle }d\stackrel{}{r}\{\varphi _i(\stackrel{}{r},s)a(\stackrel{}{r},s)`$ (5)
$`+\phi _i(\stackrel{}{r},s)a^{}(\stackrel{}{r},s)\},`$
where $`b_i`$ is the annihilation operator of a negative-energy Bogoliubov quasi-particle with amplitudes $`\varphi _i(\stackrel{}{r},s)`$ for presence and $`\phi _i(\stackrel{}{r},s)`$ for absence. $`I`$ is the number of the basis states of the employed representation. For a three-dimensional Cartesian mesh (3D-mesh) representation,$`^{\text{8})}`$ it is the number of the mesh points (times four when spin-orbit interactions are included) and typically $`10^4`$-$`10^5`$. One should vary $`\{\varphi _i,\phi _i\}_{i=1,\mathrm{},I}`$ under orthonormality conditions
$`{\displaystyle \underset{s}{}}{\displaystyle 𝑑\stackrel{}{r}\left\{\varphi _i^{}(\stackrel{}{r},s)\varphi _j(\stackrel{}{r},s)+\phi _i^{}(\stackrel{}{r},s)\phi _j(\stackrel{}{r},s)\right\}}`$
$`=`$ $`\delta _{ij}(1ijI),`$ (6)
and a constraint on the expectation value of the number of nucleons,
$`\mathrm{\Psi }|\widehat{N}|\mathrm{\Psi }`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{I}{}}}{\displaystyle \underset{s}{}}{\displaystyle 𝑑\stackrel{}{r}|\phi _i(\stackrel{}{r},s)|^2}=N,`$ (7)
$`\widehat{N}`$ $`=`$ $`{\displaystyle \underset{s}{}}{\displaystyle 𝑑\stackrel{}{r}a(\stackrel{}{r},s)a^{}(\stackrel{}{r},s)}.`$ (8)
The nucleon density for state (4) is expressed as
$$\rho (\stackrel{}{r})=\underset{i=1}{\overset{I}{}}\underset{s}{}|\phi _i(\stackrel{}{r},s)|^2.$$
(9)
The essential difference between the HF and the HFB is that one has to consider only $`N`$ $`10^2`$ wavefunctions in the former while one has to treat explicitly as many single-particle wavefunctions as the number of the basis in the latter.
Owing to the Bloch-Messiah theorem,$`^{\text{9})}`$ the state (4) can be expressed (for the ground states of even-even nuclei) as,
$`|\mathrm{\Psi }`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{K}{}}}\left(u_i+v_ia_i^{}a_{\overline{ı}}^{}\right)|0,`$ (10)
$`a_i^{}`$ $`=`$ $`{\displaystyle \underset{s}{}}{\displaystyle 𝑑\stackrel{}{r}\psi _i(\stackrel{}{r},s)a^{}(\stackrel{}{r},s)},`$ (11)
where $`a_i^{}`$ and $`a_{\overline{ı}}^{}`$ create a nucleon with wavefunction $`\psi _i(\stackrel{}{r},s)`$ and $`\psi _{\overline{ı}}(\stackrel{}{r},s)`$, respectively, which are called as the canonical basis$`^{\text{9})}`$ or the natural orbitals$`^{\text{2})}`$ of the HFB vacuum $`|\mathrm{\Psi }`$. One must use $`K=\frac{1}{2}I`$ for the exact equivalence between Eqs. (4) and (10) in general cases. When $`|\mathrm{\Psi }`$ is a time-reversal invariant state, which we assume in this paper, $`\psi _i`$ and $`\psi _{\overline{ı}}`$ are the time-reversal state of each other. In this case, only one wavefunction of each time reversal pair should be counted as independent variables of the variational procedure.
To obtain solutions in the canonical-basis framework, one should vary $`\{\psi _i,u_i,v_i\}_{i=1,\mathrm{},K}`$ under three kinds of constraints, i.e, the orthonormality conditions,
$`\psi _i|\psi _j`$ $`=`$ $`{\displaystyle \underset{s}{}}{\displaystyle 𝑑\stackrel{}{r}\psi _i^{}(\stackrel{}{r},s)\psi _j(\stackrel{}{r},s)}`$
$`=`$ $`\delta _{ij}(1ijK),`$ (12)
fixed expectation value of the number of nucleons,
$$\mathrm{\Psi }|\widehat{N}|\mathrm{\Psi }=2\underset{i=1}{\overset{K}{}}v_i^2=N,$$
(13)
and the normalization of the $`u`$-$`v`$ factors $`u_i^2+v_i^2=1`$.
The nucleon density is expressed as
$$\rho (\stackrel{}{r})=2\underset{i=1}{\overset{K}{}}\underset{s}{}v_i^2\left|\psi _i(\stackrel{}{r},s)\right|^2.$$
(14)
Reinhard et al. regarded that the advantage of the representation (10) over (4) is that one has to consider only a single set of wavefunctions $`\{\psi _i\}_{i=1,\mathrm{},I}`$ unlike a double set $`\{\varphi _i,\phi _i\}_{i=\pm 1,\mathrm{},\pm I/2}`$.$`^{\text{2})}`$ However, we expect much greater benefit from the canonical-basis representation. Namely, $`i`$ may be truncated as $`iK`$ = $`𝒪(N)`$ $``$ $`\frac{1}{2}I`$ to a very good approximation. It is because $`\psi _i`$ appearing on the right-hand side of Eq. (14) must be a localized function as $`\rho (\stackrel{}{r})`$ on the left-hand side, while the orthogonality (12) does not allow many low-energy wavefunctions to exist in the vicinity of the nucleus. For 3D-mesh representations, $`I`$ is proportional to the volume of the cavity while $`K`$ is proportional to the volume of the nucleus. The latter is $`10^1`$-$`10^3`$ times as small as the former.
Incidentally, the situation is quite different in the quasiparticle formalism. On the one hand, the localization of the density demands only the localization of $`\phi _i`$ through Eq. (9) while $`\varphi _i`$ does not have to be localized in general. On the other hand, the orthogonality condition (6) involves both $`\phi _i`$ and $`\varphi _i`$. This discrepancy enables many quasiparticle states having similar $`\phi _i`$ to be orthogonal to each other by differing their $`\varphi _i`$.
## Mean fields for zero-range interactions
Let us present the effective Hamiltonian employed in this paper. We adopt a density-dependent zero-range interaction. Zero-rangeness makes the mean-field potentials local, which is an essential advantage for coordinate-space solutions. On the other hand, the omission of momentum dependences are merely for the sake of simplicity and there will not be essential differences in the formulation if we use the full-form Skyrme force.$`^{\text{9})}`$ Our force is expressed using the parameterization of the Skyrme force as,
$`\widehat{v}(\stackrel{}{r}_1,s_1;\stackrel{}{r}_2,s_2)`$
$`=`$ $`\left(t_0+{\displaystyle \frac{1}{6}}t_3\rho \left(\stackrel{}{r}_1\right)^\alpha \right)\delta \left(\stackrel{}{r}_1\stackrel{}{r}_2\right),`$ (15)
where $`\stackrel{}{r}_i`$ and $`s_i`$ are the position vector and the spin state of the two interacting nucleons, $`i=1,2`$. Dependence on the isospin state is redundant because of the zero-rangeness. We adopt $`t_0`$ = $`983.4`$ MeV fm<sup>3</sup>, $`t_3`$ = $`13106`$ MeV fm<sup>3+3α</sup>, and $`\alpha `$ = 0.98 when the force is used to construct the mean-field (HF) potential.$`^{\text{10})}`$ We use different strengths to make the pairing potential. We express the pairing force in the parameterization of Ref.$`^{\text{11})}`$ as
$`\widehat{v}_\mathrm{p}(\stackrel{}{r}_1,s_1;\stackrel{}{r}_2,s_2)`$ (16)
$`=`$ $`v_\mathrm{p}{\displaystyle \frac{1P_\sigma }{2}}\left(1{\displaystyle \frac{\rho (\stackrel{}{r}_1)}{\rho _\mathrm{c}}}\right)\delta \left(\stackrel{}{r}_1\stackrel{}{r}_2\right).`$
$`P_\sigma `$ is the exchange operator of spin variables. $`\frac{1}{2}(1P_\sigma )`$ is a projector into spin-singlet states so that the interaction acts only between like nucleons. We use $`\rho _\mathrm{c}`$ = 0.32 fm<sup>-3</sup> (to roughly vanish the volume-changing effect$`^{\text{11})}`$) and $`v_\mathrm{p}=440`$ MeV fm<sup>3</sup>.
When one considers both protons and neutrons, the state of the nucleus is expressed as a product of two BCS type states (10) for the protons and the neutrons:
$$|\mathrm{\Psi }=\underset{\mathrm{q}=\mathrm{p}}{\overset{\mathrm{n}}{}}\underset{i=1}{\overset{K}{}}\left(u_i+v_ia_{i,\mathrm{q}}^{}a_{\overline{ı},\mathrm{q}}^{}\right)|0,$$
(17)
where q distinguishes between protons (p) and neutrons (n). $`a_{i,\mathrm{p}}^{}`$ creates a proton having a wavefunction $`\psi _{i,\mathrm{p}}(\stackrel{}{r},s)`$ while $`a_{i,\mathrm{n}}^{}`$ creates a neutron with a wavefunction $`\psi _{i,\mathrm{n}}(\stackrel{}{r},s)`$. The product form is due to the pairing interaction (16) acting only between like nucleons.
For the sake of simplicity, we treat $`N`$=$`Z`$ nuclei without Coulomb interaction in this paper. In this case, the wavefunctions are the same between protons and neutrons, i.e., $`\psi _{i,\mathrm{p}}(\stackrel{}{r},s)`$ = $`\psi _{i,\mathrm{n}}(\stackrel{}{r},s)`$, $`\psi _{\overline{ı},\mathrm{p}}(\stackrel{}{r},s)`$ = $`\psi _{\overline{ı},\mathrm{n}}(\stackrel{}{r},s)`$. Moreover, because the potentials are independent of the spin, the wavefunctions $`\psi _i(\stackrel{}{r},s)`$ can be factorized into a product of a spin wavefunction and a real function of the position, which we write $`\psi _i(\stackrel{}{r})`$ in the following. It holds $`\psi _{i,\mathrm{p}}(\stackrel{}{r})`$ = $`\psi _{\overline{ı},\mathrm{p}}(\stackrel{}{r})`$ = $`\psi _{i,\mathrm{n}}(\stackrel{}{r})`$ = $`\psi _{\overline{ı},\mathrm{n}}(\stackrel{}{r})`$.
With the interactions (15) and (16), the total energy for state (17) is written as,
$`E`$ $`=`$ $`\mathrm{\Psi }|H|\mathrm{\Psi }={\displaystyle \left(\stackrel{}{r}\right)𝑑\stackrel{}{r}},`$ (18)
$`(\stackrel{}{r})`$ $`=`$ $`{\displaystyle \frac{\mathrm{}^2}{2m}}\tau (\stackrel{}{r})+{\displaystyle \frac{3}{8}}t_0\rho (\stackrel{}{r})^2+{\displaystyle \frac{1}{16}}t_3\rho (\stackrel{}{r})^{2+\alpha }`$ (19)
$`+`$ $`{\displaystyle \frac{1}{8}}v_\mathrm{p}\left(1{\displaystyle \frac{\rho (\stackrel{}{r})}{\rho _\mathrm{c}}}\right)\stackrel{~}{\rho }(\stackrel{}{r})^2,`$
where $`m`$ is the average of the proton and the neutron masses divided by $`11/A`$ for the correction of the center-of-mass motion. $`(\stackrel{}{r})`$ is called the Hamiltonian density while function of position $`\stackrel{}{r}`$ in the right-hand side are
$`\tau (\stackrel{}{r})`$ $`=g{\displaystyle \underset{i=1}{\overset{K}{}}}v_i^2|\stackrel{}{}\psi _i(\stackrel{}{r})|^2`$ : kinetic energy density,
$`\rho (\stackrel{}{r})`$ $`=g{\displaystyle \underset{i=1}{\overset{K}{}}}v_i^2|\psi _i(\stackrel{}{r})|^2`$ : density,
$`\stackrel{~}{\rho }(\stackrel{}{r})`$ $`=g{\displaystyle \underset{i=1}{\overset{K}{}}}u_iv_i|\psi _i(\stackrel{}{r})|^2`$ : pairing density,
where $`g=4`$, which is a factor to account for the situation that a wavefunction $`\psi _i`$ takes care of four nucleons for the spin and isospin degeneracy.
The mean-field potential $`V`$ and the pairing potential $`\stackrel{~}{V}`$ are defined as
$`V`$ $`={\displaystyle \frac{}{\rho }}`$ $`={\displaystyle \frac{3}{4}}t_0\rho +{\displaystyle \frac{2+\alpha }{16}}t_3\rho ^{1+\alpha }{\displaystyle \frac{v_\mathrm{p}}{8\rho _\mathrm{c}}}\stackrel{~}{\rho }^2,`$ (20)
$`\stackrel{~}{V}`$ $`={\displaystyle \frac{}{\stackrel{~}{\rho }}}`$ $`={\displaystyle \frac{1}{4}}v_\mathrm{p}\left(1{\displaystyle \frac{\rho }{\rho _\mathrm{c}}}\right)\stackrel{~}{\rho },`$ (21)
in which $`V`$, $`\stackrel{~}{V}`$, $``$, $`\rho `$, and $`\stackrel{~}{\rho }`$ are local functions of $`\stackrel{}{r}`$ while $`t_0`$, $`t_3`$, $`\alpha `$, $`v_\mathrm{p}`$, and $`\rho _\mathrm{c}`$ are constants.
The mean-field and the pairing Hamiltonians are
$`h`$ $`=`$ $`{\displaystyle \frac{\mathrm{}^2}{2m}}\stackrel{}{}^2+V,`$ (22)
$`\stackrel{~}{h}`$ $`=`$ $`\stackrel{~}{V}.`$ (23)
The quasiparticle states are the eigenvectors of the so-called HFB super matrix composed of $`h`$ and $`\stackrel{~}{h}`$:
$$\left(\begin{array}{cc}h& \stackrel{~}{h}\\ \stackrel{~}{h}& h\end{array}\right)\left(\begin{array}{c}\varphi _\alpha \\ \psi _\alpha \end{array}\right)=ϵ_\alpha \left(\begin{array}{c}\varphi _\alpha \\ \psi _\alpha \end{array}\right).$$
(24)
This is just an eigenvalue problem of a hermitian (because of the time-reversal symmetry) matrix. The canonical orbitals are also determined by $`h`$ and $`\stackrel{~}{h}`$ but in a more complex way as described in the next section.
## Gradient method for canonical-basis HFB
In this section we describe a procedure to obtain the canonical-basis solution of the HFB equation directory, not by way of quasi-particle states. Instead of minimizing $`E`$ = $`\mathrm{\Psi }|H|\mathrm{\Psi }`$ with $`|\mathrm{\Psi }`$ given by Eq. (17) under constraints of Eqs. (12) and (13), one may introduce a Routhian $`R`$,
$`R`$ $`=`$ $`Eϵ_\mathrm{F}g{\displaystyle \underset{i=1}{\overset{K}{}}}v_i^2`$
$``$ $`g{\displaystyle \underset{i=1}{\overset{K}{}}}{\displaystyle \underset{j=1}{\overset{K}{}}}\lambda _{ij}\left\{\psi _i|\psi _j\delta _{ij}\right\},`$ (25)
and minimize it without constraints. $`ϵ_\mathrm{F}`$ is probably the most familiar Lagrange multiplier, whose physical meaning is the Fermi level. In the definition (25), $`K^2`$ Lagrange multipliers $`\lambda _{ij}`$ obeying hermiticity,
$$\lambda _{ij}=\lambda _{ji}^{},$$
(26)
are introduced instead of $`\frac{1}{2}K(K+1)`$ independent multipliers. This hermitization of $`\lambda `$ is adopted in order to make $`R`$ real so that two conditions, $`\delta R/\delta \psi _i=0`$ and $`\delta R/\delta \psi _i^{}=0`$, become equivalent and thus one has to consider only one of them. Note that $`\delta _{ij}`$ is subtracted from $`\psi _i|\psi _j`$, in contrast to Ref. $`^{\text{2})}`$, which is in order to treat $`\lambda _{ij}`$ not as constants like $`ϵ_\mathrm{F}`$ but as functionals of the wavefunctions.
Using notations,
$$ϵ_i=\psi _i|h|\psi _i,\mathrm{\Delta }_i=\psi _i|\stackrel{~}{h}|\psi _i,$$
(27)
the stationary conditions of $`R`$ result in two kinds of equations. One is $`R/v_i=0`$, which concerns the occupation amplitudes $`v_i`$ and is fulfilled by
$$v_i^2=\frac{1}{2}\pm \frac{1}{2}\frac{ϵ_iϵ_\mathrm{F}}{\sqrt{(ϵ_iϵ_\mathrm{F})^2+\mathrm{\Delta }_i^2}}.$$
(28)
The minimum of $`R`$ for the variations of $`u_i`$ and $`v_i`$ is obtained (when $`\mathrm{\Delta }_i0`$) by taking the minus sign for the double sign in the right-hand side of Eq. (28) and choosing the same sign as $`v_i`$ for $`u_i`$ = $`\pm \sqrt{1v_i^2}`$. The stationary condition for a wavefunction $`\psi _i(\stackrel{}{r},s)`$ gives the following equation:
$`{\displaystyle \frac{1}{g}}{\displaystyle \frac{\delta R}{\delta \psi _i^{}}}`$ $`=_i\psi _i{\displaystyle \underset{j=1}{\overset{K}{}}}\lambda _{ij}\psi _j`$
$``$ $`{\displaystyle \underset{j=1}{\overset{K}{}}}{\displaystyle \underset{k=1}{\overset{K}{}}}{\displaystyle \frac{\delta \lambda _{jk}}{\delta \psi _i^{}}}\left\{\psi _j|\psi _k\delta _{jk}\right\}=0,`$ (29)
where
$$_i=v_i^2h+u_iv_i\stackrel{~}{h}.$$
(30)
One can regard $`_i`$ as a state-dependent single-particle Hamiltonian. This dependence on states makes the orthogonalization conditions essential to the method. For HF, the orthogonalization conditions are easily fulfilled because $`\psi _i`$ are eigenstates of the same hermite operator $`h`$ and thus are orthogonal between themselves at the solution : $`\psi _j|\psi _i(ϵ_jϵ_i)`$ = 0. The orthogonalization procedure is needed only because states satisfying the orthogonality are unstable for decaying into Pauli-forbidden configurations. On the other hand, for the canonical-basis HFB method, the orthogonalization is essential because the single-particle Hamiltonians $`_i`$ differs from state to state. Therefore, the determination of the explicit functional form of $`\lambda _{ij}`$ is the most important part of the method. Reinhard et al. have proposed
$$\lambda _{ij}=\frac{1}{2}\psi _j|\left(_i+_j\right)|\psi _i.$$
(31)
Let us reason on which grounds the above definition can be deduced. Understanding these grounds is indispensable in order to modify the definition later for faster convergences. From the requirement that Eq. (29) must hold at the solution (where $`\psi _i|\psi _j=\delta _{ij}`$), one can deduce,
$$\lambda _{ij}=\psi _j|_i|\psi _i\text{at the solution}.$$
(32)
Eqs. (31) and (32) are equivalent at the solution because $`\lambda _{ij}`$ is defined to be hermite by Eq. (26). Since this hermiticity must hold at any points to ensure the equality between the number of constraints and the number of independent multipliers, one should not adopt Eq. (32) but Eq. (31).
One can utilize the gradient method to obtain the HFB solutions in the canonical-basis formalism. The most naive implementation of the gradient method agrees with the imaginary-time evolution method in its first order approximation of the size of the imaginary-time step $`\mathrm{\Delta }\tau `$:
$$\psi _i\psi _i\frac{1}{g}\mathrm{\Delta }\tau \frac{\delta R}{\delta \psi _i^{}}.$$
(33)
We have developed a 3D-mesh canonical-basis HFB program from scratch according to the above formulation. We take an example of our calculations using the program for the ground state of <sup>40</sup>Ca. The wavefunctions are expressed with $`39\times 39\times 39`$ mesh points with mesh spacing of $`a`$=0.8 fm. We employed the 17-point finite-difference approximation to the Laplacian. Note that the requirement of precision is higher for HFB than for HF because one has to treat larger momentum components than the Fermi momentum in HFB. The vanishing boundary conditions are imposed on the boundary (the 0th and the 40th mesh points) and the wavefunctions are anti-symmetrically reflected in the boundary to apply the finite-difference formula. We considered $`K=20`$ canonical basis, which can contain $`Kg=80`$ (=$`2\times A`$, $`A=40`$) nucleons.
For the imaginary time step size $`\mathrm{\Delta }\tau `$, it must hold$`^{\text{12})}`$
$$\mathrm{\Delta }\tau \frac{2}{T_{\mathrm{max}}},T_{\mathrm{max}}=3\frac{\mathrm{}^2}{2m}\left(\frac{\pi }{a}\right)^2.$$
(34)
We took $`\mathrm{\Delta }\tau =1/T_{\mathrm{max}}`$.
In Fig. 1, the error of the second equality of Eqs. (29) neglecting the error of orthogonality, i.e., $`\mathrm{max}_{i=1,\mathrm{},K}|_i\psi _i_j\lambda _{ij}\psi _j|`$ is plotted with a solid line versus the number of evolution steps. The corresponding quantity for HF, $`\mathrm{max}_{i=1,\mathrm{},A/4}|h\psi _i\psi _i|h|\psi _i\psi _i|`$ , is also plotted with a dash line. The figure demonstrates that one can indeed obtain HFB solutions with the natural-orbital HFB method in the 3D-mesh representation. We obtained similar convergence curves for the error of the orthogonality and for the inconsistency between the potential and the densities.
The speed of the convergence is, however, about ten times as slow as the HF case. This is the subject of the next section.
We adopted an additional procedure which is not indispensable to obtain the solutions but effective to make the convergence more robust and somewhat quicker: After every 25th gradient-method steps, we orthogonalize $`\{\psi _i\}`$ with the Gram-Schmidt algorithm in the ascending order of $`ϵ_i`$, defined in Eqs. (27), and then diagonalize the super matrix of the HFB Hamiltonian (24) by expanding the quasi-particle wavefunctions $`(\varphi _i,\phi _i)`$ ($`1iK`$ and their negative-energy partners) in a $`2K`$-dimensional basis $`\{\psi _i\}`$ $``$ $`\{\psi _i\}`$, and finally transform the resulting quasi-particle wavefunctions to canonical orbitals and occupation amplitudes $`\{\psi _i,v_i\}`$ to renew them.
Incidentally, the period of 25 steps may be too frequent because the effect seems to saturate at periods around 100. We adopt the period of 25 in most calculations, however, because the increase in the computation time is only a several percent of the total time with this period.
## Acceleration of the gradient method
We show the origin of the slow convergence and present a solution of the difficulty in this section.
Steepest-descent paths, which the gradient method draws, depend on the choice of the independent variables. For example, Eq. (33) is obtained when one uses $`(\mathrm{Re}\psi _i,\mathrm{Im}\psi _i)`$ as independent variables to define the gradient vector and then express it in coordinates $`(\psi _i,\psi _i^{})`$. If one uses scale-transformed wavefunctions $`\chi _i\alpha _i^{1/2}\psi _i`$, where $`\alpha _i`$ is a scaling factor, a gradient step becomes,
$$\chi _i\chi _i\frac{1}{g}\mathrm{\Delta }\tau \frac{\delta R}{\delta \chi _i^{}},$$
(35)
which is equivalent to
$$\psi _i\psi _i\frac{1}{g}\alpha _i\mathrm{\Delta }\tau \frac{\delta R}{\delta \psi _i^{}}.$$
(36)
The change from Eq. (33) to Eq. (36) is equivalent to multiplying $`\alpha _i`$ to the single-particle Hamiltonian $`_i`$ in Eqs. (29).
When one parameterizes the scaling factor as $`\alpha _i`$ = $`v_i^{2\nu }`$, the modified single-particle Hamiltonian becomes
$`\alpha _i_i`$ $`=`$ $`v_i^{22\nu }h+v_i^{12\nu }u_i\stackrel{~}{h}`$ (40)
$`=`$ $`\{\begin{array}{cccc}\hfill v_i^2h& +& \hfill v_iu_i\stackrel{~}{h}& (\nu =0),\hfill \\ \hfill v_ih& +& \hfill u_i\stackrel{~}{h}& (\nu =\frac{1}{2}),\hfill \\ \hfill h& +& \hfill \frac{u_i}{v_i}\stackrel{~}{h}& (\nu =1).\hfill \end{array}`$
When $`\nu =0`$ (i.e., $`\alpha _i=1`$), to which the imaginary-time evolution (33) corresponds, the single-particle Hamiltonian $`\alpha _i_i`$ (=$`_i`$) can be very small for canonical orbitals whose $`ϵ_i`$ is much higher than the Fermi level (i.e., $`ϵ_i`$ $``$ $`ϵ_\mathrm{F}`$ $``$ $`\mathrm{\Delta }_i`$). This smallness makes the changes of such orbitals very slow. On the other hand, for $`\nu =1`$, the potential can be very deep for such high-lying orbitals due to the factor $`u_i/v_i`$ in front of $`\stackrel{~}{h}`$. In this case, however, the gradient step may be numerically dangerous. We usually use $`\frac{1}{2}\nu <1`$, which provides a fast and numerically stable method of solution.
When one introduces the acceleration factors (i.e., $`\alpha _i>1`$), the multipliers $`\lambda _{ij}`$ should be modified from Eq. (31) by the following reason. In the computation of a gradient vector given by the first of Eqs. (29), the last term takes much more computing time than the first two terms due to $`\delta \lambda _{jk}/\delta \psi _i^{}`$. One can forget the last term if the orthogonality relations (12) are fulfilled along the path of the steepest descent. Let’s suppose that the relations are satisfied before a gradient-method step is taken and require that they are conserved to the first order in $`\mathrm{\Delta }\tau `$ after the step, i.e.,
$$\psi _i^{}|\psi _j^{}=\delta _{ij}+𝒪\left(\left(\mathrm{\Delta }\tau \right)^2\right)\text{if}\psi _i|\psi _j=\delta _{ij},$$
(41)
with
$$\psi _i^{}=\psi _i\alpha _i\mathrm{\Delta }\tau \left(_i\psi _i\underset{j}{}\lambda _{ij}\psi _j\right).$$
(42)
Substituting Eq. (42) into Eq. (41) and requiring the hermiticity (26) result in
$$\lambda _{ij}=\frac{1}{\alpha _i+\alpha _j}\psi _j|\left(\alpha _i_i+\alpha _j_j\right)|\psi _i.$$
(43)
This form of $`\lambda _{ij}`$ fulfills the requirement that it should agree with the expression (32) at the solution as the naive form of Eq. (31) does. Two forms differ, however, before reaching the solution if one assumes $`\alpha _i\alpha _j`$ in general. Therefore, in order to conserve the orthogonality to vanish the last term in Eq. (29) one must not use Eq. (31) but Eq. (43). We have indeed suffered from large errors of orthogonality by using the naive form (31). On the other hand, by using the correct form (43), we have observed not only that the error does not grow during the evolution but also that the error decreases without performing explicit orthogonalization procedures periodically during the evolution. This decrease should originate in the second order terms in $`\mathrm{\Delta }\tau `$ in Eq. (41), whose effects we did not consider.
We compare the results of calculations between $`\nu =0`$ and $`\nu =\frac{1}{2}`$ in Fig. 2. The wavefunctions are expressed with $`19\times 19\times 19`$ mesh points with mesh spacing of $`a`$=1.0 fm. We considered $`K=16`$ canonical basis, which can contain $`Kg=64`$ (=$`2\times A`$) nucleons. The figure shows the convergence history to a HFB solution. Four quantities are plotted as functions of the number of gradient steps. They are, from the top to the bottom, the total energy $`E`$, the pairing gap $`\overline{\mathrm{\Delta }}`$ (averaged with weight $`u_iv_i`$), the size of quadrupole deformation $`\beta `$, and the size of triaxiality of deformation $`\gamma `$. The last two quantities are determined from the mass quadrupole moments.$`^{\text{12})}`$
The dot curves were obtained without accelerations, i.e., with $`\alpha _i=1`$ or $`\nu =0`$ in Eq. (40), while the solid ones were obtained with the acceleration method with $`\alpha _i=1/v_i`$ or $`\nu =\frac{1}{2}`$. One can see that the convergences of these quantities become by far faster by using the acceleration method. This result demonstrates that canonical-basis HFB can be solved without very heavy numerical computations.
## Cut-off of the pairing interaction for canonical-basis HFB
Finally let us discuss on the cut-off schemes of the paring interaction in relation to the canonical-basis HFB method.
Delta function forces without cut-off leads to a divergence of the strength of the pairing correlation.$`^{\text{13})}`$ In order to circumvent the divergence, in the conventional method of solution, one usually takes only quasiparticles whose excitation energy $`ϵ^{\mathrm{qp}}`$ is lower than some cut-off energy parameter $`ϵ_{\mathrm{cut}}`$ to construct the ground state.$`^{\text{1})}`$ Namely Eq. (4) is modified to
$$|\mathrm{\Psi }=\underset{i=1}{\overset{I}{}}\theta \left(ϵ_{\mathrm{cut}}ϵ_i^{\mathrm{qp}}\right)b_i|0$$
(44)
where $`\theta `$ is the step function.
In the canonical-basis method, the restriction on the number of canonical orbitals may prevent the divergence without introducing explicitly a cut-off energy. We have examined this idea by performing numerical calculations for various situations. Then, we noticed that observables sometimes jumps suddenly in the course of long-time evolution.
An example is shown in Fig. 3. The calculation was done for a nucleus <sup>32</sup>S on a 19$`\times `$19$`\times `$19 cubic mesh with a mesh spacing $`a`$=1 fm. We considered $`K=20`$ canonical orbitals, to which we gave harmonic oscillator wavefunctions at the beginning. An 11-point formula was employed for the Laplacian. The acceleration parameter in Eq. (40) is taken as $`\nu =0.7`$. At first we suspected that these jumps were due to the acceleration method. However, with smaller $`\nu `$, we still observed jumps; only they come later.
We investigated the origin of the jumps and found that each sudden change was due to a shrinkage of a high-lying (i.e., having large $`ϵ_i`$) canonical orbital to a mesh point. This shrinkage can decrease the total energy of the nucleus for the following reason: The contribution of a canonical orbital to the kinetic energy density is proportional to its BCS occupation probability $`v_i^2`$, while pairing density is proportional to $`u_iv_i`$. Because $`v_i1`$ and $`u_i1`$ for high-lying orbitals, it holds that $`u_iv_iv_i^2`$. Therefore the increase in the kinetic energy due to the shrinkage is easily compensated by the gain of the pairing correlation energy at the shrunken point.
The observed jumps indicate that, in most cases (or maybe all the cases), there are no potential barriers between such physically meaningless solutions and the physically reasonable one.
Incidentally, we confirmed that the lack of potential barriers was not due to the discrete approximation of the Laplacian in the kinetic energy term: The approximation was based on the Lagrange polynomial interpolation, which tend to underestimate the expectation value for high momentum components. One might suspect that more accurate treatments of the kinetic energy could restore a barrier between the physical and unphysical solutions. However, we observed similar jumps by using the Fourier transformation with periodic boundary conditions,$`^{\text{14})}`$ which gives the exact result up to $`\pi /a`$.
We have decided that it is necessary to introduce cut-off for the reliability of the method.
As the cut-off scheme, we first employed cut-off factors which are dependent on orbitals. In the BCS approximation, a smooth cut-off method$`^{\text{8})}`$ is often utilized, in which the interaction is modified as
$`\widehat{v}_{\mathrm{pair}}`$ $`=`$ $`G\left({\displaystyle \underset{i}{}}f_ia_i^{}a_{\overline{ı}}^{}\right)\left({\displaystyle \underset{j}{}}f_ja_{\overline{ȷ}}a_j\right),`$ (45)
$`f_i`$ $`=`$ $`f(ϵ_i),`$ (46)
where $`f(ϵ)`$ is a function of single-particle energy $`ϵ`$. The function takes on $`1`$ well below a chosen cut-off energy and smoothly becomes zero above it. In analogy to the smooth cut-off method, we modified the pairing density as
$$\stackrel{~}{\rho }=g\underset{i=1}{\overset{K}{}}u_iv_i|\psi _i|^2g\underset{i=1}{\overset{K}{}}f_iu_iv_i|\psi _i|^2$$
(47)
with
$$\begin{array}{c}f_i=\mathrm{exp}\left(\frac{\mu ^2}{4}k_i^2\right),\mu =1.2\text{fm},\hfill \\ \\ k_i^2=\psi _i^{}\mathrm{}\psi 𝑑\stackrel{}{r}.\hfill \end{array}$$
(48)
We made the dependence to be on the kinetic energy ($`k_i^2`$), not on the mean-field energy $`ϵ_i`$ as in Eq. (46), to avoid a highly complicated expression of the gradient for the latter case. (In BCS, such complications are just neglected.)
The result was successful to prevent the shrinkages to points. However, this cut-off scheme has an disadvantage that the HFB super matrix in Eq. (24) cannot be defined. It follows that we do not have well-defined quasiparticle states and cannot utilize them to express the HFB ground state. This drawback is rather serious because quasiparticles are useful to improve the precision of HFB solutions obtained by the canonical-basis method and to accelerate the convergence further, as described in the fourth section.
As an alternative method, we have introduced a repulsive pairing-density dependence to the pairing force, Eq. (16), in addition to the usual density dependence:
$`\widehat{v}_\mathrm{p}(\stackrel{}{r}_1,s_1;\stackrel{}{r}_2,s_2)=v_\mathrm{p}{\displaystyle \frac{1P_\sigma }{2}}`$ (49)
$`\times `$ $`\left\{1{\displaystyle \frac{\rho (\stackrel{}{r}_1)}{\rho _\mathrm{c}}}\left({\displaystyle \frac{\stackrel{~}{\rho }(\stackrel{}{r}_1)}{\stackrel{~}{\rho }_\mathrm{c}}}\right)^2\right\}\delta \left(\stackrel{}{r}_1\stackrel{}{r}_2\right).`$
A set of reasonable values of the parameters are $`v_\mathrm{p}=440`$ MeV fm<sup>3</sup>, $`\rho _\mathrm{c}=0.32`$ fm<sup>-3</sup>, and $`\stackrel{~}{\rho }_\mathrm{c}=0.3`$ fm<sup>-3</sup>.
With forces (15) and (49), the expectation value of the energy for $`N=Z`$ systems is expressed as a space integral of a Hamiltonian density:
$`(\stackrel{}{r})`$ $`=`$ $`{\displaystyle \frac{\mathrm{}^2}{2m}}\tau (\stackrel{}{r})+{\displaystyle \frac{3}{8}}t_0\rho (\stackrel{}{r})^2+{\displaystyle \frac{1}{16}}t_3\rho (\stackrel{}{r})^{2+\alpha }`$ (50)
$`+`$ $`{\displaystyle \frac{1}{8}}v_\mathrm{p}\left\{1{\displaystyle \frac{\rho (\stackrel{}{r})}{\rho _\mathrm{c}}}\left({\displaystyle \frac{\stackrel{~}{\rho }(\stackrel{}{r})}{\stackrel{~}{\rho }_\mathrm{c}}}\right)^2\right\}\stackrel{~}{\rho }(\stackrel{}{r})^2.`$
The mean-field potential $`V`$ remains the same as Eq. (20) while the pairing potential $`\stackrel{~}{V}`$ has an additional term:
$$\stackrel{~}{V}=\frac{}{\stackrel{~}{\rho }}=\frac{1}{4}v_\mathrm{p}\left\{1\frac{\rho }{\rho _\mathrm{c}}2\left(\frac{\stackrel{~}{\rho }}{\stackrel{~}{\rho }_\mathrm{c}}\right)^2\right\}\stackrel{~}{\rho }.$$
(51)
With this new type of force, the shrinkage problem is completely removed.
In Fig. 4, we show the dependence of the pairing gap (averaged with weight $`u_iv_i`$) on the parameter $`\stackrel{~}{\rho }_\mathrm{c}`$, which controls the pairing density dependence. The values are taken after 5,000 gradient steps. The set up of the calculations are the same as in Fig. 3 except for $`\stackrel{~}{\rho }_\mathrm{c}`$. One can see that the pairing gap has reasonable values with $`\stackrel{~}{\rho }_\mathrm{c}1`$ fm<sup>-3</sup>. This is a good news because the pairing-density-dependent term, $`(\stackrel{~}{\rho }/\stackrel{~}{\rho _\mathrm{c}})^2`$, can be small for the values of $`\stackrel{~}{\rho }`$ which one finds in physical solutions, in contrast to the usual density-dependent term, $`\rho /\rho _\mathrm{c}`$, which cancels roughly 50% of the density-independent term inside nucleus. This situation is illustrated in Fig. 5. The plotted quantity is the pairing Hamiltonian density $`\stackrel{~}{}`$, which is the term in the second line of Eq. (50). The introduction of the new term demands only little change of the other parameters of the force if one adopts $`\stackrel{~}{\rho }_\mathrm{c}=0.3`$ fm<sup>-3</sup>.
We show an example of the time evolution of $`ϵ_i`$ in Fig. 6. The set up of the calculation is the same as in Figs. 3 and 4 except that a 7-point approximation is used for the Laplacian to favor the emergence of unphysical solutions and make this calculation a very severe test of the cut-off scheme. The abscissa is the number of the gradient steps while the ordinate is the expectation value $`ϵ_i`$ of the mean-field Hamiltonian for all the $`K`$ canonical orbitals. The pairing density dependence parameter $`\stackrel{~}{\rho }_\mathrm{c}`$ is set at the standard value of 0.3 fm<sup>-3</sup> before step 3,000 (interval I) and after step 6,000 (interval III). Between steps 3,000 and 6,000 (interval II) $`\stackrel{~}{\rho }_\mathrm{c}`$ is increased temporally to 3 fm<sup>-3</sup>, which results in a too weak dependence to prevent the divergence. One can see that an unphysical solution emerges without the pairing-density dependence term (in interval II) and that the unphysical solution is suppressed (in interval I) and is quickly restored to a physical one (in interval III) with the presence of the term.
Before ending the section, let us mention that our next subject concerning the cut-off is whether repulsive momentum dependences of the Skyrme force help to make the cut-off scheme more natural because it has been known that such momentum dependences make well-developed plateau before the divergence sets in.$`^{\text{13})}`$
## Conclusions
We have developed a method to obtain canonical-basis HFB solutions in a coordinate-space three-dimensional (3D) Cartesian mesh representation. The features of our method are summarized as follows.
i) It is not for spherical but for deformed nuclei and thus it can treat both deformation and continuum pairing simultaneously.
ii) It is not based on the oscillator-basis expansion but described in the 3D Cartesian mesh representation and thus can treat e.g. deformed halo-like orbitals.
iii) There is a strong reason to believe that the necessary number of canonical orbitals is much smaller than the number of single-particle basis.
iv) In order to perform variations under constraint of orthogonality between the wavefunctions of the canonical orbitals, Lagrange multipliers were introduced as functionals of the wavefunctions in Ref.$`^{\text{2})}`$. We have clarified the necessary conditions for the form of these Lagrange multiplier functionals. We have modified the functionals appropriately so that they conserve the orthogonality during the course of the accelerated evolutions explained in the next item.
v) We have found that the convergence to HFB solutions is very slow when one employs a naive gradient method. On the other hand, the convergence is quite rapid for Hartree-Fock solutions, which neglects the pairing correlation. We have investigated the origin of the slow convergence and found that the time scale of the gradient evolution is different from one canonical orbital to another depending on their BCS occupation amplitudes $`v_i`$. The difference ranges over many orders of magnitude. We have introduced an orbital-dependent acceleration method of the gradient evolutions and could overcome the difficulty of the slow convergence.
vi) We have examined the effects of the cut-off of the pairing interaction. The 3D mesh mean-field methods have a practical use only for zero-range interactions like the Skyrme force presently. One needs a cut-off for zero-range pairing forces, without which the pairing correlation energy diverges to $`\mathrm{}.^{\text{13})}`$ We have found that zero-range forces need cut-off even when the number of canonical orbitals are finite. The divergence occurs through shrinking of high-energy canonical orbital(s) to (a) mesh point(s).
vii) To suppress the divergence, we have introduced a pairing-density dependent interaction as a better choice than orbital-dependent cut-off factors.
We believe that, by choosing a faster gradient path with the acceleration method and adopting the cut-off scheme in terms of the pairing-density dependence, the canonical-basis HFB method is now fully understood and has the potential to become the standard method to treat neutron-rich nuclei in HFB. The details of this work will be published soon.$`^{\text{15})}`$
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This paper has been published in RIKEN Review No. 26 (January, 2000), pp. 87-94. The issue is devoted to the proceedings of an international symposium on Models and Theories of the Nuclear Mass, RIKEN, Wako-shi, Saitama, Japan, July 19-23, 1999. PDF files of all the issues of RIKEN Review can be obtained freely from the following URL: http://www.riken.go.jp/lab-www/library/publication/review/contents.html |
warning/0002/math0002245.html | ar5iv | text | # The Flag Major Index and Group Actions on Polynomial Rings
## 1 Introduction
### 1.1 Outline
The major index, $`\mathrm{𝑚𝑎𝑗𝑜𝑟}(\pi )`$, of a permutation $`\pi `$ in the symmetric group $`S_n`$ is the sum (possibly zero) of all indices $`1i<n`$ for which $`\pi (i)>\pi (i+1)`$. The length of a permutation $`\pi `$ is the minimal number of factors in an expression of $`\pi `$ as a product of the Coxeter generators $`(i,i+1),1i<n`$. A fundamental property of the major index is its equidistribution with the length function \[MM\]; namely, the number of elements in $`S_n`$ of a given length $`k`$ is equal to the number of elements having major index $`k`$. Bijective proofs and generalizations were given in \[F, FS , Ca , GG, Go, Ro3\].
Candidates for a major index for the classical Weyl groups of type $`B`$ have been suggested by Clarke-Foata \[CF1–3\], Reiner \[Rei1–2\], Steingrimsson \[Ste\], and others. Unfortunately, unlike the case of the symmetric group, the various alternatives are not equidistributed with the length function (defined with respect to the Coxeter generators of $`B_n`$).
In this paper we present a new definition of the major index for the groups $`B_n`$, and more generally for wreath products of the form $`C_mS_n`$, where $`C_m`$ is the cyclic group of order $`m`$. This major index is shown to be equidistributed with the length function for the groups of type $`B`$ (the case $`m=2`$), and to play a crucial role in the study of the actions of these groups on polynomial rings.
The rest of the paper is organized as follows. The definition of the flag major index is presented in Subsection 1.2 (and, in more detail, in Section 2). Main results are surveyed in Subsection 1.3. Basic properties and equidistribution with length are proved in Section 2. In Section 3 we give a combinatorial interpretation of the flag major index. In Section 4 we study $`C_mS_n`$-actions on tensor powers of polynomial rings. The combinatorial interpretation of the flag major index is then applied to obtain a simple representation of the Hilbert Series of the diagonal action invariant algebra.
### 1.2 The Flag Major Index
The groups $`C_mS_n`$ are generated by $`n1`$ involutions, $`s_1,\mathrm{},s_{n1}`$, which satisfy the usual Moore-Coxeter relations of $`S_n`$, together with an exceptional generator, $`s_0`$, of order $`m`$. We consider a different set of generators:
$$t_i:=\underset{j=0}{\overset{i}{}}s_{ij}(0in1).$$
These are Coxeter elements in a distinguished flag of parabolic subgroups. See Section 2 for more details.
An element $`\pi C_mS_n`$ has a unique representation as a product
$$\pi =t_{n1}^{k_{n1}}t_{n2}^{k_{n2}}\mathrm{}t_1^{k_1}t_0^{k_0}$$
with $`0k_i<m(i+1)`$ $`(i)`$. Define the flag major index of $`\pi `$ by
$$\text{flag-major}(\pi ):=\underset{i=0}{\overset{n1}{}}k_i.$$
For $`m=1`$, this definition gives a new interpretation of a well-known parameter.
Claim 1. For $`m=1`$ (i.e., for the symmetric group $`S_n`$) the flag major index coincides with the major index.
See Claim 2.1 below.
### 1.3 Main Results
From Claim 1 together with MacMahon’s classical result it follows that the flag major index is equidistributed with length, for $`m=1`$. This property extends to $`m=2`$.
Theorem 2. For $`m=2`$ (i.e., for the hyperoctahedral group $`B_n`$) the flag major index is equidistributed with length.
Here “length” is used in the usual sense, in terms of the Coxeter generators $`s_0,s_1,\mathrm{},s_{n1}`$. See Theorem 2.2 below.
For $`m3`$, the flag major index is no longer equidistributed with length (with respect to $`s_0,\mathrm{},s_{n1}`$), but nevertheless it does play a central role in the study of naturally defined algebras of polynomials.
Let $`P_n:=𝑪[x_1,\mathrm{},x_n]`$ be the algebra of polynomials in $`n`$ indeterminates. There is a natural action of $`G:=C_mS_n`$ on $`P_n`$ (presented explicitly in Section 4). Consider now the tensor power $`P_n^t`$ with the natural tensor action $`\phi _T`$ of $`G^t:=G\times \mathrm{}\times G`$ ($`t`$ factors), and the corresponding diagonal action of $`G`$. For more details see Section 4.
The tensor invariant algebra TIA is a subalgebra of the diagonal invariant algebra DIA. Let $`F_D(\overline{q})`$, where $`\overline{q}=(q_1,\mathrm{},q_t)`$, be the multi-variate generating function (Hilbert series) for the dimensions of the homogeneous components in DIA:
$$F_D(\overline{q}):=\underset{n_1,\mathrm{},n_t𝑵}{}(\text{ dim }_𝑪\text{ DIA}_{n_1,\mathrm{},n_t}\text{ })\text{ }q_1^{n_1}\mathrm{}q_t^{n_t},$$
where $`\text{ DIA}_{n_1,\mathrm{},n_t}`$ is the homogeneous piece of multi-degree $`(n_1,\mathrm{},n_t)`$ in DIA. Define similarly $`F_T(\overline{q})`$ for TIA. Then
Theorem 3. For all $`m,n,t1`$,
$$\frac{F_D(\overline{q})}{F_T(\overline{q})}=\underset{\pi _1\mathrm{}\pi _t=1}{}\underset{i=1}{\overset{t}{}}q_i^{\mathrm{𝑓𝑙𝑎𝑔}\mathrm{𝑚𝑎𝑗𝑜𝑟}(\pi _i)}$$
where the sum extends over all $`t`$-tuples $`(\pi _1,\mathrm{},\pi _t)`$ of elements in $`G=C_mS_n`$ such that the product $`\pi _1\pi _2\mathrm{}\pi _t`$ is equal to the identity element.
See Theorem 4.1 below.
## 2 Flag Major Index and Length
The groups $`C_mS_n`$ are generated by $`n1`$ involutions, $`s_1,\mathrm{},s_{n1}`$, together with an exceptional generator, $`s_0`$, of order $`m`$. The $`n1`$ involutions satisfy the usual Moore-Coxeter relations of $`S_n`$,
$$(s_is_{i+1})^3=1(1i<n),$$
$$(s_is_j)^2=1(|ij|>1),$$
while the exceptional generator satisfies the relations:
$$(s_0s_1)^{2m}=1,$$
$$s_0s_i=s_is_0(1<i<n).$$
For $`m=1`$, $`s_0`$ is the identity element.
Consider now a different set of generators:
$$t_i:=\underset{j=0}{\overset{i}{}}s_{ij}(0in1).$$
These are Coxeter elements \[Hu, §3.16\] in a distinguished flag of parabolic subgroups
$$1<G_1<\mathrm{}<G_n=C_mS_n$$
where $`G_iC_mS_i`$ is the subgroup of $`C_mS_n`$ generated by $`s_0,s_1,\mathrm{},s_{i1}`$.
An element $`\pi C_mS_n`$ has a unique representation as a product
$$\pi =t_{n1}^{k_{n1}}t_{n2}^{k_{n2}}\mathrm{}t_1^{k_1}t_0^{k_0}$$
with $`0k_i<m(i+1)`$ $`(i)`$. Define the flag major index of $`\pi `$ by
$$\text{flag-major}(\pi ):=\underset{i=0}{\overset{n1}{}}k_i.$$
For $`m=1`$, this definition gives a new interpretation of a well-known parameter.
Claim 2.1. For $`m=1`$ (i.e., for the symmetric group $`S_n`$) the flag major index coincides with the major index.
Proof. Consider the natural action of $`S_n`$ on the letters $`1,\mathrm{},n`$, where $`s_i`$ ($`1i<n`$) acts as the transposition $`(i,i+1)`$. Here $`t_0=s_0`$ is the identity permutation, whereas for $`1r<n`$, $`t_r=s_r\mathrm{}s_1`$ acts as the cycle $`(r+1,r,\mathrm{},2,1)`$. The claim may be proved by induction on the exponents $`k_i`$. Obviously, equality holds when $`k_1=\mathrm{}=k_{n1}=0`$. It suffices to prove that for any $`\pi =t_r^{k_r}\mathrm{}t_1^{k_1}`$ with $`1r<n`$ and $`0k_r<r`$, $`\mathrm{𝑚𝑎𝑗𝑜𝑟}(t_r\pi )\mathrm{𝑚𝑎𝑗𝑜𝑟}(\pi )=1`$. In this case $`\pi (j)=j`$ for $`j>r+1`$ and $`\pi (r+1)\{2,3,\mathrm{},r+1\}`$. Observe that, for $`1jr+1`$, $`t_r\pi (j)>t_r\pi (j+1)`$ if and only if $`\pi (j)>\pi (j+1)`$, unless $`\pi (j)=1`$ or $`\pi (j+1)=1`$. If $`\pi (j)=1`$ (this holds for a unique $`1jr`$) then $`t_r\pi (j)=r+1`$, so that $`\pi (j1)>\pi (j)<\pi (j+1)`$ whereas $`t_r\pi (j1)<t_r\pi (j)>t_r\pi (j+1)`$; thus the descent at $`j1`$ in $`\pi `$ is replaced by a descent at $`j`$ in $`t_r\pi `$. The special case $`j=1`$ is similar.
$`\mathrm{}`$
In particular, by MacMahon’s classical result, the flag major index is equidistributed with length (for $`m=1`$). This property extends to $`m=2`$.
Theorem 2.2. For $`m=2`$ (i.e., for the hyperoctahedral group $`B_n`$), the flag major index is equidistributed with length.
Here “length” is used in the usual sense, in terms of the Coxeter generators $`s_0,s_1,\mathrm{},s_{n1}`$. Theorem 2.2 can be proved by an explicit bijection.
Proof. For $`0m<2n`$ define $`r_{n,m}B_n`$ by
$$r_{n,m}:=\{\begin{array}{cc}id,\hfill & \text{if }m=0\text{;}\hfill \\ _{j=nm}^{n1}s_j,\hfill & \text{if }0<mn\text{;}\hfill \\ _{j=0}^{mn1}s_{mnj}_{j=0}^{n1}s_j,\hfill & \text{if }n<m<2n\text{,}\hfill \end{array}$$
where $`id`$ is the identity element in $`B_n`$.
Note that the length of $`r_{n,m}`$ is $`m`$. The set $`\{r_{n,m}|0m<2n\}`$ forms a complete set of representatives of minimal length for the left cosets of $`B_{n1}`$ in $`B_n`$. It follows that every element $`\pi B_n`$ has a unique representation as a product $`\pi =_{i=1}^nr_{n+1i,m_{n+1i}}`$, where $`0m_j<2j`$ for every $`j`$, and then $`\mathrm{}(\pi )=_{j=1}^nm_j`$.
On the other hand, every element in $`B_n`$ has a unique representation as a product of the form $`_{i=1}^nt_{ni}^{k_{ni}}`$, where $`0k_j<2(j+1)`$ for every $`j`$. By definition, $`\text{flag-major}(\pi )=_{j=1}^nk_{j1}`$.
It follows that the map $`\varphi :B_nB_n`$ defined by
$$\varphi (\underset{i=1}{\overset{n}{}}r_{n+1i,m_{n+1i}}):=\underset{i=1}{\overset{n}{}}t_{ni}^{m_{n+1i}}$$
is bijective and sends the length function to the flag major index.
$`\mathrm{}`$
For $`m3`$ (when $`C_mS_n`$ is no longer a Coxeter group), the flag major index is no longer equidistributed with length (with respect to $`s_0,\mathrm{},s_{n1}`$); nevertheless, it does play a central role in the study of naturally defined algebras of polynomials, as will be shown in Section 4.
## 3 A Combinatorial Interpretation
The standard major index for permutations has a natural generalization to any finite sequence of letters from a linearly ordered alphabet (see, e.g., \[F\]); namely, for a finite sequence $`a=(a_1,a_2,\mathrm{},a_n)`$ of letters from a linearly ordered alphabet, define $`\mathrm{𝑚𝑎𝑗𝑜𝑟}(a)`$ to be the sum (possibly zero) of all indices $`1i<n`$ for which $`a_i>a_{i+1}`$.
Let $`\omega 𝑪`$ be a primitive $`m`$-th root of unity. An element of the wreath product $`C_mS_n`$ (where $`C_m`$ is the cyclic group of order $`m`$) may be described as a generalized permutation $`\pi =(\pi (1),\pi (2),\mathrm{},\pi (n))`$, where, for every $`i`$, $`\pi (i)𝑪`$, $`\frac{\pi (i)}{|\pi (i)|}`$ is a power of $`\omega `$, and the sequence of absolute values $`|\pi |:=(|\pi (1)|,|\pi (2)|,\mathrm{},|\pi (n)|)`$ is a permutation in $`S_n`$. In this setting, the involutions $`s_i`$ ($`1i<n`$) are defined by
$$s_i(j):=\{\begin{array}{cc}i+1,\hfill & \text{if }j=i\text{;}\hfill \\ i,\hfill & \text{if }j=i+1\text{;}\hfill \\ j,\hfill & \text{otherwise,}\hfill \end{array}$$
whereas the exceptional generator $`s_0`$ is defined by
$$s_0(j):=\{\begin{array}{cc}\omega 1,\hfill & \text{if }j=1\text{;}\hfill \\ j,\hfill & \text{otherwise.}\hfill \end{array}$$
In particular, $`C_2S_n`$ is the group of signed permutations, also known as the hyperoctahedral group, or the classical Weyl group of type $`B`$.
Consider now the linearly ordered alphabet
$$1\omega ^{m1}<\mathrm{}<n\omega ^{m1}<\mathrm{}\mathrm{}<1\omega ^1<\mathrm{}<n\omega ^1<1\omega ^0<\mathrm{}<n\omega ^0.$$
In this section we prove
Theorem 3.1. For any $`\pi C_mS_n`$ ,
$$\text{flag-major}(\pi )=m\mathrm{𝑚𝑎𝑗𝑜𝑟}(\pi )+\underset{j=0}{\overset{m1}{}}j\mathrm{\#}\{i:\frac{\pi (i)}{|\pi (i)|}=\omega ^j\},$$
$`(3.1)`$
where $`\mathrm{𝑚𝑎𝑗𝑜𝑟}(\pi )`$ is defined with respect to the above order.
To simplify the proof some notations are needed. For any generalized permutation $`\pi C_mS_n`$ and $`1in`$ define
$$\mathrm{Log}_\omega \pi (i):=\mathrm{min}\left\{d0:\omega ^d=\frac{\pi (i)}{|\pi (i)|}\right\}.$$
$`(3.2)`$
Then clearly
$$\underset{i=1}{\overset{n}{}}\mathrm{Log}_\omega \pi (i)=\underset{j=0}{\overset{m1}{}}j\mathrm{\#}\{i:\frac{\pi (i)}{|\pi (i)|}=\omega ^j\}.$$
$`(3.3)`$
Denote the right hand side of (3.1) by $`\mathrm{𝑚𝑎𝑗𝑜𝑟}_{m,n}(\pi )`$. By (3.3),
$$\mathrm{𝑚𝑎𝑗𝑜𝑟}_{m,n}(\pi )=m\mathrm{𝑚𝑎𝑗𝑜𝑟}(\pi )+\underset{i=1}{\overset{n}{}}\mathrm{Log}_\omega \pi (i).$$
$`(3.4)`$
Lemma 3.2. For any $`\pi C_mS_n`$ with $`\pi (n)\omega ^{m1}`$,
$$\mathrm{𝑚𝑎𝑗𝑜𝑟}_{m,n}(t_{n1}\pi )\mathrm{𝑚𝑎𝑗𝑜𝑟}_{m,n}(\pi )=1.$$
Proof of Lemma 3.2. By definition,
$$t_{n1}(i)=\{\begin{array}{cc}i1,\hfill & \text{if }i1\text{;}\hfill \\ \omega n,\hfill & \text{if }i=1\text{.}\hfill \end{array}$$
Hence, for any $`\pi C_mS_n`$,
$$t_{n1}\pi (j)=\{\begin{array}{cc}\omega ^{\mathrm{Log}_\omega \pi (j)}(\pi (j)1),\hfill & \text{if }|\pi (j)|1\text{;}\hfill \\ \omega ^{\mathrm{Log}_\omega \pi (j)+1}n,\hfill & \text{if }|\pi (j)|=1\text{.}\hfill \end{array}$$
$`(3.5)`$
Let $`1i_0n`$ be the unique index for which $`|\pi (i_0)|=1`$.
Case (a). $`\mathrm{Log}_\omega \pi (i_0)<m1`$.
* It follows from (3.5) that if $`\mathrm{Log}_\omega \pi (i_0)<m1`$ then
$$\mathrm{𝑚𝑎𝑗𝑜𝑟}(t_{n1}\pi )=\mathrm{𝑚𝑎𝑗𝑜𝑟}(\pi )$$
with respect to the linear order defined before Theorem 3.1. Hence, in this case
$$\mathrm{𝑚𝑎𝑗𝑜𝑟}_{m,n}(t_{n1}\pi )=m\mathrm{𝑚𝑎𝑗𝑜𝑟}(t_{n1}\pi )+\underset{i=1}{\overset{n}{}}\mathrm{Log}_\omega t_{n1}\pi (i)=$$
$$=m\mathrm{𝑚𝑎𝑗𝑜𝑟}(\pi )+\underset{i=1}{\overset{n}{}}\mathrm{Log}_\omega t_{n1}\pi (i)=$$
$$=m\mathrm{𝑚𝑎𝑗𝑜𝑟}(\pi )+\underset{ii_0}{}\mathrm{Log}_\omega \pi (i)+(\mathrm{Log}_\omega \pi (i_0)+1)=\mathrm{𝑚𝑎𝑗𝑜𝑟}_{m,n}(\pi )+1.$$
Case (b). $`\mathrm{Log}_\omega \pi (i_0)=m1`$.
* In this case $`\pi (i_0)=1\omega ^{m1}`$ and $`t_{n1}\pi (i_0)=n`$. By assumption $`i_0n`$, so that $`\pi `$ has a descent at $`i_01`$ (unless $`i_0=1`$), whereas $`t_{n1}\pi `$ has a descent at $`i_0`$. Thus
$$\mathrm{𝑚𝑎𝑗𝑜𝑟}(t_{n1}\pi )\mathrm{𝑚𝑎𝑗𝑜𝑟}(\pi )=1$$
and $`\mathrm{Log}_\omega t_{n1}\pi (i_0)=0`$. Hence, in this case,
$$\mathrm{𝑚𝑎𝑗𝑜𝑟}_{m,n}(t_{n1}\pi )=m\mathrm{𝑚𝑎𝑗𝑜𝑟}(t_{n1}\pi )+\underset{i=1}{\overset{n}{}}\mathrm{Log}_\omega t_{n1}\pi (i)=$$
$$=m(\mathrm{𝑚𝑎𝑗𝑜𝑟}(\pi )+1)+\underset{ii_0}{}\mathrm{Log}_\omega \pi (i)=\mathrm{𝑚𝑎𝑗𝑜𝑟}_{m,n}(\pi )+m(m1).$$
$`\mathrm{}`$
Proof of Theorem 3.1. By induction on $`n`$. Obviously, Theorem 3.1 holds in the group $`C_mS_1=s_0`$. Assume that the theorem holds in the group $`C_mS_n`$, for some $`n1`$. Any element of $`C_mS_{n+1}`$ has the form $`t_n^{k_n}\pi `$, where $`\pi C_mS_n`$ and $`0k_n<m(n+1)`$. By definition,
$$\text{flag-major}(t_n^{k_n}\pi )=k_n+\text{flag-major}(\pi )$$
and
$$\mathrm{𝑚𝑎𝑗𝑜𝑟}_{m,n+1}(\pi )=\mathrm{𝑚𝑎𝑗𝑜𝑟}_{m,n}(\pi ).$$
Hence, by the induction hypothesis, it suffices to prove that for any $`\pi C_mS_n`$ and $`0k_n<m(n+1)`$
$$\mathrm{𝑚𝑎𝑗𝑜𝑟}_{m,n+1}(t_n^{k_n}\pi )\mathrm{𝑚𝑎𝑗𝑜𝑟}_{m,n+1}(\pi )=k_n.$$
This equality readily follows from iterations of Lemma 3.2, thereby completing the proof.
$`\mathrm{}`$
## 4 Diagonal Action on Tensor Powers
Let $`P_n:=𝑪[x_1,\mathrm{},x_n]`$ be the algebra of polynomials in $`n`$ indeterminates. There is a natural action of $`G:=C_mS_n`$ on $`P_n`$, $`\phi :G\text{Aut}(P_n)`$, defined on generators by
$$\phi (s_0)(x_j)=\{\begin{array}{cc}\omega x_j,\hfill & \text{if }j=1\text{;}\hfill \\ x_j,\hfill & \text{otherwise,}\hfill \end{array}(\omega :=\text{exp}(2\pi i/m)𝑪)$$
$$\phi (s_i)(x_j)=\{\begin{array}{cc}x_{i+1},\hfill & \text{if }j=i\text{;}\hfill \\ x_i,\hfill & \text{if }j=i+1\text{;}\hfill \\ x_j,\hfill & \text{otherwise,}\hfill \end{array}(1in1)$$
where each $`\phi (s_i),0in1`$, is extended to an algebra automorphism of $`P_n`$. Equivalently, in terms of generalized permutations,
$$\phi (\pi )(x_j)=\frac{\pi (j)}{|\pi (j)|}x_{|\pi (j)|}(\pi G,1jn)$$
extended multiplicatively to monomials and additively to all of $`P_n`$.
Consider now the tensor power $`P_n^t:=P_n\mathrm{}P_n`$ ($`t`$ factors) with the natural tensor action $`\phi _T`$ of $`G^t:=G\times \mathrm{}\times G`$ ($`t`$ factors). The diagonal embedding
$$d:GG^t$$
defined by
$$g(g,\mathrm{},g)G^t(gG)$$
defines the diagonal action of $`G`$ on $`P_n^t`$:
$$\phi _D:=\phi _Td.$$
The tensor invariant algebra
$$\text{TIA}:=\{\overline{p}P_n^t|\phi _T(\overline{g})(\overline{p})=\overline{p},\overline{g}G^t\}$$
is a subalgebra of the diagonal invariant algebra
$$\text{DIA}:=\{\overline{p}P_n^t|\phi _D(g)(\overline{p})=\overline{p},gG\}.$$
Note that $`\text{TIA}=(P_n^G)^t`$, where $`P_n^G`$ is the subalgebra of $`P_n`$ invariant under $`\phi (G)`$.
The algebra $`P_n^t`$ is $`𝑵^t`$-graded by multi-degree, where $`𝑵:=\{0,1,2,\mathrm{}\}`$. Let $`F_D(\overline{q})`$, where $`\overline{q}=(q_1,\mathrm{},q_t)`$, be the multivariate generating function (Hilbert series) for the dimensions of the homogeneous components in DIA:
$$F_D(\overline{q}):=\underset{n_1,\mathrm{},n_t𝑵}{}(\text{ dim }_𝑪\text{ DIA}_{n_1,\mathrm{},n_t}\text{ })\text{ }q_1^{n_1}\mathrm{}q_t^{n_t},$$
where $`\text{ DIA}_{n_1,\mathrm{},n_t}`$ is the homogeneous piece of multi-degree $`(n_1,\mathrm{},n_t)`$ in DIA. Define similarly $`F_T(\overline{q})`$ for TIA. Then
Theorem 4.1. For all $`m,n,t1`$,
$$\frac{F_D(\overline{q})}{F_T(\overline{q})}=\underset{\pi _1\mathrm{}\pi _t=1}{}\underset{i=1}{\overset{t}{}}q_i^{\mathrm{𝑓𝑙𝑎𝑔}\mathrm{𝑚𝑎𝑗𝑜𝑟}(\pi _i)}$$
where the sum extends over all $`t`$-tuples $`(\pi _1,\mathrm{},\pi _t)`$ of elements in $`G=C_mS_n`$ such that the product $`\pi _1\pi _2\mathrm{}\pi _t`$ is equal to the identity element.
## 5 Proof of Theorem 4.1
### 5.1 Preliminaries
To prove Theorem 4.1 we need two theorems of Garsia and Gessel.
A $`t`$-partite partition with $`n`$ parts (in the sense of Gordon \[G\] and Garsia-Gessel \[GG\]) is a sequence $`f=(f_1,\mathrm{},f_t)`$ of non-negative-integer valued functions $`f_i:\{1,2,\mathrm{},n\}𝑵`$, $`1it`$, satisfying the condition:
If $`f_i(j)=f_i(j+1)`$ for all $`i<i_0`$, then $`f_{i_0}(j)f_{i_0}(j+1)`$ $`(i_0,j)`$.
In particular, for $`i_0=1`$:
$$f_1(1)f_1(2)\mathrm{}f_1(n)0,$$
i.e., $`f_1`$ is a partition with (at most) $`n`$ parts.
Example. $`n=4,t=2`$ ; $`f_1=(1,1,0,0)`$, $`f_2=(1,0,2,2)`$.
Let $`B_{t,n}`$ be the set of all $`t`$-partite partitions $`f=(f_1,\mathrm{},f_t)`$ with $`n`$ parts. Denote the sum $`\underset{j=1}{\overset{n}{}}f_i(j)`$ by $`|f_i|`$. The following theorem was first proved by Garsia and Gessel.
Theorem GG1 \[GG, Theorem 2.2 and Remark 2.2\]
$$\underset{fB_{t,n}}{}q_1^{|f_1|}\mathrm{}q_t^{|f_t|}=\frac{\underset{\pi _1\mathrm{}\pi _t=1}{}\underset{i=1}{\overset{t}{}}q_i^{\mathrm{𝑚𝑎𝑗𝑜𝑟}(\pi _i)}}{_{i=1}^t_{j=1}^n(1q_i^j)}$$
where the sum in the numerator of the right-hand side extends over all $`t`$-tuples $`(\pi _1,\mathrm{},\pi _t)`$ of permutations in $`S_n`$ such that the product $`\pi _1\pi _2\mathrm{}\pi _t`$ is equal to the identity permutation.
Let $`n_1,\mathrm{},n_r`$ be non-negative integers such that $`_{i=1}^rn_i=n`$. Recall that the $`q`$-multinomial coefficient $`\left[\genfrac{}{}{0pt}{}{n}{n_1,\mathrm{},n_r}\right]_q`$ is defined by:
$$[0]!_q:=1,$$
$$[n]!_q:=[n1]!_q(1+q+\mathrm{}+q^{n1})(n1),$$
$$\left[\genfrac{}{}{0pt}{}{n}{n_1\mathrm{}n_r}\right]_q:=\frac{[n]!_q}{[n_1]!_q\mathrm{}[n_r]!_q}.$$
Now let $`(N_1,\mathrm{},N_r)`$ be a partition of the set $`N:=\{1,\mathrm{},n\}`$. For each $`1ir`$ let $`\pi _i`$ be a permutation on the elements of $`N_i`$. Recall that a permutation $`\sigma S_n`$ is a shuffle of $`\pi _1,\pi _2,\mathrm{},\pi _r`$ if, for every $`i`$, the letters of $`N_i`$ appear in $`\sigma `$ in the same order as the corresponding letters appear in $`\pi _i`$.
Example. $`N_1=\{1,2,4\}`$, $`N_2=\{3,5\}`$; $`\pi _1=241`$, $`\pi _2=35`$. Here $`\sigma =32541`$ is a shuffle of $`\pi _1`$ and $`\pi _2`$.
Theorem GG2 \[GG, Theorem 3.1\] Let $`\mathrm{\Omega }(\pi _1,\mathrm{},\pi _r)`$ be the collection of all shuffles of given permutations $`\pi _1,\pi _2,\mathrm{},\pi _r`$. Then
$$\underset{\sigma \mathrm{\Omega }(\pi _1,\mathrm{},\pi _r)}{}q^{\mathrm{𝑚𝑎𝑗𝑜𝑟}(\sigma )}=\left[\genfrac{}{}{0pt}{}{n}{n_1\mathrm{}n_r}\right]_qq^{\mathrm{𝑚𝑎𝑗𝑜𝑟}(\pi _1)+\mathrm{}+\mathrm{𝑚𝑎𝑗𝑜𝑟}(\pi _r)},$$
where $`n_i`$ is the number of elements acted upon by $`\pi _i`$ ($`1ir`$).
### 5.2 The Case $`m=1`$
Before we get to the actual proof of Theorem 4.1, let us sketch the proof of the case $`m=1`$ (i.e., $`G=S_n`$). This is done as an indication of the structure of the general case; full details will be given in the actual proof.
The tensor invariant algebra TIA is equal to $`(P_n^G)^t`$, and $`P_n^G`$ is freely generated (as an algebra) by the $`n`$ elementary symmetric functions in $`n`$ indeterminates $`x_1,\mathrm{},x_n`$. Thus
$$F_T(\overline{q})=\frac{1}{_{i=1}^t_{j=1}^n(1q_i^j)}.$$
The diagonal invariant algebra DIA is linearly spanned by the polynomials
$$\underset{gS_n}{}\phi _D(g)(\overline{x}^f)$$
where $`\overline{x}^f`$ runs through all the monomials in $`P_n^t`$. Now $`\{\overline{x}^f|fB_{t,n}\}`$ is a complete set of representatives for the orbits of monomials in $`P_n^t`$ under the diagonal action $`\phi _D`$ of $`S_n`$. It follows that a basis for DIA is
$$\{\underset{gS_n}{}\phi _D(g)(\overline{x}^f)|fB_{t,n}\}.$$
Therefore, the generating function $`F_D(\overline{q})`$ is given by the left-hand side of the equation in Theorem GG1. By Claim 2.1, this proves Theorem 4.1 for $`m=1`$. (Theorem GG2 is not needed in this case.)
In the next subsection we shall construct an explicit basis for the diagonal invariant algebra DIA, for general $`m`$. In Subsection 5.5 we shall compute its generating function.
### 5.3 A Basis for DIA
Consider now the general case of $`G=C_mS_n`$. A linear basis for $`P_n^t`$ consists of the (tensor) monomials
$$\overline{x}^f:=\underset{i=1}{\overset{t}{}}\underset{j=1}{\overset{n}{}}x_j^{f_i(j)},$$
where $`f_i(j)`$ are non-negative integers ($`i,j`$). Let $`F`$ denote the set of all such multi-powers $`f`$. The canonical projection $`\pi :P_n^t\text{DIA}`$ is defined by
$$\pi (\overline{p}):=\underset{gG}{}\phi _D(g)(\overline{p})(\overline{p}P_n^t),$$
so that
$$\text{DIA}=\text{ span }\{\pi (\overline{x}^f)|fF\}.$$
We are looking for a subset $`B`$ of $`F`$ such that the set $`\{\pi (\overline{x}^f)|fB\}`$ is a basis for DIA.
First, let
$$F_0:=\{fF|\underset{i=1}{\overset{t}{}}f_i(j)0(\text{mod}m)(j)\}.$$
Claim 5.1. For $`fF`$,
$$\pi (\overline{x}^f)0fF_0.$$
Proof.
$$\phi _D(s_0)(\overline{x}^f)=\omega ^{\alpha (f)}\overline{x}^f,$$
where
$$\alpha (f):=\underset{i=1}{\overset{t}{}}f_i(1).$$
Now, for any $`\alpha 𝒁`$:
$$\underset{k=0}{\overset{m1}{}}(\omega ^\alpha )^k=\{\begin{array}{cc}m,\hfill & \text{if }\alpha 0(\text{mod}m)\text{;}\hfill \\ 0,\hfill & \text{otherwise.}\hfill \end{array}$$
Therefore, if $`C`$ is any left coset in $`G`$ of the subgroup generated by $`s_0`$, then
$$\underset{gC}{}\phi _D(g)(\overline{x}^f)=0$$
unless
$$\underset{i=1}{\overset{t}{}}f_i(1)0(\text{mod}m).$$
Replacing $`s_0`$ by its $`n1`$ conjugates, it follows that
$$\pi (\overline{x}^f)0fF_0.$$
For the converse, let $`H`$ be the normal commutative subgroup of $`G`$ generated by $`s_0`$ and its conjugates: $`H`$ consists of all “generalized identity permutations” $`hG`$ satisfying $`|h(i)|=i`$ $`(i)`$. Let $`\widehat{G}`$ be the subgroup of $`G`$ consisting of all “unsigned permutations” $`\widehat{g}G`$, satisfying $`\mathrm{Log}_\omega \widehat{g}(i)=0`$ $`(i)`$. Then $`\widehat{G}S_n`$ and $`G`$ is the semidirect product of $`\widehat{G}`$ and $`H`$, hence each $`gG`$ has a unique representation as a product
$$g=\widehat{g}h(\widehat{g}\widehat{G},\text{ }hH).$$
For any $`fF_0`$, $`\phi _D(h)(\overline{x}^f)=\overline{x}^f`$ for all $`hH`$. Hence, for $`fF_0`$,
$$\underset{g\widehat{g}H}{}\phi _D(g)(\overline{x}^f)=|H|\phi _D(\widehat{g})(\overline{x}^f)(\widehat{g}\widehat{G}).$$
In other words, the sum over a coset of $`H`$ in $`G`$ is a monomial multiplied by a positive integer. It follows that
$$fF_0\pi (\overline{x}^f)0.$$
$`\mathrm{}`$
For $`\overline{p}=\underset{fF}{}c_f\overline{x}^fP_n^t`$ (a finite sum), define the support
$$\mathrm{𝑠𝑢𝑝𝑝}(\overline{p}):=\{fF\text{ }|\text{ }c_f0\}.$$
The following result is a consequence of the proof of Claim 5.1.
Claim 5.2. For $`f,hF_0`$, the following are equivalent:
$`(i)`$ $`\mathrm{𝑠𝑢𝑝𝑝}(\pi (\overline{x}^f))\mathrm{𝑠𝑢𝑝𝑝}(\pi (\overline{x}^h))\mathrm{}`$;
$`(ii)`$ $`\pi (\overline{x}^f)=\pi (\overline{x}^h)`$;
$`(iii)`$ $`\sigma S_n`$, such that
$$h_i(j)=f_i(\sigma (j))(i,j).$$
Let us now describe explicitly a subset $`B`$ of $`F_0`$ such that $`\{\pi (\overline{x}^f)|fB\}`$ is a basis for DIA. This subset will, of course, be a complete set of representatives for the orbits of all monomials $`\{\overline{x}^f|fF_0\}`$ under the action of $`S_n`$ described in Claim 5.2(iii).
Intuitively, we classify the exponents $`f_i(j)`$ according to their vector of residues modulo $`m`$, and attach a $`t`$-partite partition to each possible vector of residues. Formally, define
$$R:=\{(r_1,\mathrm{},r_t)𝒁^t|\mathrm{\hspace{0.17em}0}r_i<m(i)\text{ and }\underset{i=1}{\overset{t}{}}r_i0(\text{mod}m)\}$$
and choose an arbitrary linear order $`_R`$ on $`R`$.
Define a bijection
$$\theta :F_0F\times R^n$$
by
$$\theta (f):=(h,r),$$
where $`h=(h_i(j))_{i,j}F`$ and $`r=(r_i(j))_{i,j}R^n`$ are defined as the quotients and remainders, respectively, obtained when the entries of $`f=(f_i(j))_{i,j}F_0`$ are divided by $`m`$:
$$f_i(j)=mh_i(j)+r_i(j)(i,j).$$
Now let $`B`$ be the set of all $`fF_0`$ such that $`\theta (f)=(h,r)`$ satisfies:
* $`r_{}(1)_R\mathrm{}_Rr_{}(n)`$, where $`r_{}(j):=(r_1(j),\mathrm{},r_t(j))R`$ ($`j`$);
* If $`r_{}(j)=r_{}(j+1)`$, and also $`h_i(j)=h_i(j+1)`$ for all $`i<i_0`$, then
$$h_{i_0}(j)h_{i_0}(j+1).$$
We can interpret $`fB`$ as follows: For each $`r_{}R`$ there is a range (possibly empty) of indices $`j_1<jj_2`$ for which $`r_{}(j)=r_{}`$, and a sequence of vectors $`(h_{}(j))_{j=j_1+1}^{j_2}`$ which forms a $`t`$-partite partition with $`j_2j_1`$ parts (as defined in Subsection 5.1 above). The total number of vectors (for all $`r_{}R`$) is $`n`$.
Clearly, $`B`$ is a complete system of representatives for the orbits of $`F_0`$ under the action of $`S_n`$ defined in Claim 5.2(iii). By Claims 5.1 and 5.2 we conclude
Lemma 5.3. The set
$$\{\pi (\overline{x}^f)|fB\}$$
is a homogeneous basis for DIA.
Corollary 5.4. The generating function for DIA is
$$F_D(\overline{q})=\underset{fB}{}q_1^{|f_1|}\mathrm{}q_t^{|f_t|}=\underset{(n_r)}{}\left[\underset{rR}{}(q_1^{r_1}\mathrm{}q_t^{r_t})^{n_r}F_{t,n_r}(q_1^m,\mathrm{},q_t^m)\right],$$
where $`F_{t,n}(q_1,\mathrm{},q_t)`$ is the generating function for $`t`$-partite partitions described in Theorem GG1, and the sum is over all partitions of $`n`$ into non-negative integers $`(n_r)`$ indexed by $`rR`$.
Example. $`m=t=2`$.
* In this case $`R=\{(1,1),(0,0)\}`$, so that $`f=(f_1(j),f_2(j))_{1jn}`$ belongs to $`B`$ if and only if there exists an integer $`0kn`$ such that:
+ For $`1jk`$, $`f_1(j)`$ and $`f_2(j)`$ are both odd (i.e., $`r_{}(j)=(1,1)`$); and for $`k<jn`$ they are both even (i.e., $`r_{}(j)=(0,0)`$).
+ The vector pairs $`(h_1(j),h_2(j))_{1jk}`$ and $`(h_1(j),h_2(j))_{k<jn}`$ are both $`2`$-partite partitions.
### 5.4 Constructing Sequences of Generalized Permutations
In this subsection we construct a bijection between $`t`$-tuples of generalized permutations in $`G`$ with product equal to the identity element, and data consisting of numbers, sets and permutations defined below. This bijection will allow us to apply Theorem GG2 in the computation of the generating function, to be carried out in Subsection 5.5.
Every sequence $`(\pi _1,\mathrm{},\pi _t)G^t`$ such that $`\pi _t\mathrm{}\pi _1=1`$ may be constructed using the following steps:
Choose non-negative integers $`(n_r)_{rR}`$ such that
$$\underset{rR}{}n_r=n.$$
For each $`1it`$ choose a set-partition $`(N_r^{(i)})_{rR}`$ of $`N:=\{1,\mathrm{},n\}`$ such that
$$\mathrm{\#}N_r^{(i)}=n_r(rR),$$
where $`\mathrm{\#}S`$ denotes the size of the set $`S`$.
For each $`1it1`$ and $`rR`$ choose a bijective function (“permutation”) $`\widehat{\pi }_r^{(i)}:N_r^{(i)}N_r^{(i+1)}`$ (essentially, $`\widehat{\pi }_r^{(i)}S_{n_r}`$); and define $`\widehat{\pi }_r^{(t)}:N_r^{(t)}N_r^{(1)}`$ such that $`\widehat{\pi }_r^{(t)}\mathrm{}\widehat{\pi }_r^{(1)}=1`$.
For each $`1it`$ define a permutation $`\widehat{\pi }_iS_n`$ and a generalized permutation $`\pi _iG`$ as follows: if $`1jn`$ and $`jN_r^{(i)}`$ then
$$\widehat{\pi }_i(x_j):=\widehat{\pi }_r^{(i)}(x_j)$$
and
$$\pi _i(x_j):=\omega ^{r_i}\widehat{\pi }_r^{(i)}(x_j).$$
Conversely, let $`(\pi _1,\mathrm{},\pi _t)G^t`$ satisfy $`\pi _t\mathrm{}\pi _1=1`$. For each $`1jn`$ and $`1it`$ let
$$r_i(j):=\mathrm{Log}_\omega \pi _i(|\pi _{i1}\mathrm{}\pi _1(j)|).$$
Then
$$r_{}(j):=(r_1(j),\mathrm{},r_t(j))R,$$
since $`\pi _t\mathrm{}\pi _1=1`$.
Now define $`(n_r)_{rR}`$ and $`(N_r^{(i)})_{rR}`$ by:
$$N_r^{(1)}:=\{1jn|r_{}(j)=r\}(rR),$$
$$n_r:=\mathrm{\#}N_r^{(1)}(rR),$$
$$N_r^{(i)}:=\{|\pi _{i1}\mathrm{}\pi _1(j)|:jN_r^{(1)}\}(rR,2it),$$
and
$$\pi _r^{(i)}:=\pi _i|_{N_r^{(i)}}(rR,1it).$$
To sum up, there is a bijection
$$(\pi _1,\mathrm{},\pi _t)((n_r),(N_r^{(i)}),(\widehat{\pi }_r^{(i)}))$$
between $`t`$-tuples of generalized permutations in $`G`$ with product equal to the identity element, and data consisting of numbers, sets and permutations as above.
### 5.5 Computation of the Generating Function
In this subsection we prove the claim of Theorem 4.1, namely that
$$\frac{F_D(\overline{q})}{F_T(\overline{q})}=\underset{\genfrac{}{}{0pt}{}{\pi _1,\mathrm{},\pi _tG}{\pi _t\mathrm{}\pi _1=1}}{}\underset{i=1}{\overset{t}{}}q_i^{\mathrm{𝑓𝑙𝑎𝑔}\mathrm{𝑚𝑎𝑗𝑜𝑟}(\pi _i)}.$$
$`(5.1)`$
From Theorem 3.1 and from the bijection in the previous subsection it follows that the right-hand side of (5.1) is equal to
$$\text{RHS}=\underset{\genfrac{}{}{0pt}{}{\pi _1,\mathrm{},\pi _tG}{\pi _t\mathrm{}\pi _1=1}}{}\underset{i=1}{\overset{t}{}}q_i^{\mathrm{𝑓𝑙𝑎𝑔}\mathrm{𝑚𝑎𝑗𝑜𝑟}(\pi _i)}=$$
$`(5.2)`$
$$=\underset{(n_r)}{}\underset{(N_r^{(i)})}{}\underset{(\widehat{\pi }_r^{(i)})}{}\underset{i=1}{\overset{t}{}}q_i^{m\mathrm{𝑚𝑎𝑗𝑜𝑟}(\pi _i)+_{rR}n_rr_i}=$$
$$=\underset{(n_r)}{}\left[\underset{rR}{}(q_1^{r_1}\mathrm{}q_t^{r_t})^{n_r}\underset{(N_r^{(i)})}{}\underset{(\widehat{\pi }_r^{(i)})}{}\underset{i=1}{\overset{t}{}}q_i^{m\mathrm{𝑚𝑎𝑗𝑜𝑟}(\pi _i)}\right].$$
Now, $`\text{TIA}=(P_n^G)^t`$, and $`P_n^G`$ consists of all the symmetric polynomials in $`x_1^m,\mathrm{},x_n^m`$. Thus
$$F_T(\overline{q})=\frac{1}{_{i=1}^t_{j=1}^n(1q_i^{mj})},$$
$`(5.3)`$
Combining (5.3) with Corollary 5.4 (from the end of Subsection 5.3), we obtain that the left-hand side of (5.1) is equal to
$$\text{LHS}=\frac{F_D(\overline{q})}{F_T(\overline{q})}=\underset{i=1}{\overset{t}{}}\underset{j=1}{\overset{n}{}}(1q_i^{mj})\underset{(n_r)}{}\left[\underset{rR}{}(q_1^{r_1}\mathrm{}q_t^{r_t})^{n_r}F_{t,n_r}(q_1^m,\mathrm{},q_t^m)\right].$$
$`(5.4)`$
Comparing (5.2) with (5.4) we conclude that it suffices to show that, for every choice of $`(n_r)_{rR}`$,
$$\underset{i=1}{\overset{t}{}}\underset{j=1}{\overset{n}{}}(1q_i^{mj})\underset{rR}{}F_{t,n_r}(q_1^m,\mathrm{},q_t^m)=\underset{(N_r^{(i)})}{}\underset{(\widehat{\pi }_r^{(i)})}{}\underset{i=1}{\overset{t}{}}q_i^{m\mathrm{𝑚𝑎𝑗𝑜𝑟}(\pi _i)}.$$
$`(5.5)`$
Theorem GG1 above gives an explicit expression for $`F_{t,n_r}`$:
$$F_{t,n_r}(q_1^m,\mathrm{},q_t^m)=\left[\underset{i=1}{\overset{t}{}}\underset{j=1}{\overset{n_r}{}}(1q_i^{mj})\right]^1\underset{\widehat{\pi }_r^{(t)}\mathrm{}\widehat{\pi }_r^{(1)}=1}{}\underset{i=1}{\overset{t}{}}q_i^{m\mathrm{𝑚𝑎𝑗𝑜𝑟}(\widehat{\pi }_r^{(i)})},$$
where $`\widehat{\pi }_r^{(i)}S_{n_r}`$ are unsigned permutations.
Denoting $`q:=q_i^m`$, the definition of $`q`$-multinomial coefficients gives
$$\underset{j=1}{\overset{n}{}}(1q^j)\left[\underset{rR}{}\underset{j=1}{\overset{n_r}{}}(1q^j)\right]^1=\left[\genfrac{}{}{0pt}{}{n}{(n_r)_{rR}}\right]_q,$$
so that the left-hand side of (5.5) is equal to
$$\underset{i=1}{\overset{t}{}}\left[\genfrac{}{}{0pt}{}{n}{(n_r)_{rR}}\right]_{q_i^m}\underset{(\widehat{\pi }_r^{(i)})}{}\underset{rR}{}\underset{i=1}{\overset{t}{}}q_i^{m\mathrm{𝑚𝑎𝑗𝑜𝑟}(\widehat{\pi }_r^{(i)})}.$$
Here the sum is over all choices of $`\widehat{\pi }_r^{(i)}S_{n_r}`$ ($`rR`$, $`1it`$) such that $`\widehat{\pi }_r^{(t)}\mathrm{}\widehat{\pi }_r^{(1)}=1`$ ($`r`$).
Thus, all we need to prove is that, for every choice of $`(n_r)`$ and $`(\widehat{\pi }_r^{(i)})`$,
$$\underset{i=1}{\overset{t}{}}\left[\genfrac{}{}{0pt}{}{n}{(n_r)}\right]_{q_i^m}\underset{rR}{}\underset{i=1}{\overset{t}{}}q_i^{m\mathrm{𝑚𝑎𝑗𝑜𝑟}(\widehat{\pi }_r^{(i)})}=\underset{(N_r^{(i)})}{}\underset{i=1}{\overset{t}{}}q_i^{m\mathrm{𝑚𝑎𝑗𝑜𝑟}(\pi _i)}.$$
$`(5.6)`$
Since the choice of the partition $`(N_r^{(i)})_{rR}`$ of $`N`$ can be made independently for each value of $`i`$, we can “interchange” the sum and product in the right-hand side of (5.6), and thus it suffices to show that (again, denoting $`q:=q_i^m`$):
$$\left[\genfrac{}{}{0pt}{}{n}{(n_r)}\right]_q\underset{rR}{}q^{\mathrm{𝑚𝑎𝑗𝑜𝑟}(\widehat{\pi }_r^{(i)})}=\underset{(N_r^{(i)})_{rR}}{}q^{\mathrm{𝑚𝑎𝑗𝑜𝑟}(\pi _i)}(i,(n_r),(\widehat{\pi }_r^{(i)})).$$
$`(5.7)`$
Now this amounts exactly to the statement of Theorem GG2, since the different choices of $`(N_r^{(i)})_{rR}`$ correspond to the various shuffles of the permutations $`(\widehat{\pi }_r^{(i)})_{rR}`$, each yielding a different $`\pi _i`$.
$`\mathrm{}`$
If $`t=2`$, a somewhat simpler argument may be given. We demonstrate it in the special case $`m=t=2`$.
Example: $`m=t=2`$ .
In this case, by Corollary 5.4, Theorem GG1 and (5.4) (using simplified notation)
$$\frac{F_D(\overline{q})}{F_T(\overline{q})}=\underset{i=1}{\overset{2}{}}\underset{j=1}{\overset{n}{}}(1q_i^{2j})\{\underset{k=0}{\overset{n}{}}(q_1q_2)^k\frac{\underset{\pi _1S_k}{}q_1^{2\mathrm{𝑚𝑎𝑗𝑜𝑟}(\pi _1)}q_2^{2\mathrm{𝑚𝑎𝑗𝑜𝑟}(\pi _1^1)}}{\underset{i=1}{\overset{2}{}}\underset{j=1}{\overset{k}{}}(1q_i^{2j})}\}$$
$$\left\{\frac{\underset{\pi _1S_{nk}}{}q_1^{2\mathrm{𝑚𝑎𝑗𝑜𝑟}(\pi _1)}q_2^{2\mathrm{𝑚𝑎𝑗𝑜𝑟}(\pi _1^1)}}{\underset{i=1}{\overset{2}{}}\underset{j=1}{\overset{nk}{}}(1q_i^{2j})}\right\}=$$
$$=\underset{k=0}{\overset{n}{}}q_1^kq_2^k\underset{i=1}{\overset{2}{}}\left[\genfrac{}{}{0pt}{}{n}{k}\right]_{q_i^2}\mathrm{\Sigma }_{k,i}$$
where
$$\mathrm{\Sigma }_{k,i}=\left(\underset{\pi _1S_k}{}q_1^{2\mathrm{𝑚𝑎𝑗𝑜𝑟}(\pi _1)}q_2^{2\mathrm{𝑚𝑎𝑗𝑜𝑟}(\pi _1^1)}\right)\left(\underset{\pi _1S_{nk}}{}q_1^{2\mathrm{𝑚𝑎𝑗𝑜𝑟}(\pi _1)}q_2^{2\mathrm{𝑚𝑎𝑗𝑜𝑟}(\pi _1^1)}\right)=$$
$$\underset{\pi _1S_k\pi _2S_{nk}}{}q_1^{2(\mathrm{𝑚𝑎𝑗𝑜𝑟}(\pi _1)+\mathrm{𝑚𝑎𝑗𝑜𝑟}(\pi _2))}q_2^{2(\mathrm{𝑚𝑎𝑗𝑜𝑟}(\pi _1^1)+\mathrm{𝑚𝑎𝑗𝑜𝑟}(\pi _2^1))}.$$
Now, observe that every signed permutation in $`B_n=C_2S_n`$ may be constructed in the following way:
Choose a number $`0kn`$.
Choose $`k`$ digits from the set $`\{1,\mathrm{},n\}`$ and mark them “negative” (the rest will be “positive”).
Choose two permutations: $`\pi _1S_k`$ (on the “negative” digits), and $`\pi _2S_{nk}`$ (on the “positive” ones).
Choose a shuffle of $`\pi _1`$ and $`\pi _2`$.
It should be noted that the major index (in the sense defined at the beginning of Section 3) of the resulting signed permutation is independent of the choice in Step 2. In fact, it is exactly the major index of the shuffle (chosen in Step 4) of $`\pi _1`$ and $`\pi _2`$. On the other hand, the major index of the inverse of the resulting signed permutation does not depend on the choice in Step 4. In fact, it is exactly the major index of the shuffle (defined by Step 2) of $`\pi _1^1`$ and $`\pi _2^1`$.
Combining these facts with Theorem GG2, it follows that for any choosen integer $`0kn`$ and any given pair of permutations $`\pi _1S_k`$ and $`\pi _2S_{nk}`$
$$\underset{i=1}{\overset{2}{}}\left[\genfrac{}{}{0pt}{}{n}{k}\right]_{q_i^2}q_1^{2(\mathrm{𝑚𝑎𝑗𝑜𝑟}(\pi _1)+\mathrm{𝑚𝑎𝑗𝑜𝑟}(\pi _2))}q_2^{2(\mathrm{𝑚𝑎𝑗𝑜𝑟}(\pi _1^1)+\mathrm{𝑚𝑎𝑗𝑜𝑟}(\pi _2^1))}=$$
$$=\underset{\sigma B(\pi _1,\pi _2)}{}q_1^{2\mathrm{𝑚𝑎𝑗𝑜𝑟}(\sigma )}q_2^{2\mathrm{𝑚𝑎𝑗𝑜𝑟}(\sigma ^1)},$$
where $`B(\pi _1,\pi _2)`$ is the set of all signed permutations constructed by choosing $`\pi _1`$ and $`\pi _2`$ at the third step.
Thus,
$$\underset{i=1}{\overset{2}{}}\left[\genfrac{}{}{0pt}{}{n}{k}\right]_{q_i^2}\underset{\pi _1S_k\pi _2S_{nk}}{}q_1^{2(\mathrm{𝑚𝑎𝑗𝑜𝑟}(\pi _1)+\mathrm{𝑚𝑎𝑗𝑜𝑟}(\pi _2))}q_2^{2(\mathrm{𝑚𝑎𝑗𝑜𝑟}(\pi _1^1)+\mathrm{𝑚𝑎𝑗𝑜𝑟}(\pi _2^1))}=$$
$$=\underset{\sigma B_n(k)}{}q_1^{2\mathrm{𝑚𝑎𝑗𝑜𝑟}(\sigma )}q_2^{2\mathrm{𝑚𝑎𝑗𝑜𝑟}(\sigma ^1)},$$
where $`B_n(k):=\{\sigma B_n|\sigma \text{ has }k\text{ “negative” digits}\}`$. We conclude that
$$\frac{F_D(\overline{q})}{F_T(\overline{q})}=\underset{k=0}{\overset{n}{}}q_1^kq_2^k\underset{\sigma B_n(k)}{}q_1^{2\mathrm{𝑚𝑎𝑗𝑜𝑟}(\sigma )}q_2^{2\mathrm{𝑚𝑎𝑗𝑜𝑟}(\sigma ^1)}=$$
$$=\underset{\sigma B_n}{}q_1^{2\mathrm{𝑚𝑎𝑗𝑜𝑟}(\sigma )+k(\sigma )}q_2^{2\mathrm{𝑚𝑎𝑗𝑜𝑟}(\sigma ^1)+k(\sigma )},$$
where $`k(\sigma )`$ is the number of “negative” digits in $`\sigma `$. Note that $`k(\sigma )=k(\sigma ^1)`$. Theorem 3.1 completes the proof of the desired result (for $`m=t=2`$).
## 6 Final Remarks
* The flag major index may be defined on dihedral groups in an analogous way (with respect to the Coxeter generators). An exact analogue of Theorem 4.1 may be derived. This will be proved elsewhere.
* Using Theorem 3.1 it is possible to connect the flag major index with multiplicities of irreducible representations in homogeneous components of the coinvariant algebra of the group $`C_mS_n`$. These multiplicities were calculated by Stembridge \[Stem\] and involve major indices of skew standard Young tableaux. For more details see \[AR, Section 5\].
Acknowledgments. Thanks to Richard Stanley, Dominique Foata, Christian Krattenthaler, Igor Pak and Victor Reiner for their comments. |
warning/0002/astro-ph0002476.html | ar5iv | text | # Some statistical properties of spiral galaxies
## 1 Introduction
The inclination of a spiral galaxy ( i.e., the angle between the galactic plane and the tangent plane of the celestial sphere) is not only an important parameter, but also difficult to determine. A spiral galaxy consists of a thin disk, a bulge and spiral arms that are thought to be situated in the disk. If we assume that the thickness of the spiral plane is rather negligible in comparision to its extension, and that when a spiral galaxy is inclined moderately to the plane of sky, the thickness of the nucleus can be omitted, the inclination $`(\gamma )`$ can be obtained by:
$$\gamma =\mathrm{arccos}(\frac{d}{D}),$$
(1)
where $`D`$ and $`d`$ are the apparent major and minor isophotal diameters respectively. When a spiral galaxy is seen edge-on, it is not possible to consider the thickness of the nuclear part as negligible. Thus, Eq. (1) cannot be used to calculate the inclination. The reason for this is that the apparent minor isophotal-diameter consists of two parts. One is attributed by the disk and, another by the bulge, the latter of which will decrease the real value of the inclination.
Considering that the disk is not infinitely thin, Aaronson et al. (1980) corrected Eq. (1) by
$$\gamma =\mathrm{arccos}\sqrt{1.042(\frac{d}{D})^20.042}+3^{}.$$
(2)
The constant of $`3^{}`$ is added in accordance with an empirical recipe. A more elaborate specification of the axial ratio for an edge-on system that depends on the Hubble type could be justified. The thinnest galaxies are Sc spirals’, earlier types have larger bulges. Giovanelli et al. (1997) provided an example to justify why they assumed the axial ratio of Sc galaxies to be 0.13. A smaller value of the axial ratio for an edge-on system results in smaller derived inclinations, where the spirals are more face-on. Besides, if the values $`D`$ and $`d`$ are approximations due to errors, the inclination obtained by Eq. (2) is not an exact value. Ma et al. (1997, 1998) proposed a method to determine the inclination of a spiral galaxy by fitting a spiral arm with a logarithmic spiral form with constant pitch angle. They obtained the inclinations of 72 northern spiral galaxies.
The question of the mathematical form of spiral arms was recognized at the beginning of this century (von der Pahlen, 1911; Groot, 1925). Then, Danver (1942), Kennicutt (1981) and Kennicutt & Hodge (1982) systematically studied the shapes of spiral arms. Using the method of the least squares and as many points as possible situated on the spiral arm in question, Danver (1942) studied a sample of 98 nearby spirals by drawing the projected images on white paper and then, by copying it on the paper to be used for the measurement thanks to transparent illumination. Kennicutt (1981) measured the shapes of spiral arms in 113 nearby Sa-Sc galaxies by disposing directly of photographic enlargement and using an iterative procedure to correct for inclination effects. He gave an initial estimate of the inclination and pitch angle to orient the spiral to a face-on geometry, and then used any residual sinusoidal deviations in the arm shapes to make small corrections to the derived orientation. Using the IRAF software, Ma et al. (1997, 1998) fitted the shapes of spiral arms on the images, so that they could show clearly whether the fitting was good or not. The DISPLAY program of IRAF software can enlarge the image and change its grey scale to minimize any personal prejudice about the regularity and prominence of arms. But we must emphasize that the DISPLAY program of IRAF has many variables, so the results are not always objective. In our program, we modify z1 (minimum greylevel to be displayed) and z2 (maximum greylevel to be displayed) in the DISPLAY program in order to display the images clearly. In the procedure of fitting, we emphasize the global spiral structure, where, except for the small-scale distortions, the arms can be represented by the logarithmic spiral forms.
There has been much interest concerning the separation of disk and bulge components in the observed surface brightness distribution of spiral galaxies. de Vaucouleurs (1958), for instance, established an isophotic map of M 31 in blue light by means of direct photoelectric scans, spaced at 10 intervals in declination from +3931 to 4230. From photoelectric photometry, he determined that the thickness of the flat component is about 0.8 kpc. By assuming that a galaxy has an infinitesimally thin disk, Freeman (1970) and Sandage et al. (1970) collected and studied the radial distribution of the surface brightness $`I(r)`$ for thirty-six S0 and spiral galaxies, and showed that $`I(r)`$ distribution for these galaxies can be presented by two main components: an inner spheroidal component which follows the law of
$$\mathrm{log}I(r)r^{1/4}$$
(3)
and an outer exponential component (disk), with
$$\mathrm{log}I(r)=I_0e^{r/h_r},$$
(4)
where $`h_r`$ is defined as a radial scale length. Van der Kruit and Searle (1981a) proposed a model for the light distribution in the disks of edge-on spiral galaxies, assuming that a galaxy has a locally isothermal, self-gravitating and truncated exponential disk. This model has the feature of being isothermal in $`z`$ at all radii with a scale parameter $`z_0`$ and has an exponential dependence of surface brightness upon $`r`$ with a scale length $`h_r`$. The space-luminosity of this model can be described by
$$L(r,z)=L_0e^{r/h_r}\mathrm{sech}^2(z/z_0).$$
(5)
With this model, van der Kruit & Searle (1981a, 1981b, 1982a, 1982b) determined $`h_r`$ and $`z_0`$ for seven disk-dominated and one spheroid-dominated spiral galaxies by using three-color surface photometry. Peng et al (1979) investigated three-dimensional disk galaxies, based on the fundamental assumption by Parenago that the density distribution along $`z`$-direction for a finite thickness disk is
$$\rho (r,\varphi ,z)=\frac{1}{H_z}\sigma (r,\varphi )e^{|z|/h},$$
(6)
where $`h`$ is defined as an exponential scale height, $`H_z`$ is defined as a thickness of disk and equals $`2h`$, and $`\sigma (r,\varphi )`$ is the surface density. By solving Poission’s equation for a logarithmic density perturbation, Peng et al. (1979) obtained a criterion for density waves to appear, which is
$$r>r_0=\frac{H_z\sqrt{m^2+\mathrm{\Lambda }^2}}{2},$$
(7)
where $`(r_0,\varphi _0)`$ is the polar coordinate of the starting point from which arms of a galaxy stretch outward on the galactic plane, and $`m`$ is the number of the arms in a spiral galaxy. Based on this criterion, Peng (1988) proposed a method for estimating the thickness of a non-edge-on spiral, and derived the thicknesses of four galaxies.
Guthrie (1992) derived the axial ratios $`R`$ of disc components for 262 edge-on spiral galaxies on print copies of the blue Palomar Sky Survey plates by using a microscope fitted with a micrometer eyepiece. He then analyzed the distribution of isophotal axial ratios for 888 diameter-limited normal Sa-Sc galaxies to give information on the true axial ratios $`R_0`$, and at last presented the mean value of $`\mathrm{log}R_0`$ is $`0.95\pm 0.03`$.
Ma et al. (1997, 1998) derived the thicknesses of 72 spirals by using Peng’s proposal (Peng, 1988) and presented some statistical correlations between thickness or flatness and other parameters.
The value of $`R_0`$ derived by Guthrie (1992) is based on observational work and, should be much more reliable. So, it is important to compare Ma et al.’s results (1997, 1998) to Guthrie’s (1992). In Ma et al. (1997, 1998), the flatnesses for 72 galactic disks ($`H_z/D_0`$) <sup>1</sup><sup>1</sup>1$`D_0`$, which is measured at or reduced down to the surface brightness level $`\mu _B=25.0B`$ magnitudes per square arcsecond, and corrected to the “face-on” ($`\gamma =0^{}`$). For the Galactic extinction, but not for redshift, $`D_0`$ is from the Third Catalogue of Bright Galaxies by de Vaucouleurs et al. (hereafter RC3). were given, and the mean value is $`0.033\pm 0.002`$. Suppose that the value of ratio of radial scale length ($`h_r`$) over exponential scale height ($`h`$) for an average exponential galactic disk is equal to $`R_0`$ from Guthrie (1992), which is 9. From Freeman (1970), we can derive $`D_0/(2h_r)5`$. Thus, we obtain $`H_z/D_00.023`$, which is in relative agreement with the mean value of Ma et al. (1997, 1998). At the same time, we calculate the values of $`\overline{\mathrm{log}R_0}`$ for spirals of various types T and list them in Table 1. T, from the RC3, is morphological types, and $`\sigma `$ is the dispersion. From this table and Table 3 of Guthrie (1992), we can see that our results are in agreement with Guthrie’s (1992). Except for the data concerning Scd galaxies, our results also agree with de Grijs’ (1998).
The structure of this paper is as following: in Sect. 2, we outline the principles of obtaining an inclination and a pitch angle; Sect. 3 presents some statistical properties; and conclusions will be shown in Sect. 4.
## 2 Principles of obtaining an inclination and a pitch angle
As we know, when the line of intersection (namely, the major axis of the image) between the galactic plane and tangent plane is taken as the polar axis, it is easily proved that
$$r=\rho \sqrt{1+\mathrm{tan}^2\gamma \mathrm{sin}^2\theta }$$
(8)
and
$$\mathrm{tan}\varphi =\frac{\mathrm{tan}\theta }{\mathrm{cos}\gamma },$$
(9)
where $`r`$ and $`\varphi `$ are the polar co-ordinates in the spiral plane and $`\rho `$, $`\theta `$ are the corresponding co-ordinates in the tangent-plane, and $`\gamma `$ is the inclination. If it is possible to represent arms by equiangular spirals as
$$r=r_0e^{\lambda (\varphi \varphi _0)}$$
(10)
and
$$\mu =\mathrm{arctan}\lambda ,$$
(11)
where $`r`$ and $`\varphi `$ are the polar co-ordinates on the spiral arm in the spiral plane, and $`\mu `$, defined as a pitch angle, is the angle between the tangent to the arm and the circle with the constant $`r`$. The mathematical form of Eq. (10) in the tangent plane of the celestial sphere is
$$\rho (\theta ,\gamma )=\rho _0\frac{f(\theta _0,\gamma )}{f(\theta ,\gamma )}e^{\lambda B(\theta ,\gamma )},$$
(12)
here
$$f(\theta ,\gamma )=\sqrt{\mathrm{sin}^2\theta +\mathrm{cos}^2\theta \mathrm{cos}^2\gamma }$$
(13)
and
$$B(\theta ,\gamma )=g(\theta ,\gamma )g(\theta _0,\gamma ),$$
(14)
where
$`g(\theta ,\gamma )=`$
$$\{\begin{array}{cc}\mathrm{arctan}(\mathrm{tan}\theta /\mathrm{cos}\gamma ),\hfill & k\pi \frac{\pi }{2}<\theta <k\pi +\frac{\pi }{2}\hfill \\ \pi +\mathrm{arctan}(\mathrm{tan}\theta /\mathrm{cos}\gamma ),\hfill & k\pi +\frac{\pi }{2}<\theta <k\pi +\frac{3\pi }{2},\hfill \end{array}$$
(15)
where $`k`$ is an integer. Supposing that ($`\rho _i`$, $`\theta _i`$) are the measured co-ordinates of points of the arms, and making use of the least squares method, we set
$$\underset{i=1}{\overset{n}{}}[\rho _i\rho (\theta _i,\gamma )]=minimum.$$
(16)
By direct differentiation with respect to $`\lambda `$, we obtain the equation for determining the parameter $`\lambda `$ of
$$\underset{i=1}{\overset{n}{}}B(\theta _i,\gamma )\rho (\theta _i,\gamma )[\rho _i\rho (\theta _i,\gamma )]=0.$$
(17)
This is a transcendental equation. Using the least squares method, we can obtain $`\gamma `$ and $`\mu `$. Details have been shown in Ma et al. (1997, 1998).
## 3 Statistical properties
### 3.1 Sample
Our statistical sample contains 72 northern spiral galaxies for which the values of disk thickness and inclination and pitch angle of individual spiral arms are from Ma et al. (1997, 1998). Except for M 31, the images, which have distinguishable spiral arms, are from the Digitized Palomar Sky Survey. The image of M 31 is digitized from POSSII (field No. 295) via PDS measurement at Purple Mountain Observatory of China. Most galaxies in this sample have the grand design structure. All the galaxies have the total color indexes ($`(BV)_T^0`$), which are corrected for the differential Galactic and internal extinction to “face-on”, and for redshift between $`B`$ and $`V`$ bands. The mean numerical Hubble stage indexes (T) of the galaxies, which are quoted in the RC3, are from 2 to 6, and 2, 3, 4, 5 and 6 correspond to Sab, Sb, Sbc, Sc and Scd respectively. This is an inclination-limited sample, i.e., the values of $`\mathrm{log}(D_{25}/d_{25}`$) for these galaxies are smaller than the limited-value (0.76), where $`D_{25}`$ and $`d_{25}`$, taken from the RC3, are the apparent major and minor isophotal diameters measured at or reduced down to the surface brightness level $`\mu _B=25.0B`$ magnitudes per square arcsecond. Clearly, this is not a diameter-limited sample.
### 3.2 Thickness vs absolute magnitude
Ma et al. (1998) have presented some statistical correlation between thickness or flatness and other parameters. In this section, we will investigate some other statistical correlations.
Fig. 1 plots the correlation between the thickness and the total (“ face-on”) absolute magnitude in the B band. The latter is derived by using the formula
$$M_B^o=B_T^o+55\mathrm{log}d,$$
(18)
here $`B_T^o`$ is the total “face-on” magnitude from the RC3, corrected for the Galactic and internal extinction, as well as for redshift, d the distance, in units of pc, from Ma et al. (1997, 1998). One can find that there is a significant correlation between thickness and B-band absolute magnitude, i.e., thicker galaxies are brighter.
Based on the careful studies of Binggeli et al. (1985) about the Virgo cluster, Binggeli et al. (1988) have thoroughly reviewed what is currently known about the luminosity function $`\mathrm{\Phi }(L)`$. They have derived the luminosity function $`\mathrm{\Phi }(L,T)`$ for each morphological type separately in the Virgo cluster. Furthermore, according to the detailed analysis of the member galaxies for the Virgo cluster (Zhao & Shao, 1994; Shao & Zhao, 1996), Shu et al. (1995) also investigated its LF for individual morphological types. Combining with the luminosity function $`\mathrm{\Phi }(L,T)`$ of the Virgo cluster (Binggeli et al., 1988; Shu, et al., 1995), Fig. 1 automatically implies that early-type spirals, which are on average brighter than late-type ones, are also thicker.
### 3.3 Dependence of $`H_z/H_r`$ on morphological type, luminosity and color
In this section, we will present some statistical correlations between $`H_z/H_r`$ and morphological type, luminosity and color for spiral galaxies.
Ma et al. (1997, 1998) derived the flatnesses ($`H_z/D_0`$) for 72 spiral galaxies. When we take $`D_0/H_r5`$ ($`H_r`$ is defined as $`2h_r`$), we can derive the values of $`H_z/H_r`$ for these spirals. Fig. 2 shows $`H_z/H_r`$ as a function of different Hubble types. It can be seen that $`H_z/H_r`$ of spiral galaxies decreases smoothly an average along the Hubble type, but the dispersion in $`H_z/H_r`$ among galaxies of the same Hubble sequence is very large. This, perhaps, reflects that the intrinsic flattening of spiral disks is smaller for later type galaxies. de Grijs (1998) also shows that galaxies become systematically thinner when going from type S0 to Sc.
Fig. 3 plots $`H_z/H_r`$ as a function of total (“face-on”) absolute magnitude in the B band. In Fig. 3, we can see a clear trend reflecting an increase in the values of $`H_z/H_r`$ as the galaxies become brighter for all but a few galaxies. Combining this data with the luminosity function $`\mathrm{\Phi }(L,T)`$ of the Virgo cluster (Binggeli et al., 1988; Shu et al., 1995), we find that early-type spirals, which are brighter on average than late-type ones, have larger values of $`H_z/H_r`$.
Fig. 4 plots $`H_z/H_r`$ as a function of total galactic color index ($`(BV)_T^0`$) in the RC3 and shows a strong trend which suggests that a bluer galaxy has a smaller value of $`H_z/H_r`$. Roberts & Haynes (1994) studied the physical parameters along the Hubble sequence systematically by making use of two primary catalogues, the RC3 and a private catalogue maintained by R. Giovanclli & M. Haynes. He presented the well-established trend between morphology and mean color which reveals that the E and S0 galaxies are clearly redder than spirals and that late-type spirals are bluer on average than early-type ones. So, Fig. 4 also implies that later-type galaxies have smaller values of $`H_z/H_r`$.
### 3.4 Statistical property of inclination
Tully & Fisher (1977) showed that if a spiral structure is well defined, the opening of the spiral structure could be used to define the inclination of the disk. It has been shown by Danver (1942), Kennicutt (1981) and Kennicutt & Hodge (1982) that the spiral arm can be represented by a logarithmic spiral form with a constant pitch angle. Grosb$`\varphi `$l (1985) studied 605 galaxies with inclination angles inferior to $`56^{}`$, and estimated the position and inclination angles of these galaxies based on the bisymmetric intensity distribution in their outer parts applying a one dimensional Fourier transform method. Ma et al. (1997, 1998) obtained the inclinations of 72 northern spiral galaxies by fitting the arms on the images. The RC3 lists the values of $`D_{25}/d_{25}`$ and their errors, $`D_{25}`$ and $`d_{25}`$ are the apparent major and minor isophotal diameters measured at or reduced down to the surface brightness level $`\mu _B=25.0`$ B magnitudes per square arcsecond. We calculate the inclinations by use of Eq. (2). Fig. 5 presents the correlation between them and those obtained by Ma et al. (1997, 1998). This figure shows that, when we take into account the errors in the inclination due to errors in $`D_{25}/d_{25}`$, and although there is some scatter, Ma et al.’s result (1997, 1998) is consistent with what is derived from Eq. (2).
### 3.5 Mean value of the pitch angle along the Hubble sequences in the RC3
Hubble (1926, 1936) introduced an early scheme to classify galaxies. Its concepts are still in use, a sequence starting from elliptical to spiral galaxies, and including lenticular. This scheme has been extended by some astronomers (Holmberg, 1958; de Vaucouleurs, 1956, 1959; Morgan, 1958, 1959; van den Bergh, 1960a, b, 1976; Sandage, 1961; Sandage & Tammann, 1981, 1987; Sandage & Bedke, 1993) over the years, who tried to employ multiple classification criteria. The two main systems commonly used are derived from Hubble’s original classification criteria. One is the Hubble system as explained in detail by Sandage (1961), Sandage & Tammann (1981, 1987) and Sandage & Bedke (1993), and another, developed by de Vaucouleurs (1956, 1959), adds more detailed descriptions to the notation and makes a division and extension of the Sc and SBc families by introducing the Scd, Sd, Sdm, Sm and Im subdivisions.
Galaxy morphological classification is still mainly done visually by dedicated astronomers, based on the Hubble’s original scheme and its modification. It is possible for each observer to give slightly different weights to various criteria, although the criteria for classification are accepted generally by them. Lahav et al. (1995) and Naim et al. (1995) investigated the consistency of visual morphological classifications of galaxies from six independent observers. They found that individual observers agree with one another with combined rms dispersions of between 1.3 and 2.3 type units, typically about 1.8 units of the revised Hubble numerical type index T, although there are some cases where the range of types given to it was superior to 4 units in width.
The tightness of the spiral arm, in addition to the degree of resolution in the arm and the relative size of the unresolved nuclear region, are the fundamental criteria in Hubble’s (1926, 1936) classification of spiral galaxies. In order to better understand the nature and origin of the Hubble sequence and evaluate the difference of classification, Kennicutt (1981) compared the arm pitch angles as measured by him with the Hubble type as determined by Sandage & Tammann (1981, hereafter ST) and Yerkes class by Morgan (1958, 1959). From Figs. 7 and 8 in Kennicutt (1981), we can see that, although the ST’s classification is based almost solely upon disk resolution and the Morgan’s is based entirely on the central concentration of the galaxies, the trends between the arm pitch angles and the two types are almost the same. In the classifications by de Vaucouleurs et al. (1976, RC2; 1991, RC3), all these three criteria are considered.
A most important feature of spiral galaxies is that the stellar content between the spiral arms and the spheroidal component is very different, and the central part is redder than the arm region. The Yerkes system by Morgan (1958, 1959) contains the information on the state of stellar evolution in the central region of galaxies and it complements the information in the Hubble sequence by de Vaucouleurs (1956, 1959), which, in cases of conflicting criteria, emphasizes more strongly the strength of population in the arm. There are other sets of classifications in which disk-to-bulge is the primary discriminant (van den Bergh, 1976) or the ratio of disk-to-bulge to spiral arm morphology is primarily considered (Dressler, 1980; Wilkerson, 1980). Among these three criteria, perhaps, spiral arm morphology may be distance independent.
Until now, there are only a few papers dealing with pitch angles. So, the mean value of pitch angle for the different Hubble types, which is an important parameter in the decision of whether the WKB approximation can be satisfied (Binney & Tremaine, 1987), has not been presented. As we know, the WKB approximation is an indispensable tool for understanding the origin and evolution of a density wave in spiral galaxies.
Before deriving the mean value of pitch angle for the different types in the RC3, it is useful to compare the pitch angles of the common galaxies derived by Kennicutt (1981) and Ma et al. (1998). 22 of the spirals measured by Ma et al. (1998) were previously done by Kennicutt (1981), and the average pitch angles are compared in Fig. 6. Except for 5 galaxies, the results are consistent between Kennicutt (1981) and Ma et al. (1998). Some of the discrepancies can be attributed to our measured error, and the others may be expected from the galaxies themselves; for example, when a spiral galaxy has two asymmetric arms: different authors (Kennicutt, 1981; Ma et al. 1998) did not measure the same arm in the same galaxy.
Table 2 presents the mean values of pitch angle ($`\mu `$) from the ones measured by Kennicutt (1981) and Ma et al. (1997, 1998) for the different Hubble sequences in the RC3, where $`\sigma `$ is the dispersion and N the number of galaxies. The important result from Table 2 is that the mean pitch angles of Sa-Sc are not larger than $`15.5^{}`$, so that cotan$`\mu 3.6`$. Thus the WKB approximation is satisfied at the mean pitch angles from Sa-Sc, but not by very much.
Fig. 7 presents the mean values versus the Hubble types in the RC3, and it shows that, from the early to late Hubble types, the mean values of the pitch angle increase, although the value for S0/a’s is larger than that for Sa’s. In fact, it is difficult to distinguish between S0/a and Sa.
Besides Hubble classification systems, Elmegreen & Elmegreen (1982a, b; 1987) introduced another classification system which is designed to emphasize arm continuity, length and symmetry, and it is related to the presence or strength of density waves. This system is not dependent on the galaxy Hubble sequence and contains 12 distinct arm classes, corresponding to a systematic change from the ragged and patchy arms in ‘flocculent’ galaxies to the two symmetric and continuous arms in ‘grand design’ ones. Intermediate arm classes show characteristics of both the ‘flocculent’ and ‘grand design’ types. Based on their spiral arm classification system, Elmegreen & Elmegreen (1982a, b; 1987) classified the spiral galaxies in the field, in binary systems, in groups and in clusters, and suggested that bars tend to correlate with spiral density waves, companions may influence (or generate) symmetric density waves, grand design galaxies are physically larger than flocculent ones by a factor of $``$ 1.5, and grand design galaxies are also preferentially in dense groups.
## 4 Conclusions
In this paper, we investigate some statistical correlations about the thickness of galactic disks and pitch angles of spiral arms. The main conclusions are: (1). Early-type spirals, that are brighter on the average, are thicker; (2). The axis ratio ($`H_z/H_r`$, here $`H_r`$ and $`H_z`$ are defined as two times the radial scale length ($`h_r`$) and exponential scale height ($`h`$)) of galactic disk tends to be smaller along the Hubble sequence; (3). Except for a few galaxies, early-type spiral galaxies have larger values of $`H_z/H_r`$; (4). $`H_z/H_r`$ correlates strongly with galaxy color; (5). The inclinations obtained by fitting the pattern of spiral structure with a logarithmic spiral form are nearly the same as those obtained by using the formulas of Aaronson et al. (1980); (6). The mean measured pitch angles for different Hubble sequences in the RC3 are derived, and the results show that the mean pitch angles of Sa-Sc’s are not larger than $`15.5^{}`$, so that cotan$`\mu 3.6`$. Thus the WKB approximation can be satisfied at the mean pitch angles from Sa-Sc’s, but not by very much; (7). From early to late Hubble types, the mean value of pitch angle increases, despite some scattering.
Although the method, proposed by Peng (1988) for deriving the thickness of a face-on disc, is effective and simple, it relies on a spiral structure theory that predicts that a spiral arm has to end somewhere in the disk. However, this theory might not be completely right. For example, Zaritsky et al. (1993) has presented K-band (2.2 $`\mu m`$) images of M 51, which reveal remarkable dynamical structure not visible in the conventional optical observations, and show that the spiral arms extend significantly further towards the galaxy’s center than previously observed. In the optical images, the spiral arms begins at a radius of about $`30^{^{\prime \prime }}`$ from the center. But, the K-band residual image showed the spiral arms wind through an additional $`540^{}`$ beyond that seen in the B-band images of the entire galaxy, and end at about $`10^{^{\prime \prime }}`$ from the center of the galaxy. If spiral galaxies whose spiral arms do not stop anywhere in the disk exist, the parameters $`\rho _0`$ and $`r_0`$ (Peng, 1988; Ma et al., 1997, 1998) cannot be found. In the K-band images, which minimize the effect of dust and maximize sensitivity to the dominant stellar population, we can derive our reliable values of $`\rho _0`$ and $`r_0`$. As an example, we apply our method to the image of M 51 in the K-band<sup>2</sup><sup>2</sup>2The image of M 51 in K-band is provided by Prof. Zaritsky. and derive the flatness of disk ($`H_z/D_0`$), which is $`0.010\pm 26.4\%`$. Comparing this value with Peng’s (1988) ($`H_z/D_0=0.013\pm 16.4\%`$), we can find that the value of the flatness based on the K-band image is smaller. The reason is that, in the K-band image, the effect of dust can be minimized and the values of $`\rho _0`$ and $`r_0`$ may be reliably derived. We also emphasize that the images, which Ma et al. (1998) used, are from the Digitized Palomar Sky Survey, in which many images have burnt-out centers. There are some pictures that have burnt-out centers in Ma et al. (1998), but we can change the parameters of DISPLAY program to minimize the effect.
###### Acknowledgements.
We are indebted to the anonymous referee for many critical comments and helpful suggestions that have greatly improved our paper, and for correcting our English. We are grateful to Prof. Zaritsky for providing us the image of M 51 in the K-band. J. Ma gratefully acknowledges the support of K. C. Wong Education Foundation, Hong Kong. This work is supported partly by the Chinese National Science Foundation, No. 19603003. |
warning/0002/cond-mat0002010.html | ar5iv | text | # Enhancement of Kondo Effect in Quantum Dots with an Even Number of Electrons
## Abstract
We investigate the Kondo effect in a quantum dot with almost degenerate spin-singlet and triplet states for an even number of electrons. We show that the Kondo temperature as a function of the energy difference between the states $`\mathrm{\Delta }`$ reaches its maximum around $`\mathrm{\Delta }=0`$ and decreases with increasing $`\mathrm{\Delta }`$. The Kondo effect is thus enhanced by competition between singlet and triplet states. Our results explain recent experimental findings. We evaluate the linear conductance in the perturbative regime.
The Kondo effect takes place when a localized spin $`S`$ is brought in contact with electron Fermi sea. The Kondo effect gives rise to a new many-body ground state that has a lesser spin. Recently the Kondo effect has been observed in semiconductor quantum dots connected to external leads by tunnel junctions . In this case the localized spin is formed by electrons in the dot. The number of electrons $`N`$ is fixed by Coulomb blockade to integer values and can be tuned by gate voltage. Despite the Coulomb blockade, the ground state of the dot is usually similar to what one obtains disregarding the interaction. The discrete spin-degenerate levels in the dot are consecutively occupied, and the total spin is zero or 1/2 for an even and odd number of electrons, respectively. Then the Kondo effect takes place only in the latter case .
Significant deviations from this plain picture were recently observed in so-called “vertical” quantum dots . The strength of the electron-electron Coulomb interaction in such dots is comparable with the spacing of discrete levels, and this may give rise to a complicated ground state. For example, if two electrons are put into nearly degenerate levels, the exchange interaction favors a spin-triplet state. This state can be changed to a spin singlet by applying a magnetic field since the magnetic field increases the level spacing .
This gives a unique possibility to change the spin of the ground state during an experiment and even obtain extra degenerate states by tuning the energies of different spin configurations to the same value. Such a possibility hardly exists in the traditional solid state context. The Kondo effect in multilevel quantum dots has been investigated by several groups . In this letter we examine a novel effect that stems from the competition between spin-singlet and threefold spin-triplet states for an even number of electrons in a dot. Since the energy difference between the states $`\mathrm{\Delta }`$ can be controlled experimentally, we elucidate $`\mathrm{\Delta }`$ dependence of Kondo temperature as a typical energy scale for the Kondo effect and linear conductance through the dot. This enables direct comparison between our calculations and recent experimental results .
At large positive $`\mathrm{\Delta }`$ the dot is in a triplet ground state and an extra singlet state can be disregarded. The Kondo effect follows a usual $`S=1`$ scenario. At large negative $`\mathrm{\Delta }`$ the ground state of the dot is a spin singlet and the Kondo effect ceases to exist. From this one could suggest that the Kondo temperature decreases as $`\mathrm{\Delta }0`$ from the positive side. Our results show just the opposite. The Kondo temperature $`T_K(\mathrm{\Delta })`$ is enhanced at small $`\mathrm{\Delta }`$ and reaches its maximum at $`\mathrm{\Delta }T_K^{max}`$. At $`\mathrm{\Delta }T_K^{max}`$ the Kondo temperature decreases with increasing $`\mathrm{\Delta }`$ obeying a power law $`T_K(\mathrm{\Delta })1/\mathrm{\Delta }^\gamma `$. The exponent $`\gamma `$ is not universal but depends on model parameters. Our results clearly demonstrate the importance of one of the basic principles of Kondo physics: although the Kondo effect occurs at small energy scale $`T_K`$, the value of this scale is determined by all energies from $`T_K`$ up to the upper cutoff. In our case, the energies from $`\mathrm{\Delta }`$ to the upper cutoff would feel 4-fold degeneracy of the dot states, which enhances the Kondo temperature. At $`\mathrm{\Delta }<\mathrm{\Delta }_c`$ ($`|\mathrm{\Delta }_c|T_K^{max}`$), the Kondo effect is not relevant. We stress the difference between this mechanism and one proposed in Ref. , where the Kondo effect arises from extra degeneracy between one component of the spin-triplet state and a singlet state, which is brought by the Zeeman splitting.
To model the situation, it is sufficient to consider two extra electrons in a quantum dot at the background of a singlet state of all other $`N2`$ electrons, which we will regard as the vacuum. These two extra electrons occupy two levels of different orbital symmetry . The energies of the levels are $`\epsilon _1,\epsilon _2`$. Possible two-electron states are (i) the threefold spin-triplet state, (ii) the spin-singlet state of the same orbital symmetry as the triplet state, $`1/\sqrt{2}(d_1^{}d_2^{}d_1^{}d_2^{})|0`$, and (iii) two singlets of different orbital symmetry, $`d_1^{}d_1^{}|0`$, $`d_2^{}d_2^{}|0`$. Among the singlet states, we only consider a state of the lowest energy, which belongs to the group (iii). Thus we restrict our attention to four states, $`|SM`$:
$`|11`$ $`=`$ $`d_1^{}d_2^{}|0`$ (1)
$`|10`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}(d_1^{}d_2^{}+d_1^{}d_2^{})|0`$ (2)
$`|11`$ $`=`$ $`d_1^{}d_2^{}|0`$ (3)
$`|00`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}(C_1d_1^{}d_1^{}C_2d_2^{}d_2^{})|0,`$ (4)
where $`d_{i\sigma }^{}`$ creates an electron with spin $`\sigma `$ in level $`i`$. Their energies, $`E_{S=1}`$, $`E_{S=0}`$, and coefficients in the singlet state, $`C_1`$, $`C_2`$ ($`|C_1|^2+|C_2|^2=2`$), are determined by the electron-electron interaction and one-electron level spacing $`\delta =\epsilon _2\epsilon _1`$. At the moment, we set $`C_1=C_2=1`$ and show that this is the general case afterwards. The energy difference, $`\mathrm{\Delta }=E_{S=0}E_{S=1}`$, is changed by applying a magnetic field $`B`$ as shown in Fig. 1(a). We disregard the Zeeman splitting of spin states since this is a much smaller energy scale than the orbital effect of the magnetic field in semiconductor heterostructures in use . The exact condition for this is $`E_{\mathrm{Zeeman}}T_K`$.
The dot is connected to two external leads $`L`$, $`R`$ with free electrons being described by
$$H_{\mathrm{leads}}=\underset{\alpha =L,R}{}\underset{k\sigma i}{}\epsilon _{\alpha k}^{(i)}c_{\alpha ,k\sigma }^{(i)}c_{\alpha ,k\sigma }^{(i)}.$$
(5)
The tunneling between the dot and the leads is written as
$$H_T=\underset{\alpha =L,R}{}\underset{k\sigma i}{}(V_{\alpha ,i}c_{\alpha ,k\sigma }^{(i)}d_{i\sigma }+\mathrm{H}.\mathrm{c}.).$$
(6)
Here $`c_{\alpha ,k\sigma }^{(i)}`$ is the creation operator of an electron in lead $`\alpha `$ with momentum $`k`$, spin $`\sigma `$, and orbital symmetry $`i`$ $`(=1,2)`$. We assume that the orbital symmetry is conserved in the tunneling processes. Therefore we have two electron “channels” in each lead.
We assume that the state of the dot with $`N`$ electrons is stable, so that addition/extraction energies, $`E^\pm E(N\pm 1)E(N)\mu `$ where $`\mu `$ is the Fermi energy in the leads, are positive. We are interested in the case when $`E^\pm |\mathrm{\Delta }|`$, $`\delta `$ and also exceed the level broadening $`\mathrm{\Gamma }_\alpha ^i=\pi \nu |V_{\alpha ,i}|^2`$ ($`\nu `$ being density of states in the leads) and temperature $`T`$ (Coulomb blockade region). In this case we can integrate out the states with one or three extra electrons. This is equivalent to Schrieffer-Wolf transformation which is used to obtain the conventional Kondo model . We obtain the following effective low-energy Hamiltonian
$$H_{\mathrm{eff}}=H^{S=1}+H^{S=10}+H_{\mathrm{eff}}^{}+H_{\mathrm{dot}}.$$
(7)
The first term involves components of the spin-triplet state and resembles a conventional Kondo Hamiltonian for $`S=1`$.
$`H^{S=1}`$ $`=`$ $`{\displaystyle \underset{kk^{}}{}}{\displaystyle \underset{\alpha \beta =L,R}{}}{\displaystyle \underset{i=1,2}{}}J_{\alpha \beta }^{(i)}\left[\widehat{S}_+c_{\alpha k^{}}^{(i)}c_{\beta k}^{(i)}+\widehat{S}_{}c_{\alpha ^{}k}^{(i)}c_{\beta k}^{(i)}+\widehat{S}_z(c_{\alpha ^{}k}^{(i)}c_{\beta k}^{(i)}c_{\alpha ^{}k}^{(i)}c_{\beta k}^{(i)})\right]`$ (8)
$`=`$ $`{\displaystyle \underset{kk^{}}{}}{\displaystyle \underset{\alpha \beta =L,R}{}}{\displaystyle \underset{i=1,2}{}}J_{\alpha \beta }^{(i)}[\sqrt{2}(f_{11}^{}f_{10}+f_{10}^{}f_{11})c_{\alpha k^{}}^{(i)}c_{\beta k}^{(i)}+\sqrt{2}(f_{10}^{}f_{11}+f_{11}^{}f_{10})c_{\alpha k^{}}^{(i)}c_{\beta k}^{(i)}`$ (10)
$`+(f_{11}^{}f_{11}f_{11}^{}f_{11})(c_{\alpha k^{}}^{(i)}c_{\beta k}^{(i)}c_{\alpha k^{}}^{(i)}c_{\beta k}^{(i)})].`$
Here we have introduced pseudo-fermion operators $`f_{SM}^{}`$ ($`f_{SM}`$) which create (annihilate) the state $`|SM`$. It is required that $`_{SM}f_{SM}^{}f_{SM}=1`$. The second term in $`H_{\mathrm{eff}}`$ describes the conversion between the spin-triplet and singlet states accompanied by interchannel scattering of conduction electrons
$`H^{S=10}={\displaystyle \underset{kk^{}}{}}{\displaystyle \underset{\alpha \beta =L,R}{}}\{\stackrel{~}{J}_{\alpha \beta }[\sqrt{2}(f_{11}^{}f_{00}f_{00}^{}f_{11})c_{\alpha k^{}}^{(1)}c_{\beta k}^{(2)}+\sqrt{2}(f_{00}^{}f_{11}f_{11}^{}f_{00})c_{\alpha k^{}}^{(1)}c_{\beta k}^{(2)}`$ (11)
$`(f_{10}^{}f_{00}+f_{00}^{}f_{10})(c_{\alpha k^{}}^{(1)}c_{\beta k}^{(2)}c_{\alpha k^{}}^{(1)}c_{\beta k}^{(2)})]+\stackrel{~}{J}_{\alpha \beta }^{}[12]\}.`$ (12)
The third term $`H_{\mathrm{eff}}^{}`$ represents the scattering processes without change of the dot state and is not relevant for the current discussion. The coupling constants are given by $`J_{\alpha \beta }^{(i)}=V_{\alpha ,i}V_{\beta ,i}^{}/(2E_c),\stackrel{~}{J}_{\alpha \beta }=V_{\alpha ,1}V_{\beta ,2}^{}/(2E_c),`$ where $`1/E_c=1/E^++1/E^{}`$. The Hamiltonian of the dot reads
$$H_{\mathrm{dot}}=\underset{S,M}{}E_Sf_{SM}^{}f_{SM}.$$
(13)
To avoid the complication due to the fact that there are two leads $`\alpha =L,R`$, we perform a unitary transformation for electron modes in the leads along the lines of Ref. ; $`c_{k\sigma }^{(i)}=(V_{L,i}c_{L,k\sigma }+V_{R,i}c_{R,k\sigma })/V_i`$, $`\overline{c}_{k\sigma }^{(i)}=(V_{R,i}^{}c_{L,k\sigma }V_{L,i}^{}c_{R,k\sigma })/V_i`$, with $`V_i=\sqrt{|V_{L,i}|^2+|V_{R,i}|^2}`$. The modes $`\overline{c}_{k\sigma }^{(i)}`$ are not coupled to the quantum dot and shall be disregarded. The coupling constants for modes $`c_{k\sigma }^{(i)}`$ become
$`J^{(i)}`$ $`=`$ $`{\displaystyle \frac{|V_{Li}|^2+|V_{Ri}|^2}{2E_c}},`$ (14)
$`\stackrel{~}{J}`$ $`=`$ $`{\displaystyle \frac{1}{2E_c}}{\displaystyle \frac{(V_{L1}^2+V_{R1}^2)(V_{L2}^2+V_{R2}^2)}{\sqrt{(|V_{L1}|^2+|V_{R1}|^2)(|V_{L2}|^2+|V_{R2}|^2)}}}.`$ (15)
The spin-flip processes included in our model are shown in the inset of Fig. 1(b).
We calculate the Kondo temperature $`T_K`$ with the poor man’s scaling technique . By this method, we can properly consider the energies from $`T_K`$ to the upper cutoff. We concentrate on evaluating the exponential part of $`T_K`$. We assume constant density of states in the leads $`\nu `$ in the energy band of $`[D,D]`$. By changing the energy scale from $`D`$ to $`D|dD|`$, we obtain a closed form of the scaling equations for $`J^{(1)}`$, $`J^{(2)}`$, and $`\stackrel{~}{J}`$ in two limits.
In the first limit, the energy difference $`|\mathrm{\Delta }|`$ is negligible ($`|\mathrm{\Delta }|D`$) and $`H_{\mathrm{dot}}`$ can be safely disregarded. The scaling equations are best presented in the following matrix form:
$$\frac{d}{d\mathrm{ln}D}\left(\begin{array}{cc}J^{(1)}& \stackrel{~}{J}\\ \stackrel{~}{J}^{}& J^{(2)}\end{array}\right)=2\nu \left(\begin{array}{cc}J^{(1)}& \stackrel{~}{J}\\ \stackrel{~}{J}^{}& J^{(2)}\end{array}\right)^2.$$
(16)
The equations can be readily rewritten for eigenvalues of the matrix, $`J_\pm =(J^{(1)}+J^{(2)})/2\pm \sqrt{(J^{(1)}J^{(2)})^2/4+|\stackrel{~}{J}|^2}`$. The larger one, $`J_+`$, diverges faster upon decreasing the bandwidth $`D`$ and hence determines $`T_K`$. If the equations remain valid till the scaling ends ($`|\mathrm{\Delta }|T_K`$), the Kondo temperature is $`T_K(0)=D_0\mathrm{exp}[1/2\nu J_+]`$. Here $`D_0`$ is the initial bandwidth given by $`\sqrt{E^+E^{}}`$ . In another limiting case, $`\mathrm{\Delta }D`$. In this case the ground state of the dot is spin triplet and the singlet state can be disregarded. $`J^{(1)}`$ and $`J^{(2)}`$ evolve independently
$$\frac{d}{d\mathrm{ln}D}J^{(i)}=2\nu J^{(i)2},$$
(17)
and $`\stackrel{~}{J}`$ does not change. If these equations remain valid in the whole scaling region ($`\mathrm{\Delta }>D_0`$), it yields $`T_K(\mathrm{})=D_0\mathrm{exp}[1/2\nu J^{(1)}]`$. Here we assume $`J^{(1)}J^{(2)}`$. This is the Kondo temperature for spin-triplet localized spins .
To determine $`T_K`$ in the intermediate region, $`T_K(0)\mathrm{\Delta }D_0`$, we match the solutions of Eqs. (16) and (17) at $`D\mathrm{\Delta }`$. $`\stackrel{~}{J}`$ saturates at this point while $`J^{(1)}`$ and $`J^{(2)}`$ continue to grow with decreasing $`D`$. This yields power law dependence on $`\mathrm{\Delta }`$
$$T_K(\mathrm{\Delta })=T_K(0)\left(T_K(0)/\mathrm{\Delta }\right)^\gamma ,$$
(18)
where $`\sqrt{\gamma }=|\stackrel{~}{J}|/[\sqrt{(J^{(1)}J^{(2)})^2/4+|\stackrel{~}{J}|^2}+|J^{(1)}J^{(2)}|/2]`$. The exponent $`\gamma `$ appears to be nonuniversal, depending on a ratio of the initial coupling constants. In general, $`0<\gamma 1`$. For $`\mathrm{\Delta }<0`$, all the coupling constants saturate and no Kondo effect is expected, provided $`|\mathrm{\Delta }|T_K(0)`$.
In a simple case of the identical couplings, $`J^{(1)}=J^{(2)}=\stackrel{~}{J}`$ $`(J)`$, $`T_K(0)=D_0\mathrm{exp}[1/4\nu J]`$. For $`\mathrm{\Delta }>0`$, $`T_K`$ decreases with increasing $`\mathrm{\Delta }`$ as $`T_K(\mathrm{\Delta })=T_K(0)^2/\mathrm{\Delta }`$ ($`\gamma =1`$ in Eq. (18)) and finally converges to $`D_0\mathrm{exp}[1/2\nu J]=T_K(0)^2/D_0`$. For $`\mathrm{\Delta }<0`$, $`T_K`$ drops to zero suddenly at $`|\mathrm{\Delta }|T_K(0)`$. The dependence of the Kondo temperature on $`\mathrm{\Delta }`$ is schematically shown in Fig. 1(b).
We have discussed so far the case of $`C_1=C_2=1`$ in Eq. (4) for the spin-singlet state. This is not required by symmetry and $`C_1C_2`$ in general. To justify the assumption we made, let us consider the renormalization equations for $`C_1C_2`$. The coupling constants $`\stackrel{~}{J}_1=C_1\stackrel{~}{J}`$ and $`\stackrel{~}{J}_2=C_2\stackrel{~}{J}`$ are renormalized now in a different way, involving the scattering processes without spin flip in the dot
$$H_{\mathrm{eff}}^{}=\underset{kk^{}\sigma }{}\underset{i=1,2}{}\left[J^{(i)}c_{k^{}\sigma }^{(i)}c_{k\sigma }^{(i)}\underset{M}{}f_{1M}^{}f_{1M}+J^{\prime \prime (i)}c_{k^{}\sigma }^{(i)}c_{k\sigma }^{(i)}f_{00}^{}f_{00}\right].$$
(19)
General scaling equations for $`|\mathrm{\Delta }|D`$ are given by
$`d\stackrel{~}{J}_1/d\mathrm{ln}D`$ $`=`$ $`2\nu (J^{(1)}+J^{(2)})\stackrel{~}{J}_1+\nu J^{}\stackrel{~}{J}_1`$ (20)
$`d\stackrel{~}{J}_2/d\mathrm{ln}D`$ $`=`$ $`2\nu (J^{(1)}+J^{(2)})\stackrel{~}{J}_2\nu J^{}\stackrel{~}{J}_2`$ (21)
$`dJ^{}/d\mathrm{ln}D`$ $`=`$ $`8\nu (|\stackrel{~}{J}_1|^2|\stackrel{~}{J}_2|^2)`$ (22)
$`dJ^{(i)}/d\mathrm{ln}D`$ $`=`$ $`2\nu \left[J^{(i)2}+(|\stackrel{~}{J}_1|^2+|\stackrel{~}{J}_2|^2)/2\right]`$ (23)
where $`J^{}=J^{(1)}J^{(2)}J^{\prime \prime (1)}+J^{\prime \prime (2)}`$. When $`\mathrm{\Delta }D`$, the equations are identical to Eq. (17). Our point is that if we concentrate on the most rapidly divergent solutions of Eqs. (20)-(23), which are proportional to $`1/\mathrm{ln}D`$, $`\stackrel{~}{J}_1`$ and $`\stackrel{~}{J}_2`$ appear to be the same. To this leading order, the renormalization is the same as given by Eq. (16). Consequently the Kondo temperature is the same as that in the case of $`C_1=C_2=1`$, apart from a prefactor.
We calculate perturbation corrections to conductance when $`T_KTE_c`$. The third order perturbations in $`J`$’s yield the logarithmic corrections $`G_K`$ typical for the Kondo effect . At $`T|\mathrm{\Delta }|`$,
$`G_K/(2e^2/h)={\displaystyle \underset{i=1,2}{}}{\displaystyle \frac{4\mathrm{\Gamma }_L^i\mathrm{\Gamma }_R^i}{(\mathrm{\Gamma }_L^i+\mathrm{\Gamma }_R^i)^2}}6\pi ^2J^{(i)}\nu \left[(J^{(i)}\nu )^2+|\stackrel{~}{J}\nu |^2\right]\mathrm{ln}{\displaystyle \frac{D_0}{T}}`$ (24)
$`+{\displaystyle \frac{2(\mathrm{\Gamma }_L^1\mathrm{\Gamma }_R^2+\mathrm{\Gamma }_L^2\mathrm{\Gamma }_R^1)}{(\mathrm{\Gamma }_L^1+\mathrm{\Gamma }_R^1)(\mathrm{\Gamma }_L^2+\mathrm{\Gamma }_R^2)}}12\pi ^2\left[J^{(1)}\nu +J^{(2)}\nu \right]|\stackrel{~}{J}\nu |^2\mathrm{ln}{\displaystyle \frac{D_0}{T}}.`$ (25)
The interplay between the spin-singlet and triplet states largely enhances the conductance. In the opposite case, $`T|\mathrm{\Delta }|`$, the interplay becomes less effective and the logarithmic corrections become smaller
$$G_K/(2e^2/h)=\underset{i=1,2}{}\frac{4\mathrm{\Gamma }_L^i\mathrm{\Gamma }_R^i}{(\mathrm{\Gamma }_L^i+\mathrm{\Gamma }_R^i)^2}8\pi ^2J^{(i)}\nu \left[(J^{(i)}\nu )^2\mathrm{ln}\frac{D_0}{T}+|\stackrel{~}{J}\nu |^2\mathrm{ln}\frac{D_0}{\mathrm{\Delta }}\right]$$
(26)
for $`\mathrm{\Delta }>0`$ and disappear for $`\mathrm{\Delta }<0`$. If one varies $`\mathrm{\Delta }`$ at fixed $`T`$, one sees the enhanced conductance for $`|\mathrm{\Delta }|T`$.
We believe that the conductance approaches the unitary limit $`2e^2/h`$ at $`TT_K`$ in our model. However, not all the terms in $`G_K`$ can be renormalized to a universal function $`G(T/T_K)`$. Because of the multichannel nature of our model, one should expect nonuniversal logarithmic terms along with the universal ones. Those, however, will be much smaller than $`2e^2/h`$.
Recent experiment has shown a significant enhancement of conductance to the values of the order of $`e^2/h`$, around the crossing point between spin-singlet and triplet states, in vertical quantum dots. The conductance remains low in the stability domains of singlet or triplet states. Our results provide a possible theoretical explanation for the results. The conductance increase can be attributed to the Kondo effect enhanced by the competition between the two states, near the crossing point. The Kondo temperature elsewhere is probably low in comparison with the actual temperature, so that no conductance increase is seen.
In conclusion, we have shown that the competition between spin-singlet and triplet states enhances the Kondo effect. This may explain recent experimental observations. We have predicted power law dependence of the Kondo temperature on the energy difference between the states.
The authors are indebted to L. P. Kouwenhoven for suggesting the topic of the research presented, S. De Franceschi, J. M. Elzerman, K. Maijala, S. Sasaki, W. G. van der Wiel, Y. Tokura, L. I. Glazman, M. Pustilnik, and G. E. W. Bauer for valuable discussions. The authors acknowledge financial support from the “Netherlandse Organisatie voor Wetenschappelijk Onderzoek” (NWO). M. E. is also grateful for financial support from the Japan Society for the Promotion of Science for his stay at Delft University of Technology. |
warning/0002/nlin0002020.html | ar5iv | text | # Approximation theorem for the self-focusing Nonlinear Schrödinger Equation and for the periodic curves in 𝐑³.
## Abstract
It is shown, that any sufficiently smooth periodic solution of the self-focusing Nonlinear Schrödinger equation can be approximated by finite-gap ones with an arbitrary small error. As a corollary an analogous result for the motion of closed curves in $`R^3`$ guided by the Filament equation is proved. This equation describes the dynamics of very thin filament vortices in a fluid.
thanks: This work was supported by the INTAS grant No. 96-770 and by the Russian Foundation for Basic Research grant No. 98-01-01161.
Dedicated to V.E.Zakharov’s 60th birthday
One of the basic questions of the finite-gap theory is the following: are the finite-gap solutions of a given equation sufficiently generic or they belong to some special subclass? To answer this question it is reasonable to check if arbitrary periodic in spatial variables solution can be approximated by finite-gap ones.
The study of simplest examples shows, that the uniform approximation for all $`x`$ and $`t`$ is impossible, because to do it we have to keep all space and time frequencies simultaneously, and we have too many constraints on the deformations to fulfill all of them.
Therefore it is natural to ask the following questions:
1) Let us have a smooth periodic in spatial variables solution and an arbitrary compact domain $`U`$ in the $`(x,t)`$ space. Is it possible to approximate our solution on $`U`$ by finite-gap ones with arbitrary small error? The approximating solutions are allowed to be non-periodic. We shall call such approximations local.
2) Let us have a smooth periodic in spatial variables solution and an arbitrary compact domain $`U`$ in the $`t`$ space. Is it possible to approximate our solution by periodic in $`x`$ finite-gap ones with the same $`x`$-periods for all $`x`$ and $`tU`$? We shall call such approximations periodic.
Of course any periodic approximation is automatically local, but the characterization of periodic finite-gap solutions is usually sufficiently complicated, therefore the transition from local approximations to periodic ones may be rather nontrivial.
In our text we construct periodic approximations for the self-focusing Nonlinear Schrödinger equation (SfNLS)
$$iq_t+q_{xx}+2q^2\overline{q}=0,\text{where}$$
(1)
$`q=q(x,t)`$ is a complex-valued function of two real variables, and for the Filament equation
$$\frac{\stackrel{}{\gamma }(s,t)}{t}=k(x,t)\stackrel{}{b}(s,t),\text{where}$$
(2)
$`\stackrel{}{\gamma }(s,t)`$ is a $`t`$-dependent family of smooth curves in $`R^3`$, $`s`$ is the natural parameter (i.e. $`|_s\stackrel{}{\gamma }(s,t)|1`$), $`\stackrel{}{b}(s,t)`$, $`k(s,t)`$ denote the binormal vector from the Frenet reper and the curvature function respectively.
The first periodic approximation theorem was proved in 1975 by Marchenko and Ostrovskii for the real Korteveg de Vries (KdV) equation. The method of is based on the theory of conformal maps and it can be naturally extended to some soliton systems including the defocusing NLS. But for other systems the question is still open and the answer depends on the equation. For example, from results of Krichever it is rather clear that any periodic Kadomthsev-Petviashvili II (KP II) solution allows finite-gap approximations, but it is likely that for KP I it is not so.
Direct attempts to generalize the approach of to SfNLS meet the following problems
1. The space of spectral curves corresponding to real periodic $`g`$-gap KdV solutions is topologically $`R\left(R^+\right)^g`$. But in the SfNLS case this space it the real part of some ramified covering of $`C^{g+1}`$, and the structure of the Marchenko-Ostrovskii conformal map is essentially more complicated. To avoid a detailed study of the parameters space we use the method of isoperiodic deformations suggested by M.Schmidt and author in .
2. In the KdV case the characterization of admissible divisors is very simple: the $`n`$-th point of divisor is located at an arbitrary point of the $`n`$-th compact real oval. In the SfNLS case the characterization is less explicit (see below) and we have to describe how we vary the divisor after perturbing the spectral curve to preserve the admissibility.
3. Dubrovin equations for the real KdV are non-singular therefore a small change of parameters and starting point slightly affects the solution. But in the SfNLS case Dubrovin equations may have singularities, and the solutions may have branch points (nevertheless the corresponding SfNLS potential is smooth). Therefore we have to check that our variations of spectral data do not change the solution too much. To do it we introduce some new “symmetric” variables.
Let us recall some basic facts from the SfNLS theory. The scattering transform for NLS was found in 1971 by Zakharov and Shabat . Finite-gap NLS solutions were first constructed in 1976 by Its and Kotljarov . The characterization of SfNLS admissible divisors as well as a proof that all solutions with reduction (6) are automatically nonsingular were obtained by Cherednik . Infinite-gap periodic problem for matrix operators including the NLS $`L`$-operator was studied by M.Schmidt . A lot of useful information about the NLS theory including the Hamiltonian theory is contained in the book by Faddeev and Takhtadjan. Finite-gap NLS theory is discussed in details in the article by Previato. Effictivisation of low genus formulas by NLS was studied by Kamchatnov .
The zero-curvature representation for SfNLS reads as:
$$\frac{\mathrm{\Psi }(\lambda ,x,t)}{x}=U(\lambda ,x,t)\mathrm{\Psi }(\lambda ,x,t),\frac{\mathrm{\Psi }(\lambda ,x,t)}{t}=V(\lambda ,x,t)\mathrm{\Psi }(\lambda ,x,t),$$
(3)
where $`\mathrm{\Psi }(\lambda ,x,t)`$ is a 2-component vector,
$$\mathrm{\Psi }(\lambda ,x,t)=\left[\begin{array}{c}\psi _1(\lambda ,x,t)\\ \psi _2(\lambda ,x,t)\end{array}\right],$$
(4)
$`U(\lambda ,x,t)`$, $`V(\lambda ,x,t)`$ are the following $`2\times 2`$ matrices:
$$U(\lambda ,x,t)=\left[\begin{array}{cc}i\lambda & iq(x,t)\\ ir(x,t)& i\lambda \end{array}\right],V(\lambda ,x,t)=2\lambda U(\lambda ,x,t)+\left[\begin{array}{cc}iqr& q_x\\ r_x& iqr\end{array}\right],$$
(5)
$$r(x,t)=\overline{q}(x,t).$$
(6)
We shall assume that $`q(x,t)`$ is periodic in $`x`$ with the period 1
$$q(x+1,t)q(x,t).$$
(7)
We shall fix our attention on the spectral transform for a fixed $`t=t_0`$, therefore starting from this moment we shall omit $`t`$ in our notations.
The Bloch eigenfunction $`\mathrm{\Psi }(\lambda ,x)`$ is by definition the common eigenfunction of $`L=_xU(\lambda ,x)`$ and the shift operator
$$L\mathrm{\Psi }(\lambda ,x)=0,\mathrm{\Psi }(\lambda ,x+1)=e^{ip(\lambda )}\mathrm{\Psi }(\lambda ,x).$$
(8)
Equation (8) defines $`\mathrm{\Psi }(\lambda ,x)`$ up to a constant factor. We fix it assuming
$$\mathrm{\Phi }(\lambda ,0)1,\text{where}\mathrm{\Phi }(\lambda ,x)=\psi _1(\lambda ,x)+\psi _2(\lambda ,x).$$
(9)
The function $`p(\lambda )`$ is defined up to adding $`2\pi n`$, $`nZ`$. It is called the quasimomentum.
To calculate the Bloch function we have to diagonalize the $`2\times 2`$ monodromy matrix $`T(\lambda )`$, which is an entire function of $`\lambda `$. The eigenfunctions of $`T(\lambda )`$ lie on a two-sheeted covering $`\mathrm{\Gamma }`$ of the $`\lambda `$-plane. $`\mathrm{\Gamma }`$ is called spectral curve. Denote the permutation of sheets of $`\mathrm{\Gamma }`$ by $`\sigma `$. $`detT(\lambda )1`$, therefore $`p(\gamma )+p(\sigma \gamma )0(\text{mod}2\pi )`$. (We shall denote points of $`\mathrm{\Gamma }`$ by $`\gamma `$ and the projection $`\mathrm{\Gamma }C`$ by $`𝒫`$, $`\lambda =𝒫\gamma `$). $`p(\gamma )`$ is a locally holomorphic multivalued function on $`\mathrm{\Gamma }`$, $`dp=(dp(\lambda )/d\lambda )d\lambda `$ is a holomorphic differential on the finite part of $`\mathrm{\Gamma }`$, $`\sigma (dp)=dp`$. $`p(\gamma )\pm \lambda `$ as $`\lambda \mathrm{}`$ (as an asymptotic series), therefore $`\mathrm{\Gamma }`$ is compactified by 2 infinite points $`\mathrm{}_+`$, $`\mathrm{}_{}`$, $`\sigma \mathrm{}_+=\mathrm{}_{}`$, $`\sigma \mathrm{}_{}=\mathrm{}_+`$, $`p(\gamma )\pm \lambda `$ as $`\gamma \mathrm{}_\pm `$ respectively.
A point $`\lambda C`$ is called regular if $`p(\gamma )0(\text{mod}\pi )`$, where $`𝒫\gamma =\lambda `$ and irregular otherwise. Let $`\lambda _k`$ be an irregular point. The Tailor expansion of $`trT(\lambda )`$ reads as: $`trT(\lambda )=\pm 2+T_k^n(\lambda \lambda _k)^n+\mathrm{}`$. Let us call $`n`$ the order of the point $`\lambda _k`$, $`n=ord_q(\lambda _k)`$. ($`q`$ means that the order is defined in terms of the quasimomentum function). $`\lambda _k`$ is a branch point of $`\mathrm{\Gamma }`$ if $`n`$ is odd and a double point of $`\mathrm{\Gamma }`$ if $`n`$ is even. A branch point is called simple if $`n=1`$. If the opposite is not stated explicitly we have one point of $`\mathrm{\Gamma }`$ over each double point. To have an uniform representation for our equations we shall treat an irregular point of order $`n`$ as the result of fusion $`n`$ simple branch point.
If $`\mathrm{\Psi }(\lambda ,x)`$ is a Bloch solution of (8), then
$$\mathrm{\Psi }^+(\overline{\lambda },x)=\left[\begin{array}{c}\overline{\psi }_2(\lambda ,x)\\ \overline{\psi }_1(\lambda ,x)\end{array}\right],$$
(10)
is also a Bloch solution of (8) with the quasimomentum $`p^+(\overline{\lambda })=\overline{p}(\lambda )`$. Therefore $`\mathrm{\Gamma }`$ has the following antiholomorphic involutions $`\gamma \sigma \overline{\gamma }`$ and $`\gamma \overline{\gamma }`$. (We assume that $`p(\overline{\gamma })=\overline{p}(\gamma )`$, $`p(\sigma \overline{\gamma })=\overline{p}(\gamma )`$). For real $`\lambda `$ $`T(\lambda )`$ is an unitary matrix, $`p(\lambda )R`$, and $`\mathrm{\Gamma }`$ has no real branch points (but may have real double points).
Let $`\{E_k\}`$, $`\{E_k^+\}`$ be the lists of all irregular points, the index $`k`$ takes all integer values. We assume that
1. $`ImE_k0`$.
2. $`E_k^+=\overline{E_k}`$.
3. If $`ord_q(\lambda _k)=n`$ the point $`\lambda _k`$ has exactly $`n`$ entries in our lists. For example if $`ord_q(\lambda _k)=4`$ and $`\lambda _kR`$ then we have exactly 2 integers $`k_1`$, $`k_2`$ such that $`E_{k_1}=E_{k_2}=E_{k_1}^+=E_{k_2}^+`$.
###### Lemma 1
It is possible to enumerate the irregular points so, that for sufficiently large $`|k|`$
1. $`E_k=(\pi sgnk)\sqrt{k^2I_1(q)}+o\left(\frac{1}{k}\right)`$ where $`I_1(q)=_0^1q(x)\overline{q}(x)𝑑x`$.
2. $`E_kE_k^+0`$ faster than any degree of $`k^1`$ as $`|k|\mathrm{}`$ (We assume $`q(x)`$ to be smooth).
The next important object for us is the set of zeroes of the quasimomentum differential $`dp`$. They are invariant under the involution $`\sigma `$ therefore we shall consider their projections to the $`\lambda `$-plane instead. Denote them by $`\alpha _k`$ where $`k`$ takes all integer values. As above we use the following agreement
1. If $`\lambda _k`$ is a regular points it has $`n`$ entries in the list $`\{\alpha _k\}`$ where $`n`$ is the order of zero of $`dp`$ at one sheet.
2. If $`\lambda _k`$ is an irregular points it has $`\left[\frac{ord_q(\lambda _k)}{2}\right]`$ entries to the list $`\{\alpha _k\}`$ where $`[]`$ denotes the integer part.
###### Lemma 2
It is possible to enumerate the points $`\alpha _k`$ so, that for sufficiently large $`|k|`$
1. $`Im\alpha _k=0`$.
2. $`\alpha _k=ReE_k+o(ImE_k)`$.
Let us define now the “second part” of the spectral data – the divisor of poles of the Bloch function.
Let $`\stackrel{~}{\mathrm{\Psi }}(\gamma ,x)`$ denote a Bloch eigenfunction of $`L`$ with some non-singular locally holomorphic normalisation (of course $`\mathrm{\Psi }(\gamma ,x)=\stackrel{~}{\mathrm{\Psi }}(\gamma ,x)/\stackrel{~}{\mathrm{\Phi }}(\gamma ,0)`$ where $`\stackrel{~}{\mathrm{\Phi }}(\gamma ,x)=\stackrel{~}{\psi }_1(\gamma ,x)+\stackrel{~}{\psi }_2(\gamma ,x)`$. Consider the Wronskian of the Bloch functions $`\stackrel{~}{W}(\gamma )=\stackrel{~}{\psi }_1(\gamma ,x)\stackrel{~}{\psi }_2(\sigma \gamma ,x)\stackrel{~}{\psi }_2(\gamma ,x)\stackrel{~}{\psi }_1(\sigma \gamma ,x)`$. It is defined up to a non-zero holomorphic multiplier and does not wanish at regular points. Let $`\lambda _k`$ be an irregular point. We have $`W(\lambda )=\pm w_k^m(\lambda \lambda _k)^{m/2}(1+o(1))`$ where $`m`$ is even if $`\lambda _k`$ is a double point and odd if $`\lambda _k`$ is a branch point, $`m0`$. Denote $`m`$ by $`ord_b(\lambda _k)`$. It is easy to check that $`ord_b(\lambda _k)ord_q(\lambda _k)`$. A double point $`\lambda _k`$ is called removable if $`ord_b(\lambda _k)=0`$. It is well-known that removable double points can be treated as regular points and we can forget about them.
The divisor of Bloch function zeroes is a list of points of $`\mathrm{\Gamma }`$ $`\{\gamma _k(x)\}`$ where $`k`$ takes all integer values such that each zero of $`\stackrel{~}{\mathrm{\Phi }}(\gamma ,x)`$ generates $`l`$ entries to this list if $`l`$ is the multiplicity of it and each irregular point generates $`(ord_qord_b)/2`$ entries. The divisor of Bloch function poles $`\{\gamma _k\}`$ coinsides with the divisor of Bloch function zeroes taken at the point $`x=0`$.
###### Lemma 3
1. The spectral curve $`\mathrm{\Gamma }`$ has only finite number of non-removable double points and degenerate branch points.
2. All real double points are removable.
###### Lemma 4
It is possible to enumerate the points $`\gamma _k`$ so, that for sufficiently large $`|k|`$ $`𝒫\gamma _k=ReE_k+O(ImE_k)`$.
It is well-known, that the spectral curve and the divisor of poles completely define the potential $`q(x)`$. To reconstruct the potential we can use Dubrovin equations
$$\frac{}{x}\lambda _j(x)=2i\left[\lambda _j(x)+\underset{k=\mathrm{}}{\overset{\mathrm{}}{}}\left(\frac{E_k+E_k^+}{2}\lambda _k(x)\right)\right]\nu _j(x),$$
(11)
where
$$\nu _j(x)=\sqrt{(\lambda _j(x)E_j)(\lambda _j(x)E_j^+)}\underset{kj}{}\frac{\sqrt{(\lambda _j(x)E_k)(\lambda _j(x)E_k^+)}}{\lambda _j(x)\lambda _k(x)},$$
(12)
and the reconstruction formula
$$q(x)=\underset{k=\mathrm{}}{\overset{\mathrm{}}{}}\left(\frac{E_k+E_k^+}{2}\lambda _k(x)\right)+\underset{j=\mathrm{}}{\overset{\mathrm{}}{}}\nu _j(x).$$
(13)
The infinite sums and products in the formulas above perfectly converge.
Let us recall the characterization of divisors corresponding to operators with the reduction (6). Consider the following 1-form on $`\mathrm{\Gamma }`$: $`\mathrm{\Omega }(\gamma ,x)=\omega (\gamma ,x)d\lambda `$, where
$$\omega (\gamma ,x)=\frac{(\stackrel{~}{\psi }_1(\gamma ,x)+\stackrel{~}{\psi }_2(\gamma ,x))(\stackrel{~}{\psi }_2(\sigma \gamma ,x)\stackrel{~}{\psi }_1(\sigma \gamma ,x))}{\stackrel{~}{\psi }_1(\gamma ,x)\stackrel{~}{\psi }_2(\sigma \gamma ,x)\stackrel{~}{\psi }_2(\gamma ,x)\stackrel{~}{\psi }_1(\sigma \gamma ,x)}.$$
(14)
Denote by $`U(R)`$ be the domain $`|\lambda |<R`$ in $`\mathrm{\Gamma }`$. Consider the following function $`\omega _R(\gamma ,x)=\omega (\gamma ,x)\underset{k:|E_k|<R}{}\sqrt{(\lambda E_k)(\lambda E_k^+)}`$. Denote by $`D(\omega _R,x)`$ the divisor of zeroes of $`\omega _R(\gamma ,x)`$ and by $`D(\omega ,x)`$ the limit of $`D(\omega _R,x)`$ as $`R\mathrm{}`$.
###### Lemma 5
$`D(\omega ,x)`$ coincide with the set $`\{\gamma _k(x),\overline{\gamma }_k(x)\}`$, where $`\{\gamma _k(x)\}`$ is the divisor of Bloch function zeroes.
Let $`\delta _j`$ be an arbitratry collection of pairwise distinct real points such that for sufficiently large $`|j|`$ $`\delta _j=ReE_j`$.
###### Lemma 6
1. The form $`\mathrm{\Omega }`$ reads as $`\mathrm{\Omega }=\left[1\stackrel{~}{\kappa }(\gamma ,x)\right]d\lambda `$
$$\stackrel{~}{\kappa }(\gamma ,x)=\underset{j=\mathrm{}}{\overset{\mathrm{}}{}}\frac{\kappa _j(x)}{\sqrt{(\lambda E_j)(\lambda E_j^+)}}\underset{kj}{}\frac{\lambda \delta _k}{\sqrt{(\lambda E_k)(\lambda E_k^+)}},$$
(15)
where $`\kappa _j(x)`$ are some real functions of $`x`$.
2. $`|\kappa (\gamma ,x)|1`$ for all $`xR`$, $`\gamma R`$.
In particular $`\kappa (\delta _j,x)|1`$ for all $`j`$. It gives us the following estimate on the functions $`\kappa _j(x)`$:
$$|\kappa _j(x)|\left|\sqrt{(\delta _jE_j)(\delta _jE_j^+)}\underset{kj}{}\frac{\sqrt{(\delta _jE_k)(\delta _jE_k^+)}}{\delta _j\delta _k}\right|\left|\delta _jE_j\right|C_\mathrm{\Gamma },$$
(16)
where $`C_\mathrm{\Gamma }`$ is a positive constant, depending only on the spectral curve. In particular, if we have a removable double point $`E_k=E_k^+=\delta _k`$, then $`\kappa _k(x)0`$.
###### Lemma 7
Let $`\kappa _k(0)`$ be a collection of real numbers such, that $`|\kappa (\gamma ,0)|1`$ for all $`\gamma \mathrm{\Gamma }`$, where $`\kappa (\gamma ,0)`$ is defined by (15), $`D(\omega ,0)`$ be the corresponding divisor, $`\{\gamma _j(0)\}`$ be any set of points such that $`D(\omega ,0)=\{\gamma _k(0),\overline{\gamma }_k(0)\}`$. Then the corresponding operator $`L`$-operator satisfy (6), and the potential $`q(x)`$ is nonsingular.
For us the following definition will be convenient: potential $`q(x)`$ is called finite-gap if $`E_j=E_j^+`$ for all $`|j|J_0`$. Then all points $`E_j=E_j^+`$ are removable double points, $`\alpha _j=\gamma _j(x)=E_j`$ for all $`|j|J_0`$ and $`\mathrm{\Gamma }`$ has only finite number of branch points and non-removable double points. Finite-gap solutions of soliton equations were first introduced by Novikov in 1974 for KdV . The corresponding solutions can be written explicitly in terms of Riemann $`\theta `$-functions.
The first step of the approximation procedure is to construct a finite-gap deformation of $`\mathrm{\Gamma }`$ generating solutions with the same period. To do it we need the following lemma proved by M.Schmidt and the author in .
###### Lemma 8
Let $`\alpha _kR`$ be the projection of a zero of the quasimomentum such, that $`\alpha _k\alpha _j`$ for $`jk`$, $`\alpha _kE_j`$, $`\alpha _kE_j^+`$, for all $`j`$. Consider the following system of ODE’s, associated with the point $`\alpha _k`$:
$$\frac{E_j}{\tau }=\frac{c_k(\tau )}{E_j\alpha _k},\frac{E_j^+}{\tau }=\frac{c_k(\tau )}{E_j^+\alpha _k},$$
$$\frac{\alpha _j}{\tau }=\frac{c_k(\tau )}{\alpha _j\alpha _k}\text{ for }jk,$$
(17)
$$\frac{\alpha _k}{\tau }=c_k(\tau )\left[\underset{jk}{}\frac{1}{\alpha _j\alpha _k}\frac{1}{2}\underset{j=\mathrm{}}{\overset{\mathrm{}}{}}\left(\frac{1}{E_j\alpha _k}+\frac{1}{E_j^+\alpha _k}\right)\right],$$
where $`c_k(\tau )`$ is an arbitrary real function of $`\tau `$.
Denote by $`\mathrm{\Gamma }(\tau )`$ the solution of (17) with the initial value $`\mathrm{\Gamma }(0)=\mathrm{\Gamma }`$, where the spectral curve $`\mathrm{\Gamma }`$ corresponds to a periodic with the period 1 potential $`q(x)`$ (of course this solution is defined only in some neighborhood of zero $`U(0)`$). Then for all $`\tau U(0)`$ the curve $`\mathrm{\Gamma }(\tau )`$ generates periodic with the period 1 potentials (the $`x`$-quasifrequencies of the potentials do not depend on the divisor).
Let $`|k|`$ be sufficiently large. Then using this deformation we can merge the pair $`E_k`$, $`E_k^+`$ to a removable double point, and the corresponding shift of all points $`E_j`$, $`E_j^+`$, $`\alpha _j`$ with $`jk`$ is of order $`o(ImE_k)`$. Therefore applying this deformation to all $`k`$ such that $`|k|K`$, where $`K>0`$ is a sufficiently large integer, we obtain a finite-gap spectral curve $`\mathrm{\Gamma }_K`$ (it is almost evident that the superposition of infinitely many deformations perfectly converges).
###### Lemma 9
For any $`ϵ>0`$ there exists a $`K`$ such, that
1. $`|E_j\stackrel{~}{E}_j|<ϵ`$, $`|E_j^+\stackrel{~}{E}_j^+|<ϵ`$, $`|\alpha _j\stackrel{~}{\alpha }_j|<ϵ`$, for all $`j`$.
2. $`|Im(E_j\stackrel{~}{E}_j)|<ϵ|ImE_j|`$, for all $`j`$ such, that $`|j|<K`$.
where $`\stackrel{~}{E}_j`$, $`\stackrel{~}{E}_j^+`$, $`\stackrel{~}{\alpha }_j`$ are the branch points of the curve $`\mathrm{\Gamma }_K`$ and the quasimomentum zeroes respectively.
We have constructed a family of finite-gap curves $`\mathrm{\Gamma }_K`$ approximating the curve $`\mathrm{\Gamma }`$. Let us discuss now the admissible divisors.
###### Lemma 10
There exists a pair of positive integer constants $`K_1`$, $`K_2`$ such, that for all $`KK_2`$ the points of any admissible divisor $`\gamma _k`$ on $`\mathrm{\Gamma }_K`$ can be enumerated so, that
1. For all $`k`$ such that $`|k|K_1`$ $`|\lambda _k|<K_1+1/10`$.
2. For all $`k`$ such that $`|k|>K_1`$ $`|(\lambda _k\stackrel{~}{\delta }_k)|Im\stackrel{~}{E}_k`$, and $`|\stackrel{~}{\delta }_k|>K_1+11/10`$.
The proof follows from the characterization of admissible divisors given by Lemmas 5-7.
Equations (11)-(12) have singularities at the right-hand side. To simplify the structure of Dubrovin equation it is convenient to introduce the following new variables:
1. $`s_k(x)`$, $`q_k(x)`$, $`1k2K_1+1`$ – the first $`2K_1+1`$ expansion coefficients at $`\mathrm{}`$ of the function
$$\mathrm{\Xi }(\gamma ,x)=\left(1+\stackrel{~}{\kappa }(\gamma ,x)\right)\underset{k}{}\frac{\sqrt{(\lambda E_k)(\lambda E_k^+)}}{(\lambda \lambda _k(x))}$$
(18)
$$\mathrm{\Xi }(\lambda ,x)=\pm \left(1+\underset{k>0}{}\frac{s_k(x)}{\lambda ^k}\right)+\underset{k>0}{}\frac{q_k(x)}{\lambda ^k}\text{as}\gamma \pm \mathrm{}.$$
(19)
2. $`\stackrel{~}{\lambda }_k(x)=\lambda _k(x)\delta _k`$, $`|k|>K_1`$.
3. $`\stackrel{~}{\nu }_k(x)=\sqrt{(\lambda _k(x)E_k)(\lambda _k(x)E_k^+)}`$, $`|k|>K_1`$.
These variables are dependent. Denote this set of variables by $`𝒮`$.
Consider the following norm:
$$𝒮_n=\sqrt{\underset{|k|K_1}{}\left(|s_k|^2+|q_k|^2\right)+\underset{|k|>K_1}{}|k|^n\left(|\stackrel{~}{\lambda }_k|^2+|\stackrel{~}{\nu }_k|^2\right)}$$
(20)
This norm is bounded on the space of admissible divisors for any positive $`n`$.
###### Lemma 11
For any sufficiently large $`n`$ there exists a constant $`C_n(\mathrm{\Gamma })`$ such, that for any admissible pair $`𝒮_1(x)`$, $`𝒮_2(x)`$ of solutions of Dubrovin equations we have the following estimate
$$\frac{}{x}\left(𝒮_1(x)𝒮_2(x)\right)_nC_n(\mathrm{\Gamma })𝒮_1(x)𝒮_2(x)_n.$$
(21)
It is easy to check, that for any $`ϵ_1>0`$ there exists a constant $`K_3(n)`$ such that
$$𝒮_n^{(2)}<ϵ_1,\text{where}𝒮_n^{(2)}=\sqrt{\underset{|k|>K_3(n)}{}|k|^n\left(|\stackrel{~}{\lambda }_k|^2+|\stackrel{~}{\nu }_k|^2\right)}.$$
(22)
We need also the following semi-norm
$$𝒮_n^{(1)}=\sqrt{\underset{|k|K_1}{}\left(|s_k|^2+|q_k|^2\right)+\underset{K_1<|k|<K_3(n)}{}|k|^n\left(|\stackrel{~}{\lambda }_k|^2+|\stackrel{~}{\nu }_k|^2\right)}.$$
(23)
It is evident, that $`𝒮_n^{}𝒮_n^{(1)}+𝒮_n^{(2)}`$ and $`𝒮_n^{(1)}𝒮_n`$.
Combining (21) and (22) we obtain the following estimate
$$𝒮_1(x)𝒮_2(x)_ne^{C_n(\mathrm{\Gamma })|x|}\left(𝒮_1(0)𝒮_2(0)_n^{(1)}+ϵ_1\right).$$
(24)
Therefore in approximate calculations we can truncate the Dubrovin system to a finite-dimensional one, removing the variables $`\stackrel{~}{\lambda }_k(x)`$, $`\stackrel{~}{\nu }_k(x)`$, with the $`|k|>K_3(n)`$. For the truncated system small variations of the curve and of the starting point result in small variations of the solution. To complete the proof it is sufficient to check, that choosing $`K`$ sufficiently large we can make the admissible variation of divisor arbitrary small. But it follows from the characterization of admissible divisors presented above.
These arguments can be applied also for the Dubrovin equations, describing the $`t_l`$ -evolution of the divisor, where $`t_l`$ denotes the $`l`$-s time from the NLS hierarchy (in these notations $`t=t_1`$). Taking into account that the first $`k`$ $`x`$-derivatives of $`q(x)`$ are continuous functionals in the norm $`_n`$ for sufficiently large $`n`$ we obtain:
###### Theorem 1
Let $`q(x,t)`$ be an arbitrary SfNLS solution with smooth $`x`$-periodic Cauchy data $`q(x,0)=q_0(x)`$. Then for any $`ϵ>0`$, $`N>0`$ and $`𝒯>0`$ there exists a finite-gap SfNLS solution $`q^F(x,t)`$ such, that
$$\left|\frac{^n}{x^n}\left(q^F(x,t)q(x,t)\right)\right|<ϵ\text{for all}xR,|t|<𝒯,nN.$$
(25)
At the end let us say a few words how to prove an analogous theorem for the Filament equation. Equivalence between the SfNLS and the Filament equation is given by the Hasimoto map . ($`\theta `$-functional solutions of Filament equations were studied by Sym in ).
$$q(s,t)=\frac{1}{2}k(s,t)e^{i^s\kappa (\stackrel{~}{s},t)𝑑\stackrel{~}{s}},$$
(26)
where $`k(s,t)`$, $`\kappa (s,t)`$ are the curvature and the torsion functions respectively. In it was shown, that equations (17) (except one corresponding to $`\alpha _0=0`$) preserve the periodicity in $`s`$ of the Filament equation solution. Therefore the technique developed above can be applied without changes.
Acknowledgments. The author is grateful to Martin Schmidt for explaining some results from the periodic NLS theory. |
warning/0002/cond-mat0002295.html | ar5iv | text | # I Introduction
## I Introduction
Electric double layers form spontaneously whenever surfaces carrying ionizable groups are suspended in a polar solvent, most frequently water. The high dielectric constant $`ϵ`$ of the latter favours the dissociation of the functional surface groups, so that the surface acquires a net charge per unit area, $`\sigma `$. Oppositely charged counterions are released into the solvent, which generally also contains a finite concentration of microscopic anions and cations (microions), providing thus additional coions and counterions. The electric double-layer results from the build-up of a charge density (or ”cloud”) of opposite sign to that of the surface charge, which tends to screen the electrostatic potential due to the latter. The width of the double-layer, which is a measure of its capacity, is determined by the competition between the thermal motion of the microions, which tends to spread out, or homogenize their distribution, in order to increase their entropy, and the electrostatic interactions, which attract the counterions towards the surface, while repelling the coions.
The present review deals with electric double-layers around mesoscopic charged particles, which will be referred to as polyions (frequently also called macroions), and more specifically with the effective interactions between electric double-layers associated with different polyions. Such double-layers are ubiquitous in many physical, geophysical, chemical and biological systems, including complex fluids, clay soils, polyelectrolytes (e.g. DNA) or cell membranes. The main focus of the review will be on colloidal dispersions of spherical, rod-like or lamellar polyions, or micelles resulting from the self-assembly of ionic surfactants . Advanced experimental techniques, like surface force machines or video microscopy in combination with optical tweezers, allow a direct measurement of the forces acting between charged surfaces or polyions. Such direct measurements contrast with indirect determinations of effective forces via small angle X-ray or neutron scattering measurements of static structure factors . For the theoretician, the main challenge is the highly asymmetric nature of dispersions involving mesoscopic polyions, and microscopic solvent molecules and ions. Clearly, at least three widely different length scales are involved, namely the characteristic size (radius) $`a`$ of the polyions, a typical width of the electric double-layer, of the order of the Debye screening length $`\lambda _D`$, and a typical correlation length $`l`$ of the solvent, of the order of a few molecular diameters. Under most physical conditions the double inequality $`a\lambda _Dl`$ holds, so that a coarse-grained statistical description of the suspension is clearly warranted. This strategy may be cast in the unifying framework of the density functional theory (DFT) of inhomogeneous fluids , as explained in the following section. A number of key issues, which are presently the object of intense experimental and theoretical scrutiny, will be addressed in subsequent sections of this review. They include the following questions and topics :
* The limitations of the standard Poisson-Boltzmann or mean-field theory of electric double- layers and the importance of spatial correlations between microions. A particularly important question relates to the possibility of an attractive component of the effective interaction between polyions, induced by microion correlations.
* The notion of effective polyion charge (or charge renormalization), as controlled by dissociation equilibria and counterion adsorption (or condensation ), and of charge regulation.
* The relation between effective interactions between polyions, and the phase behaviour of their dispersions.
* The effect of confinement on the effective interactions between colloids, and more generally the influence of electrostatic boundary conditions at interfaces.
* The introduction of solvent granularity into the statistical description of double-layers.
This review will emphasise the more recent experimental and theoretical developments in this very active field, covering mostly work published during the past decade. Earlier work is adequately covered in some previous reviews of the subject . Electrostatic CGS units will be used throughout.
## II Multi-component versus effective one-component description
Consider a suspension of $`N`$ polyions, of radius $`a`$ and charge $`Z`$, in a polar solvent with co and counterions, of radii $`a_\pm `$ and charges $`z_\pm e`$ ; the radius of the solvent molecules is comparable to the microion radii $`a_\pm `$ , i.e. of the order of $`0.1`$ to 0.2 nm. In micellar solutions, 1 nm $`<a<10`$ nm, and $`10<|Z|<100`$, while in most colloidal dispersions, 10 nm $`<a<10^3`$ nm, and $`10^2<|Z|<10^4`$. A statistical description of the highly asymmetric multi-component system is a considerable challenge, which requires some degree of coarse-graining. In most theoretical investigations of micellar or colloidal systems, the solvent is regarded as a mere continuum, characterized by its macroscopic dielectric constant $`ϵ`$. This amounts to a ”primitive model” (PM) level of description, where polyions and microions are assumed to be charged hard spheres (non-spherical polyions will be considered later in this review) interacting via the Coulomb potential $`e^2/ϵr`$, outside the excluded volume range of inter-particle distances $`r`$. It is implicitly assumed that the mesoscopic polyions have the same macroscopic dielectric constant $`ϵ`$ as the solvent, thus avoiding complications due to dielectric discontinuities (e.g. image charges).
Even within the PM, the theoretician is still faced with the polyion-microion asymmetry. In micellar solutions, where the size and charge ratios $`a/a_\pm `$ and $`|Z|/|z_\pm |`$ are roughly a factor of 10, the asymmetry may still be handled within the multicomponent PM level of description, using the theoretical techniques of the theory of classical fluids . In particular, the partial pair distribution functions $`g_{\alpha \beta }(r)`$ (where $`\alpha `$ and $`\beta `$ are species indices) may be calculated from the usual fluid integral equations, including the hypernetted chain (HNC) and mean spherical approximations (MSA) , or their variants and hybrid combinations . Alternatively these correlation functions may be obtained from Monte Carlo (MC) or Molecular Dynamics (MD) simulations of the asymmetric PM. The multicomponent PM point of view is, in practice, limited to micellar systems with $`|Z|<10^2`$, because integral equation closures are increasingly unreliable for asymmetric systems, and their numerical solution tends to become unstable. Similarly simulations become more and more difficult, and are at present limited to $`|Z|60`$ . The multi-component PM description becomes untractable in the much more asymmetric colloidal range, which requires a coarse-grained description, based on effective interactions between polyions, to be discussed below.
Important simplifications occur within the PM in a number of limiting situations.
1. In the limit of infinitely low concentration $`n`$ (number per unit volume) of polyions, only a pair of these, within a solution of microions, need to be considered. The effective interaction between the two polyions (p) is then exactly given by the potential of mean force, $`v_{\mathrm{pp}}(r)=k_\mathrm{B}T\mathrm{log}[g_{\mathrm{pp}}(r)]`$ . The multicomponent problem reduces to that of an electrolyte of microions in the external field due to two fixed polyions, which leads to considerable simplifications in simulations or the numerical solution of integral equations.
2. The problem may be further simplified by taking the limit $`a\mathrm{}`$ ; the initial dispersion then reduces to the much simpler system of an electrolyte confined between two parallel planes carrying a, generally uniform, surface charge $`\sigma `$. The classic problem of two interacting planar double-layers goes back to the work of Gouy , and has been the object of considerable theoretical and experimental work, some of which will be discussed in a subsequent section.
3. In the opposite limit of high concentration, each polyion is, at least temporarily, trapped in the ”cage” formed by neighbouring polyions. This regime is reasonably described by a one-body model, where a single polyion is located at the centre of a Wigner-Seitz cell, together with an inhomogeneous distribution of microions, such that the total charge inside the cell is zero. For practical purposes the cell is chosen of a simple geometry reflecting the shape of the polyion (e.g. a spherically symmetric cell around a spherical polyion), and the electrostatic boundary conditions are chosen such as to mimic the average effect of the surrounding polyions . Despite the considerable simplifications which they imply, cell models allow valuable insight into concentrated dispersions, and will be discussed further on.
A unifying framework for the statistical description of interacting electric double-layers, or, more generally, of inhomogeneous fluids and interfaces, is provided by density functional theory (DFT) .
### A DFT and effective interactions
In view of the considerable polyion-microion asymmetry, it seems natural to combine a discrete representation of the former with a field description of the latter. Let $`\{𝐑_i\}(1iN)`$ denote the positions of the $`N`$ polyions, assumed here to be spherical, and let $`V_N^{\mathrm{dir}}\{𝐑_i\}`$ be their direct interaction energy for a given configuration; $`V_N^{\mathrm{dir}}`$ is, to a good approximation, pairwise additive, with a pair potential $`v_{\mathrm{pp}}(R)`$ including a short-range excluded volume repulsion, the long- range Coulomb repulsion $`Z^2e^2/ϵR`$ and a van der Waals attraction (dispersion force), of intensity characterized by a Hamaker constant . The inhomogenous distributions of co- and counterions, in the ”external” field of the polyions, are characterized by the local densities (or concentrations) $`\rho _\alpha (𝐫)(\alpha =+,)`$. The equilibrium densities satisfy the variational principle :
$`{\displaystyle \frac{\delta \mathrm{\Omega }[\rho _+^{}(𝐫),\rho _{}^{}(𝐫)]}{\delta \rho _\alpha ^{}(𝐫)}}|_{\rho _\alpha ^{}=\rho _\alpha }=0;\alpha =+,`$ (1)
where $`\mathrm{\Omega }`$ is the grand potential functional of the trial densities , given by :
$`\mathrm{\Omega }[\rho _+^{},\rho _{}^{}]=F[\rho _+^{},\rho _{}^{}]{\displaystyle \underset{\alpha }{}}{\displaystyle \mathrm{\Phi }_\alpha (𝐫)\rho _\alpha ^{}(𝐫)d𝐫}`$ (2)
In (2), $`F`$ denotes the intrinsic free energy functional, while $`\mathrm{\Phi }_\alpha (𝐫)=\mu _\alpha \mathrm{\Phi }_\alpha ^{\mathrm{ext}}(𝐫)`$, with $`\mu _\alpha `$, the chemical potential of species $`\alpha `$ ,and $`\mathrm{\Phi }_\alpha ^{\mathrm{ext}}`$, the external potential acting on ions of this species, equal to the sum of the interactions with the $`N`$ polyions :
$`\mathrm{\Phi }_\alpha ^{\mathrm{ext}}(𝐫)={\displaystyle \underset{i=1}{\overset{N}{}}}v_{\mathrm{p}\alpha }(𝐫𝐑_i)`$ (3)
Note that in view of (3), the equilibrium density profiles $`\rho _\alpha (𝐫)`$, and the resulting grand potential $`\mathrm{\Omega }`$, depend parametrically on the polyion configuration. The chemical potentials $`\mu _\alpha `$ are either fixed at some reservoir value (semi-grand canonical ensemble), or determined a posteriori by the canonical ensemble constraints:
$`{\displaystyle \frac{1}{V}}{\displaystyle _V}\rho _\alpha (𝐫)dr=n_\alpha `$ (4)
where $`V`$ is the total volume of the dispersion and $`n_\alpha `$ is the mean (macroscopic) concentration of $`\alpha `$-ions.
Explicit calculations of the density profiles require some approximation for the generally unknown free energy functional $`F`$, which is traditionally split into ideal and excess parts, $`F=F_{\mathrm{id}}+F_{\mathrm{ex}}`$; $`F_{\mathrm{id}}`$ is known exactly :
$`F_{\mathrm{id}}[\rho _+,\rho _{}]=k_\mathrm{B}T{\displaystyle \underset{\alpha =+,}{}}{\displaystyle _V}\rho _\alpha (𝐫)\left[\mathrm{log}\left(\mathrm{\Lambda }_\alpha ^3\rho _\alpha (𝐫)1\right)\right]d𝐫`$ (5)
If $`n_\alpha ^{(0)}`$ denotes the concentrations of some homogeneous reference state (i.e. in the absence of polyions), and $`\rho _\alpha (𝐫;\xi )=n_\alpha ^{(0)}+\xi \mathrm{\Delta }\rho _\alpha (𝐫)(0\xi 1)`$ denotes a continuous set of density profiles, such that $`\rho _\alpha (𝐫;\xi =1)`$ leads back to $`\rho _\alpha (𝐫)`$, the equilibrium density profiles in the presence of polyions, then $`F_{\mathrm{ex}}`$ is given by the formally exact expression:
$`F_{\mathrm{ex}}[\rho _+,\rho _{}]`$ $`=`$ $`F_{\mathrm{ex}}(n_+^{(0)},n_{}^{(0)})+{\displaystyle \underset{\alpha }{}}\mu _\alpha ^{\mathrm{ex}}{\displaystyle \mathrm{\Delta }\rho _\alpha (𝐫)d𝐫}`$ (6)
$``$ $`k_\mathrm{B}T{\displaystyle \underset{\alpha }{}}{\displaystyle \underset{\beta }{}}{\displaystyle _0^1}𝑑\xi (\xi 1)`$ (7)
$`\times `$ $`{\displaystyle 𝑑𝐫d𝐫^{^{}}\mathrm{\Delta }\rho _\alpha (𝐫)c_{\alpha \beta }[\left\{\rho _\alpha (\xi )\right\};𝐫,𝐫^{^{}}]\mathrm{\Delta }\rho _\beta (𝐫^{^{}})}`$ (8)
$`\mu _\alpha ^{\mathrm{ex}}`$ is the excess part of the chemical potential, while the $`c_{\alpha \beta }\left[\left\{\rho _\alpha (\xi )\right\}\right]`$ are the set of direct correlation functions for the inhomogeneous electrolyte with local densities $`\rho _\alpha (𝐫;\xi )`$. For Coulombic systems, these decay asymptotically as $`v_{\alpha \beta }(r)/k_\mathrm{B}T=z_\alpha z_\beta l_\mathrm{B}/r`$ (where $`l_\mathrm{B}=e^2/ϵk_\mathrm{B}T`$ is the Bjerrum length). It proves convenient to subtract this from the $`c_{\alpha \beta }`$, leaving the short- ranged part $`c_{\alpha \beta }^{\mathrm{sr}}`$ . The excess free energy function then splits into a mean-field Coulombic contribution, and a correlation part :
$`F_{\mathrm{ex}}[\rho _+,\rho _{}]`$ $`=`$ $`F_{\mathrm{Coul}}+F_{\mathrm{corr}}`$ (9)
$`=`$ $`{\displaystyle \frac{e^2}{2ϵ}}{\displaystyle d𝐫d𝐫^{^{}}\frac{\rho (𝐫)\rho (𝐫^{^{}})}{|𝐫𝐫^{^{}}|}}+F_{\mathrm{corr}}`$ (10)
where $`\rho (𝐫)=z_+\rho _+(𝐫)+z_{}\rho _{}(𝐫)`$ is the microion charge density. The correlation term is formally given by (8), with $`c_{\alpha \beta }^{\mathrm{sr}}`$ replacing $`c_{\alpha \beta }`$ .
Starting from these exact expressions, there are basically two strategies to proceed with approximations.
1. The most common strategy focusses on the density profiles, and the resulting free energies. $`F_{corr}`$ is approximated by a local density ansatz (LDA) or by some weighted density approximation (WDA) which generally requires the direct correlation functions in a homogeneous reference state as input. Once the variational problem (1) has been solved for the approximate free energy functional, the effective interaction energy between polyions is given by :
$`V_N\left(\{𝐑_i\}\right)=V_N^{\mathrm{dir}}\left(\{𝐑_i\}\right)+\mathrm{\Omega }\left(\{𝐑_i\}\right)`$ (11)
where the equilibrium grand potential accounts for their indirect interactions, induced by the double-layer-forming microions, which have effectively been traced out. It is very important to realize that this contribution depends on the thermodynamic state of the dispersion, and is generally not pair-wise additive.
If correlations are neglected altogether, $`F_{\mathrm{ex}}`$ reduces to the mean-field Coulombic part, and (1) immediately leads to the equilibrium profiles
$`\rho _\alpha (𝐫)=\zeta _\alpha \mathrm{exp}\left\{\left[\mathrm{\Phi }_\alpha ^{\mathrm{dir}}(𝐫)+z_\alpha e\psi (𝐫)\right]/k_\mathrm{B}T\right\}`$ (12)
where $`\zeta _\alpha `$ is the fugacity of species $`\alpha `$ (equal to its reservoir concentration in the absence of correlations), $`\mathrm{\Phi }_\alpha ^{\mathrm{sr}}`$ is the short-range (excluded volume) part of the interaction of an ion $`\alpha `$ with the $`N`$ polyions, and $`\psi (𝐫)`$ is the total electrostatic potential at $`𝐫`$, which satisfies Poisson’s equation :
$`^2\psi (𝐫)={\displaystyle \frac{4\pi e}{ϵ}}\left[\rho ^{\mathrm{ext}}(𝐫)+\rho (𝐫)\right]`$ (13)
with $`\rho ^{\mathrm{ext}}(𝐫)`$ the ”external” charge density carried by the $`N`$ polyions. Eqs (12) and (13) constitute the multi-centre version of mean-field non-linear Poisson-Boltzmann (PB) theory. Despite their apparent simplicity, numerical solution of the PB equations in the presence of $`N`$ polyions is a formidable task. In practice, eq. (13) is often solved in the domain outside the polyions, and the contribution of the latter to the electrostatics is treated as a boundary value problem.
2. A second, more ambitious strategy is to seek information on both the density profiles, and the pair correlations between microions within the double-layers. This may be achieved by relating the direct correlation functions, appearing in (8), to their functional inverses, the total correlation functions $`h_{\alpha \beta }`$ , via the Ornstein-Zernike relations for inhomogeneous fluids , supplemented by an approximate closure relation between the h and c functions. This strategy leads to coupled equations for the density profiles and the pair correlation functions, as in the widely used inhomogeneous hypernetted chain (IHNC) theory . In view of its numerical complexity, this strategy can only be applied to very simple geometries. In practice it is limited to planar geometry and will be discussed in a subsequent section.
### B Cell model
As stated earlier, cell models, involving a single polyion, prove useful to study concentrated dispersions, including colloidal crystals. The paradigm is provided by the much studied case of a globular polyion of radius $`a`$, placed at the centre of a spherical cavity of radius $`R`$ determined by the polyion concentration, i.e. $`R=(3/4\pi n)^{1/3}`$, containing counterions and salt which ensure overall charge neutrality. The density profiles $`\rho _+(r)`$ and $`\rho _{}(r)`$ depend only on the distance $`r`$ from the centre (1d problem). Gauss’s law implies the boundary conditions (b.c.) that the electric field vanish on the surface of the cell, i.e. $`d\psi /dr|_{r=R}=0`$. The field on the surface of the particle (i.e. at $`r=a`$) is determined by the surface charge density $`\sigma `$ . Finally, due to the b.c. at $`r=R`$, the osmotic pressure $`P`$ of the microions is given by :
$`P=k_\mathrm{B}T[\rho _+(R)+\rho _{}(R)]`$ (14)
PB theory, embodied in eqs. (12) and (13), reduces here to a spherically symmetric one-centre problem, giving rise to a simple second order non-linear differential equation for the local potential $`\psi (r)`$ , which is easily solved numerically . Improvements over the mean-field approximation are achieved by including correlations within the LDA or the WDA . The resulting osmotic pressures are consistently lower than those predicted by PB theory, by a factor of three or more for highly charged polyions. This trend is confirmed by MC simulations of the microions in the cell, subjected to the electrostatic potential of the polyion .
Cell model calculations have been extended to other geometries, e.g. rods or platelets in cylindrical cells. They prove very useful in the determination of effective polyion charges, as discussed in the next section.
## III Charge regulation and renormalization
Most published work on electric double-layers is based on the assumption of constant charge on the polyions, e.g. in the form of a constant and uniform surface charge. This is clearly an oversimplification, as would be the other extreme assumption of constant surface potential, for two reasons : the strong coupling between electrostatics and chemical dissociation equilibrium at the surface (charge regulation), and the adsorption and strong physical binding of counterions to the surface, often referred to as counterion condensation, which leads to a reduction of the apparent polyion charge seen at larger distances (charge renormalization).
### A Charge regulation
The bare (or structural) surface charge of most polyions results from the dissociation of functional groups like sulphate or carboxylic groups, the number of which may be measured by titration . Generally the ionization is only partial, and the dissociation equilibrium, governed by a law of mass action, depends on the local ionic environment, e.g. on local salt concentration or pH. Any variation of this environment, linked to the relative motion of neighbouring polyions, will lead to a fluctuation of the surface charge, but for given macroscopic conditions, one may define some average bare charge. From a theoretical point of view, the coupling of the surface chemistry to the local inhomogeneities induced by electrostatics poses a difficult challenge. Early attempts were based on a combination of the law of mass action with PB theory . A more microscopic point of view is adopted in the ”charge regulation primitive model” (CRPM) (see also Ref ), which adds a strong attraction of chemical origin to the short-range part $`v_{\mathrm{p}\alpha }^{sr}`$ of the polyion-microion pair potential, such that :
$`\mathrm{exp}\left\{v_{\mathrm{p}\alpha }^{sr}(r)/k_\mathrm{B}T\right\}`$ $`=`$ $`V_\alpha \delta (rd_\alpha );r<{\displaystyle \frac{1}{2}}(a+a_\alpha )`$ (15)
$`=`$ $`1;r>{\displaystyle \frac{1}{2}}(a+a_\alpha )`$ (16)
This corresponds to an infinitely deep and narrow potential well localized on a sphere of radius $`d_\alpha `$ from the centre of the polyion; in practice $`d_\alpha `$ must be chosen to be smaller than the polyion radius $`a`$, to prevent the chemical binding of the same microion to two polyions.
The pair structure of the model has been analysed by diagrammatic expansion, and by numerical solutions of the HNC equation for small polyions corresponding to mineral oxide particles . The charge regulation strongly affects the effective interaction between spherical particles, to be discussed later.
### B Charge renormalization
The bare (or structural) charge resulting from the dissociation equilibrium is frequently very large, typically $`|Z|10^4`$ for polyion radii $`a10^2`$ nm. To describe electric double-layers near such highly charged polyions, the traditional phenomenological approach is to divide the counterions into two populations. The first includes ions that are tightly bound (or adsorbed) to the surface by the strong electric field $`E=4\pi \sigma /ϵ`$, thus forming a so-called Stern layer of ”condensed” counterions. A rough estimate of the thickness $`\mathrm{\Delta }`$ of a Stern layer is obtained by balancing the electrostatic work $`E\times ze\times \mathrm{\Delta }`$ against the thermal energy $`k_\mathrm{B}T`$. The resulting $`\mathrm{\Delta }=ϵk_\mathrm{B}T/4\pi \sigma ze`$ is of the order of a few Å in water at room temperature for typical surface charge densities and monovalent counterions; thus $`\mathrm{\Delta }`$ is comparable to the size of microions, so that the ”condensate” may be expected to be roughly a monolayer. This monolayer, of opposite sign to the surface charge, strongly compensates the latter, reducing the total polyion charge to an effective charge $`Z^{}`$ significantly smaller than the bare charge $`Z`$. The remaining counterions feel a much reduced ”external” potential and form the ”diffuse” part of the double-layer, which can often be treated within linearized PB theory (LPB).
The phenomenological approach is to consider $`Z^{}`$ as a parameter adjusted to experimental data, e.g. polyion structure factors determined by light scattering , but much recent theoretical effort has gone into determining the effective charge from first principles. Most schemes are based on the observation that the asymptotic behaviour of the potential or density profiles is correctly described by the simple exponential screening predicted by LPB theory. Charge renormalization should account both for non-linearities in the mean-field PB approach and for microion correlations. Early attempts focused on the single polyion problem within the cell model . Within LPB theory the electrostatic potential in a spherical cell is easily calculated to be of the form :
$`\psi (r)=C+{\displaystyle \frac{Ze}{ϵr}}\left(Ae^{\kappa _Dr}+Be^{\kappa _Dr}\right)`$ (17)
where $`\kappa _D=1/\lambda _D`$ is the inverse Debye screening length :
$`\kappa _D=\left(4\pi l_\mathrm{B}{\displaystyle \underset{\alpha }{}}n_\alpha z_\alpha ^2\right)^{1/2}`$ (18)
The integration constants $`A`$ and $`B`$ are determined by the b.c. at $`r=a`$ and $`r=R`$; $`C`$ is conveniently chosen such that $`\psi (r=R)=0`$ . The validity of the LPB approximation is expected to be reasonable on the cell surface, but to be very poor near the polyion surface if the latter is highly charged. The bare charge $`Z`$ appearing in the LPB potential (17) is then adjusted to a (lower) effective value $`Z^{}`$ by requiring that the resulting microion charge density on the cell surface, $`\rho (R)`$, match the corresponding density obtained from a numerical integration of the PB equation under the same b.c. . In view of the latter, and of Poisson’s equation, this means that the resulting LPB potential matches the non-linear PB potential at the surface up to the third derivative. Note that a renormalization of $`Z`$ also implies a renormalization of $`\kappa _D`$ at low (or vanishing) concentration of added salt, due to the overall charge neutrality requirement. In the absence of salt (counterions only), the matching condition also implies equality of osmotic pressures. Since LPB theory underestimates the charge density near the polyion surface, it overestimates $`\rho (r=R)`$ , compared to the PB charge profile, so that $`Z^{}`$ is always lower than $`Z`$. $`Z^{}Z`$ for low polyion charge, then increases with $`Z`$, but saturates for high values of the bare charge, as a consequence of counterion condensation; the degree of counterion condensation is not greatly affected by salt concentration . The saturation value of $`Z^{}`$ is roughly proportional to the polyion radius $`a`$.
The effects of microion correlations on the values of $`Z^{}`$ have been investigated with LDA and WDA versions of DFT, and by MC simulations . Correlations tend to lead to a further decrease of $`Z^{}`$ compared to $`Z`$, and even to a maximum value, followed by an actual decrease of $`Z^{}`$ as a function of $`Z`$ . If the effective charge is obtained by matching the LPB osmotic pressure to ”exact” MC estimates, the resulting $`Z^{}`$ turns out to depend only weakly on the ratio $`a/R`$, i.e. on the concentration of polyions.
Other semi-phenomenological criteria for determining the effective charge are discussed in ref. . The Debye-Hückel-Bjerrum theory, that has proved very successful to describe the phase behaviour of PM electrolytes has been extended to charged colloids . The Bjerrum pairs of simple electrolytes are replaced by clusters involving each a polyion and a variable number of adsorbed counterions, which are in chemical equilibrium with the free, non-adsorbed counterions, assuming a simple form for the internal partition function of a cluster. The distribution of cluster sizes turns out to be sharply peaked, allowing a rather clear-cut definition of $`Z^{}`$ .
A more fundamental approach to a definition of effective charges is based on a careful analysis of the asymptotic decay of density profiles and correlation functions. This may be achieved within the ”dressed ion” reformulation of linear response theory, i.e. of the polarization of an electrolyte by the average electrostatic potential . Within the PM, the point charges of the microions are replaced by a local, spread-out charge distribution incorporating the short-range fraction of the polarization charge around each bare ion, which is thus replaced by a dressed ion. The short-range fraction of the polarization charge is defined in terms of short-range contributions $`h_{\alpha \beta }^{\mathrm{sr}}`$ to the pair correlation functions $`h_{\alpha \beta }(r)=g_{\alpha \beta }(r)1`$ ; the $`h_{\alpha \beta }^{\mathrm{sr}}`$ are related to the short-range parts $`c_{\alpha \beta }^{\mathrm{sr}}`$ of the direct correlation functions, introduced after eq. (8), via coupled Ornstein-Zernike (OZ) relations . This leads to an exact, non-local generalization of the LPB equations, and a residue analysis of the corresponding Fourier transform of the potential $`\psi _i(r)`$ around an ion i leads to the following asymptotic decay for large $`r`$ :
$`\psi _i(r){\displaystyle \frac{z_i^{}e\mathrm{exp}(\kappa r)}{ϵ^{}r}}`$ (19)
This is precisely of the LPB (or Debye-Hückel) form, but with renormalized values $`z_i^{},\kappa `$ and $`ϵ^{}`$ of the valence, inverse screening length and dielectric constant of the solution (as opposed to the pure solvent), which can all be expressed in terms of the $`h_{\alpha \beta }^{\mathrm{sr}}`$ ; the latter may be calculated by supplementing the OZ relations with some approximate closure, e.g. HNC.
In the weak-coupling limit of very low ionic bulk concentration, $`z_i^{}z_i`$, $`\kappa \kappa _\mathrm{D}`$ and $`ϵ^{}ϵ`$. Simple Stillinger-Lovett sum rule considerations show that proper inclusion of the finite size $`a_\alpha `$ of the ions leads to an enhancement of screening, i.e. $`\kappa >\kappa _\mathrm{D}`$ . For sufficiently strong coupling (i.e. at high ionic concentrations), the exponential decay (19) goes over to a damped oscillatory behaviour characterized by a pair of complex conjugate inverse decay lengths $`\kappa `$ and $`\kappa ^{}`$ . The location of the cross-over from monotonous to oscillatory decay depends on the approximate closure.
The ”dressed ion” reformulation for the bulk PM may be extended to the case where one ionic species consists of polyions, and hence $`z_i^{}`$ is to be identified with the effective charge $`Z^{}`$ of the latter . A similar treatment has been applied to planar electric double layers, to determine effective surface charge densities . The latter tend to saturate for monovalent counterions, increasingly so as their concentration increases, while for concentrated solutions of divalent ions, the effective surface charge goes through a maximum before decreasing as the bare surface charge increases. This behaviour is reminiscent of that observed in a spherical cell .
The effective polyion charge $`Z^{}`$ is not easily accessible to experiment. Indirect determinations are based on the assumption of the validity of some simple functional form of the effective pair interaction between polyions, generally a screened Coulomb (or DLVO) potential. The thermodynamic properties and pair structure of a system of polyions interacting via this potential can be determined accurately from fluid integral equations or simulations . The effective charge $`Z^{}`$ is then adjusted to provide the best theoretical fit to experimental data, e.g. the polyion pair structure as measured by scattering experiments, or the freezing line as a function of salt concentration and bare particle charge . The latter measurements confirm the saturation predicted by the PB cell model , but should be re-analysed since the underlying simulation data do not take into account an important ”volume” contribution to the free energies, which strongly affects the phase diagram at low salt concentration .
Recent experiments, carried out at very low electrolyte concentration, on silica particles, with weakly dissociating silanol surface groups, and on latex particles, with strongly dissociating sulfonic acid groups, were able to measure simultaneously the bare polyion charges by conductometric titration, and the effective charges by conductivity measurements . The bare charge was controlled by varying the amount of added NaOH. Under these experimental conditions, $`Z^{}`$ was not observed to saturate, but to follow very nearly a square root law $`Z^{}\sqrt{Z}`$ for both colloidal systems.
## IV Planar geometry
The simplest, and most widely studied double-layer geometry is that of uniformly charged, infinite planes in an ionic solution . These planes may represent stacks of thin charged lamellae, like clay platelets or rigid membranes, or correspond to the surfaces of spherical colloidal particles separated by a distance $`ha`$, so that the curvature of the facing surfaces may be neglected in first approximation. In the former situation, one may consider the simplified model of a single charged plane within a slab of thickness h equal to the mean spacing between lamellae in the stack; this would be the 1d equivalent of the cell model. The two situations are shown schematically in Figure 1; $`\sigma `$ denotes the surface charge density, $`ϵ`$ is the macroscopic dielectric constant of the solvent (within a PM representation) and $`ϵ^{^{}}`$ is the dielectric constant of the colloidal particles. When $`ϵ^{^{}}ϵ`$, which is the rule rather than the exception, the dielectric discontinuity at $`z=\pm h/2`$ must be properly incorporated in the boundary condition, e.g. by the introduction of electrostatic image charges . In most published theoretical work the assumption $`ϵ^{^{}}=ϵ`$ is made for the sake of simplicity.
Microion density profiles $`\rho _+(z)`$ and $`\rho _{}(z)`$ depend only on $`z`$, and the intrinsic free energy functional per unit area perpendicular to the $`z`$ axis is still given by eqs (5)-(10), with the volume integrals replaced by 1d integrals over the interval $`[\frac{h}{2},\frac{h}{2}]`$ . Fixed charge (or Neumann) boundary conditions imply that the potential satisfy :
$`{\displaystyle \frac{d\psi }{dz}}|_{z=\pm \frac{h}{2}}=\pm {\displaystyle \frac{4\pi \sigma }{ϵ}}`$ (20)
The key physical quantity, which is in principle measurable using a surface force apparatus , is the force per unit area acting between two charged planes and their associated electric double-layers. This force per unit area is the net osmotic pressure (or disjoining pressure) $`\mathrm{\Delta }P=PP_{\mathrm{bulk}}`$, where $`P`$ is the osmotic pressure exerted by the microions between the charged plates, and $`P_{\mathrm{bulk}}`$ is the pressure of the electrolyte in the reservoir which fixes the chemical potentials of the microions. $`P`$ may be calculated via the mechanical route, by averaging the local stress tensor, or by the thermodynamic route, by differentiation of the total free energy, or grand potential, with respect to $`h`$. Both routes lead to a natural separation of the pressure $`P`$ into kinetic, electrostatic and collisional parts :
$`P=P_{\mathrm{kin}}+P_{\mathrm{el}}+P_{\mathrm{coll}}`$ (21)
where
$`P_{\mathrm{kin}}=k_\mathrm{B}T{\displaystyle \underset{\alpha }{}}\rho _\alpha (z)`$ (22)
and explicit expressions for $`P_{\mathrm{el}}`$ and $`P_{\mathrm{coll}}`$ in terms of density profiles and inhomogeneous pair correlation functions are given in ref. . Note that, while each of the terms in eq. (21) depends on the coordinate $`z`$, their sum must be independent of z for mechanical equilibrium. In practice, simplified expressions are obtained by calculating $`P`$ in mid-plane ($`z=0`$ in Fig. 1) or at the surface of the charged planes. In the latter case the expression for $`P`$ reduces to the simple contact theorem :
$`P=k_\mathrm{B}T{\displaystyle \underset{\alpha }{}}\rho _\alpha \left(z=\pm (h/2a_\alpha )\right){\displaystyle \frac{2\pi \sigma ^2}{ϵ}}`$ (23)
This expression is not very convenient for simulation purposes, since $`P`$ appears as a difference of two large numbers, the first of which involves the contact value of the counterion density profile, affected by large numerical uncertainties. The calculation at midplane is then preferable, since the $`\rho _\alpha `$ are expected to be close to their bulk values there, but the average force between microions on both sides of this plane must be evaluated .
### A Poisson-Boltzmann approximation
The mean-field approximation, where correlations between microions, due to excluded volume and Coulomb interactions, are neglected, is well understood since the pioneering work of Gouy ; density profiles and forces are known analytically, or given by simple quadratures . The pressure is given throughout by :
$`P=k_\mathrm{B}T{\displaystyle \underset{\alpha }{}}\rho _\alpha (z){\displaystyle \frac{ϵ}{8\pi }}\left[E(z)\right]^2`$ (24)
where $`E(z)=d\psi (z)/dz`$ is the mean local electric field. This is most easily evaluated in the mid- plane, where $`E=0`$ by symmetry, i.e. $`P=k_\mathrm{B}T_\alpha \rho _\alpha (0)`$ , showing that the force between equally charged plates is always repulsive within PB theory.
In the limit of low surface charge and electrolyte concentration, PB theory may be linearized, and the disjoining pressure reduces to :
$`\mathrm{\Delta }P=2\sigma ^2e^{\kappa _\mathrm{D}h}`$ (25)
A charge renormalization to account for non-linearities, very similar to the procedure within the cell model described earlier, leads to a much reduced effective surface charge $`\sigma ^{}`$, and to a saturation for large bare charges $`\sigma `$ .
### B The effect of microion correlations
While PB theory always predicts a purely repulsive interaction between planar electric double-layers, it was suggested by Oosawa as early as 1968 that the force might turn attractive if microion correlations were taken into account. The first convincing evidence for double-layer attraction came from MC simulations which indicated large deviations from the predictions of PB theory in the presence of divalent counterions . This early prediction of an attractive minimum in the effective force between planar double-layers at short separations h was confirmed by subsequent simulations in the presence of salt ($`1:2`$ and $`2:2`$ electrolytes) and polyvalent counterions , by inhomogeneous HNC calculations , and an improved version based on the reference HNC (RHNC) closure , and by LDA and WDA versions of DFT .
The reduction of the double-layer repulsion and eventual correlation-induced attraction may be understood by considering the contact theorem (23), or the expression for the pressure in midplane. At contact the correlations keep the counterions in the first layer apart, thus limiting the piling up following from the uncorrelated mean field treatment, and reducing the contact value of the counterion density; according to eq. (23) this leads to a lowering of the pressure, which can even become negative (corresponding to an effective attraction), due to the negative Coulombic contribution. The lowering of the counterion concentration immediately at contact leads to an enhanced attraction of the next layers of counterions to the surface, so that correlations enhance the counterion density in that region. This enhancement in turn entails a depletion of the counterion density in the midplane, while electrostatic correlations make a negative contribution ($`P_{\mathrm{el}}`$ in eq. (21)) which generally dominates the smaller, positive collisional part. Overall these various effects lead to a much reduced $`\mathrm{\Delta }P`$ compared to the repulsive PB interaction. The attraction observed for very short separations may be explained by a 2d lattice-like structuring of the adsorbed counterions, when the counterion patterns on opposite surfaces are shifted relative to each other to minimize the electrostatic interactions .
The reason why an attraction is in general observed only with divalent (e.g. $`\mathrm{Ca}^{++}`$) counterions is that the entropic cost associated with their ”condensation” near the surface is smaller than for monovalent ions, since only half as many counterions are needed to provide the same electrostatic shielding. Another related aspect of microion correlations is the phenomenon of charge reversal (or inversion). The apparent charge density of the surface placed at $`z=h/2`$, seen at the abscissa $`z`$ is :
$`\sigma ^{}(z)=\sigma +e{\displaystyle _{\frac{h}{2}}^z}{\displaystyle \underset{\alpha }{}}z_\alpha \rho _\alpha (z^{^{}})dz^{^{}}`$ (26)
Within PB theory $`\sigma ^{}`$ never changes sign, but in the presence of correlations this may occur at a critical z, beyond which $`\sigma ^{}`$ is of opposite sign to $`\sigma `$ , due to an over-compensation of the bare surface charge by the counterions. Thus a test-charge will be attracted, rather than repelled, by a plane carrying a surface charge $`\sigma `$ of the same sign, beyond a critical separation. Charge inversion can only occur when the microion density profiles have a non-monotonic behaviour. The over- compensation of the surface charges of two parallel plates is enhanced when the counterions are linked together to form polyelectrolyte chains, due once more to the reduction of entropic repulsion .
Charge inversion and effective attraction between platelets may even be observed in the case of monovalent counterions, provided the coions are larger than the counterions .
Attractive forces between charged surfaces immersed in aqueous solutions of $`\mathrm{CaCl}_2`$ and $`\mathrm{Ca}(\mathrm{NO}_3)_2`$ have been measured directly with a surface force apparatus . The existence of correlation-induced attraction between charged surfaces of equal sign has important consequences on colloid stability, and provides a plausible explanation for the common observation that the addition of multivalent ions to solutions or suspensions of polyions often leads to precipitation.
## V Spherical polyions
Aqueous solutions and dispersions of spherical or quasi-spherical polyions, including globular proteins and charge-stabilized colloidal particles, like silica mineral particles or polymer latex spheres, have been thoroughly studied experimentally and theoretically for many decades. Much of the interest stems from the realization that such colloidal systems exhibit a phase behaviour reminiscent of that of simple molecular systems, albeit on very different scales. A particularly attractive feature of colloidal dispersions is that the effective interactions between particles may be tuned by the experimentalist, e.g. by varying the concentration of added electrolyte, thus providing an additional handle on phase behaviour. The beautiful early microscopy observations by Hachisu et al. unambiguously showed the existence, at low volume fractions, of colloidal crystals and crystalline alloys, which Bragg-reflect visible light and are, in particular, responsible for the iridescence of opals.
The coexistence of ordered crystalline and disordered fluid phases has been monitored by very careful experiments using confocal laser scanning microscopy, which allows the direct observations of ordered and disordered configurations, and ultra-small-angle X-ray scattering (USAXS), which provides statistically averaged information on ordering of colloids via the static structure factor $`S(k)`$ . These measurements provide three-dimensional phase diagrams by varying the colloid concentration $`n`$ or packing fraction $`\eta =4\pi na^3/3`$, the effective colloid charge $`Z^{}`$ (as determined by conductivity measurements), and the monovalent salt concentration $`n_s`$, equal to the concentration $`n_{}`$ of coions, typically in the range $`10^610^5`$M. The observed phase diagrams exhibit a striking re-entrant behaviour, e.g. for fixed $`n`$ and $`n_s`$, the colloidal dispersion crystallizes into an ordered BCC lattice upon increasing $`Z^{}`$, as one might expect, but remelts into a disordered fluid phase upon further increase of $`Z^{}`$, . This re-entrant behaviour may be qualitatively understood in terms of an increase of the total ionic strength, linked to an increase of the counterion concentration $`n_+`$ with $`Z^{}`$, which enhances the screening power of the microions and hence reduces the range of the screened Coulomb repulsion between polyions. However, the standard DLVO representation of the latter fails to provide a quantitative explanation of the experimental data .
Another striking feature of low concentration dispersions $`(\eta <0.05)`$ , reported by the Kyoto School, is the appearance of strongly inhomogeneous patterns in highly deionized samples. USAXS measurements and confocal laser scanning microscopy show clear evidence of regions of relatively high colloid concentration, which may be crystalline , amorphous or fluid in character, coexisting with regions of extremely low concentration, or voids. Apart from direct visual observation, the existence of voids, and the fraction of the total volume occupied by them, may be inferred from a comparison between the mean inter-particle spacing, $`d_nn^{1/3}`$, calculated under the assumption of a homogeneous dispersion, and the spacing $`d_x`$ deduced from the position of the main X-ray diffraction peak in the structure factor $`S(k)`$. The observation that $`d_x`$ is significantly less than $`d`$ points to the existence of voids occupying a fraction $`f=1(d_x/d)^3`$ of the total sample volume. The considerable literature on the subject is summarized in ref . The existence of voids, and the related, but controversial , observation of a complete separation between a high concentration colloidal ”liquid” phase, and a much more dilute ”gas” phase , are claimed to be evidence for a long range attractive component of the effective pair interaction between colloidal particles carrying charges of equal sign. The same attraction is also invoked to explain the re-entrant liquid-solid coexistence described earlier .
However recent direct measurements of the effective interaction between pairs of colloidal particles show no evidence of an attraction, at least at low concentration. Such direct measurements are based on video microscopy. In the simplest method, a large number of instantaneous configurations of a dilute suspension are recorded, and the colloid-colloid pair distribution function $`g(r)`$ is calculated from the measured distribution of inter-particle distances . The effective pair potential coincides with the potential of mean force in the low concentration limit and is hence given by :
$`v(r)=k_\mathrm{B}T\mathrm{log}[g(r)]`$ (27)
An alternative method follows the relative Brownian motion of a pair of colloidal particles released from initial positions where they were localized by optical tweezers . Both sets of measurements confirm that the effective pair potential is purely repulsive, and may be fitted to a screened Coulomb (DLVO) form by adjusting the effective charge $`Z^{}`$ and the inverse screening length $`\kappa _\mathrm{D}`$, compatible with estimated salt concentrations. However this same DLVO potential, used within a harmonic approximation, seems to be unable to reproduce the measured bulk modulus of a colloidal crystal at much higher concentrations , although the difference between experiment and theory may be linked to the omission of the ”volume” contribution to the effective interaction energy between polyions, to be discussed in the following section.
### A DLVO theory revisited
The effective pair potential between spherical polyions, first derived by Derjaguin and Landau, and Verwey and Overbeek , is easily recovered within the framework of DFT. Upon tracing out the microion degrees of freedom, the effective interaction energy between polyions is given by eq. (11), once the grand potential $`\mathrm{\Omega }`$ has been determined from the variational principle (1). The required intrinsic free energy functional is defined by eqs. (5) and (10). Neglecting $`F_{\mathrm{corr}}`$ in the latter amounts to mean-field theory, and the Euler-Lagrange equation associated with the variational principle (1) leads back to the PB equation (13) for the local electrostatic potential $`\psi (𝐫)`$ . Since the external charge associated with $`N`$ polyions is a multi-centre distribution, the numerical task of solving the PB equation for $`N`$ interacting polyions is a formidable one, which can only be handled numerically using advanced optimization techniques, coupling Molecular Dynamics (MD) simulations with DFT .
A more phenomenological approach may be adopted, whereby the bare polyion charge is reduced to a considerably lower effective value by a Stern layer of tightly bound counterions. The coupling between the polyions and the remaining microions (forming the so-called ”diffuse” double-layer) is accordingly strongly reduced, so that the corresponding density profiles $`\rho _\alpha (𝐫)`$ vary much more smoothly in the vicinity of the polyion surfaces. In that case the integrands of the ideal contributions to F in eq. (10) may be expanded to second order in the deviations $`\mathrm{\Delta }\rho _\alpha (𝐫)=\rho _\alpha (𝐫)n_\alpha `$ of the local densities from their bulk values, i.e. :
$`F_{\mathrm{id}}{\displaystyle \underset{\alpha =\pm }{}}\left\{F_{\mathrm{id}}(V,T,n_\alpha )+{\displaystyle \frac{k_\mathrm{B}T}{2n_\alpha }}{\displaystyle _V}[\mathrm{\Delta }\rho _\alpha (𝐫)]^2d𝐫\right\}`$ (28)
where $`F_{\mathrm{id}}(V,T,n_\alpha )`$ is the Helmholtz free energy of an ideal gas of density $`n_\alpha `$ . The free energy functional defined by (5) and (8), neglecting $`F_{\mathrm{corr}}`$, is now a quadratic functional of the local densities, and the resulting Euler-Lagrange equation reduces to a linear multi-centre PB (LPB) equation for the local potential $`\psi (𝐫)`$ , in the familiar form :
$`\left(^2\kappa _\mathrm{D}^2\right)\psi (𝐫)={\displaystyle \frac{4\pi }{ϵ}}\rho ^{\mathrm{ext}}(𝐫)`$ (29)
This linear equation is easily solved by Fourier transformation in the bulk, with the boundary condition that the potential and its gradient vanish at infinity. Since $`\rho ^{\mathrm{ext}}`$ is a linear superposition of contributions from the $`N`$ polyions, the same is true of $`\psi (𝐫)`$ and the resulting $`\rho _\alpha (𝐫)`$ . Since microions cannot penetrate the spherical polyions, the excluded volume condition, $`\rho _\alpha (𝐫)=0;|𝐫𝐑_i|<a(1iN)`$ , must be imposed via a constraint, or by the use of a polyion- microion pseudo potential . Note that, strictly speaking, the linear superposition of densities only holds provided the electric double-layer associated with a given polyion does not overlap a neighbouring polyion (”weak overlap approximation”). This is not a significant limitation except at high polyion packing fraction. The excluded volume condition leads to an additional polyion charge renormalization such that the effective charge $`Z^{}`$ is multiplied by the DLVO factor $`e^{\kappa _\mathrm{D}}/(1+\kappa _\mathrm{D}a)`$; the product will henceforth be designated by $`Z^{}`$. The local potential is finally of the form :
$`\psi (𝐫)={\displaystyle \underset{i=1}{\overset{N}{}}}\psi ^{(i)}(𝐫)={\displaystyle \underset{i=1}{\overset{N}{}}}{\displaystyle \frac{Z^{}e}{ϵ}}{\displaystyle \frac{\mathrm{exp}\left(\kappa _\mathrm{D}|𝐫𝐑_i|\right)}{|𝐫𝐑_i|}}`$ (30)
When (30) and the resulting $`\rho _\alpha (𝐫)`$ are substituted into the quadratic free energy functional, the effective interaction energy (11) between the $`N`$ polyions reduces to :
$`V_N(\{𝐑_i\})=V_0+{\displaystyle \underset{i<j}{}}v\left(|𝐑_i𝐑_j|\right)`$ (31)
where $`V_0(T,n,n_+,n_{})`$ is a state-dependent ”volume” term , the detailed expression of which is given in and . This term has been generally overlooked, but it strongly influences the calculated phase behaviour of charge-stabilized colloidal dispersions, as discussed later. The main contribution to the volume term is the sum of cohesive electrostatic free energies arising from the electrostatic attraction between each polyion and its associated double layer of opposite charge. The physical interpretation of $`V_0`$ is discussed in detail in ref. . The effective pair potential is precisely of the well known DLVO form , namely:
$`v(R)={\displaystyle \frac{Z^2e^2}{ϵ}}{\displaystyle \frac{e^{\kappa _\mathrm{D}R}}{R}}`$ (32)
Note that the pair-wise additivity in eq. (31) is a direct consequence of the quadratic nature of the approximate free energy functional; the exact ideal contribution (5) to the free energy would lead to many-body effective interactions.
There are alternative contraction procedures to reduce the initial asymmetric, multicomponent PM to an effective one-component system, based on the OZ equations and approximate closures, like the MSA, for the partial pair distribution functions $`g_{\alpha \beta }(r)`$ . The simple assumption of the asymptotic behaviour $`c_{\alpha \beta }(r)v_{\alpha \beta }(r)/k_\mathrm{B}T`$ of the direct correlation functions is sufficient to show that the effective pair interaction is asymptotically of the DLVO form (32), and reduces to (32) in the appropriate limits. Corrections to the DLVO potential at short range may be obtained from numerical solutions of the fluid integral equations, and confirm the purely repulsive nature of the effective pair potential at low colloid concentration, at least when all microions are monovalent .
The phase diagram of a monodisperse system of particles interacting via the DLVO pair potential (32) has been determined by extensive MC simulations as a function of colloid and salt concentrations, the latter determining the inverse Debye screening length $`\kappa _\mathrm{D}`$ . As expected from the purely repulsive nature of the effective pair potential, the phase diagram exhibits a first- order transition between a single (disordered) fluid phase and (ordered) FCC or BCC crystalline phases, BCC being the stable phase in the $`\kappa _\mathrm{D}0`$ limit, which corresponds to the widely studied ”one-component plasma” (OCP) model . However the coarse-graining procedure leading from the initial multicomponent polyion/microion ”mixture” to the effective one-component system also introduces the volume term $`V_0`$ into the effective interaction energy between polyions. Since $`V_0`$ is a non-linear function of the polyion density $`n`$, due to the overall charge neutrality condition $`nZ^{}+n_+z_++n_{}z_{}=0`$, this term must be included in the free energy for a proper determination of the phase behaviour . At high concentrations of added salt, the variation of $`V_0`$ with $`n`$ is sufficiently slow not to affect the phase behaviour significantly. At salt concentrations lower than $`10^5`$M, the volume term drives a van der Waals-like instability of the fluid, which separates into two phases of very different colloid concentrations, reminiscent of the ”gas” and ”liquid” phases of ordinary molecular fluids. Depending on the effective colloid surface charge, this phase separation is completely disconnected from the freezing transition driven by the repulsive pair interaction (32), and exhibits upper and lower critical points (re-entrant behaviour) in the $`nn_s`$ plane (see Figure 2); or it merges with the freezing line, leading to a considerable broadening of the fluid-solid coexistence region, and the possibility of upper and lower triple points,as seen in Figure 3. The possibility of a fluid miscibility gap in polyelectrolytes had been conjectured already in 1938 by Langmuir, who referred to it as ”unipolar coacervation” , and the complex phase scenario of references has recently been confirmed by an extension of Debye-Hückel theory to the highly asymmetric PM . These calculations provide a natural explanation for the observed phase behaviour of charge-stabilized colloids, including the formation of voids , without the assumption of a long-range attractive component in the effective pair interaction which is frequently made , but lacks a firm theoretical basis.
### B Beyond DLVO
It was shown earlier that the effective force between two uniformly charged plates may turn attractive at short distances if correlations between microions are properly accounted for. A similar effect is expected to hold between two spherical polyions, as first conjectured for two isolated spheres by Patey on the basis of numerical solutions of the HNC equation. When two spheres are sufficiently close, such that the shortest separation h between their surfaces is much less than their radius $`a`$, the problem reduces essentially to that of two charged planes discussed earlier : a short-range correlation- induced attraction may be expected for divalent counterions. A collective polarization mechanism of adsorbed counterions operates in the opposite limit $`ha`$, and leads to a ”doubly screened” attractive component, proportional to $`e^{2\kappa _\mathrm{D}R}`$ , which is always dominated by the DLVO repulsion . It is important to stress that effective attractions between spherical polyions require a finite concentration of the latter in the absence of salt, i.e. when only counterions are present . No such restriction holds at finite salt concentration, where the problem of two isolated polyions is a meaningful limit.
A similar ”doubly screened” attraction between spherical polyions occurs in the charge regulation PM introduced earlier (c.f. eq. (12)); explicit HNC calculations have been carried out for small mineral oxide particles . Note that the effective interaction between two isolated colloids is purely attractive (for $`R>2a`$) in the case of a symmetrical adsorption of co and counterions ($`V_+=V_{}`$ in eq. (12)), because the colloidal particles are then, on average, neutral.
### C Computer simulations
Density-functional molecular dynamics simulations were performed for monovalent counterions without and with added salt. The effective interaction between the polyions was found to be repulsive, and in quantitative agreement with DLVO theory, provided the charge was renormalized according to the cell model prescription incorporating the excluded volume correction . This approach includes counterion correlations, but is not completely equivalent to a full simulation of the PM since an approximate polyion-counterion pseudopotential is used and the exact density functional of the electrolyte is unknown. Due to the high charge asymmetry between poly- and microions, a full simulation of the PM requires the inclusion of many counter- and salt ions and is not feasible on present-day computers. There are two ways to escape from this: either by considering only a small number of polyions or by reducing the charge asymmetry significantly. By considering only two polyions within the PM, the effective pair potential can be calculated within a simulation by averaging over the microscopic ions, while keeping the polyion positions fixed. For monovalent microions, the simulation yields repulsive forces correctly described by DLVO theory both with and without added salt. For divalent counterions, attractive forces are obtained, which may be attributed to counterion correlations and Coulomb depletion , resulting from a depletion zone of counterions between two nearly touching polyions due to the strong counterion repulsion. A typical snapshot of the counterions around two fixed polyions and the averaged force versus polyion distance are shown in Figure 4. A system of three polyions was simulated in order to extract effective triplet interactions which were found, however, to be small with respect to the pairwise part .
On the other hand, there are quite a number of full simulations of the PM with a reduced charge asymmetry corresponding to micellar rather than colloidal polyions. Most of the work is summarized in a review by Vlachy . As for more recent work, the role of salt valency has been investigated in some detail . Moreover Linse and Lobaskin pushed the charge asymmetry to $`Z:z_+=60:1,\mathrm{\hspace{0.33em}60}:2,\mathrm{\hspace{0.33em}60}:3`$ using efficient cluster move algorithms and treating 80 polyions simultaneously . For the $`60:3`$ asymmetry clear indications for a phase separation were reported which can be attributed to an effective attraction between polyions.
## VI Polyions in confined geometries
The bulk behaviour of polyionic dispersions is expected to change significantly in the presence of neutral or charged interfaces confining the dispersion to a restricted volume. Such confinement occurs naturally in the vicinity of the sample container, which may be approximated locally by an infinite plane restricting the suspension to a half space. Many experiments are carried out in a slit geometry, with the suspension confined between parallel planes (e.g. glass plates); if the interplanar spacing is comparable to the diameter $`2a`$ of the polyions, the suspension behaves like a quasi-two- dimensional (2d) many-body system exhibiting interesting phase behaviour .
The electrostatic effects of confinement are essentially threefold. (i) The interface between two different dielectric media, e.g. the suspension and the glass of the container wall, must be characterized by appropriate boundary conditions; in particular the dielectric discontinuity implies the presence of electrostatic image charges. (ii). The confining surfaces lead to a reduction in screening power of the electrolyte, which ceases to be exponential in directions parallel to the surface . (iii) If the surfaces are charged, they attract or repel the polyions and microions, and release additional counterions into the suspension.
These effects will modify the electrostatic interactions between double-layers associated with polyions close to the confining walls. This has been clearly demonstrated by a number of direct measurements of the effective forces between spherical polyions confined near a glass wall, or in narrow slits, using digital video microscopy techniques similar to those mentioned earlier for measurements in the bulk. Measurements of the colloid pair distribution function $`g(r)`$ carried out at sufficiently low concentration allow the effective pair potential to be extracted directly from eq. (23) , while data collected from more concentrated samples, where the effective potential $`v(r)`$ is expected to differ significantly from the potential of mean force defined by eq. (27), require an elaborate inversion procedure based on fluid integral equations (properly adapted to a 2d geometry) , or on an iterative method involving computer simulations . These experiments leave little doubt for the existence of an attractive component of the effective polyion pair potential at low ionic strength (i.e. for salt concentrations estimated to be below about $`10^5`$M), when the polyions are highly confined, as achieved in narrow slits , or by localizing them close to a single plane by optical tweezers . The observed attractive well in $`v(r)`$ is relatively shallow (depth of the order of $`k_\mathrm{B}T`$), and long-ranged; typically $`v(r)`$ becomes negative for interparticle distances considerably larger than the particle diameter $`2a`$. The experiments carried out for two particles close to a charged plate with the centre-to-centre vector parallel to the plate, clearly show that the attractive well disappears when the particles are moved away from the plate; the measured $`v(r)`$ returns to its bulk DLVO form (32) when the distance to the plate is several particle diameters . An additional twist for suspensions of colloidal particles confined to a slit is provided by the introduction of quenched obstacles, in the form of larger, charged colloids, which form a disordered porous 2d matrix. As the concentration of such obstacles is increased, the effective interaction potential between the smaller colloids is observed to develop a second attractive component extending to larger inter-particle separations . Closely related observations of metastable colloidal crystallites suggest that facets of such crystals behave very much like charged plates, inducing effective attractions between nearby polyions due to many-body interactions.
The DFT formulation of DLVO theory in the bulk has been extended to the case of confined particles. The electrostatic potential due to a point polyion near a planar surface separating two media of dielectric constants $`ϵ`$ and $`ϵ^{^{}}`$ was determined within LPB theory by Stillinger as early as 1961 . The case of two spherical polyions of finite radius a near an uncharged planar surface was examined within LPB theory in . Since microions cannot leak beyond the dividing surface, the spherical symmetry of the electric double layers around each of the polyions is broken. The effective repulsion between the polyions is enhanced, and decays like $`1/R^3`$, where $`R`$ is the separation parallel to the surface, as expected from the considerations in .
If the confining surface carries an electric charge density $`\sigma `$ , as would be the case for a glass plate in contact with an aqueous dispersion, the planar electric double-layer building up near the surface will modify the local co and counterion concentrations in the vicinity of nearby polyions, and hence influence their mutual interaction. Considering the charged wall as an additional colloidal particle of infinite radius, it is clear from the inherent pair-wise additivity that the quadratic free energy functional, defined by eqs (28) and (10) (with $`F_{\mathrm{corr}}=0`$), which underlies LPB theory, will not affect the effective interaction between two polyions near the wall. Clearly the intrinsic three-body nature of the problem must be taken into account. The minimal theory to achieve this is to carry the expansion (28) one order further, allowing for a direct coupling between the three double-layers associated with the polyions and surface. Application of DFT perturbation theory leads to the prediction of a long-range attraction between polyions, decaying like $`1/R^3`$ along the surface, and exponentially normal to the surface, for finite concentrations of polyions, to allow for an imbalance between co and counterion concentrations . Both the location and depth (typically of the order of $`k_\mathrm{B}T`$) of the predicted potential well are in semi-quantitative agreement with direct experimental measurements . The attraction predicted within mean-field theory occurs when the surfaces of the spherical polyions are several Debye lengths apart, and has nothing to do with the correlation-induced-attraction at very short range predicted in the bulk, and which is significant only for multivalent counterion .
The generalization of DLVO theory may also be extended to slit and pore geometries. The quasi-1d confinement in the latter case leads to an even slower decay of the long- range attraction. The existence of an effective attraction between two like-charged polyions confined by a charged cylindrical surface was also predicted on the basis of numerical solutions of the non-linear PB equation , but it has since been proved rigorously that PB theory can only lead to purely repulsive effective interaction between two identical polyions confined to a cylinder of arbitrary cross-section , whatever the boundary conditions at the charged surfaces. This proof does not invalidate the results of the perturbation theory , since charged renormalization (which accounts for short-range correlations), followed by an LPB perturbation treatment of the non-adsorbed counterions is not a mere approximation to non-linear PB theory, which is, moreover, unphysical near highly charged surfaces.
A direct simulation of the PM for confined colloids is only feasible for low surface charges and small particle sizes . In this regime, the attraction between colloids is strongly enhanced by the presence of a charged wall with respect to that in the bulk . Also the wall-particle interaction was found to be attractive for divalent counterions.
## VII Rod-like polyions
There is a large variety of rod-like polyions ranging from synthetic suspensions to biological macromolecules. These include colloidal $`\beta `$-FeOOH , imogolite , boehmite , polytetrafluoroethylene , ellipsoidal polystyrene latex particles , cylindrical micellar aggregates as well as virus solutions from tobacco-mosaic (TMV) , or bacterial fd and Pf1 viruses . Another motivation to study rod-like colloids comes from the self-assembly of charged stiff biopolymers such as DNA strands , F-actin fibers and microtubules which constitute rod-like polyions on a supramolecular rather than on a colloidal length scale.
The simplest model system is the electric double layer around a single infinitely long cylindrical polyion of radius $`a`$ which is homogeneously charged with line charge density $`\lambda `$. LPB theory leads to the electric potential $`\mathrm{\Psi }(r)=\lambda K_0(\kappa _\mathrm{D}r)/aK_1(\kappa _\mathrm{D}a)ϵ`$ where $`r`$ is the distance from the rod axis and $`K_0(x),K_1(x)`$ are Bessel functions of imaginary argument. Consequently, LPB theory predicts the effective interaction $`v(r)`$ between two parallel rods per unit length to be of the following form:
$`v(r)={\displaystyle \frac{\lambda ^2}{(\kappa _\mathrm{D}aK_1(\kappa _\mathrm{D}a))^2ϵ}}K_0(\kappa _\mathrm{D}r)`$ (33)
where the factor $`1/[\kappa _\mathrm{D}aK_1(\kappa _\mathrm{D}a)]`$ is the excluded volume correction. Contrarily to the three-dimension al case, PB theory can be solved analytically within a cylindrical cell around the charged rod for vanishing salt concentration . The analytical solution for the counterion density field around the colloidal rod exhibits a strong peak at contact providing a theoretical framework of Manning counterion condensation . As for spherical macroions, one cannot extract an effective interaction from the PB solution for a single double-layer directly. One possibility is to proceed as for the spherical PB cell model. Matching the LPB theory at the cell boundary , one derives the effective interaction potential $`v(r)=\lambda ^2/ϵK_0(\kappa r)`$, where $`\kappa `$ and $`\lambda ^{}`$ are renormalized according to the cell model prescription.
A more direct theory, which requires a larger numerical effort, is to solve the two-dimensional PB theory for two parallel rods in a slit geometry and extract the effective force from there. In a recent numerical study no attraction was found, consistent with the exact result in three dimensions.
Beyond the PB level of description, computer simulations for the effective interaction between two parallel, homogeneously charged rods have been performed for divalent counterions and no added salt . An effective attraction between the rods was found. The question is whether the attraction is due to correlations or to fluctuations . If the former is true, the attraction should increase with decreasing temperature and may even persist for zero temperature, while fluctuation- induced attraction should increase with temperature. On the basis of a simple model , Levin and coworkers gave significant support for a correlation picture. This can be intuitively understood as the gain in electrostatic energy upon bringing together staggered arrays of adsorbed counterions . Fluctuation and polarization effects may still play an important additional role. For parallel rods on a triangular lattice, a negative pressure was found by computer simulation with divalent counterions , which can be attributed to some attractive component in the effective interactions.
Qualitatively, the situation closely resembles the case of spherical macroions. But there are also differences. First, Manning condensation theory applies only to rod-like polyions. A possible mechanism of attraction of nearly touching rods by sharing Manning-condensed counterions was proposed, which is only possible for rods . Secondly, the Coulomb depletion mechanism for attraction was found to be irrelevant for parallel rods.
For finite rod lengths and arbitray orientations of the rods, the effective interaction between two rods will depend both on the rod orientations and on their center-of-mass separation. Starting from point charges distributed along the rods, LPB theory results in a Yukawa-segment model which was confirmed by density-functional Molecular-Dynamics simulations for monovalent counterions , provided the charge and the screening constant $`\kappa `$ are renormalized according to the cell model prescription. The associated phase diagram for parameters suitable for the TMV suspension was calculated, involving different liquid crystalline phases .
A full simulation of the PM was performed for stiff polyelectrolytes , resulting in bundle formation, which is a possible sign of an effective attraction. Clearly, the effective pair potential picture is insufficient for bundles and many-body interactions play a significant role .
## VIII Lamellar polyions
Lamellar polyions may be schematically looked upon as rigid or flexible charged surfaces or platelets, providing a 2d counterpart of charged rods or polyelectrolytes. Examples are the geologically and technologically important smectite clays and self-assembled bilayers of ionic surfactants, which constitute the proto-type of biological membranes .
Consideration will first be given to electric double-layers around infinitely thin, uniformly charged circular or square platelets, which constitute a reasonable model of smectite clay particles, and in particular of the widely studied synthetic Laponite mineral particles. To a good approximation the latter are rigid, thin discs, of thickness $`d1`$nm, radius $`a15`$nm, carrying a structural surface charge $`\sigma e/(\mathrm{nm})^2`$. Natural montmorrilonite clays are silicate mineral platelets of similar chemical composition and crystal structure, but of irregular and polydisperse shape, and of much larger lateral dimension, implying some degree of bending flexibility.
Dry powders of clay will swell upon the addition of water, releasing counterions into the interlamellar volume, and building up interacting electric double-layers; the swelling is essentially driven by the double-layer repulsion of mostly entropic origin. During the initial stages of swelling, the spacing $`h`$ between platelets remains small compared to their lateral extension, so that a moderately swollen lamellar phase may be reasonably modelled by a stack of infinite charged planes . The results for planar geometry, discussed earlier, apply then directly to the swollen phase. In particular, limited swelling, often observed in the presence of divalent counterions, may be related to the cohesive behaviour (effective attraction between planes, or negative disjoining pressure) due to correlation effects within the PM . An important issue in the understanding of the swelling behaviour is the competitive condensation of counterions of different valence and/or size . The PM model is, however, expected to be inadequate when the interlamellar spacing h is only of the order of a few microion diameters. Under those conditions the molecular nature of the solvent (water) must be explicitly taken into account. Constant pressure MC simulations clearly show the importance of counterion hydration in determining the swelling behaviour, and point to the role of $`\mathrm{K}^+`$ ions as a clay swelling inhibitor . When swelling proceeds until the interlamellar spacing $`h`$ becomes comparable to the platelet radius $`a`$, finite size (edge) effects become important. The force acting on a platelet follows from the integration of the stress tensor $`\underset{¯}{\underset{¯}{\mathrm{\Pi }}}`$ over the two faces $`\mathrm{\Sigma }_+`$ and $`\mathrm{\Sigma }_{}`$ of the platelet:
$`𝐅={\displaystyle _{\mathrm{\Sigma }_+,\mathrm{\Sigma }_{}}}\underset{¯}{\underset{¯}{\mathrm{\Pi }}}d𝐬`$ (34)
where $`\mathrm{d}𝐬=\widehat{𝐧}\mathrm{d}s`$ is the surface element oriented along the outward normal $`\widehat{𝐧}`$, and the stress tensor has the standard form :
$`\underset{¯}{\underset{¯}{\mathrm{\Pi }}}=\left[P(𝐫)+{\displaystyle \frac{ϵ}{8\pi }}|𝐄|^2\right]𝐈{\displaystyle \frac{ϵ}{4\pi }}𝐄(𝐫)𝐄(𝐫)`$ (35)
which generalizes the uniaxial expression (24). Within LPB theory, eq. (34) leads to the following expression for the force between two coaxial discs separated by $`h`$, immersed in an ionic solution of inverse Debye length $`\kappa _\mathrm{D}`$ :
$`F_z(h)=(\pi a^2)\times {\displaystyle \frac{4\pi \sigma ^2}{ϵ}}{\displaystyle _0^{\mathrm{}}}J_1^2(x){\displaystyle \frac{1}{x}}\mathrm{exp}\left\{{\displaystyle \frac{h}{a}}\sqrt{x^2+\kappa _D^2a^2}\right\}dx`$ (36)
Numerical solutions of the non-linear PB equation, for the same coaxial geometry, show that the LPB expression (36) strongly overestimates the force , as in the spherical case, thus requiring a proper renormalization of the surface charge density $`\sigma `$. Recent MC simulations show that the force may become attractive for divalent counterions, as in the case of infinite charged planes, but that the finite size of the discs leads to significant differences in the relative weight of electrostatic and contact contributions .
Finite concentration effects in highly swollen stacks of parallel platelets of finite size may be examined within a cell model, compatible with the shape of the platelet. A coaxial cylindrical cell is chosen for a disc-shaped platelet, while a parallelepipedic cell is better adapted to square platelets . Note that the platelet concentration determines the cell volume, but not the aspect ratio of the cell, e.g. the ratio $`R/h`$ in the case of a cylindrical cell of radius $`R`$ and height $`h`$ equal to the inter-lamellar spacing; the optimum aspect ratio for given volume and electrolyte concentration is determined by minimizing the free energy of the microions in the cell. The latter has been calculated within LPB theory which can be solved analytically , and within non-linear PB theory, which requires numerical solution; the latter is greatly simplified within a Green’s function formulation .
The osmotic pressures calculated as a function of Laponite concentration within PB theory agree reasonably well with experimental data , for reservoir salt concentrations of the order of $`10^3`$M or larger, but differ dramatically from the predictions of 1d PB theory for stacks of infinite platelets, pointing to the importance of edge effects in highly swollen clays .
In very dilute dispersions of clay platelets, such that the distance between the centres of neighbouring polyions is significantly larger than the particle radius, the platelets can rotate more or less freely (sol phase). As the concentration increases, a gel phase is formed, depending on salt concentration. Several recent experiments have attempted to establish a link between the mesoscopic fractal structure of the gel, and its rheological properties , but no clear-cut scenario has yet emerged, due to metastability and ageing of the dispersions .
A theoretical description of the sol-gel transition of clay dispersions hinges on a knowledge of the highly anisotropic effective interaction between charged platelets in an electrolyte. The charged segment, or site-site model, introduced earlier for charged rods has been generalized to charged circular platelets, and used in MD simulations of Laponite dispersions .
For well separated platelets, the effective pair potential reduces to a sum of screened interactions between the electrostatic multipoles associated with the anisotropic electric double layers around each platelet . A simplified version of such an effective interaction, involving infinitely thin discs carrying an unscreened quadrupole moment, has been used in MC simulations which predict a reversible sol-gel transition .
Lipid bilayers and membranes constitute another class of lamellar polyions, which are flexible. Electrostatic interactions renormalize the bending rigidity of these flexible membranes , a subject of ongoing work beyond the scope of this review. In relation to the correlation-induced attraction between like-charged planes discussed earlier, the fluidity of membranes provides an additional mechanism for attraction between membranes, resulting from the lateral charge fluctuations within the planes of the latter .
## IX Discrete solvent effects
So far all theoretical considerations of interacting electric double-layers were based on the PM, which ignores the molecular nature of the solvent. The PM thus neglects excluded volume effects of the latter, as well as hydration of ions and the expected reduction of the local dielectric constant near highly charged surfaces, due to polarization of nearby water molecules. Neglect of these and other solvent effects is expected to be particularly inadequate in situations where the spacing between two charged polyion surfaces is only of the order of a few molecular diameters.
To go beyond the PM level of description, simple models have been used to describe the solvent molecules (generally water). The crudest model is to represent the latter by neutral hard spheres of appropriate diameter, while keeping a macroscopic dielectric constant $`ϵ`$ in the Coulombic interactions between ions; this model accounts only for excluded volume effects. The next refinement is to consider hard spheres with embedded point dipoles, to account for the highly polar nature of the solvent , while a reasonable local coordination of the solvent molecules can only be achieved by adding higher order multipoles . In MC or MD simulations, much more sophisticated pair potentials may be used, involving three or more interaction sites on each water molecule, as in the widely used SPC/E potential .
A purely HS solvent is implicit in the modified PB formulation of Kralj-Iglic and Iglic and others . As expected, the modified PB equation leads to a saturation of the counterion density at contact, for high surface charge densities. The force between two charged surfaces in a mixture of charged and neutral hard spheres was calculated by Tang et al. within a DFT generalizing the earlier theory of the same authors for a PM electrolyte ; the force is a strongly oscillatory function of the spacing between the plates, due to HS layering, and is rather insensitive to surface charge and salt concentration. The latter observation, which holds when solvent, anions and cations are of the same size, does not carry over to the more realistic case of unequal diameters, which has been investigated within IHNC theory . At low surface charge, the total force between platelets is reasonably approximated by a superposition of the pure HS solvent ”hydration” force, and the electrostatic contribution of the ions, as calculated within the PM, in the case of equal diameters .
The dipolar HS model for the solvent was first used to determine the structure of the electric double layer near a single charged wall within the MSA . More accurate and complete results on the planar electric double layer in a dipolar solvent were obtained from numerical solutions of the coupled RHNC equations for the three density profiles . These calculations give Statistical Mechanics evidence for the reduction of the local dielectric constant near the charged surface, and for significant electrostriction. A generic density functional, based on Rosenfeld’s successful ”fundamental measure” theory for hard core fluids has been put forward, which can be adapted to any polyion geometry ; when applied to a single planar double layer, this theory shows a considerable enhancement of the counterion density at contact, when a dipolar, rather than bare HS solvent is used.
The effective solvation force between two charged plates immersed in an ionic solution of charged and dipolar HS was calculated within a quadratic free energy functional of the local density, charge density and polarization, generalizing eqs. (28) and (8) (with direct correlation functions approximated by bulk MSA solutions) . This calculation provided the first convincing evidence of the strong influence of a granular (as opposed to continuous) solvent on the solvation forces at short range.
The most complete investigation so far of the potential of mean force between two spherical polyions immersed in an ionic solution with a ”realistic” solvent involving hard spheres with point dipoles and tetrahedral quadrupoles , appears to be the work by Kinoshita et al , who solved the RHNC equations for this highly asymmetric multicomponent ”mixture”, for ratios $`10a/a_s30`$, where $`a`$ and $`a_s`$ are the polyion and solvent molecule radii. While maintaining $`a_s`$ fixed at 0.14nm, a value appropriate for water, these authors varied the counterion radius. Their most striking finding is that larger counterions are more strongly adsorbed to the polyion in the presence of a molecular solvent, leading to a greater reduction of the Coulomb repulsion between the polyions, and even to the possibility of an effective attraction for monovalent counterions, when their radius exceeds $`a_s`$ by more than $`20\%`$. The trend towards a reduction of effective repulsion with increasing counterion size is the exact opposite of the prediction of the PM, a clear illustration of the pitfalls of the latter in describing the short-range behaviour of solvation forces! Obviously more work on discrete solvent models is needed.
## X Outlook
Although the understanding of effective interactions between electric double layers has clearly advanced in the last decade through a combination of new theoretical approaches, quantitative measurements and large-scale computer simulations, there are still many open questions. Future research should focus on the influence of image charges due to dielectric discontinuities between the solvent and the container walls or the polyions themselves. Real samples possess, moreover, an intrinsic polydispersity in size, charge and shape which becomes relevant for a quantitative comparison between theory and experimental data. Another rapidly growing field concerns flexible polyelectrolyte chains whose stiffness is governed by their persistence length . A theoretical approach needs input from both nonlinear screening and polymer theories. Colloids can also be stabilized by polyelectrolytes, and the resulting effective interaction becomes qualitatively different from that between charged spheres, particularly near the overlap concentration. The microscopic incorporation of the solvent is still in its early stage and full molecular theories describing hydration forces and hydrogen bonding are highly desirable for aqueous suspensions. Lastly alternative approaches, such as recent field theoretical formulations , might lead to additional insight.
## Acknowledgements
We thank E.Allahyarov, D.Goulding, Y. Levin, P. Linse, V. Lobaskin, P. Pincus, and R. van Roij for helpful comments.
Figure captions
Figures 1a, 1b: Two different plate geometries.
Figure 2: Phase diagram for a charged suspensions in the $`\eta n_s`$ plane with $`\eta =4\pi na^3/3`$ and $`Z=350`$, $`a=380`$nm , $`ϵ=13.9`$ (penthanol) and monovalent counterions at room temperature. Possible phases are fluid (F), gas (G), liquid (L), and face-centered-cubic (FCC) crystals. Tie-lines denote coexistence conditions.
Figure 3: Same as Figure 2, but now for $`Z=1000`$, $`a=350`$nm, $`ϵ=25.3`$ (ethanol). There is a triple point of coexisting gas, liquid and crystal.
Figure 4: Effective force $`F`$ in units of $`F_0=\frac{Z^2e^2}{a^2}\times 10^4`$ between two polyions versus reduced distance $`r/a`$. The squares are for aqueous suspensions. The force is repulsive in good agreement with DLVO-theory (dashed line). The circles are for a solvent with a strongly reduced dielectric constant $`ϵ`$ and show attraction. The parameters are given in detail in Ref. . The inset shows a counterion configuration around two fixed polyions. |
warning/0002/hep-th0002139.html | ar5iv | text | # IMSC/2000/02/07hep-th/0002139 Loop Variables and Gauge Invariant Interactions - I This is a detailed description of an approach, outlined in a talk at the Puri Workshop in 1996, to use loop variables to string interactions.
## 1 Introduction
The loop variable approach introduced in (hereafter I) (see also ) is an attempt to write down gauge invariant equations of motion for both massive and massless modes. This method being rooted in the sigma model approach ,the computations are expected to be simpler and the gauge transformation laws more transparent. This hope was borne out at the free level and also to a certain extent in the interacting case ( hereafter II). The gauge transformations at the free level can be summarized by the equation
$$k(t)k(t)\lambda (t)$$
(1.0.1)
Here $`k(t)`$ is the generalized momentum Fourier-conjugate to $`X`$ and $`\lambda `$ is the gauge parameter. This clearly has the form of a rescaling and one can speculate on the space-time interpretation of the string symmetries as has been done for instance in I.
In II the interacting case was discussed. It was shown that the leading interactions could be obtained by the simple trick of introducing an additional parameter ‘$`\sigma `$’ as $`k(t)k(t,\sigma )`$, parametrizing different interacting strings. Thus, for instance, $`k_1^\mu (\sigma _1)k_1^\nu (\sigma _2)`$ could stand for two massless photons when $`\sigma _1\sigma _2`$, but when $`\sigma _1=\sigma _2`$ it would represent a massive “spin 2” excitation of one string. The gauge transformations admit a corresponding generalization
$$k(t,\sigma )k(t,\sigma )𝑑\sigma _1\lambda (t,\sigma _1)$$
(1.0.2)
It was shown, however that this prescription introduces only the leading interaction terms.<sup>1</sup><sup>1</sup>1Our notation, unfortunately, is perverse: The variable $`t`$ originally used in the free theory lies along the string, and $`\sigma `$ introduced in II labelling as it does the number of interaction vertices, parametrizes evolution.
In a talk some years ago (hereafter III) we showed that there is a natural generalization of this construction to include the full set of interactions that one expects based on the operator product expansion (OPE) of vertex operators. It was shown that this construction gives gauge invariant equations. The generalization involves introducing $`\sigma `$-dependence in the $`X`$ coordinates also. Gauge invariance at the level of loop variables is very easy to see. What was not clear at the time was whether there was a consistent map to space-time fields. Here we show that this is in fact the case. It crucially involves keeping a finite cutoff on the world sheet and also making the space-time fields a function of $`x_n`$. Keeping a finite cutoff is required when going off shell . In the presence of a finite cutoff there are problems with gauge invariance as discussed in . It was shown there that to lowest order these problems could be resolved by adding a massive mode with an appropriate transformation law. It was also speculated that maintaining exact gauge invariance would be possible if all the modes are kept. This is shown to be true in the present work. We have the full gauge invariance and it is not violated by a finite cutoff and the construction necessarily requires all the modes.
Another feature that emerges from the present work is that the space-time fields have to be functions of the gauge coordinates $`x_n`$. This is forced on us when we require that it be possible to define the gauge transformation laws for the space time fields in a consistent way. This does not introduce any new physical degrees since these can be gauge fixed. Nevertheless it is amusing to note that space-time has effectively become infinite dimensional.
In order to make precise contact with string theory one has to perform another step that we do not discuss in this paper. It involves generalizing to the interacting case the dimensional reduction that was done in I (for the free case). Given that the basic technique involves calculation of correlators of vertex operators on the world sheet we are more or less guaranteed that we will reproduce bosonic string amplitudes. What needs to be shown is that the dimensional reduction does not violate the gauge invariance. We reserve this issue for a future publication .
This paper is organized as follows. In section II we give a short review of and elaborate on the role of the parameter $`\sigma `$. In section III we describe the generalization outlined in . In section IV we discuss the gauge invariance of the Loop Variable. Section V contains some examples of equations of motion. Section VI discusses how one obtains the gauge transformation laws of the space-time fields. Section VII discusses the consistency issue and shows why the fields have to be functions of $`x_n`$. Section VIII contains a summary and some concluding remarks. An Appendix contains some details of a covariant Taylor expansion.
## 2 Review
### 2.1 Free theory
In I the following expression was the starting point to obtain the equations of motion at the free level:
$$e^A=e^{k_0^2\mathrm{\Sigma }+_{n>0}k_n.k_0\frac{}{x_n}\mathrm{\Sigma }+_{n,m>0}k_n.k_m(\frac{^2}{x_nx_m}\frac{}{x_{n+m}})\mathrm{\Sigma }+ik_nY_n}$$
(2.1.1)
The prescription was to vary w.r.t $`\mathrm{\Sigma }`$ and evaluate at $`\mathrm{\Sigma }=0`$ to get the equations of motion. Here, $`2\mathrm{\Sigma }<Y(z)Y(z)>`$ and $`Y=_n\alpha _n\frac{^nX}{(n1)!}_n\alpha _n\stackrel{~}{Y}_n`$. $`\alpha _n`$ are the modes of the einbein $`\alpha (t)`$ used in defining the loop variable
$$e^{i_c\alpha (t)k(t)_zX(z+t)dt+ik_0X}=e^{i_nk_nY_n}$$
(2.1.2)
$$\alpha (t)=\underset{n0}{}\alpha _nt^n$$
$$k(t)=\underset{n0}{}k_nt^n$$
One can also show easily that $`Y_n=\frac{Y}{x_n}`$. $`\mathrm{\Sigma }`$ is thus a generalization of the Liouville mode, and what we have is a generalization of the Weyl invariance condition on vertex operators.
There is an alternative way to obtain the $`\mathrm{\Sigma }`$ dependence . This is to perform a general conformal transformation on a vertex operator by acting on it with $`e^{_n\lambda _nL_{+n}}`$ using the relation <sup>2</sup><sup>2</sup>2This relation is only true to lowest order in $`\lambda `$. The exact expression is given in :
$$e^{_n\lambda _nL_{+n}}e^{iK_m\stackrel{~}{Y}_m}=e^{K_n.K_m\lambda _{nm}+\stackrel{~}{Y_n}\stackrel{~}{Y_m}\lambda _{+n+m}+imK_n\stackrel{~}{Y_m}\lambda _{n+m}}e^{iK_m\stackrel{~}{Y_m}}$$
(2.1.3)
The anomalous term is $`K_n.K_m\lambda _{nm}`$ and the classical term is $`mK_n\stackrel{~}{Y_m}\lambda _{n+m}`$. We will ignore the classical piece: this can be rewritten as a $`(mass)^2`$ term, which will be reproduced by performing a dimensional reduction, and other pieces involving derivatives of $`\mathrm{\Sigma }`$ (defined below) that correspond to field redefinitions . We can apply (2.1.3) to the loop variable (5.3) by setting $`K_m=_nk_{mn}\alpha _n`$. Defining
$$\mathrm{\Sigma }=\underset{p,q}{}\alpha _p\alpha _q\lambda _{pq}$$
(2.1.4)
we recover (2.1.1). It is the approach described above that generalizes more easily to the interacting case.
The equations thus obtained are invariant under
$$k_n\underset{m}{}k_{nm}\lambda _m$$
(2.1.5)
which is just the mode expansion of (1.0.1).
That this is an invariance of the equations of motion derived from (2.1.1) follows essentially from the fact that the transformation (2.1.5), applied to (2.1.1) changes it by a total derivative.
$$\delta A=\underset{n}{}\lambda _n\frac{}{x_n}[A]$$
(2.1.6)
The equations are obtained by the operation $`\frac{\delta }{\delta \mathrm{\Sigma }}A_{\mathrm{\Sigma }=0}`$. Thus consider the gauge variation of this:
$$\delta _{gauge}\frac{\delta }{\delta \mathrm{\Sigma }}A=\frac{\delta }{\delta \mathrm{\Sigma }}\delta _{gauge}A$$
$$=\frac{\delta }{\delta \mathrm{\Sigma }}\lambda _n\frac{}{x_n}A$$
Now $`A`$ being linear in $`\mathrm{\Sigma }`$ and its derivatives can always be expressed after integration by parts as $`\mathrm{\Sigma }B`$ for some $`B`$. Thus we have
$$=\frac{\delta }{\delta \mathrm{\Sigma }}\lambda _n\frac{}{x_n}(\mathrm{\Sigma }B)=\lambda _n(\frac{}{x_n}B+\frac{}{x_n}B)=0$$
Thus the equations obtained from (2.1.1) are invariant.
The connection between these variables and transformation laws and the usual fields and gauge transformations was described in I. Briefly, the fields were defined by
$$S_{n,m,\mathrm{}}^{\mu \nu ..}(k_0)=<k_n^\mu k_m^\nu \mathrm{}>=[\underset{n}{}dk_nd\lambda _n]k_n^\mu k_m^\nu \mathrm{}\mathrm{\Psi }[k_0,k_1,k_2,\mathrm{},k_n,\mathrm{}\lambda _m\mathrm{}]$$
(2.1.7)
where $`\mathrm{\Psi }`$ is some “string field” that describes a given configuration.
And the gauge parameters $`\mathrm{\Lambda }_{p,n,m..}^{\mu ,\nu }(k_0)`$ were defined by a similar equation involving one power of $`\lambda _p`$, $`p=1,2,\mathrm{}`$, and arbitrary numbers of $`k_n,k_m..`$ in the integrand.
However there are some caveats. In proving (2.1.6) one needs to use equations such as
$$\frac{}{x_1}(\frac{^2}{x_1^2}\frac{}{x_2})\mathrm{\Sigma }=(\frac{^3}{x_1^3}\frac{^2}{x_1x_2})\mathrm{\Sigma }=2(\frac{^2}{x_1x_2}\frac{}{x_3})\mathrm{\Sigma }$$
(2.1.8)
which follow from the basic definitions . This implies that equations of motion obtained by varying $`\mathrm{\Sigma }`$ will not be invariant. To see this consider the following expression:
$$2(\frac{^2}{x_1x_2}\frac{}{x_3})\mathrm{\Sigma }A+(\frac{^2}{x_1^2}\frac{}{x_2})\mathrm{\Sigma }\frac{A}{x_1}$$
(2.1.9)
Using (2.1.8) we get
$$=\frac{}{x_1}[(\frac{^2}{x_1^2}\frac{}{x_2})\mathrm{\Sigma }A]$$
(2.1.10)
which is a total derivative. However if we vary (2.1.9) w.r.t. $`\mathrm{\Sigma }`$, one gets
$$2\delta \mathrm{\Sigma }(\frac{^2}{x_1x_2}+\frac{}{x_3})A+\delta \mathrm{\Sigma }(\frac{^2}{x_1^2}+\frac{}{x_2})\frac{A}{x_1}$$
(2.1.11)
which is not zero. On the other hand if we rewrite (2.1.9) as (using (2.1.8))
$$(\frac{^3}{x_1^3}\frac{^2}{x_1x_2})\mathrm{\Sigma }A+(\frac{^2}{x_1^2}\frac{}{x_2})\mathrm{\Sigma }\frac{A}{x_1}$$
(2.1.12)
and vary w.r.t $`\mathrm{\Sigma }`$ we get
$$\delta \mathrm{\Sigma }(\frac{^3}{x_1^3}\frac{^2}{x_1x_2})A+\delta \mathrm{\Sigma }(\frac{^3}{x_1^3}+\frac{^2}{x_1x_2})A$$
(2.1.13)
which is zero.
Thus one has to be careful about varying w.r.t $`\mathrm{\Sigma }`$ indiscriminately. Let us review the solution to this as we will face the same issue in the interacting case discussed in the next section. Consider the variation of the exponent $`A`$ (2.1.1), reproduced below, due to $`\lambda _p`$:
$$e^A=e^{k_0^2\mathrm{\Sigma }+_{n>0}k_n.k_0\frac{}{x_n}\mathrm{\Sigma }+_{n,m>0}k_n.k_m(\frac{^2}{x_nx_m}\frac{}{x_{n+m}})\mathrm{\Sigma }+ik_nY_n}$$
(2.1.14)
The change is
$$\lambda _p(\underset{n}{}k_{np}.k_0\frac{}{x_n}\mathrm{\Sigma }+\underset{n,mp}{}k_{np}.k_m(\frac{^2}{x_nx_m}\frac{}{x_{n+m}})\mathrm{\Sigma }+$$
$$+\underset{m}{}k_m.k_0(\frac{^2}{x_mx_p}\frac{}{x_{m+p}})\mathrm{\Sigma }+i\underset{n}{}k_{np}Y_n)$$
(2.1.15)
If we assume tracelessness of the gauge parameter so that any term of the form $`\lambda _pk_n.k_m`$ is zero then the second sum in (2.1.15) vanishes and using the fact that the first sum cancels the second term in the last sum we can rewrite the variation of $`A`$ as
$$\lambda _p\frac{}{x_p}\{\underset{m}{}k_m.k_0\frac{}{x_m}\mathrm{\Sigma }+\underset{n}{}ik_{np}Y_{np}+\underset{n,m}{}k_n.k_m(\frac{^2}{x_nx_m}\frac{}{x_{n+m}})\mathrm{\Sigma }\}$$
$$=\lambda _p\frac{}{x_p}A$$
(2.1.16)
Note that in the first line of this equation we have added a term that vanishes by the tracelessness constraint, viz terms involving $`\lambda _pk_n.k_m`$ . But it is important that we have not used identities of the type given in (2.1.8). Thus tracelessness of the gauge parameters ensures the gauge invariance of the equations.
### 2.2 Interactions
In II this approach was generalized to include some interactions. The basic idea was to introduce a new parameter $`\sigma :0\sigma 1`$ to label different strings and to replace each $`k_n`$ in the free equation by $`_0^1𝑑\sigma k_n(\sigma )`$. The next step was to assume that
$$<k_1^\mu (\sigma _1)k_1^\nu (\sigma _2)>=S^{\mu \nu }\delta (\sigma _1\sigma _2)+A^\mu A^\nu $$
(2.2.17)
where $`<\mathrm{}>`$ denotes $`𝒟k(\sigma )\mathrm{}\mathrm{\Psi }[k(\sigma )]`$, $`\mathrm{\Psi }`$ being the“string field” defined in I.<sup>3</sup><sup>3</sup>3No special property of $`\mathrm{\Psi }`$ is assumed other than this. This corresponds to saying that when $`\sigma _1=\sigma _2`$, both the $`k_1`$’s belong to the same string and otherwise to different strings where they represent two photons at an interaction point.<sup>4</sup><sup>4</sup>4It will be seen that (2.2.17)has to be generalized by replacing the $`\delta `$-function on the RHS by something else, when we attempt to reproduce string amplitudes . However in this paper we will not do so. The gauge transformation is replaced by (1.0.2). This is easily seen to give interacting interacting equations. However the fact is that this is only a leading term in the infinite set of interaction vertices.
As a prelude to generalizing this construction, let us explain more precisely the nature of the replacement $`k_n_0^1𝑑\sigma k_n(\sigma )`$. Let us split the interval $`(0,1)`$ into $`N`$ bits of width $`a=\frac{1}{N}`$. We will assume that when $`\sigma `$ satisfies $`\frac{n}{N}\sigma \frac{n+1}{N}`$ it represents the $`(n+1)`$th string. Let us also define a function
$`D(\sigma _1,\sigma _2)`$ $`=`$ $`1if\sigma _1,\sigma _2belongtothesameinterval`$ (2.2.18)
$`=`$ $`0if\sigma _1,\sigma _2belongtodifferentintervals.`$
Thus $`_0^1𝑑\sigma _1D(\sigma _1\sigma _2)=a=_0^1𝑑\sigma _1_0^1𝑑\sigma _2D(\sigma _1\sigma _2)`$.
Then we set
$$<k^\mu (\sigma _1)k^\nu (\sigma _2)>=\frac{D(\sigma _1,\sigma _2)}{a}S^{\mu \nu }+A^\mu A^\nu $$
(2.2.19)
In the limit $`N\mathrm{},a0,\frac{D(\sigma _1,\sigma _2)}{a}\delta (\sigma _1\sigma _2)`$ and we recover (2.2.17).
In effect (2.1.1) has been modified to
$$e^{_0^1_0^1𝑑\sigma _1𝑑\sigma _2[k_0(\sigma _1)k_0(\sigma _2)\mathrm{\Sigma }+k_n(\sigma _1).k_0(\sigma _2)\frac{}{x_n}\mathrm{\Sigma }+_{n,m}(\frac{^2}{x_nx_m}\frac{}{x_{n+m}})\mathrm{\Sigma }]+_0^1𝑑\sigma k_n(\sigma )Y_n}$$
(2.2.20)
The final step (which is also necessary in the free case), is to dimensionally reduce to obtain the massive equations. For details we refer the reader to I.
The modification (2.2.17), that replaces $`S^{\mu \nu }`$ by $`S^{\mu \nu }+A^\mu A^\nu `$ can be understood in terms of the OPE. Consider a correlation function involving two vector vertex operators and any other set of operators, that we represent as
$$𝒜=<V_1V_2\mathrm{}V_N:k_1^\mu _zX^\mu e^{i{\scriptscriptstyle k_0Y}}:q_1^\nu _wX^\nu e^{iq_0Y}>$$
(2.2.21)
The OPE of $`:k_1^\mu _zX^\mu (z)e^{i{\scriptscriptstyle k_0Y}}:`$ and $`:q_1^\nu _wX^\nu (w)e^{iq_0Y}:`$
is given by
$$:k_1^\mu _zX^\mu (z)e^{i{\scriptscriptstyle k_0Y}}::q_1^\nu _wX^\nu (w)e^{iq_0Y}:=$$
$$:k_1^\mu q_1^\nu _zX^\mu _wX^\nu e^{i(k_0X(z)+q_0X(w))}:+termsinvolvingcontractions.$$
(2.2.22)
We can Taylor expand
$$X(w)=X(z)+(wz)_zX+O(wz)+\mathrm{}$$
(2.2.23)
This gives for the leading term in (2.2.21)
$$𝒜=<V_1V_2\mathrm{}V_N:k_1^\mu q_1^\nu _zX^\mu _zX^\nu e^{i(k_0X(z)+q_0X(w))}:>$$
(2.2.24)
Compare this with the correlation involving $`S^{\mu \nu }`$:
$$𝒜^{}=<V_1V_2\mathrm{}V_N:k_1^\mu k_1^\nu _zX^\mu _zX^\nu e^{i{\scriptscriptstyle k_0Y}}:>$$
(2.2.25)
We see that $`𝒜`$ and $`𝒜^{}`$ give identical terms except that $`S^{\mu \nu }`$ is replaced by $`A^\mu A^\nu `$. It is in this sense that the substitution given in II, gives the leading term in the OPE. The crucial point is that, while in (2.2.20) we have introduced the parameter $`\sigma `$ in the $`k_n`$’s we have not done so for the $`Y_n`$’s. This is equivalent to approximating $`X(w)`$ by $`X(z)`$ in (2.2.23). Clearly, the generalization required to get all the terms is to introduce the parameter $`\sigma `$ in $`Y`$ also. We turn to this in the next section.
## 3 Interactions
### 3.1 Introducing $`\sigma `$-dependence in the loop variable
We will introduce the parameter $`\sigma `$ in all the variables keeping in mind the basic motivation that $`\sigma `$ labels different vertex operators. Thus all the variables that are required to define a vertex operator become $`\sigma `$ dependent. Thus
$$X^\mu (z)X^\mu (z(\sigma ))$$
(3.1.1)
$$x_nx_n(\sigma )$$
(3.1.2)
in addition to
$$k_n^\mu k_n^\mu (\sigma )$$
(3.1.3)
The $`\sigma `$-dependence of $`x_n`$ in eqn. 3.1.2 is only an intermediate step. At the end of the day (but before any integration by parts is done) we will set all the $`x_n`$’s to be the same. One can think of this merely as a device for keeping track of which term is being differentiated.
(3.1.1) and (3.1.2) imply that
$$\frac{}{x_n}Y\frac{}{x_n(\sigma )}Y(z(\sigma ),x_n(\sigma ))$$
(3.1.4)
Note that $`X`$ need not be an explicit function of $`\sigma `$ since at a given location $`z`$, on the world sheet there can only be one $`X(z)`$. As an example of the above consider the case when we have regions $`(0,1/2)`$ and $`(1/2,1)`$. When $`0\sigma 1/2`$ one has $`z(\sigma )z`$ and for $`1/2\sigma 1`$ one has $`z(\sigma )w`$. Similarly $`x_n(\sigma )`$ could be called $`x_n,y_n`$ in the two regions and $`k_n(\sigma )`$ could be called $`k_n,p_n`$ in the two regions. Thus in this example the vertex operator $`k_n(\sigma )Y_n(z(\sigma ),x_n(\sigma ))e^{ik_0(\sigma )Y(\sigma )}`$ stands for $`k_n\frac{Y}{x_n}(z,x_i)e^{ik_0Y(z,x_n)}`$ and $`p_n\frac{Y}{y_n}(w,y_i)e^{ip_0Y(w,y_n)}`$ in the two regions.
Now we have to clarify what we mean by a derivative w.r.t $`x_n(\sigma )`$: In (3.1.4) we have $`\frac{Y(z(\sigma ),x_i(\sigma ))}{x_n(\sigma )}`$ : One has to specify the meaning of $`\frac{x_n(\sigma )}{x_n(\sigma ^{})}`$. Clearly what we want is: If $`\sigma ,\sigma ^{}`$ belong to the same interval, then $`\frac{x_n(\sigma )}{x_n(\sigma ^{})}=\mathrm{\hspace{0.33em}1}`$ and zero otherwise. Thus using (2.2.18)
$$\frac{x_n(\sigma )}{x_n(\sigma ^{})}=D(\sigma ,\sigma ^{})$$
(3.1.5)
or more generally
$$\frac{x_n(\sigma )}{x_m(\sigma ^{})}=\delta _{nm}D(\sigma ,\sigma ^{})$$
(3.1.6)
Note that this is not the same as the conventional functional derivative. However we can define
$$\frac{\delta x_n(\sigma )}{\delta x_n(\sigma ^{})}\frac{D(\sigma ,\sigma ^{})}{a}$$
(3.1.7)
which, in the limit $`a0`$ becomes the usual functional derivative. Thus
$$𝑑\sigma ^{}\frac{\delta Y(\sigma )}{\delta x_n(\sigma ^{})}=\frac{Y(\sigma )}{x_n(\sigma )}$$
(3.1.8)
We can now write down the generalization of (2.1.1)
$$exp\{d\sigma _1d\sigma _2\{k_0(\sigma _1).k_0(\sigma _2)[\stackrel{~}{\mathrm{\Sigma }}(\sigma _1,\sigma _2)+\stackrel{~}{G}(\sigma _1,\sigma _2)]$$
$$+𝑑\sigma _3𝑑\sigma _4\underset{n,m0}{}k_n(\sigma _1).k_m(\sigma _2)$$
$$\frac{1}{2}[\frac{\delta ^2}{\delta x_n(\sigma _1)\delta x_m(\sigma _2)}\delta (\sigma _1\sigma _2)\frac{\delta }{\delta x_{n+m}(\sigma _1)}][\stackrel{~}{\mathrm{\Sigma }}(\sigma _3,\sigma _4)+\stackrel{~}{G}(\sigma _3,\sigma _4)]\}\}$$
$$exp\{i𝑑\sigma k_n(\sigma )Y_n(\sigma )\}$$
(3.1.9)
For convenience of notation have assumed the following:
$$\frac{\delta \alpha _n(\sigma _1)}{\delta x_0(\sigma _2)}=\frac{D(\sigma _1\sigma _2)}{a}\alpha _n(\sigma _1)$$
This saves us the trouble of writing separately the case $`n=0`$ in the sum in (3.1.9).
In (3.1.9) $`G(\sigma _1,\sigma _2)=\stackrel{~}{G}(z(\sigma _1),z(\sigma _2))`$ =$`<Y(z(\sigma _1))Y(z(\sigma _2))>`$ is the Green function which starts out as $`ln(z_1z_2)`$. We have suppressed the Lorentz indices. One might expect by Lorentz invariance $`\stackrel{~}{G}^{\mu \nu }=\delta ^{\mu \nu }\stackrel{~}{G}`$. However in I it was seen that the $`D+1`$th coordinate has a special role and is like the ghost coordinate of bosonic string theory. So there is no reason to expect the full $`SO(D+1)`$ invariance. In fact we will have to assume some specific properties for $`\stackrel{~}{G}^{D+1,D+1}`$ in order to reproduce string amplitudes.
More precisely, if we define: \[Using the notation $`z_i=z(\sigma _i)`$\]
$$D_{z_1}=D_{z(\sigma _1)}1+\alpha _1(\sigma _1)\frac{}{z(\sigma _1)}+\alpha _2\frac{^2}{z^2(\sigma _1)}+\mathrm{}$$
(3.1.10)
so that
$$Y(z(\sigma ))=D_{z(\sigma )}X(z(\sigma ))$$
(3.1.11)
then,
$$\stackrel{~}{G}(z_1,z_2)=D_{z_1}D_{z_2}G(z_1,z_2)$$
(3.1.12)
$$\stackrel{~}{\mathrm{\Sigma }}(\sigma _1,\sigma _2)=D_{z_1}D_{z_2}\rho (\sigma _1,\sigma _2)$$
(3.1.13)
where
$$\rho (\sigma _1,\sigma _2)=\frac{\lambda (z(\sigma _1))\lambda (z(\sigma _2))}{z(\sigma _1)z(\sigma _2)}$$
(3.1.14)
is the generalization of the usual Liouville mode $`\rho (\sigma )`$ which is equal to $`\frac{d\lambda }{dz}`$. The $`\stackrel{~}{\mathrm{\Sigma }}`$ dependence in (3.1.9) is obtained by the following step:
$$e^{:\frac{1}{2}{\scriptscriptstyle 𝑑u\lambda (u)[_zX(z+u)]^2}:}e^{ik_n\frac{}{x_n}D_{z_1}X}e^{ip_m\frac{}{x_m}D_{z_2}X}$$
(3.1.15)
defines the action of the Virasoro generators on the two sets of vertex operators.
$$=e^{ik_n.p_m_{x_n}_{y_m}D_{z_1}D_{z_2}{\scriptscriptstyle 𝑑u{\scriptscriptstyle \frac{\lambda (u)}{z_1z_2}}[{\scriptscriptstyle \frac{1}{z_1u}}{\scriptscriptstyle \frac{1}{z_2u}}]}}$$
(3.1.16)
$$=e^{ik_n.p_m_{x_n}_{y_m}\stackrel{~}{\mathrm{\Sigma }}}$$
(3.1.17)
This expression is only valid to lowest order in $`\lambda `$ which is all we need here.<sup>5</sup><sup>5</sup>5The exact expression is given in . The expression
$$𝑑\sigma _1𝑑\sigma _2\frac{1}{2}[\frac{\delta ^2}{\delta x_n(\sigma _1)\delta x_m(\sigma _2)}\delta (\sigma _1\sigma _2)\frac{\delta }{\delta x_{n+m}(\sigma _1)}][\stackrel{~}{\mathrm{\Sigma }}(\sigma _3,\sigma _4)+\stackrel{~}{G}(\sigma _3,\sigma _4)]$$
(3.1.18)
can easily be seen to be equal to
$$\frac{^2}{x_n(\sigma _3)x_m(\sigma _3)}\stackrel{~}{\mathrm{\Sigma }}(\sigma _3,\sigma _4)$$
(3.1.19)
In the limit $`\sigma _3=\sigma _4=\sigma `$ this is just equal to $`1/2[\frac{^2}{x_n(\sigma )x_m(\sigma )}\frac{}{x_{m+n}(\sigma )}]\stackrel{~}{\mathrm{\Sigma }}(\sigma ,\sigma )`$ and reduces to the free field case described by (2.2.20)(provided the limit is taken after differentiation).
Let us show that the gauge transformation (1.0.2) changes (3.1.9) by a total derivative
$$\delta A=𝑑\sigma ^{}\lambda (\sigma ^{})𝑑\sigma \frac{\delta }{\delta x_n(\sigma )}A$$
(3.1.20)
## 4 Invariance of the Loop Variable
Our starting point is the loop variable $`e^A`$ given by:
$$e^{i\{{\scriptscriptstyle 𝑑\sigma k_0(\sigma )Y(\sigma )}+i_{n>0}k_n(\sigma )\frac{Y(\sigma )}{x_n(\sigma )}\}}$$
$$e^{{\scriptscriptstyle }{\scriptscriptstyle }d\sigma _1d\sigma _2\{k_0(\sigma _1)k_0(\sigma _2)[\stackrel{~}{\mathrm{\Sigma }}+\stackrel{~}{G}](\sigma _1,\sigma _2)+(_{n>0}k_n(\sigma _1).k_0(\sigma _2)\frac{[\stackrel{~}{\mathrm{\Sigma }}+\stackrel{~}{G}](\sigma _1,\sigma _2)}{x_n(\sigma _1)}+\sigma _1\sigma _2)\}}$$
$$e^{{\scriptscriptstyle 𝑑\sigma _1𝑑\sigma _2}\{_{n,m>0}k_n(\sigma _1).k_m(\sigma _2)\frac{^2[\stackrel{~}{\mathrm{\Sigma }}+\stackrel{~}{G}](\sigma _1,\sigma _2)}{x_1(\sigma _1)x_1(\sigma _2)}\}}$$
(4.0.1)
Under a gauge transformation:
$$k_n(\sigma _1)𝑑\sigma \lambda _p(\sigma )k_{np}(\sigma _1)$$
(4.0.2)
Let us consider $`p=1`$.
$$k_1(\sigma _1).k_0(\sigma _2)\frac{}{x_1(\sigma _1)}[\stackrel{~}{\mathrm{\Sigma }}+\stackrel{~}{G}](\sigma _1,\sigma _2)𝑑\sigma \lambda _1(\sigma )k_0(\sigma _1).k_0(\sigma _2)\frac{}{x_1(\sigma _1)}[\stackrel{~}{\mathrm{\Sigma }}+\stackrel{~}{G}](\sigma _1,\sigma _2)$$
(4.0.3)
$$k_0(\sigma _1).k_1(\sigma _2)\frac{}{x_1(\sigma _1)}[\stackrel{~}{\mathrm{\Sigma }}+\stackrel{~}{G}](\sigma _1,\sigma _2)𝑑\sigma \lambda _1(\sigma )k_0(\sigma _1).k_0(\sigma _2)\frac{}{x_1(\sigma _2)}[\stackrel{~}{\mathrm{\Sigma }}+\stackrel{~}{G}](\sigma _1,\sigma _2)$$
(4.0.4)
Adding the two we get:
$$𝑑\sigma \lambda _1(\sigma )[\frac{}{x_1(\sigma _1)}+\frac{}{x_1(\sigma _2)}]k_0(\sigma _1).k_0(\sigma _2)[\stackrel{~}{\mathrm{\Sigma }}+\stackrel{~}{G}](\sigma _1,\sigma _2)$$
(4.0.5)
Similarly,
$$k_2(\sigma _1).k_0(\sigma _2)\frac{}{x_2(\sigma _1)}[\stackrel{~}{\mathrm{\Sigma }}+\stackrel{~}{G}](\sigma _1,\sigma _2)𝑑\sigma \lambda _1(\sigma )k_1(\sigma _1).k_0(\sigma _2)\frac{}{x_2(\sigma _1)}[\stackrel{~}{\mathrm{\Sigma }}+\stackrel{~}{G}](\sigma _1,\sigma _2)$$
(4.0.6)
$$k_0(\sigma _1).k_2(\sigma _2)\frac{}{x_2(\sigma _2)}[\stackrel{~}{\mathrm{\Sigma }}+\stackrel{~}{G}](\sigma _1,\sigma _2)𝑑\sigma \lambda _1(\sigma )k_0(\sigma _1).k_1(\sigma _2)\frac{}{x_2(\sigma _2)}[\stackrel{~}{\mathrm{\Sigma }}+\stackrel{~}{G}](\sigma _1,\sigma _2)$$
(4.0.7)
$$k_1(\sigma _1).k_1(\sigma _2)\frac{^2[\stackrel{~}{\mathrm{\Sigma }}+\stackrel{~}{G}](\sigma _1,\sigma _2)}{x_1(\sigma _1)x_1(\sigma _2)}d\sigma \lambda _1(\sigma )(k_1(\sigma _1).k_0(\sigma _2)+k_0(\sigma _1).k_1(\sigma _2))\frac{^2[\stackrel{~}{\mathrm{\Sigma }}+\stackrel{~}{G}](\sigma _1,\sigma _2)}{x_1(\sigma _1)x_1(\sigma _2)}$$
(4.0.8)
Adding we get,
$$d\sigma \lambda _1(\sigma )\{[\frac{}{x_1(\sigma _1)}+\frac{}{x_1(\sigma _2)}]k_1(\sigma _1).k_0(\sigma _2)\frac{}{x_1(\sigma _1)}[\stackrel{~}{\mathrm{\Sigma }}+\stackrel{~}{G}]+$$
$$[\frac{}{x_1(\sigma _1)}+\frac{}{x_1(\sigma _2)}]k_0(\sigma _1).k_1(\sigma _2)\frac{}{x_1(\sigma _2)}[\stackrel{~}{\mathrm{\Sigma }}+\stackrel{~}{G}]\}$$
(4.0.9)
$$=d\sigma \lambda _1(\sigma )[\frac{}{x_1(\sigma _1)}+\frac{}{x_1(\sigma _2)}]\{k_1(\sigma _1).k_0(\sigma _2)\frac{}{x_1(\sigma _1)}[\stackrel{~}{\mathrm{\Sigma }}+\stackrel{~}{G}]+\sigma _1\sigma _2\}$$
(4.0.10)
Now consider
$$k_2(\sigma _1).k_1(\sigma _2)\frac{^2[\stackrel{~}{\mathrm{\Sigma }}+\stackrel{~}{G}]}{x_2(\sigma _1)x_1(\sigma _2)}+k_1(\sigma _1).k_2(\sigma _2)\frac{^2[\stackrel{~}{\mathrm{\Sigma }}+\stackrel{~}{G}]}{x_1(\sigma _1)x_2(\sigma _2)}$$
(4.0.11)
$$𝑑\sigma \lambda _1(\sigma )k_1(\sigma _1).k_1(\sigma _2)\frac{^2[\stackrel{~}{\mathrm{\Sigma }}+\stackrel{~}{G}]}{x_2(\sigma _1)x_1(\sigma _2)}+𝑑\sigma \lambda _1(\sigma )k_1(\sigma _1).k_1(\sigma _2)\frac{^2[\stackrel{~}{\mathrm{\Sigma }}+\stackrel{~}{G}]}{x_1(\sigma _1)x_2(\sigma _2)}$$
(4.0.12)
$$=𝑑\sigma \lambda _1(\sigma )[\frac{}{x_1(\sigma _1)}+\frac{}{x_1(\sigma _2)}]k_1(\sigma _1).k_1(\sigma _2)\frac{^2[\stackrel{~}{\mathrm{\Sigma }}+\stackrel{~}{G}](\sigma _1,\sigma _2)}{x_1(\sigma _1)x_1(\sigma _2)}$$
(4.0.13)
From the above it is clear that we get the following:
$$\delta A=𝑑\sigma \lambda _1(\sigma )[\frac{}{x_1(\sigma _1)}+\frac{}{x_1(\sigma _2)}]A$$
(4.0.14)
On setting $`x_n(\sigma _1)=x_n(\sigma _2)=x_n`$ we get
$$\delta A=\lambda _1(\sigma )\frac{}{x_1}A$$
(4.0.15)
Thus to lowest order in $`\lambda _1`$, $`e^A`$ changes by a total derivative in $`x_1`$. This is obviously true for $`\lambda _p`$ also.
Thus the equations obtained by varying w.r.t $`\stackrel{~}{\mathrm{\Sigma }}(z(\sigma _1),z(\sigma _2),x_n(\sigma _1),x_n(\sigma _2))`$ are invariant. However $`\stackrel{~}{\mathrm{\Sigma }}`$ is not a local field on the world sheet. $`A`$ has terms of the form $`[\frac{^2}{xy}\mathrm{\Sigma }(w,z,x,y)]_{x=y}\frac{^2}{x^2}[\mathrm{\Sigma }(w,z,x,y)_{x=y}]`$. Thus $`A`$ cannot be expressed in terms of $`x_n`$-derivatives of a field. We would have to use both $`x_n`$ and $`Y_n`$. But we cannot integrate by parts on both $`x_n,y_n`$ \- there is no such gauge invariance. So we first Taylor expand it in powers of $`z(\sigma _2)z(\sigma _1)`$ the coefficients of which are derivatives of a local field $`\mathrm{\Sigma }(z,x)\overline{\mathrm{\Sigma }}(z,x,y)_{x=y}`$, where $`\overline{\mathrm{\Sigma }}(v,x,y)=\stackrel{~}{\mathrm{\Sigma }}(v,v,x,y)`$. Below we have used the letter $`v`$ to denote $`\frac{z(\sigma _1)+z(\sigma _2)}{2}`$ and $`x(\sigma _1)=x,x(\sigma _2)=y`$.
$$\stackrel{~}{\mathrm{\Sigma }}(\sigma _1,\sigma _2)=\overline{\mathrm{\Sigma }}(v)+aD_1(x,y)\overline{\mathrm{\Sigma }}(v)+a^2D_2(x,y)\overline{\mathrm{\Sigma }}(v)+\mathrm{}$$
(4.0.16)
$$=\overline{\mathrm{\Sigma }}(v)+a\underset{r}{}(\gamma _r^0\frac{\overline{\mathrm{\Sigma }}}{y_{r+1}}\gamma _r^0\frac{\overline{\mathrm{\Sigma }}}{x_{r+1}})+$$
(4.0.17)
$$\frac{a^2}{2!}[\underset{s}{}(\gamma _s^1\frac{\overline{\mathrm{\Sigma }}}{y_{s+1}}+\gamma _s^1\frac{\overline{\mathrm{\Sigma }}}{x_{s+1}})2\underset{r,s}{}\gamma _r^0\gamma _s^0\frac{^2\overline{\mathrm{\Sigma }}}{x_{r+1}y_{s+1}}]+\mathrm{}$$
$`D_k`$ and $`\gamma `$ are defined in the Appendix. Very explicitly, the first few terms of the Taylor expansion are :
$$\stackrel{~}{\mathrm{\Sigma }}(\sigma _1,\sigma _2)=\overline{\mathrm{\Sigma }}(v)+a\underset{D_1(x,y)}{\underset{}{[\frac{\overline{\mathrm{\Sigma }}}{y_1}\frac{\overline{\mathrm{\Sigma }}}{x_1}+(\frac{y_1^2}{2}+y_2)\frac{\overline{\mathrm{\Sigma }}}{y_3}(\frac{x_1^2}{2}+x_2)\frac{\overline{\mathrm{\Sigma }}}{x_3}]}}+$$
$$\frac{a^2}{2}\underset{D_2(x,y)}{\underset{}{[\frac{\overline{\mathrm{\Sigma }}}{x_2}+\frac{\overline{\mathrm{\Sigma }}}{y_2}2\frac{^2\overline{\mathrm{\Sigma }}}{x_1y_1}+x_1\frac{\overline{\mathrm{\Sigma }}}{x_3}+y_1\frac{\overline{\mathrm{\Sigma }}}{y_3}2(\frac{x_1^2}{2}+x_2)\frac{^2\overline{\mathrm{\Sigma }}}{x_3y_1}2(\frac{y_1^2}{2}+y_2)\frac{^2\overline{\mathrm{\Sigma }}}{y_3x_1}]}}+\mathrm{}$$
(4.0.18)
Once you Taylor expand $`\mathrm{\Sigma }`$ we have the following problem that we encounter also in the free case. The problem is that when the gauge variation does not produce $`k_0`$ we need constraints:
$$k_2.k_1[\frac{^2\mathrm{\Sigma }}{x_2x_1}\frac{\mathrm{\Sigma }}{x_3}]\lambda _1k_1.k_1[\frac{^2\mathrm{\Sigma }}{x_2x_1}\frac{\mathrm{\Sigma }}{x_3}]$$
$$\mathrm{?}=\mathrm{?}\lambda _1k_1.k_1\frac{}{x_1}[\frac{^2\mathrm{\Sigma }}{x_1^2}\frac{\mathrm{\Sigma }}{x_2}]$$
(4.0.19)
The last equation is not an identity and follows only because certain constraints are obeyed by $`\mathrm{\Sigma }`$. This in turn requires the imposition of the constraints on the gauge parameters - $`\lambda _1k_1.k_1=0`$.
This problem is not there for the $`\lambda _2`$ variation as can be seen in the following:
$$k_2.k_1[\frac{^2\mathrm{\Sigma }}{x_2x_1}\frac{\mathrm{\Sigma }}{x_3}]\lambda _2k_0.k_1[\frac{^2\mathrm{\Sigma }}{x_2x_1}\frac{\mathrm{\Sigma }}{x_3}]$$
$$\frac{k_3.k_0}{2}\frac{\mathrm{\Sigma }}{x_3}\frac{\lambda _2k_1.k_0}{2}\frac{\mathrm{\Sigma }}{x_3}$$
They add up to:
$$\lambda _2k_2.k_0\frac{^2\mathrm{\Sigma }}{x_2x_1}=\lambda _2\frac{}{x_2}[\frac{k_1.k_0}{2}\frac{\mathrm{\Sigma }}{x_1}]$$
For the above argument to go through in the interacting case we need the following property for the Taylor expansion coefficients:
$$\frac{^2}{x_nx_m}[D_k(x,y)\overline{\mathrm{\Sigma }}]=\frac{}{x_{n+m}}[D_k(x,y)\overline{\mathrm{\Sigma }}]$$
(4.0.20)
(and the same obviously for $`Y_n`$). It is demonstrated in the Appendix that this is in fact true.
Thus in general consider:
$$k_n(\sigma _1).k_m(\sigma _2)[\frac{^2D_k(x,y)\overline{\mathrm{\Sigma }}}{x_ny_m}]_{x=y}+k_{n+m}(\sigma _2).k_0(\sigma _1)[\frac{D_k(x,y)\overline{\mathrm{\Sigma }}}{y_{n+m}}]_{x=y}$$
(4.0.21)
$$𝑑\sigma \lambda _n(\sigma )k_0(\sigma _1).k_m(\sigma _2)[\frac{^2D_k(x,y)\overline{\mathrm{\Sigma }}}{x_ny_m}]_{x=y}+𝑑\sigma \lambda _n(\sigma )k_m(\sigma _2).k_0(\sigma _1)[\frac{D_k(x,y)\overline{\mathrm{\Sigma }}}{y_{n+m}}]_{x=y}$$
(4.0.22)
$$=𝑑\sigma \lambda _n(\sigma )k_0(\sigma _1).k_m(\sigma _2)[(\frac{}{x_n}+\frac{}{y_n})\frac{D_k(x,y)\overline{\mathrm{\Sigma }}}{y_m}]_{x=y}$$
(4.0.23)
$$=d(\sigma )\lambda _n(\sigma )\frac{}{x_n}[k_0(\sigma _1).k_m(\sigma _2)\frac{D_k(x,y)\overline{\mathrm{\Sigma }}}{y_m}_{x=y}]$$
(4.0.24)
Similarly,
$$k_n(\sigma _2).k_m(\sigma _1)[\frac{^2D_k(x,y)\overline{\mathrm{\Sigma }}}{x_my_n}]_{x=y}+k_{n+m}(\sigma _1).k_0(\sigma _2)[\frac{D_k(x,y)\overline{\mathrm{\Sigma }}}{x_{n+m}}]_{x=y}$$
(4.0.25)
$$𝑑\sigma \lambda _n(\sigma )k_0(\sigma _2).k_m(\sigma _1)[\frac{^2D_k(x,y)\overline{\mathrm{\Sigma }}}{x_my_n}]_{x=y}+𝑑\sigma \lambda _n(\sigma )k_m(\sigma _1).k_0(\sigma _2)[\frac{D_k(x,y)\overline{\mathrm{\Sigma }}}{x_{n+m}}]_{x=y}$$
(4.0.26)
$$=𝑑\sigma \lambda _n(\sigma )k_0(\sigma _2).k_m(\sigma _1)[(\frac{}{x_n}+\frac{}{y_n})\frac{D_k(x,y)\overline{\mathrm{\Sigma }}}{x_m}]_{x=y}$$
(4.0.27)
$$=d(\sigma )\lambda _n(\sigma )\frac{}{x_n}[k_0(\sigma _2).k_m(\sigma _1)\frac{D_k(x,y)\overline{\mathrm{\Sigma }}}{x_m}_{x=y}]$$
(4.0.28)
Adding the two we find that the $`\lambda _n`$ variation is a total derivative in $`x_n`$ of $`A`$ even after Taylor expanding.
Similarly the tracelessness constraint of the free theory generalizes to
$$<d(\sigma )\lambda _p(\sigma )[k_n(\sigma _1).k_m(\sigma _2)]\mathrm{}.>=0$$
(4.0.29)
in the equation of motion. All the above guarantees that the variation of $`e^A`$ is a total derivative and therefore the equations of motion obtained by varying w.r.t $`\mathrm{\Sigma }`$ are invariant.
## 5 Examples
### 5.1 Vector $`k_1`$ Contribution to $`Y_1^\mu `$
Our starting point is $`e^A`$ given by
$$e^{i\{{\scriptscriptstyle 𝑑\sigma k_0(\sigma )Y(\sigma )}+i_{n>0}k_n(\sigma )\frac{Y(\sigma )}{x_n(\sigma )}\}}$$
$$e^{{\scriptscriptstyle }{\scriptscriptstyle }d\sigma _1d\sigma _2\{k_0(\sigma _1)k_0(\sigma _2)[\stackrel{~}{\mathrm{\Sigma }}+\stackrel{~}{G}](\sigma _1,\sigma _2)+(_{n>0}k_n(\sigma _1).k_0(\sigma _2)\frac{[\stackrel{~}{\mathrm{\Sigma }}+\stackrel{~}{G}](\sigma _1,\sigma _2)}{x_n(\sigma _1)}+\sigma _1\sigma _2)\}}$$
$$e^{{\scriptscriptstyle 𝑑\sigma _1𝑑\sigma _2}\{_{n,m>0}k_n(\sigma _1).k_m(\sigma _2)\frac{^2[\stackrel{~}{\mathrm{\Sigma }}+\stackrel{~}{G}](\sigma _1,\sigma _2)}{x_1(\sigma _1)x_1(\sigma _2)}\}}$$
(5.1.1)
We keep terms with one $`k_1`$ only. There are three terms that contribute.
$$e^{{\scriptscriptstyle k_0(\sigma _5)}.k_0(\sigma _6)[\stackrel{~}{\mathrm{\Sigma }}+\stackrel{~}{G}]}k_1(\sigma _1).k_0(\sigma _2)\frac{}{x_1(\sigma _1)}[\stackrel{~}{\mathrm{\Sigma }}+\stackrel{~}{G}]e^{i{\scriptscriptstyle k_0Y}}+e^{{\scriptscriptstyle k_0(\sigma _5)}.k_0(\sigma _6)[\stackrel{~}{\mathrm{\Sigma }}+\stackrel{~}{G}]}ik_1Y_1e^{i{\scriptscriptstyle k_0Y}}$$
(5.1.2)
In leading order we have:
$$\frac{}{x_1(\sigma _1)}\stackrel{~}{\mathrm{\Sigma }}(\sigma _1,\sigma _2)=\frac{\overline{\mathrm{\Sigma }}}{x_1}=\frac{1}{2}\frac{\mathrm{\Sigma }}{x_1}$$
$$\stackrel{~}{\mathrm{\Sigma }}=\overline{\mathrm{\Sigma }}=\mathrm{\Sigma }$$
(5.1.3)
We get
$$k_1(\sigma _1).k_0(\sigma _2)\frac{\overline{\mathrm{\Sigma }}}{x_1}e^{i{\scriptscriptstyle k_0Y}}+k_1(\sigma _2).k_0(\sigma _1)\frac{\overline{\mathrm{\Sigma }}}{y_1}e^{i{\scriptscriptstyle k_0Y}}+k_0(\sigma _1).k_0(\sigma _2)\overline{\mathrm{\Sigma }}ik_1Y_1e^{i{\scriptscriptstyle k_0Y}}$$
(5.1.4)
which on setting $`x_n=y_n`$ becomes
$$=\frac{k_1(\sigma _1).k_0(\sigma _2)+k_1(\sigma _2).k_0(\sigma _1)]}{2}\frac{\mathrm{\Sigma }}{x_1}e^{i{\scriptscriptstyle k_0Y}}+k_0(\sigma _1).k_0(\sigma _2)\mathrm{\Sigma }ik_1Y_1$$
(5.1.5)
Setting $`\frac{\delta }{\delta \mathrm{\Sigma }}`$ of this expression to zero we get the equation (we can set all the $`\sigma `$ ’s to be equal)
$$k_1(\sigma _1)k_0(\sigma _1)ik_0^\mu Y_1^\mu +k_0(\sigma _1).k_0(\sigma _1)ik_1^\mu Y_1^\mu =0$$
(5.1.6)
Converting to space-time fields the coefficient of $`Y_1^\mu `$ is:
$$=_\mu ^\nu A^\mu +_\mu ^\mu A^\nu =_\mu F^{\mu \nu }=0$$
(5.1.7)
which is Maxwell’s equation. (5.1.6) is clearly invariant under $`k_1(\sigma _1)k_1(\sigma _1)+𝑑\sigma \lambda _1(\sigma )k_0(\sigma _1)`$, which in terms of space-time fields is $`A_\mu A_\mu +_\mu \mathrm{\Lambda }`$.
### 5.2 $`k_1k_1`$ and $`k_2`$ Contribution to $`Y_1^\mu `$
(i)
$$\frac{1}{2!}\{k_1(\sigma _1).k_0(\sigma _2)\frac{}{x_1(\sigma _1)}[\stackrel{~}{\mathrm{\Sigma }}+\stackrel{~}{G}]+\sigma _1\sigma _2\}$$
$$\{k_1(\sigma _3).k_0(\sigma _4)\frac{}{x_1(\sigma _3)}[\stackrel{~}{\mathrm{\Sigma }}+\stackrel{~}{G}]+\sigma _3\sigma _4\}e^{{\scriptscriptstyle k_0(\sigma _5)}.k_0(\sigma _6)[\stackrel{~}{\mathrm{\Sigma }}+\stackrel{~}{G}]}e^{i{\scriptscriptstyle k_0Y}}$$
(5.2.8)
(ii)
$$k_1(\sigma _1).k_1(\sigma _2)\frac{^2[\stackrel{~}{\mathrm{\Sigma }}+\stackrel{~}{G}](\sigma _1,\sigma _2)}{x_1(\sigma _1)x_1(\sigma _2)}e^{{\scriptscriptstyle k_0(\sigma _5)}.k_0(\sigma _6)[\stackrel{~}{\mathrm{\Sigma }}+\stackrel{~}{G}]}e^{i{\scriptscriptstyle k_0Y}}$$
(5.2.9)
(iii)
$$e^{{\scriptscriptstyle k_0(\sigma _5)}.k_0(\sigma _6)[\stackrel{~}{\mathrm{\Sigma }}+\stackrel{~}{G}]}\{k_1(\sigma _1).k_0(\sigma _2)\frac{}{x_1(\sigma _1)}[\stackrel{~}{\mathrm{\Sigma }}+\stackrel{~}{G}]+\sigma _1\sigma _2\}ik_1Y_1e^{i{\scriptscriptstyle k_0Y}}$$
(5.2.10)
Let us consider each in turn:
(i)
Using the leading order expressions given in (5.1.3) we get
$$2\times \frac{1}{2!}e^{{\scriptscriptstyle k_0(\sigma _5)}.k_0(\sigma _6)[\stackrel{~}{G}]}k_1(\sigma _1).k_0(\sigma _2)\frac{\mathrm{\Sigma }}{x_1}[2k_1(\sigma _3).k_0(\sigma _4)\frac{\stackrel{~}{G}(\sigma _3,\sigma _4)}{x_1(\sigma _3)}]e^{i{\scriptscriptstyle k_0Y}}$$
(5.2.11)
Varying w.r.t $`\mathrm{\Sigma }`$ gives
$$=2e^{{\scriptscriptstyle k_0(\sigma _5)}.k_0(\sigma _6)[\stackrel{~}{G}]}k_1(\sigma _1).k_0(\sigma _2)k_1(\sigma _3).k_0(\sigma _4)\frac{\stackrel{~}{G}(\sigma _3,\sigma _4)}{x_1(\sigma _3)}ik_0\frac{Y}{x_1}e^{i{\scriptscriptstyle k_0Y}}$$
(5.2.12)
Now we have to consider all possible contractions of the $`k_n`$’s. In order to keep track of the possibilities we separate them into two cases: those involving only one point on the world sheet (i.e. only one vertex operator) and those involving two distinct points (two vertex operators).
One Vertex Operator
In the first case we have to be careful about regularizing. Let us refer to the point as $`\sigma _A`$ with $`z(\sigma _A)=z`$ as the location of the vertex operator. When there is a need for regularizing we will let $`\sigma _B`$ be the second point with $`z(\sigma _B)z(\sigma _A)=ϵ`$. We now let the various $`\sigma _i`$ be equal to $`\sigma _A`$ or $`\sigma _B`$ in all possible combinations, but we divide by 2 since these are actually the same point.
Following this procedure we see that regularization is required for $`\sigma _3,\sigma _4`$ when they stand for the same point (and also for $`\sigma _5,\sigma _6`$, which we ignore for the moment). Thus we can let
$$\sigma _3=\sigma _A,\sigma _4=\sigma _B$$
and $`\sigma _1,\sigma _2`$ can be anything. This gives
$$2k_1(\sigma _A).k_0(\sigma _A)k_1(\sigma _A).k_0(\sigma _B)\frac{1}{z_Az_B}$$
(5.2.13)
Now we let $`z_Az_B`$ and $`\sigma _A\sigma _B`$ and use:
$$<k_1^\mu (\sigma _A)k_1^\nu (\sigma _A)>=<k_1^\mu (\sigma _A)k_1^\nu (\sigma _B)>=S_{1,1}^{\mu \nu }(k_0)$$
(5.2.14)
This gives
$$2S_{1,1}^{\mu \nu }k_0^\mu k_0^\nu \frac{1}{ϵ}.$$
(5.2.15)
The other possibility is
$$\sigma _3=\sigma _B,\sigma _4=\sigma _A$$
and again $`\sigma _1,\sigma _2`$ can be anything, which gives
$$2k_1(\sigma _A).k_0(\sigma _A)k_1(\sigma _B).k_0(\sigma _A)\frac{1}{z_Bz_A}$$
(5.2.16)
Using (5.2.14) gives
$$2S_{1,1}^{\mu \nu }k_0^\mu k_0^\nu \frac{1}{ϵ}$$
(5.2.17)
Adding the two ((5.2.15) and (5.2.17)) we get zero.
Two Vertex Operators
Now we go to the second possibility viz. there are two distinct points. Let us call them $`\sigma _I`$ and $`\sigma _{II}`$ and let $`z(\sigma _I)=z`$ and $`z(\sigma _{II})=w`$. Thus we can have a)$`\sigma _1=\sigma _I`$ and $`\sigma _3=\sigma _{II}`$ or b) vice versa.
Consider a):
First we consider the case that does not require regularization.
Non-singular Case
$$2k_1(\sigma _I).k_0(\sigma _2)k_1(\sigma _{II}).k_0(\sigma _4)\frac{1}{wz(\sigma _4)}$$
(5.2.18)
Let $`\sigma _4=\sigma _I`$ and $`\sigma _2=\sigma _Ior\sigma _{II}`$
We now use the notation $`k_0(\sigma _I)=p`$ and $`k_0(\sigma _{II})=q`$. Thus
$$<k_1^\mu (\sigma _I)>=A^\mu (p)$$
$$<k_1^\mu (\sigma _{II})>=A^\mu (q)$$
(5.2.19)
and we get as contribution to the equation of motion:
$$𝑑z𝑑wA^\mu (p)(p+q)^\mu A^\nu (q)p^\nu \frac{1}{wz}$$
(5.2.20)
We have explicitly written out the integrals over $`z`$ and $`w`$ to emphasize the symmetry. Thus by antisymmetry of the integrand in $`z,w`$ this is zero.
Singular Case
If we let $`\sigma _4=\sigma _{II}`$ (and $`\sigma _2=\sigma _Ior\sigma _{II}`$) we have to regularize. So we split $`\sigma _{II}`$ into $`\sigma _A`$ and $`\sigma _B`$ as before. Again either
$$\sigma _3=\sigma _Aand\sigma _4=\sigma _B$$
which gives
$$2k_1(\sigma _I).k_0(\sigma _2)k_1(\sigma _A).k_0(\sigma _B)\frac{1}{w_Aw_B}$$
(5.2.21)
and using (5.2.19)
$$=A^\mu (p)(p+q)^\mu A^\nu (q)q^\nu (\frac{1}{ϵ})$$
(5.2.22)
or
$$\sigma _3=\sigma _Band\sigma _4=\sigma _A$$
which gives
$$2k_1(\sigma _I).k_0(\sigma _2)k_1(\sigma _B).k_0(\sigma _A)\frac{1}{w_Bw_A}$$
(5.2.23)
and using (5.2.19)
$$=A^\mu (p)(p+q)^\mu A^\nu (q)q^\nu (\frac{+1}{ϵ})$$
(5.2.24)
Adding the two contributions again gives zero.
We have also to look at possibility b) which was $`\sigma _1=\sigma _{II}`$ and $`\sigma _3=\sigma _{II}`$ Analogous to (5.2.20) one gets
$$𝑑z𝑑wA^\mu (q)(p+q)^\mu A^\nu (p)p^\nu \frac{1}{zw}$$
(5.2.25)
Note that this is (upto a sign) (5.2.20) with the labels $`p,q`$ interchanged. But $`p,q`$ being integration variables we get back (5.2.20) but the overall sign being opposite, they cancel. The integration $`𝑑z𝑑w`$ also ensures the vanishing of each term , viz (5.2.20) and (5.2.25), individually. This is also as it should be since interchanging $`z`$ with $`w`$ is equivalent to interchanging momenta.
(ii)
$$k_1(\sigma _1).k_1(\sigma _2)\frac{^2\stackrel{~}{\mathrm{\Sigma }}(\sigma _1,\sigma _2)}{x_1(\sigma _1)x_1(\sigma _2)}e^{{\scriptscriptstyle k_0(\sigma _5)}.k_0(\sigma _6)[\stackrel{~}{G}]}e^{i{\scriptscriptstyle k_0Y}}$$
(5.2.26)
Use
$$\frac{^2\stackrel{~}{\mathrm{\Sigma }}(\sigma _1,\sigma _2)}{x_1(\sigma _1)x_1(\sigma _2)}=\frac{^2\overline{\mathrm{\Sigma }}}{x_1y_1}+\mathrm{}$$
On setting $`x_n=y_n`$,
$$=\frac{1}{2}(\frac{^2\mathrm{\Sigma }}{x_1^2}\frac{\mathrm{\Sigma }}{x_2})+\mathrm{}$$
(5.2.27)
Only the first term can give, on integration by parts, a contribution to $`Y_1`$:
$$e^{{\scriptscriptstyle k_0(\sigma _5)}.k_0(\sigma _6)[\stackrel{~}{G}]}k_0(\sigma _7).k_0(\sigma _8)\frac{\stackrel{~}{G}(\sigma _7,\sigma _8)}{x_1}k_1(\sigma _1).k_1(\sigma _2)ik_0Y_1e^{i{\scriptscriptstyle k_0Y}}$$
(5.2.28)
$$\stackrel{~}{G}(\sigma _7,\sigma _8)=ln|z(\sigma _7)z(\sigma _8)|+\frac{\alpha _1}{z(\sigma _7)z(\sigma _8)}\frac{\beta _1}{z(\sigma _7)z(\sigma _8)}+\mathrm{}$$
(5.2.29)
So
$$\frac{\stackrel{~}{G}}{x_1}=0+higherorderinx_n$$
We do not get any contribution.
(iii)
$$e^{{\scriptscriptstyle k_0(\sigma _5)}.k_0(\sigma _6)[\stackrel{~}{G}]}k_1(\sigma _1).k_0(\sigma _2)\frac{\stackrel{~}{\mathrm{\Sigma }}}{x_1(\sigma _1)}+\sigma _1\sigma _2+\stackrel{~}{\mathrm{\Sigma }}\stackrel{~}{G}$$
(5.2.30)
Using $`\frac{\stackrel{~}{G}}{x_1}=0`$ the leading order contribution is zero.
Thus combining (i),(ii) and (iii) we conclude that there are no corrections to $`_\mu F^{\mu \nu }=0`$ to this order.
### 5.3 $`k_1k_1,k_2`$ Contributions to $`Y_2^\mu `$
We start with, as usual, $`e^A`$ given by
$$e^{i\{{\scriptscriptstyle 𝑑\sigma k_0(\sigma )Y(\sigma )}+i_{n>0}k_n(\sigma )\frac{Y(\sigma )}{x_n(\sigma )}\}}$$
$$e^{{\scriptscriptstyle }{\scriptscriptstyle }d\sigma _1d\sigma _2\{k_0(\sigma _1)k_0(\sigma _2)[\stackrel{~}{\mathrm{\Sigma }}+\stackrel{~}{G}](\sigma _1,\sigma _2)+(_{n>0}k_n(\sigma _1).k_0(\sigma _2)\frac{[\stackrel{~}{\mathrm{\Sigma }}+\stackrel{~}{G}](\sigma _1,\sigma _2)}{x_n(\sigma _1)}+\sigma _1\sigma _2)\}}$$
$$e^{{\scriptscriptstyle 𝑑\sigma _1𝑑\sigma _2}\{_{n,m>0}k_n(\sigma _1).k_m(\sigma _2)\frac{^2[\stackrel{~}{\mathrm{\Sigma }}+\stackrel{~}{G}](\sigma _1,\sigma _2)}{x_1(\sigma _1)x_1(\sigma _2)}\}}$$
(5.3.31)
Pick out the terms that contribute to $`Y_2^\mu `$ involving $`k_1k_1`$ and $`k_2`$.
There are four terms:
(i)
$$e^{{\scriptscriptstyle k_0(\sigma _3)}.k_0(\sigma _4)[\stackrel{~}{G}]}\{k_1(\sigma _1).k_0(\sigma _2)\frac{\stackrel{~}{\mathrm{\Sigma }}}{x_1(\sigma _1)}+\sigma _1\sigma _2\}ik_1Y_1e^{i{\scriptscriptstyle k_0Y}}$$
(5.3.32)
(ii)
$$e^{{\scriptscriptstyle k_0(\sigma _3)}.k_0(\sigma _4)[\stackrel{~}{G}]}\{k_2(\sigma _1).k_0(\sigma _2)\frac{\stackrel{~}{\mathrm{\Sigma }}}{x_2(\sigma _1)}+\sigma _1\sigma _2\}e^{i{\scriptscriptstyle k_0Y}}$$
(5.3.33)
(iii)
$$e^{{\scriptscriptstyle k_0(\sigma _3)}.k_0(\sigma _4)[\stackrel{~}{G}]}\{k_1(\sigma _1).k_1(\sigma _2)\frac{^2[\stackrel{~}{\mathrm{\Sigma }}+\stackrel{~}{G}](\sigma _1,\sigma _2)}{x_1(\sigma _1)x_1(\sigma _2)}\}e^{i{\scriptscriptstyle k_0Y}}$$
(5.3.34)
(iv)
$$e^{{\scriptscriptstyle k_0(\sigma _3)}.k_0(\sigma _4)[\stackrel{~}{G}+\stackrel{~}{\mathrm{\Sigma }}]}ik_2Y_2e^{i{\scriptscriptstyle k_0Y}}$$
(5.3.35)
We expand $`\stackrel{~}{\mathrm{\Sigma }}`$ using (5.1.3) and vary w.r.t. $`\mathrm{\Sigma }`$:
(i)
$$e^{{\scriptscriptstyle k_0(\sigma _3)}.k_0(\sigma _4)[\stackrel{~}{G}]}\{k_1(\sigma _1).k_0(\sigma _2)\frac{\overline{\mathrm{\Sigma }}}{x_1}+\sigma _1\sigma _2\}ik_1Y_1e^{i{\scriptscriptstyle k_0Y}}$$
$$=e^{{\scriptscriptstyle k_0(\sigma _3)}.k_0(\sigma _4)[\stackrel{~}{G}]}k_1(\sigma _1).k_0(\sigma _2)\frac{\mathrm{\Sigma }}{x_1}ik_1(\sigma )Y_1e^{i{\scriptscriptstyle k_0Y}}$$
$`\frac{\delta }{\delta \mathrm{\Sigma }}`$ gives:
$$e^{{\scriptscriptstyle k_0(\sigma _3)}.k_0(\sigma _4)[\stackrel{~}{G}]}k_1(\sigma _1).k_0(\sigma _2)ik_1(\sigma )Y_2e^{i{\scriptscriptstyle k_0Y}}$$
(5.3.36)
We have to make contractions of the $`k_n`$’s. Let $`z(\sigma _1)=z`$ and assign the momentum $`p`$ to this point. Let $`z(\sigma )=w`$ and assign momentum $`q`$. Now $`\sigma _2=\sigma _1`$ or $`\sigma _2=\sigma `$ are two possibilities and for each of these we can have $`\sigma _3=\sigma _1,\sigma _4=\sigma `$ or $`\sigma _3=\sigma ,\sigma _4=\sigma _1`$. None of the above need regularization. We can also include the following two possibilities that need a regulator:$`\sigma _3=\sigma _4=\sigma `$ or $`\sigma _3=\sigma _4=\sigma _1`$. For these cases we will let $`\sigma _A`$ and $`\sigma _B`$ be the “point splitting” of $`\sigma `$ . Thus ($`\sigma _3=\sigma _A`$ and $`\sigma _4=\sigma _B`$) or ($`\sigma _3=\sigma _B`$ and $`\sigma _4=\sigma _A`$). We will weight these with a factor of 1/2. This gives $`q^2lnϵ`$. Similarly point splitting $`\sigma _1`$ gives $`p^2lnϵ`$. Putting all the above together and using (5.2.19)we get:
$$A(p).(p+q)iA^\mu (q)|zw|^{2p.q}(ϵ)^{p^2+q^2}Y_2^\mu e^{i(p+q)Y}$$
(5.3.37)
If we multiply and divide by $`ϵ^{2p.q}`$ this becomes
$$A(p).(p+q)iA^\mu (q)|\frac{zw}{ϵ}|^{2p.q}(ϵ)^{(p+q)^2}Y_2^\mu e^{i(p+q)Y}$$
(5.3.38)
Interchanging the role of $`\sigma `$ and $`\sigma _1`$ i.e. setting $`z(\sigma _1)=w`$ and $`z(\sigma )=z`$, we get the same expression with $`p,q`$ interchanged. Since these are just dummy variables we can combine the two if we allow both $`p`$ and $`q`$ to vary over the full range of values.
$$A(p).(p+q)iA^\mu (q)|\frac{zw}{ϵ}|^{2p.q}(ϵ)^{(p+q)^2}Y_2^\mu e^{i(p+q)Y}$$
(5.3.39)
However
There is also the possibility that $`\sigma =\sigma _1`$. In this case we point split and let $`\sigma =\sigma _A,\sigma _1=\sigma _B`$ or vice versa (with weight 1/2 to each). Then using (5.2.14) we get
$$S_{1,1}^{\mu \nu }k_0^\nu iY_2^\mu e^{i{\scriptscriptstyle k_0Y}}(ϵ)^{k_0^2}$$
(5.3.40)
(ii)
$$e^{{\scriptscriptstyle k_0(\sigma _5)}.k_0(\sigma _6)[\stackrel{~}{G}]}\{k_2(\sigma _1).k_0(\sigma _2)\frac{\stackrel{~}{\mathrm{\Sigma }}}{x_2(\sigma _1)}+\sigma _1\sigma _2\}e^{i{\scriptscriptstyle k_0Y}}$$
(5.3.41)
Using the approximations:
$$\frac{}{x_2(\sigma _1)}\stackrel{~}{\mathrm{\Sigma }}(\sigma _1,\sigma _2)\frac{\overline{\mathrm{\Sigma }}}{x_2}=\frac{1}{2}\frac{\mathrm{\Sigma }}{x_2}$$
$$\stackrel{~}{\mathrm{\Sigma }}\overline{\mathrm{\Sigma }}=\mathrm{\Sigma }$$
(5.3.42)
$$e^{{\scriptscriptstyle k_0(\sigma _5)}.k_0(\sigma _6)[\stackrel{~}{G}]}k_2(\sigma _1).k_0(\sigma _2)\frac{\mathrm{\Sigma }}{x_2}e^{i{\scriptscriptstyle k_0Y}}$$
(5.3.43)
$`\frac{\delta }{\delta \mathrm{\Sigma }}`$ gives
$$e^{{\scriptscriptstyle k_0(\sigma _5)}.k_0(\sigma _6)[\stackrel{~}{G}]}k_2(\sigma _1).k_0(\sigma _2)ik_0^\mu Y_2^\mu e^{i{\scriptscriptstyle k_0Y}}$$
(5.3.44)
$$=(ϵ)^{k_0^2}S_2(k_0).k_0ik_0^\mu Y_2^\mu $$
(5.3.45)
(iii)
$$e^{{\scriptscriptstyle k_0(\sigma _5)}.k_0(\sigma _6)[\stackrel{~}{G}]}k_1(\sigma _1).k_1(\sigma _2)\frac{^2\stackrel{~}{\mathrm{\Sigma }}}{x_1(\sigma _1)x_1(\sigma _2)}$$
(5.3.46)
$$\frac{^2\stackrel{~}{\mathrm{\Sigma }}}{x_1(\sigma _1)x_1(\sigma _2)}\frac{1}{2}(\frac{^2\mathrm{\Sigma }}{x_1^2}\frac{\mathrm{\Sigma }}{x_2})+\mathrm{}$$
(5.3.47)
$$=e^{{\scriptscriptstyle k_0(\sigma _5)}.k_0(\sigma _6)[\stackrel{~}{G}]}k_1(\sigma _1).k_1(\sigma _2)\frac{1}{2}(\frac{^2\mathrm{\Sigma }}{x_1^2}\frac{\mathrm{\Sigma }}{x_2})e^{i{\scriptscriptstyle k_0Y}}$$
(5.3.48)
$`\frac{\delta }{\delta \mathrm{\Sigma }}`$
$$=e^{{\scriptscriptstyle k_0(\sigma _5)}.k_0(\sigma _6)[\stackrel{~}{G}]}k_1(\sigma _1).k_1(\sigma _2)ik_0Y_2e^{i{\scriptscriptstyle k_0Y}}$$
(5.3.49)
Following the procedures described earlier this gives (when $`\sigma _1\sigma _2`$):
$$(ϵ)^{(p+q)^2}|\frac{zw}{ϵ}|^{2p.q}A(p).A(q)i(p+q)^\mu Y_2^\mu e^{i{\scriptscriptstyle k_0Y}}$$
(5.3.50)
Both $`p`$ and $`q`$ are integrated over the entire range.
When $`\sigma _1=\sigma _2`$ one has to point split and this gives:
$$(ϵ)^{k_0^2}iS_2^\mu (k_0)Y_2^\mu e^{i{\scriptscriptstyle k_0Y}}$$
(5.3.51)
(iv)
$$=e^{{\scriptscriptstyle k_0(\sigma _5)}.k_0(\sigma _6)[\stackrel{~}{\mathrm{\Sigma }}+\stackrel{~}{G}]}ik_2^\nu Y_2^\mu $$
(5.3.52)
Using $`\stackrel{~}{\mathrm{\Sigma }}\mathrm{\Sigma }`$ and varying w.r.t.$`\mathrm{\Sigma }`$ gives
$$=e^{{\scriptscriptstyle k_0(\sigma _5)}.k_0(\sigma _6)[\stackrel{~}{G}]}k_0(\sigma _1).k_0(\sigma _2)ik_2^\nu Y_2^\mu e^{i{\scriptscriptstyle k_0Y}}$$
(5.3.53)
$$=(ϵ)^{k_0^2}k_0^2iS_2^\mu (k_0)Y_2^\mu e^{i{\scriptscriptstyle k_0Y}}$$
(5.3.54)
We can also check that the equations are invariant at the loop variable level:
$`\delta (i)`$ $`=`$ $`\lambda _1(\sigma )k_0(\sigma _1).k_0(\sigma _2)ik_1Y_2e^{i{\scriptscriptstyle k_0Y}}e^{{\scriptscriptstyle k_0(\sigma _5)}.k_0(\sigma _6)[\stackrel{~}{G}]}`$
$`\lambda _1(\sigma )k_1(\sigma _1).k_0(\sigma _2)ik_0Y_2e^{i{\scriptscriptstyle k_0Y}}e^{{\scriptscriptstyle k_0(\sigma _5)}.k_0(\sigma _6)[\stackrel{~}{G}]}`$
$`\delta (ii)`$ $`=`$ $`\lambda _1(\sigma )k_1(\sigma _1).k_0(\sigma _2)ik_0Y_2e^{i{\scriptscriptstyle k_0Y}}e^{{\scriptscriptstyle k_0(\sigma _5)}.k_0(\sigma _6)[\stackrel{~}{G}]}`$
$`\lambda _2(\sigma )k_0(\sigma _1).k_0(\sigma _2)ik_0Y_2e^{i{\scriptscriptstyle k_0Y}}e^{{\scriptscriptstyle k_0(\sigma _5)}.k_0(\sigma _6)[\stackrel{~}{G}]}`$
$`\delta (iii)`$ $`=`$ $`2\lambda _1(\sigma )k_1(\sigma _1).k_0(\sigma _2)ik_0Y_2e^{i{\scriptscriptstyle k_0Y}}e^{{\scriptscriptstyle k_0(\sigma _5)}.k_0(\sigma _6)[\stackrel{~}{G}]}`$ (5.3.57)
$`\delta (iv)`$ $`=`$ $`\lambda _1(\sigma )k_0(\sigma _1).k_0(\sigma _2)ik_1Y_2e^{i{\scriptscriptstyle k_0Y}}e^{{\scriptscriptstyle k_0(\sigma _5)}.k_0(\sigma _6)[\stackrel{~}{G}]}`$
$`+\lambda _2(\sigma )k_0(\sigma _1).k_0(\sigma _2)ik_0Y_2e^{i{\scriptscriptstyle k_0Y}}e^{{\scriptscriptstyle k_0(\sigma _5)}.k_0(\sigma _6)[\stackrel{~}{G}]}`$
Clearly the variations add up to zero. In Section 6 we will discuss the gauge transformation law for space-time fields.
### 5.4 $`k_1k_1k_1,k_2k_1,k_3`$ Contribution to $`Y_1^\mu `$
There are many terms that contribute. We will consider only the following term to illustrate the technique being used.
(i)
$$\frac{1}{2!}e^{k_0(\sigma _7).k_0(\sigma _8)\stackrel{~}{G}}[k_1(\sigma _1).k_0(\sigma _2)\frac{\stackrel{~}{\mathrm{\Sigma }}}{x_1(\sigma _1)}+\sigma _1\sigma _2]$$
$$[k_1(\sigma _3).k_0(\sigma _4)\frac{\stackrel{~}{G}}{x_1(\sigma _3)}+\sigma _3\sigma _4][k_1(\sigma _5).k_0(\sigma _6)\frac{\stackrel{~}{G}}{x_1(\sigma _5)}+\sigma _5\sigma _6]e^{i{\scriptscriptstyle k_0Y}}=0.$$
(5.4.59)
We start with contractions that do not involve any regularization (“non singular case”). We will treat the cases that require regularization (“singular case”) separately. In each case depending on the number of distinct points we have different contributions.
Non Singular Case:
Three Vertex Operators:
We first consider contractions involving three distinct verex operators. Let us designate $`\sigma _I,\sigma _{II}`$ and $`\sigma _{III}`$ as the labels of the vertex operators with locations $`z(\sigma _I)=u,z(\sigma _{II})=w,z(\sigma _{III})=z`$ and momenta $`p,q`$ and $`k`$ respectively being associated with these vertex operators.
We use the approximation that $`\stackrel{~}{\mathrm{\Sigma }}\mathrm{\Sigma }`$ as before and integrate by parts on $`x_1`$ to get:
$$\frac{1}{2!}e^{k_0(\sigma _7).k_0(\sigma _8)\stackrel{~}{G}}[k_1(\sigma _1).k_0(\sigma _2)\frac{\mathrm{\Sigma }}{x_1(\sigma _1)}+\sigma _1\sigma _2]$$
$$[k_1(\sigma _3).k_0(\sigma _4)\frac{\stackrel{~}{G}}{x_1(\sigma _3)}+\sigma _3\sigma _4][k_1(\sigma _5).k_0(\sigma _6)\frac{\stackrel{~}{G}}{x_1(\sigma _5)}+\sigma _5\sigma 6]e^{i{\scriptscriptstyle k_0Y}}$$
(5.4.60)
$$=\frac{1}{2!}e^{k_0(\sigma _7).k_0(\sigma _8)\stackrel{~}{G}}k_1(\sigma _1).k_0(\sigma _2)[k_1(\sigma _3).k_0(\sigma _4)\frac{\stackrel{~}{G}}{x_1(\sigma _3)}+\sigma _3\sigma _4]$$
$$[k_1(\sigma _5).k_0(\sigma _6)\frac{\stackrel{~}{G}}{x_1(\sigma _5)}+\sigma _5\sigma _6]ik_0Y_1e^{i{\scriptscriptstyle k_0Y}}$$
(5.4.61)
Let us first consider contractions that do not involve regularization. Thus consider the assignments
$`pz(\sigma _I)=u`$ $`\sigma _1=\sigma _I`$ $`\sigma _2=\sigma _I,\sigma _{II},\sigma _{III}`$
$`qz(\sigma _{II})=w`$ $`\sigma _3=\sigma _{II}`$ $`\sigma _4=\sigma _I,\sigma _{III}`$
$`kz(\sigma _{III})=z`$ $`\sigma _5=\sigma _{III}`$ $`\sigma _6=\sigma _I,\sigma _{II}`$ (5.4.62)
This gives (using (5.2.19)
$$|\frac{zw}{ϵ}|^{2k.q}|\frac{wu}{ϵ}|^{2p.q}|\frac{zw}{ϵ}|^{2p.k}(ϵ)^{(p+q+k)^2}$$
$$A(p).(p+q+k)[\frac{A(q).p}{wu}+\frac{A(q).k}{wz}][\frac{A(k).p}{zu}+\frac{A(k).q}{zw}]i(p+q+k)^\mu Y_1^\mu e^{i(p+q+k)}$$
(5.4.63)
Two Vertex Operators
Now we come to the case where there are two vertex operators. The assignment that does not need regularization is:
(I)
$`pz(\sigma _I)=w`$ $`\sigma _1=\sigma _I`$ $`\sigma _2=\sigma _I,\sigma _{II}`$
$`\sigma _3=\sigma _I`$ $`\sigma _4=\sigma _{II}`$
$`qz(\sigma _{II})=z`$ $`\sigma _5=\sigma _{II}`$ $`\sigma _6=\sigma _I`$ (5.4.64)
This gives:
$$|\frac{zw}{ϵ}|^{2p.q}(ϵ)^{(p+q)^2}S_{1,1}^{\mu \nu }(p)(p+q)^\mu q^\nu \frac{1}{wz}A^\rho (q)p^\rho \frac{1}{zw}i(p+q)^\sigma Y_1^\sigma e^{i(p+q)Y}$$
(5.4.65)
The other possible assignment is:
(II)
$`pz(\sigma _I)=w`$ $`\sigma _1=\sigma _I`$ $`\sigma _2=\sigma _I,\sigma _{II}`$
$`qz(\sigma _{II})=z`$ $`\sigma _3=\sigma _{II}`$ $`\sigma _4=\sigma _I`$ (5.4.66)
$`\sigma _5=\sigma _{II}`$ $`\sigma _6=\sigma _I`$
which gives
$$|\frac{zw}{ϵ}|^{2p.q}(ϵ)^{(p+q)^2}A^\mu (p)(p+q)^\mu S_{1,1}^{\nu \rho }(q)\frac{p^\nu p^\rho }{(wz)^2}i(p+q)^\sigma Y_1^\sigma e^{i(p+q)Y}$$
(5.4.67)
Singular Cases:
Three Vertex Operators:
We now consider assignments that require regularization. For the three vertex operator case we have:
$`pz(\sigma _I)=u`$ $`\sigma _1=\sigma _I`$ $`\sigma _2=\sigma _I,\sigma _{II},\sigma _{III}`$
$`qz(\sigma _{II})=w`$ $`\sigma _3=\sigma _{II_A}`$ $`\sigma _4=\sigma _{II_B},`$
$`kz(\sigma _{III})=z`$ $`\sigma _5=\sigma _{III}`$ $`\sigma _6=\sigma _I,\sigma _{II}`$ (5.4.68)
$$|\frac{zw}{ϵ}|^{2k.q}|\frac{wu}{ϵ}|^{2p.q}|\frac{zw}{ϵ}|^{2p.k}(ϵ)^{(p+q+k)^2}k_1(\sigma _A).[k_0(\sigma _I)+k_0(\sigma _{II})+k_0(\sigma _{III})]$$
$$\underset{=0}{\underset{}{[\frac{k_1(\sigma _A).k_0(\sigma _B)}{w_Aw_B}+\frac{k_1(\sigma _B).k_0(\sigma _A)}{w_Bw_A}]}}([\frac{k_1(\sigma _{III}.k_0(\sigma _I)}{zu}+\frac{k_1(\sigma _{III}.k_0(\sigma _I)}{zw}]$$
(5.4.69)
We have used $`\sigma _A=\sigma _{II_A}`$ and $`\sigma _B=\sigma _{II_B}`$. In the limit $`\sigma _A\sigma _B`$, $`<k_1(\sigma _A)>=<k_1(\sigma _B)>=A(q)`$ and $`<k_0(\sigma _A)>=<k_0(\sigma _B)>=q`$ which is why the expression in square brackets vanishes.
Two Vertex Operators:
We turn to the case with two vertex operators.
(I)
$`pz(\sigma _I)=w`$ $`\sigma _1=\sigma _{I_A}`$ $`\sigma _2=\sigma _I,\sigma _{II}`$
$`\sigma _3=\sigma _{I_B}`$ $`\sigma _4=\sigma _{I_A}`$
$`qz(\sigma _{II})=z`$ $`\sigma _5=\sigma _{II}`$ $`\sigma _6=\sigma _I`$ (5.4.70)
(We also have to consider the assignment with $`\sigma _A`$ and $`\sigma _B`$ interchanged.)
This gives:
$$k_1(\sigma _A).[k_0(\sigma _A)+k_0(\sigma _{II})][\frac{k_1(\sigma _B).k_0(\sigma _A)}{w_Bw_A}(+\sigma _3\sigma _4isnotallowed)][\frac{k_1(\sigma _{II}).k_0(\sigma _I)}{zw}]$$
(5.4.71)
(As before we are using $`\sigma _{A,B}`$ to denote $`\sigma _{I_{A,B}}`$.).
$`\sigma _3`$ cannot be set to $`\sigma _A`$ because $`\sigma _1=\sigma _A`$. This is why the exchange term involving $`\sigma _3`$ and $`\sigma _4`$ is not allowed.
If we interchange $`AB`$ in (5.4.71) we find
$$k_1(\sigma _B).[k_0(\sigma _B)+k_0(\sigma _{II})][\frac{k_1(\sigma _A).k_0(\sigma _B)}{w_Aw_B}(+\sigma _3\sigma _4isnotallowed)][\frac{k_1(\sigma _{II}).k_0(\sigma _I)}{zw}]$$
(5.4.72)
In the limit $`\sigma _A\sigma _B`$,
$`<k_1^\mu (\sigma _A)k_1^\nu (\sigma _B)>=S_{1,1}^{\mu \nu }`$ and $`k_0(\sigma _A)=k_0(\sigma _B)=k_0(\sigma _I)=p`$ and thus (5.4.71) and (5.4.72) add up to zero.
II
$`pz(\sigma _I)=w`$ $`\sigma _1=\sigma _I`$ $`\sigma _2=\sigma _I,\sigma _{II}`$
$`qz(\sigma _{II})=z`$ $`\sigma _3=\sigma _{II_A}`$ $`\sigma _4=\sigma _{II_B}`$ (5.4.73)
$`\sigma _5=\sigma _{II_B}`$ $`\sigma _6=\sigma _{II_A}`$
This gives:
$$k_1(\sigma _I).[k_0(\sigma _I)+k_0(\sigma _{II})][\frac{k_1(\sigma _A).k_0(\sigma _B)}{z_Az_B}\frac{k_1(\sigma _B).k_0(\sigma _A)}{z_Bz_A}]i(p+q)^\mu Y_1^\mu e^{i(p+q)Y}$$
$$+k_1(\sigma _I).[k_0(\sigma _I)+k_0(\sigma _{II})][\frac{k_1(\sigma _B).k_0(\sigma _A)}{z_Bz_A}\frac{k_1(\sigma _A).k_0(\sigma _B)}{z_Az_B}]i(p+q)^\mu Y_1^\mu e^{i(p+q)Y}$$
(5.4.74)
We have used the same shorthand notation as in previous examples. Note that just as in previous cases interchanging $`\sigma _3`$ and $`\sigma _4`$ is not allowed. Note also that the two terms do not cancel. They add to give:
$$A(p).(p+q)\frac{S_{1,1}^{\mu \nu }(q)q^\mu q^\nu }{ϵ^2}$$
(5.4.75)
We have weighted it by a factor $`\frac{1}{2}`$ as in (LABEL:..).
The final result is
$$A(p).(p+q)\frac{S_{1,1}^{\mu \nu }(q)q^\mu q^\nu }{ϵ^2}|\frac{zw}{ϵ}|^{2p.q}(ϵ)^{(p+q)^2}$$
(5.4.76)
Integrals over $`w`$ and $`p,q`$ are implicit.
One Vertex Operator
Finally we have to consider the assignment where there is only one vertex operator and this clearly is singular and needs regularization. We will also observe a serious dependence on the prescription. This is not necessarily unacceptable. Presumably different prescriptions involve field redefinitions. If we impose physical state conditions on the fields these dependences should disappear.
$`z(\sigma _A)=z_A`$ $`\sigma _1=\sigma _A`$ $`\sigma _2=\sigma `$
$`z(\sigma _B)=z_B`$ $`\sigma _3=\sigma _B`$ $`\sigma _4=\sigma _A,\sigma _C`$
$`z(\sigma _C)=z_C`$ $`\sigma _5=\sigma _C`$ $`\sigma _6=\sigma _A,\sigma _B`$ (5.4.77)
We assign the momentum $`k`$ to the vertex operator. There are a total of 3! ways of assigning labels. So we weight each possibility by $`\frac{1}{3!}`$. What is given above is only one of the possibilities.
It gives
$$k_1(\sigma _A).[k_0(\sigma _A)+k_0(\sigma _B)+k_0(\sigma _C)][\frac{k_1(\sigma _B).k_0(\sigma _A)}{z_Bz_A}+\frac{k_1(\sigma _B).k_0(\sigma _C)}{z_Bz_C}]$$
$$[\frac{k_1(\sigma _C).k_0(\sigma _B)}{z_Cz_B}+\frac{k_1(\sigma _C).k_0(\sigma _A)}{z_Cz_A}]$$
(5.4.78)
Now we take the three points to be equidistant (this is a prescription) and this implies that
$`z_Bz_A`$ $`=`$ $`ϵ`$
$`z_Cz_B`$ $`=`$ $`ϵ`$
$`z_Cz_A`$ $`=`$ $`2ϵ`$ (5.4.79)
The expression in the second pair of square brackets vanishes and this term is zero. The term obtained by interchanging $`\sigma _A`$ and $`\sigma _C`$ also vanishes. Similarly the two terms that have $`\sigma _5=\sigma _B`$ also vanish. Thus four of the six possibilities give zero. The remaining two are given by the assignment
$`z(\sigma _A)=z_A`$ $`\sigma _1=\sigma _B`$ $`\sigma _2=\sigma `$
$`z(\sigma _B)=z_B`$ $`\sigma _3=\sigma _A`$ $`\sigma _4=\sigma _A,\sigma _C`$
$`z(\sigma _C)=z_C`$ $`\sigma _5=\sigma _C`$ $`\sigma _6=\sigma _A,\sigma _B`$ (5.4.80)
and the one obtained by interchanging $`\sigma _A`$ and $`\sigma _C`$ in this. This gives:
$$k_1(\sigma _B).[k_0(\sigma _A)+k_0(\sigma _B)+k_0(\sigma _C)][\frac{k_1(\sigma _A).k_0(\sigma _B)}{z_Az_B}+\frac{k_1(\sigma _A).k_0(\sigma _C)}{z_Az_C}]$$
$$[\frac{k_1(\sigma _C).k_0(\sigma _A)}{z_Cz_A}+\frac{k_1(\sigma _C).k_0(\sigma _B)}{z_Cz_B}]$$
(5.4.81)
Plugging the space time fields and the rest of the factors we get
$$S_{1,1,1}^{\mu \nu \rho }(k_0)k_0^\mu [\frac{k_0^\nu }{(ϵ)}+\frac{k_0^\nu }{(2ϵ)}][\frac{k_0^\rho }{(2ϵ)}+\frac{k_0^\rho }{(ϵ)}](ϵ)^{k_0^2}ik_0^\sigma Y_1^\sigma $$
(5.4.82)
$$=\frac{3}{4}S_{1,1,1}^{\mu \nu \rho }k_0^\mu k_0^\nu k_0^\rho (ϵ)^{k_0^22}ik_0^\sigma Y_1^\sigma $$
(5.4.83)
We have multiplied the answer by $`\frac{2}{3!}`$ as the weight for this term.
We have thus calculated the contribution from the first term to the the equation of motion.
What is to be noted is that the field $`S_{1,1,1}`$ is present as a result of the fact that we did not throw away the singular (normal ordering) pieces. This term will be indispensable in defining gauge transformations because there will be terms that cannot be assigned to any other field - in fact the presence of this term therefore guarantees that a gauge transformation can always be defined.
## 6 Space-Time Fields and their Transformations
Now we proceed to define fields. In the first approximation they were defined by the following equations:
$`<k_1^\mu >`$ $`=`$ $`A_1^\mu `$
$`<k_1^\mu (\sigma _1)k_1^\nu (\sigma _2)>`$ $`=`$ $`{\displaystyle \frac{D(\sigma _1\sigma _2)}{a}}S_{1,1}^{\mu \nu }+A_1^\mu A_1^\nu `$
$`<k_2^\nu >`$ $`=`$ $`S_2^\mu `$ (6.0.1)
We will define the gauge transformation laws for the space time fields by comparing the variations of the loop variable expression with the field expression. Thus consider expression (i)in Section 5.4 in both forms:
A(i) given in (5.3.37) and (5.3.40):
$$A(p).(p+q)iA^\mu (q)|\frac{zw}{ϵ}|^{2p.q}(ϵ)^{(p+q)^2}Y_2^\mu e^{i(p+q)Y}S_{1,1}^{\mu \nu }k_0^\nu iY_2^\mu e^{i{\scriptscriptstyle k_0Y}}(ϵ)^{k_0^2}$$
(6.0.2)
B(i) given in (5.3.36):
$$e^{{\scriptscriptstyle k_0(\sigma _5)}.k_0(\sigma _6)[\stackrel{~}{G}]}k_1(\sigma _1).k_0(\sigma _2)ik_1(\sigma )Y_2e^{i{\scriptscriptstyle k_0Y}}$$
(6.0.3)
Integrals over $`z,w`$ and $`p,q,k`$ are implicit in all the above. Thus the integral
$$𝑑z𝑑w|\frac{zw}{ϵ}|^{2p.q}=𝑑zF(p,q)$$
(6.0.4)
$`F(p,q)`$ is defined only after suitable regularization. The actual evaluation of this function will be done later.
Now we consider the variation of B(i):
$`\delta B(i)`$ $`=`$ $`[i\lambda _1(\sigma )k_0(\sigma _1).k_0(\sigma _2)k_1^\mu (\sigma _3)Y_2^\mu `$
$`i\lambda _1(\sigma )k_1(\sigma _1).k_0(\sigma _2)ik_0^\mu (\sigma _3)Y_2^\mu ]e^{k_0(\sigma _5).k_0(\sigma _6)\stackrel{~}{G}}e^{i{\scriptscriptstyle k_0Y}}`$
We convert this to space -time fields :
$`\delta B(i)`$ $`=`$ $`[i\mathrm{\Lambda }_{1,1}^\mu (k)k_0^2(ϵ)^{k_0^2}e^{ik_0Y}Y_2^\mu `$
$`i\mathrm{\Lambda }_1(p)A_1^\mu (q)(p+q)^2|{\displaystyle \frac{zw}{ϵ}}|^{2p.q}ϵ^{(p+q)^2}e^{i(p+q)Y}Y_2^\mu ]`$
$`+[i\mathrm{\Lambda }_{1,1}^\nu (k)k_0^\nu k_0^\mu (ϵ)^{k_0^2}e^{ik_0Y}Y_2^\mu `$
$`i\mathrm{\Lambda }_1(q)A_1^\nu (p)(p+q)^\nu (p+q)^\mu |{\displaystyle \frac{zw}{ϵ}}|^{2p.q}(ϵ)^{(p+q)^2}e^{i(p+q)}Y_2^\mu ]`$
This is to be compared with $`\delta A(i)`$: We will write
$`\delta S_{1,1}^{\mu \nu }(k)e^{ik_0Y}`$ $`=`$ $`[\mathrm{\Lambda }_{1,1}^\mu (k)k_0^\nu +\mathrm{\Lambda }_{1,1}^\nu (k)k_0^\mu +\delta _{int}S_{1,1}^{\mu \nu }]`$
$`\delta A_1^\mu (p)e^{ip_0Y}`$ $`=`$ $`p^\mu \mathrm{\Lambda }_1(p)e^{ip_0Y}`$ (6.0.7)
and determine $`\delta _{int}S_{1,1}^{\mu \nu }`$
We get
$$\delta B(i)=\delta _{free}A(i)+$$
$$+dpdq\delta (p+qk)[iq.(p+q)\mathrm{\Lambda }_1(p)A_1^\mu (q)i\mathrm{\Lambda }_1(q)p^\mu A(p).(p+q)]F(p,q)(ϵ)^{(p+q)^2}e^{i(p+q)Y}+$$
$$=\delta _{free}A(i)+\delta _{int}S_{1,1}^{\mu \nu }e^{ik_0Y}(ϵ)^{k_0^2}k_0^\nu $$
(6.0.8)
This fixes $`\delta _{int}S_{1,1}^{\mu \nu }`$ to be
$$\delta _{int}S_{1,1}^{\mu \nu }(k)=𝑑p𝑑q\delta (p+qk)[i\mathrm{\Lambda }_1(p)q^\nu A_1^\mu (q)i\mathrm{\Lambda }_1(q)p^\mu A_1(p)^\nu ]F(p,q)$$
(6.0.9)
B(ii) (5.3.44)
$$e^{k_0(\sigma _5).k_0(\sigma _6)\stackrel{~}{G}}k_2(\sigma _1).k_0(\sigma _2)ik_0^\mu (\sigma _3)e^{ik_0Y}Y_2^\mu $$
$$\delta B(ii)=e^{k_0(\sigma _5).k_0(\sigma _6)\stackrel{~}{G}}[\lambda _1(\sigma )k_1(\sigma _1).k_0(\sigma _2)+\lambda _2(\sigma )k_0(\sigma _1).k_0(\sigma _2)]ik_0^\mu (\sigma )e^{ik_0Y}Y_2^\mu $$
$$=\mathrm{\Lambda }_{1,1}^\mu (k_0)k_0^\nu ik_0^\mu (ϵ)^{k_0^2}e^{ik_0Y}$$
$$𝑑p𝑑q\delta (p+qk)𝑑w\mathrm{\Lambda }_1(p)A_1^\nu (q)(p_0+q_0)^\nu i(p_0+q_0)^\mu |\frac{zw}{ϵ}|^{2p.q}(ϵ)^{(p+q)^2}e^{i(p+q)Y}Y_2^\mu $$
$$\mathrm{\Lambda }_2(k_0)k_0^2ik_0^\mu (ϵ)^{k_0^2}Y_2^\mu e^{ik_0Y}$$
(6.0.10)
A(ii) (5.3.45)
$$A(ii)=S_2^\nu (k_0^\nu )ik_0(ϵ)^{k_0^2}Y_2^\mu e^{ik_0Y}$$
$$\delta S_2^\mu (k_0)=\mathrm{\Lambda }_2(k_0)k_0^\mu +\mathrm{\Lambda }_{1,1}^\mu (k_0)+\delta _{int}S_2^\mu $$
(6.0.11)
$$\delta A(ii)=\mathrm{\Lambda }_2(k_0)k_0^2ik_0^\mu (ϵ)^{k_0^2}Y_2^\mu e^{ik_0Y}$$
$$\mathrm{\Lambda }_{1,1}^\mu (k_0)k_0^\nu ik_0^\mu (ϵ)^{k_0^2}e^{ik_0Y}$$
$$\delta _{int}S_2^\mu k_0^\nu ik_0^\mu (ϵ)^{k_0^2}Y_2^\mu e^{ik_0Y}$$
(6.0.12)
Comparing (6.0.10) with (6.0.12) we find
$$\delta B(ii)=\delta A(ii)+\delta S_{2int}^\mu k_0^\nu ik_0^\mu (ϵ)^{k_0^2}Y_2^\mu e^{ik_0Y}$$
$$𝑑p𝑑q\delta (p+qk)\mathrm{\Lambda }_1(p)A_1^\nu (q)(p_0+q_0)^\nu i(p_0+q_0)^\mu \underset{F(p,q)}{\underset{}{𝑑w|\frac{zw}{ϵ}|^{2p.q}}}(ϵ)^{(p+q)^2}e^{i(p+q)Y}Y_2^\mu $$
(6.0.13)
From this we conclude that
$$\delta _{int}S_2^\mu (k_0)=𝑑p𝑑q\delta (p+qk)\mathrm{\Lambda }_1(p)A_1^\nu (q)F(p,q)$$
(6.0.14)
Thus we obtain the transformation rules for $`S_{1,1}`$ and $`S_2`$. Equations (iii) and (iv) are clearly consistent with this since they differ only in index structure.
It is not particularly illuminating to describe in detail the gauge transformation law for $`S_{1,1,1}`$ that one obtains in this manner since the calculation is very similar to that of $`S_{1,1}`$.
## 7 Consistency of Gauge Transformations and $`x_n`$-dependence of Fields
We examine, in this section, the question of consistency of gauge transformations of space-time fields defined in earlier sections. The question arises because there are different equations that can be used to define the gauge transformation law of $`S_{1,1}`$. For instance when one integrates by parts on $`x_1`$, different vertex operators such as $`Y_1`$ or $`Y_2`$ are obtained depending on whether one differentiates $`e^{ik_0Y}`$ twice or acts once each on $`e^{ik_0Y}`$ and $`e^{k_0.k_0\stackrel{~}{G}}`$. The dependence on $`zw`$ is thus different and one obtains instead of $`F(p,q)`$ (in (6.0.9)) some other function, and thus a different transformation law.
In fact $`F(p,q)`$ is a function of $`x_n`$ because $`\stackrel{~}{G}(zw)=ln(zw)+O(x_n)`$. Thus in principle one can ask what the result of differentiating (6.0.9) by $`x_n`$ is. The RHS of (6.0.9) is non-zero on differentiating and one reaches an inconsistency unless one assumes that the LHS also is non-zero - i.e. it must be a function of $`x_n`$ as well. This leads inexorably to the conclusion that the space-time fields such as $`S_{1,1}`$ must be functions of $`x_n`$, $`S_{1,1}(k_0,x_n)`$.
In the equation defining $`S`$ (5.2.14) there is a natural way to introduce this dependence, and this is to make the “string field” $`\mathrm{\Psi }`$ a function of $`x_n`$. Thus:
$$𝑑k_n𝑑\lambda _nk_1^\mu k_1^\nu \mathrm{\Psi }[k_n,k_0,\lambda _n,x_n]=S_{1,1}^{\mu \nu }(k_0,x_n)$$
(7.0.1)
Since $`\lambda _n`$ is to lowest order the shift in $`x_n`$, we can change variables to $`y_n`$, defined by
$$\underset{n}{}\lambda _nt^n=e^{_nt^ny_n}$$
and replace (7.0.1) by
$$𝑑k_n𝑑y_nk_1^\mu k_1^\nu \mathrm{\Psi }[k_n,k_0,y_n,x_n]=S_{1,1}^{\mu \nu }(k_0,x_n)$$
(7.0.2)
Both $`x_n`$ and $`y_n`$ are gauge coordinates. It is necessary therefore to understand the presence of both of these coordinates in the $`\mathrm{\Psi }`$. Our starting point is a field $`\mathrm{\Psi }(X(z),x_n+y_n)`$. The breakup of the gauge coordinate into $`x_n+y_n`$ is similar in spirit to that in the background field method in field theory. We treat $`x_n`$ as a background or reference point. Now we do a generalized Fourier transform using the loop variable and define the variable $`Y`$ which has in it only $`x_n`$. Thus the relation between the original string field and the one we have been using in this paper can be summarized in the following way: (We use the symbol $`\mathrm{\Psi }`$ for all the fields - the arguments of the fields will make clear which field we are referring to)
$$\mathrm{\Psi }(X,x_n+y_n)=\mathrm{\Psi }(Y_n,x_n,y_n)=[dk_n]\mathrm{\Psi }(k_n,x_n,y_n)e^{i_nk_nY_n}$$
Thus while the original field is only a function of $`x_n+y_n`$, once we define the variable $`Y`$ we have specified a reference point. The space-time fields obtained by (5.2.14) thus depend on this reference point. Gauge invariance is the statement that physics is independent of $`x_n+y_n`$. In terms of the new variables it becomes independence of $`x_n`$. Thus the $`k_n,y_n`$ integrals are in the nature of Fourier transformations, whereas the $`x_n`$ integral is an imposition of gauge invariance.
One can also do the integral over $`k_0`$:
$$𝑑k_0S_{1,1}(k_0,x_n)e^{ik_0Y}=S_{1,1}(Y,x_n).$$
Note $`S`$ depends explicitly on $`x_n`$ but also implicitly on $`x_n`$ through $`Y`$ because $`Y`$ depends on $`x_n`$ and all the derivatives of $`X(z)`$. Thus $`S`$ is a non-local object in that it depends on all the derivatives of $`X`$. To put it another way, specifying $`S`$ requires specifying a curve $`X(z)`$, because no two curves will produce the same value for $`S`$ for all $`x_n`$.
Thus the dependence on the infinite number of $`x_n`$ coordinates effectively makes $`S`$ non-local in $`z`$ and therefore $`X(z)`$. Of course the relevant scale here is the string scale so at low energies one can neglect the higher derivatives and effectively $`S`$ becomes an ordinary local field. It is also possible to redefine fields so that this non-locality disappears .
Note also that $`x_n`$ being along gauge directions we can fix gauge and set them to some fixed value. So there is no increase in the number of physical degrees of freedom. This is a desirable feature.
## 8 Conclusions
We have described a general construction that gives gauge invariant equations of motion, the gauge transformation prescription (in terms of loop variables) being the same as in II. This method was outlined in III but many of the details had not been worked out. One of the problems that was left unsolved was whether the map from loop variables to space time fields is unambiguous. In particular it was not obvious that there was a map that correctly reproduced the gauge transformations. The results of the present paper indicate that it is indeed possible to define space time fields and their gauge transformations consistently. There are two crucial ingredients. One is that one has to carefully keep all the singular terms that are normally discarded by “normal ordering”. We have to keep a finite cutoff in order for this procedure to make sense. This is not unexpected - we already know that in order to define off-shell Green functions in this approach, one needs a finite cutoff . As shown in , even U(1) gauge invariance of the massless vector is violated when a finite cutoff is introduced in order to go off-shell, and one needs to introduce massive modes to restore gauge invariance. In the loop variable formalism all the modes are present from the start and there is no problem. Gauge invariance is present, on or off-shell. However the exact value of the Koba-Nielsen integral will depend on the cutoff prescription. Presumably these are equivalent to field redefinitions (of the space time fields).
The second ingredient is that the string-field $`\mathrm{\Psi }`$ and thus the space-time fields are functions of the gauge coordinates $`x_n`$. This is crucial for consistency of the definitions of gauge transformations. Thus effectively “space-time” has become infinite dimensional!
At this stage we have a non-trivial interacting theory with an infinite tower of higher spin gauge fields and a large gauge invariance. By construction these modes are essentially <sup>6</sup><sup>6</sup>6“essentially” because of a technicality that is discussed in I those of the open bosonic string (including the auxiliary ones). Nevertheless we have not proved that the amplitudes of this theory are those of the bosonic string. We have to demonstrate that the procedure of “dimensional reduction” that worked for the free case goes through here also i.e. without loss to gauge invariance. If this works out then we are guaranteed that the on-shell amplitudes are those governed by the bosonic string simply because the two dimensional correlators that are being calculated here are identical to those of the bosonic string amplitude calculation. There are arguments that this is in fact the case . Furthermore as the gauge invariance does not use any on-shell conditions, these amplitudes are guaranteed to be gauge invariant off-shell also. Thus we have an off shell formulation. Further tests of the consistency of this will involve checking loop amplitudes. This is work for the future.
The main advantages are that the prescription for writing down the equations and gauge transformation laws are fairly straightforward. The gauge transformations written in terms of loop variables seem to have some geometric meaning - they look like local scale transformations. The interactions look as if they have the effect of converting a string to a membrane. The fields also appear massless in one higher dimension. These are intriguing features. Finally, assuming the above issues are resolved satisfactorily, one has to see whether this formalism provides any insight into the various other issues that have become pressing in string theory, such as duality. Some of the structure observed in may be relevant for this.
## Appendix A Appendix: Covariant Taylor Expansion
We derive the covariant Taylor expansion for $`\stackrel{~}{\mathrm{\Sigma }}(z(\sigma _1),z(\sigma _2),\sigma _1,\sigma _2)`$. We first derive a Taylor expansion for $`Y(z)`$ and then use it to obtain a Taylor expansion for $`\stackrel{~}{\mathrm{\Sigma }}`$.
### A.1 Taylor Expansion for Y
Ordinary Taylor expansion gives,
$$Y(z+a)=Y(z)+a\frac{dY}{dz}+\frac{a^2}{2}\frac{d^2Y}{dz^2}+\mathrm{}$$
(A.1)
$$Y\underset{n0}{}\alpha _n\stackrel{~}{Y}_n$$
(A.2)
where $`\stackrel{~}{Y}_n\frac{1}{(n1)!}\frac{^nX}{z^n}`$ and $`\alpha _n`$ satisfy
$$\underset{n0}{}\alpha _nt^n=e^{_{n0}x_nt^n}$$
$$\frac{\alpha _n}{x_m}=\alpha _{nm},$$
$$\frac{dY}{dz}=\stackrel{~}{Y}_1+\underset{n=1}{}n\alpha _n\stackrel{~}{Y}_{n+1},$$
$$=\stackrel{~}{Y}_1+\underset{n=1}{}nx_n\frac{}{x_{n+1}}Y.$$
(A.3)
In the above we have used $`_n[nx_n\frac{}{x_n}]\alpha _m=m\alpha _m`$.
Differentiating (A.3) gives
$$\frac{d^2Y}{dz^2}=\stackrel{~}{Y}_2+\underset{n=1}{}nx_n\frac{}{x_{n+1}}(\frac{dY}{dz})$$
Plugging in (A.3)
$$=\stackrel{~}{Y}_2+\underset{n=1}{}nx_n\frac{}{x_{n+1}}(\stackrel{~}{Y}_1+\underset{m=1}{}mx_m\frac{Y}{x_{m+1}})$$
$$\frac{d^2Y}{dz^2}=\stackrel{~}{Y}_2+\underset{n,m=1}{}nmx_nx_m\frac{Y}{x_{n+m+2}}+\underset{m=2}{}m(m1)x_{m1}\frac{Y}{x_{m+1}}$$
(A.4)
Adding (A.2),(A.3),(A.4) gives the first few terms of a Taylor series, except that we would like to express $`\stackrel{~}{Y}_i`$ in terms of $`Y_n`$ in order to make the expression covariant.
$`\stackrel{~}{Y}`$ in terms of Y:
We first write
$$\alpha (t)_zX(z+t)=\underset{n,m0}{}t^{mn}\alpha _n\stackrel{~}{Y}_{m+1}$$
(A.5)
Let $`\beta (t)=_{p0}\beta _pt^p`$. Let us evaluate $`\frac{dt}{t}\beta (t)\alpha (t)_zX(z+t)`$.
$$\frac{dt}{t}\beta (t)\alpha (t)_zX(z+t)=\underset{n,p,m=n+p}{}\beta _p\alpha _n\stackrel{~}{Y}_{m+1}$$
$$=\underset{m,p0}{}\beta _p\frac{}{x_{p+1}}\alpha _{m+1}\stackrel{~}{Y}_{m+1}$$
So
$$\frac{dt}{t}\beta (t)\alpha (t)_zX(z+t)=\underset{p0}{}\beta _p\frac{}{x_{p+1}}Y$$
(A.6)
Let us choose $`\beta (t)`$ having the property $`\beta (t)\alpha (t)=t^s;s0`$ and call it $`\beta ^s(t)`$ with the expansion
$$\beta ^s(t)=\underset{ps}{}\beta _p^st^p$$
Then (A.6) will become
$$\stackrel{~}{Y}_{s+1}=\underset{ps}{}\beta _p^sY_{p+1}$$
(A.7)
Thus if we determine $`\beta _p^s`$ we obtain the required expansion.
To determine $`\beta _p^s`$ we note that
$$\beta _p^st^p=\beta ^s(t)=t^s\alpha ^1(t)$$
$$=t^se^{_nx_nt^n}$$
$$=\underset{n0}{}\alpha _n(x_n)t^{ns}=\underset{ps}{}\alpha _{ps}(x_n)t^p$$
$$=\frac{}{(x_s)}\alpha _p(x_n,t)$$
This gives
$$\frac{\alpha _p}{x_s}_{x_nx_n}=\alpha _{ps}(x_n)=\beta _p^s$$
(A.8)
Thus for instance
$$\beta _0^0=1$$
$$\beta _1^0=x_1$$
$$\beta _2^0=\frac{x_1^2}{2}x_2$$
$$\beta _3^0=\frac{x_1^3}{6}+x_2x_1x_3$$
Therefore
$$\stackrel{~}{Y}_1=Y_1x_1Y_2+(\frac{x_1^2}{2}x_2)Y_3+\mathrm{}$$
(A.9)
It is easy to see that
$$\beta _p^s=\beta _{ps}^0$$
and so all the coefficients are easily determined.
Similarly
$$\frac{}{x_n}\beta _r^s=\beta _{rn}^s$$
Using this it is easy to see that
$$\frac{}{x_n}\stackrel{~}{Y}_s=0$$
as it should be.
Using the above results one obtains
$$\frac{dY}{dz}=\stackrel{~}{Y}_1+\underset{n=1}{}nx_n\frac{}{x_n}Y_1$$
$$=\underset{r0}{}(\beta _r^0+rx_r)Y_{r+1}$$
$$=\underset{r0}{}\gamma _r^0Y_{r+1}$$
(A.10)
Using the fact that $`\frac{^2}{x_nx_m}Y=\frac{}{x_{n+m}}Y`$ one can easily verify that $`\frac{^2}{x_nx_m}\frac{dY}{dz}=\frac{}{x_{n+m}}\frac{dY}{dz}`$. This is as it should be because the operations of differentiating w.r.t. z and w.r.t $`x_n`$ commute.
### A.2 Taylor Expansion of $`\stackrel{~}{\mathrm{\Sigma }}`$
We now use this to obtain the covariant expansion of $`\stackrel{~}{\mathrm{\Sigma }}`$. $`\stackrel{~}{\mathrm{\Sigma }}`$ was defined in section III
$$\stackrel{~}{\mathrm{\Sigma }}(z_1,z_2)=𝑑u\omega (u)<_uX(u)Y(va)><_uX(u)Y^{}(v+a)>$$
(A.11)
where $`z_2z_1=2a`$ and $`z_1+z_2=2v`$ and the contour encircles both points. We will call the gauge coordinates $`x_n`$ at $`z_1`$ and $`y_n`$ at $`z_2`$. The prime on $`Y`$ indicates indicates that it is a function of $`y_n`$.
We will use the shorthand notation $`<Y(va)Y^{}(v+a)>`$ for the above definition of $`\stackrel{~}{\mathrm{\Sigma }}`$. Thus we have the following Taylor expansion for $`\stackrel{~}{\mathrm{\Sigma }}`$:
$$<Y(va)Y^{}(v+a)>=<(Y(v)a\frac{dY}{dv}+\frac{a^2}{2}\frac{d^2Y}{dv^2}+\mathrm{})$$
$$(Y^{}(v)a\frac{dY^{}}{dv}+\frac{a^2}{2}\frac{d^2Y^{}}{dv^2}+\mathrm{})>$$
Plugging in the Taylor expansions for $`Y`$ that has been derived in this Appendix we get
$$=\overline{\mathrm{\Sigma }}(v)+a\underset{r0}{}(\gamma _r^0^{}\frac{}{y_{r+1}}\gamma _r^0\frac{}{x_{r+1}})\overline{\mathrm{\Sigma }}(v)+$$
$$a^2[\underset{r,s0}{}\gamma _r^0\gamma _s^0^{}\frac{^2}{x_{r+1}y_{s+1}}\overline{\mathrm{\Sigma }}+\frac{1}{2}\underset{r0}{}(\gamma _r^1\frac{}{x_{r+1}}+\gamma _r^{}_{}{}^{}1\frac{}{y_{r+1}})\overline{\mathrm{\Sigma }}]+O(a^3)$$
(A.12)
Here as before $`\overline{\mathrm{\Sigma }}(z,x_n,y_n)\stackrel{~}{\mathrm{\Sigma }}(z,x_n,z,y_n)`$
This is what has been used in section VI. |
warning/0002/astro-ph0002502.html | ar5iv | text | # Boltzmann equations for neutrinos with flavor mixings
## I introduction
The neutrino transport plays an important role in some astrophysical phenomena such as supernova explosions and the following proto neutron star cooling (e.g. and references therein). In their studies, the Boltzmann equation or its approximate versions are commonly employed to describe the temporal variations of neutrino distributions in the phase space. These equations are usually derived from the following assumptions: (1) the neutrinos are propagating along the geodesics for a massless particle, $`p^\mu p_\mu =0`$, and the volume in the phase space occupied by these neutrinos is not varied along their world line if there is no reaction. (2) the variation of the neutrino population due to reactions is described by the so-called collision terms obtained with the stohszahl ansatz. With the masses non-diagonal in the neutrino flavor space, the neutrino oscillation occurs among different flavors of neutrinos(e.g. and references therein). It is thus interesting from an academic point of view how this oscillation phenomenon is described by the generalized Boltzmann equations. It is also important from a practical point of view for those who are interested in the possible significant consequences the oscillation might give in astrophysical events. In collapse-driven supernova explosions, for example, this is particularly the case if the resonance of oscillation occurs near a neutrino sphere where neutrinos are interacting with other particles and thus the oscillation should be treated simultaneously with these reactions and possibly with the evolution of the matter distribution as well. The purpose of this paper is to provide the formulation which can be easily implemented in those numerical simulations.
In considering the transport equation with the oscillation, we have to rely on a more formal derivation of the Boltzmann equation. This might be done in a couple of ways. Sirera and Pérez, for instance, based their derivation on the relativistic Wigner function approach in the mean field approximation. Although they took the relativistic kinematics properly into account, they did not obtain the collision terms, since it is difficult to go beyond the mean field approximation in their formalism. Raffelt et al., on the other hand, obtained their transport equation via the density matrix approach. Although they derived the collision terms, they did not consider the spatially inhomogeneous system. In this paper, we derive the relativistic Boltzmann equation including corrections due to the oscillation both in the advection terms and the collision terms by employing the real time formalism of the nonequilibrium field theory. In this approach, the dispersion relation and the collision terms are derived on the same basis, that is, a particular approximation for the self-energy of neutrinos, which is conveniently represented with Feynman diagrams.
The paper is organized as follows. We first derive a generic form of the transport equation without specifying particular equations of motion of fields. Then, the formulation is applied to the neutrino flavor oscillations. In so doing, we ignore small corrections of the order of $`m_\nu ^2/E_\nu `$ except for the terms responsible for the flavor conversion, as is usually the case. Here $`m_\nu `$ and $`E_\nu `$ are typical mass and energy of neutrinos in the observer’s inertial frame. In this limit, as shown later, the left handed neutrinos are decoupled from the right handed ones and the difference between Majorana mass and Dirac mass never shows up in the flavor mixing. The general relativistic corrections are obtained up to the leading order of $`\lambda _\nu /R`$, where $`\lambda _\nu `$ is a typical wave length of neutrino and $`R`$ is a scale height of the background matter distribution.
## II formulation
### A general derivation of transport equations
In this section we derive general transport equations for multi-component fields based on the real time formalism of nonequilibrium field theory by Keldysh. In this formalism, we introduce path-ordered products of operators on the closed time-path, which extends from $`t=\mathrm{}`$ to $`t=+\mathrm{}`$ then back to $`t=\mathrm{}`$. In this product, the operator with a time argument which comes later on the above time-path is put to the left of other operators whose time arguments come earlier. Accordingly the path-ordered Green function is defined as
$$iG_{pij}(t_1,t_2)T_p\psi _i(t_1)\psi _j^{}(t_2).$$
(1)
Here $`T_p`$ stands for the path-ordered product of the following operators. The subscript $`ij`$ of the Green function denotes the components of the field. The bracket $`\mathrm{}`$ represents that arguments are averaged over the ensemble specified by a density operator $`\rho `$ as $`Tr\{\mathrm{}\rho \}`$, where $`Tr`$ is a trace operator. We define a generating functional of the Green function as
$$Z(J,J^{})Tr\left\{T_p\left[\mathrm{exp}\left(i\underset{i}{}_pd^4x\left(J_i^{}(x)\psi _i(x)+J_i(x)\psi _i^{}(x)\right)\right)\right]\rho \right\}\mathrm{exp}\left[iW(J,J^{})\right].$$
(2)
The Green function is obtained by the functional derivative, $`iG_{pij}(x,y)=\frac{\delta }{i\delta J_i^{}(x)}\frac{\delta }{i\delta J_j^{\text{}}(y)}Z(J,J^{})|_{J,J^{}=0}`$. The generating functional for the connected Green function is denoted as $`W(J,J^{})`$. Going to the interaction representation, we obtain
$`Z(J,J^{})`$ $`=`$ $`Tr\left\{T_p\left[\mathrm{exp}\left(i{\displaystyle \underset{i}{}}{\displaystyle _p}d^4x\left(J_i^{}(x)\psi _{\text{I}i}(x)+J_i(x)\psi _{\text{I}i}^{}(x)+_{int}(\psi _\text{I}(x),\psi _\text{I}^{}(x))\right)\right)\right]\rho _\text{I}\right\}`$ (3)
$`=`$ $`\mathrm{exp}\left[i{\displaystyle _p}d^4y_{int}({\displaystyle \frac{\delta }{i\delta J^{}}},{\displaystyle \frac{\delta }{i\delta J}})\right]Tr\left\{T_p\left[\mathrm{exp}\left(i{\displaystyle \underset{i}{}}{\displaystyle _p}d^4x\left(J_i^{}(x)\psi _{\text{I}i}(x)+J_i(x)\psi _{\text{I}i}^{}(x)\right)\right)\right]\rho _\text{I}\right\},`$ (4)
where the Lagrangian density for interactions is denoted as $`_{int}`$ and the subscript $`I`$ indicates that the variables are given in the interaction representation. The last factor of the right hand side of Eq. (4) is the generating functional for the no interaction case, $`Z_0(J,J^{})`$, and is given as
$$Z_0(J,J^{})=Z_{vac}(J,J^{})Tr\left\{\text{:}\mathrm{exp}\left(i\underset{i}{}_pd^4x\left(J_i^{}(x)\psi _{\text{I}i}(x)+J_i(x)\psi _{\text{I}i}^{}(x)\right)\right)\text{:}\rho _\text{I}\right\},$$
(5)
with the generating functional for vacuum,
$$Z_{vac}(J,J^{})=\mathrm{exp}\left(i\underset{ij}{}_pd^4xd^4yJ_j(y)G_{pij}^0(x,y)J_i^{}(x)\right).$$
(6)
Here $`G_{pij}^0(x,y)`$ is the path-ordered Green function for vacuum. The normal order product is represented by $`\text{:}\mathrm{}\text{:}`$ in Eq. (5). All the information of the ensemble is included in the last term of Eq. (5), $`N_p(J,J^{})=\mathrm{exp}[iW_p^N(J,J^{})]`$. Its connected part $`W_p^N(J,J^{})`$ is in general expanded to cumulants as
$`W_p^N(J,J^{})`$ $`=`$ $`{\displaystyle \underset{m,n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{m!n!}}{\displaystyle \mathrm{}_pd^4y_1\mathrm{}d^4y_nd^4x_1\mathrm{}d^4x_mJ(y_1)\mathrm{}J(y_n)J^{}(x_1)\mathrm{}J^{}(x_m)}`$ (8)
$`\times W_p^{Nmn}(x_1,\mathrm{},x_m|y_1,\mathrm{},y_n).`$
In the following we assume that the expansion is terminated at the quadratic order. This is true, for example, for the thermal equilibrium and the more general condition for this to be true can be found in the paper by Danielewicz. With this assumption, we can expand as usual the Green functions by the propagator which have corrections originating from a particular ensemble.
The Dyson equations are obtained by the Legendre transformations:
$$\mathrm{\Gamma }(\psi _c,\psi _c^{})=W(J,J^{})J^{}\psi _cJ\psi _c^{}$$
(9)
with $`\psi _c(x)=\frac{\delta }{\delta J^{}(x)}W`$ and $`\psi _c^{}(x)=\frac{\delta }{\delta J(x)}W`$. We use the abbreviation $`J^{}\psi _c=_i_pd^4xJ_i^{}(x)\psi _{ci}(x)`$. Then the following relations hold: $`\delta \mathrm{\Gamma }/\delta \psi _c(x)=J^{}(x)`$, $`\delta \mathrm{\Gamma }/\delta \psi _c^{}(x)=J(x)`$. The Dyson equations take the integral form on the closed time-path as
$$_pd^4zG_p^c(x,z)\mathrm{\Gamma }_p(z,y)=_pd^4z\mathrm{\Gamma }_p(x,z)G_p^c(z,y)=\delta _p(xy),$$
(10)
where the connected Green function and the vertex function are defined as $`iG_p^c(x,y)=\frac{\delta }{i\delta J^{}(x)}\frac{\delta }{i\delta J^{\text{}}(y)}W(J,J^{})`$ and $`\mathrm{\Gamma }_p(x,y)=\delta ^2\mathrm{\Gamma }/\delta \psi _c(y)\delta \psi _c^{}(x)`$, respectively. The $`\delta `$-function is extended on the closed time-path as follows: $`\delta _p(xy)=\delta (xy)`$ for $`t_x`$ and $`t_y`$ on the positive branch of the time-path extending from $`t=\mathrm{}`$ to $`t=+\mathrm{}`$ and $`\delta _p(xy)=\delta (xy)`$ for $`t_x`$ and $`t_y`$ on the negative branch of the time-path that runs from $`t=+\mathrm{}`$ to $`t=\mathrm{}`$. Introducing the matrix representations for the Green function and the vertex function as
$$\widehat{G}=\left(\begin{array}{cc}G_F& G_+\\ G_{}& G_{\overline{F}}\end{array}\right),\widehat{\mathrm{\Gamma }}=\left(\begin{array}{cc}\mathrm{\Gamma }_F& \mathrm{\Gamma }_+\\ \mathrm{\Gamma }_{}& \mathrm{\Gamma }_{\overline{F}}\end{array}\right),$$
(11)
we can recast the Dyson equation in a single time representation:
$$d^4z\widehat{G}^c(x,z)\sigma _3\widehat{\mathrm{\Gamma }}(z,y)=d^4z\widehat{\mathrm{\Gamma }}(x,z)\sigma _3\widehat{G}^c(z,y)=\sigma _3\delta (xy).$$
(12)
In the above equations, the time integration runs from $`t=\mathrm{}`$ to $`t=+\mathrm{}`$, and $`\sigma _3=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right)`$ is the Pauli matrix. The subscripts $`F`$ and $`\overline{F}`$ indicate that the time arguments are both on the positive branch and on the negative branch, respectively, while the subscript $`+`$ means that the first argument is located on the positive branch and the second on the negative branch, and the subscript $``$ represents the other way around. It is clear that $`G_F`$ is an ordinary Green function defined from the chronologically ordered product while $`G_{\overline{F}}`$ is obtained from the anti-chronological ordering. From these quantities, we further define the retarded, advanced and correlation functions as
$`G_r`$ $`=`$ $`G_FG_+,`$ (13)
$`G_a`$ $`=`$ $`G_FG_{},`$ (14)
$`G_c`$ $`=`$ $`G_F+G_{\overline{F}}.`$ (15)
The counter parts for the vertex functions are defined in an analogous way. Using the identity $`\mathrm{\Gamma }_F+\mathrm{\Gamma }_{\overline{F}}=\mathrm{\Gamma }_++\mathrm{\Gamma }_{}`$, we can express $`\mathrm{\Gamma }`$’s in general as
$`\mathrm{\Gamma }_\pm `$ $`=`$ $`i(B\pm A),`$ (16)
$`\mathrm{\Gamma }_F`$ $`=`$ $`D+iB,`$ (17)
$`\mathrm{\Gamma }_{\overline{F}}`$ $`=`$ $`D+iB,`$ (18)
$`\mathrm{\Gamma }_r`$ $`=`$ $`D+iA,`$ (19)
$`\mathrm{\Gamma }_a`$ $`=`$ $`DiA,`$ (20)
where $`A`$, $`B`$ and $`D`$ are three Hermitian matrices. Solving the Dyson equations using these quantities, we obtain the general form of the Green functions as
$`G_r`$ $`=`$ $`\left[D+iA\right]^1,`$ (21)
$`G_a`$ $`=`$ $`\left[DiA\right]^1,`$ (22)
$`G_\pm `$ $`=`$ $`i\left[D+iA\right]^1\left[B\pm A\right]\left[DiA\right]^1,`$ (23)
$`G_F`$ $`=`$ $`\left[D+iA\right]^1\left[DiB\right]\left[DiA\right]^1,`$ (24)
$`G_{\overline{F}}`$ $`=`$ $`\left[D+iA\right]^1\left[D+iB\right]\left[DiA\right]^1.`$ (25)
The dispersion relation is obtained from $`D`$ and $`A`$, while the distribution function is found from $`B`$ as shown shortly.
$`D`$, $`A`$ and $`B`$ can be represented in turn by the self-energy $`\mathrm{\Sigma }_p`$ which is defined from the two point vertex function as
$$\mathrm{\Gamma }_p=\mathrm{\Gamma }_{p0}\mathrm{\Sigma }_p,$$
(26)
where the free vertex function is $`\mathrm{\Gamma }_{p0}(xy)=S(_x)\delta _p(xy)`$. Here the derivative operator is taken from the free Lagrangian, $`_0=\psi ^{}S()\psi `$. Defining again the matrix components of the self-energy in the single time representation, we obtain
$`D`$ $`=`$ $`S(_x)\delta (xy){\displaystyle \frac{1}{2}}\left(\mathrm{\Sigma }_F\mathrm{\Sigma }_{\overline{F}}\right)`$ (27)
$`A`$ $`=`$ $`{\displaystyle \frac{1}{2}}i\left(\mathrm{\Sigma }_{}\mathrm{\Sigma }_+\right)`$ (28)
$`B`$ $`=`$ $`{\displaystyle \frac{1}{2}}i\left(\mathrm{\Sigma }_{}+\mathrm{\Sigma }_+\right).`$ (29)
The self-energy, on the other hand, is given by the relation
$$\mathrm{\Sigma }_p(x,y)=\left(iT_pj(x)j^{}(y)\delta _p(xy)\frac{\delta ^2}{\delta \psi (y)\delta \psi ^{}(x)}_{int}\right)_{1PI},$$
(30)
where the currents are defined as $`j(x)=\frac{\delta }{\delta \psi ^{}(x)}_{int}`$ and $`j^{}(x)=\delta _{int}/\delta \psi (x)`$, and the subscript $`1PI`$ means the one particle irreducible part.
Now we introduce the distribution function, $`n`$. First we define another Hermitian matrix $`N`$ from $`B`$ as
$$G_c=\mathrm{\Gamma }_r^12iB\mathrm{\Gamma }_a^1=\mathrm{\Gamma }_r^1NN\mathrm{\Gamma }_a^1.$$
(31)
Then it satisfies the following equation
$$NDDNi(NA+AN)=\mathrm{\hspace{0.17em}2}iB.$$
(32)
The matrix distribution function is finally defined as
$$N=12n,$$
(33)
where the upper and lower signs are taken for Fermion and Boson, respectively. It is easily shown that this distribution function becomes Fermi- or Bose-distribution function in the thermal equilibrium case. In that case, $`n`$ can be simultaneously diagonalized with $`D`$ and gives the distribution functions of quasi-particles. In general, however, $`n`$ has non-diagonal components even in the representation which diagonalizes $`D`$. These non-diagonal components are responsible for the flavor mixing as discussed below. Eq. (32) gives the equation satisfied by $`n`$:
$$nDDni(nA+An)=\pm i(BA)=\mathrm{\Sigma }_+.$$
(34)
Using Eq. (28), we can rewrite the above equation,
$$nDDn=\frac{1}{2}\left[(1n)(\mathrm{\Sigma }_+)+(\mathrm{\Sigma }_+)(1n)\right]\frac{1}{2}\left[n\mathrm{\Sigma }_{}+\mathrm{\Sigma }_{}n\right].$$
(35)
It is already clear that the right hand side of the above equation describes collisional processes among the quasi-particles. In fact, $`\left(\pm i\mathrm{\Sigma }_+\right)`$ and $`\left(i\mathrm{\Sigma }_{}\right)`$ can be interpreted as the emission and absorption rates of quasi-particle.
The transport equation as we know it is obtained by performing the so-called gradient expansion for the above equation. The Wigner representation of a quantity $`F(x,y)`$ is obtained by making Fourier transformation with respect to the relative coordinate as
$$F(k,X)=d^4(xy)e^{ik(xy)}F(x,y),$$
(36)
with the center of mass coordinate $`X=(x+y)/2`$. The gradient expansion is performed by taking the Wigner representation of both sides of Eq. (35) keeping only the leading order of the derivative with respect to $`X`$. Thus, we obtain the transport equation as
$`{\displaystyle \frac{1}{2}}[{\displaystyle \frac{D(k,X)}{k_\mu }}{\displaystyle \frac{n(k,X)}{X^\mu }}`$ $`+`$ $`{\displaystyle \frac{n(k,X)}{X^\mu }}{\displaystyle \frac{D(k,X)}{k_\mu }}]{\displaystyle \frac{1}{2}}[{\displaystyle \frac{D(k,X)}{X^\mu }}{\displaystyle \frac{n(k,X)}{k_\mu }}+{\displaystyle \frac{n(k,X)}{k_\mu }}{\displaystyle \frac{D(k,X)}{X^\mu }}]`$ (37)
$``$ $`i\left[D(k,X)n(k,X)n(k,X)D(k,X)\right]`$ (38)
$`=`$ $`{\displaystyle \frac{1}{2}}\left\{\left[1n(k,X)\right]\left[i\mathrm{\Sigma }_+(k,X)\right]+\left[i\mathrm{\Sigma }_+(k,X)\right]\left[1n(k,X)\right]\right\}`$ (39)
$``$ $`{\displaystyle \frac{1}{2}}\left\{n(k,X)\left[i\mathrm{\Sigma }_{}(k,X)\right]+\left[i\mathrm{\Sigma }_{}(k,X)\right]n(k,X)\right\}.`$ (40)
It is evident that the first row of the above equation represents ordinary advection terms while the right hand side stands for the collision terms. The second row, on the other hand, does not appear in the ordinary transport equation and we see below that this term causes the mixing among neutrino flavors. What remains now to do is to give the self-energy which determines not only the collision terms but also the dispersion relation, that is, $`D`$. We do this for the neutrino mixing in the next section.
### B neutrino transport equation with flavor mixings
We apply the general formulation obtained so far to the neutrino transport. The following Lagrangian density is considered:
$$=\{\begin{array}{cc}\frac{i}{2}\overline{\psi }_\text{L}\gamma ^\mu \underset{\mu }{\overset{}{}}\psi _\text{L}\frac{1}{2}\overline{\psi }\text{}_\text{L}^cM_\text{M}\psi _\text{L}\frac{1}{2}\overline{\psi }_\text{L}M_\text{M}^{}\psi _\text{L}^c+_{int}\hfill & \text{for Majorana }\nu \hfill \\ \frac{i}{2}\overline{\psi }_\text{L}\gamma ^\mu \underset{\mu }{\overset{}{}}\psi _\text{L}+\frac{i}{2}\overline{\psi }_\text{R}\gamma ^\mu \underset{\mu }{\overset{}{}}\psi _\text{R}\overline{\psi }_\text{R}M_\text{D}\psi _\text{L}\overline{\psi }_\text{L}M_\text{D}^{}\psi _\text{R}+_{int}\hfill & \text{for Dirac }\nu \hfill \end{array},$$
(41)
where the Majorana and Dirac masses are $`M_\text{M}`$ and $`M_\text{D}`$, respectively. The subscripts $`L`$ and $`R`$ stand for the spinor with left and right handed chirality, respectively, and $`\psi _\text{L}^c=C\overline{\psi }\text{}_\text{L}^\text{T}`$ with $`C`$ the charge conjugation and the superscript $`T`$ representing the transposition. The interaction Lagrangian density is denoted as $`_{int}`$. In the above equation, the indices for spinor components and neutrino flavors are suppressed. In the following, the flavor is denoted by the superscript and the spinor component by the subscript as $`\psi _i^a`$ when necessary. The matrix Green functions of interest are $`T_p\psi _\text{L}\text{}_i^a\overline{\psi }_\text{L}\text{}_j^b`$, $`T_p\psi _\text{L}\text{}_i^a\overline{\psi }_\text{R}\text{}_j^b`$, $`T_p\psi _\text{R}\text{}_i^a\overline{\psi }_\text{L}\text{}_j^b`$ and $`T_p\psi _\text{R}\text{}_i^a\overline{\psi }_\text{R}\text{}_j^b`$. Here and in the following, $`\psi _\text{R}`$ should be replaced by $`\psi _\text{L}^c`$ for the Majorana neutrino. We discuss the advection part and collision part of the Boltzmann equation separately, since we apply different approximations to the self-energies included in them.
#### 1 advection part
Following the common practice, we take the mean field approximation for the neutrino self-energy in the advection part, which is conveniently represented by a Feynman diagram shown in Fig. 1 and comes from the second term of Eq. (30). Only scattering processes contribute to this ensemble average. In the supernova core, the scatterings on free nucleons, nuclei and electrons are important. The former two of them occur only via neutral currents and as a result, the self-energies corresponding to them are proportional to the unit matrix in the flavor space:
$`_{int}^{nsc}`$ $`=`$ $`{\displaystyle \underset{a,N}{}}{\displaystyle \frac{G_F}{\sqrt{2}}}\left[\overline{\psi }_\text{L}\text{}^a\gamma ^\mu \left(1\gamma ^5\right)\psi _\text{L}^a\right]\left[\overline{\psi }_\text{N}\gamma _\mu \left(h_\text{N}^\text{V}h_\text{N}^\text{A}\gamma ^5\right)\psi _\text{N}\right],`$ (42)
$`\mathrm{\Sigma }_{F\text{LL}}^{ab}`$ $`=`$ $`\delta ^{ab}{\displaystyle \frac{G_F}{\sqrt{2}}}\gamma ^\mu \left(1\gamma ^5\right){\displaystyle \underset{N}{}}h_\text{N}^\text{V}\rho _\text{N}(X)\delta _{\mu 0},`$ (43)
$`\mathrm{\Sigma }_{F\text{RR}}^{ab}`$ $`=`$ $`\delta ^{ab}{\displaystyle \frac{G_F}{\sqrt{2}}}\gamma ^\mu \left(1+\gamma ^5\right){\displaystyle \underset{N}{}}h_\text{N}^\text{V}\rho _\text{N}(X)\delta _{\mu 0},`$ (44)
$`\mathrm{\Sigma }_{F\text{LR}}^{ab}`$ $`=`$ $`\mathrm{\Sigma }_{F\text{RL}}^{ab}=0,`$ (45)
where Eq. (44) is true only for the Majorana neutrino and $`\mathrm{\Sigma }_{F\text{RR}}^{ab}=0`$ for the Dirac neutrino. In the above equations, the subscript $`N`$ runs over neutron and proton, and $`\rho _\text{N}`$ stands for the nucleon number density. The similar equations are obtained for the scattering on nuclei. Hence, in the following, the nucleon scattering is considered.
On the other hand, the scattering on electrons gives a non-trivial structure to the self-energy in the flavor space since the process occurs not only through the neutral current but also through the charged current, and the latter is relevant only for the electron-type neutrinos in the matter in which electrons are abundant but other charged leptons are not. In that case, the interaction Lagrangian density becomes
$`_{int}^{esc}`$ $`=`$ $`{\displaystyle \frac{G_F}{\sqrt{2}}}\left[\overline{\psi }_\text{L}\text{}^{\nu _e}\gamma ^\mu \left(1\gamma ^5\right)\psi _\text{L}^{\nu _e}\right]\left[\overline{\psi }_e\gamma _\mu \left(\stackrel{~}{g}^\text{V}\stackrel{~}{g}^\text{A}\gamma ^5\right)\psi _e\right]`$ (46)
$`+`$ $`{\displaystyle \underset{a}{\overset{\nu _\mu ,\nu _\tau }{}}}{\displaystyle \frac{G_F}{\sqrt{2}}}\left[\overline{\psi }_\text{L}\text{}^a\gamma ^\mu \left(1\gamma ^5\right)\psi _\text{L}^a\right]\left[\overline{\psi }_e\gamma _\mu \left(g^\text{V}g^\text{A}\gamma ^5\right)\psi _e\right].`$ (47)
In the above equation, $`g^\text{V}=1/2+2\mathrm{sin}^2\theta _\text{W}`$ and $`g^\text{A}1/2`$ denote the vector and axial vector coupling constants of the neutral current, while the charged current is also taken into account in $`\stackrel{~}{g}^\text{V}=g^\text{V}+1`$ and $`\stackrel{~}{g}^\text{A}=g^\text{A}+1`$. The Weinberg angle is referred to as $`\theta _\text{W}`$ here. We obtain the self-energy of neutrino in the mean field approximation as
$`\mathrm{\Sigma }_{F\text{LL}}^{ab}`$ $`=`$ $`\delta ^{a\nu _e}\delta ^{b\nu _e}{\displaystyle \frac{G_F}{\sqrt{2}}}\gamma ^\mu \left(1\gamma ^5\right)\stackrel{~}{g}^\text{V}\rho _e(X)\delta _{\mu 0}`$ (48)
$`+`$ $`\delta ^{a\nu _{\mu ,\tau }}\delta ^{b\nu _{\mu ,\tau }}{\displaystyle \frac{G_F}{\sqrt{2}}}\gamma ^\mu \left(1\gamma ^5\right)g^\text{V}\rho _e(X)\delta _{\mu 0},`$ (49)
for the unpolarized electrons. Here the electron number density is denoted as $`\rho _e`$. As for the other components of the self-energy, $`\mathrm{\Sigma }_{F\text{LR}}=\mathrm{\Sigma }_{F\text{RL}}=0`$ common to both types of neutrinos, and $`\mathrm{\Sigma }_{F\text{RR}}=\mathrm{\Sigma }_{F\text{LL}}`$ with $`(1\gamma ^5)`$ replaced with $`(1+\gamma ^5)`$ for the Majorana neutrino and $`\mathrm{\Sigma }_{F\text{RR}}=0`$ for the Dirac neutrino. If the electrons are polarized in the magnetic field, the neutrino self-energy is modified to
$`\mathrm{\Sigma }_{F\text{LL}}^{ab}`$ $`=`$ $`\delta ^{a\nu _e}\delta ^{b\nu _e}{\displaystyle \frac{G_F}{\sqrt{2}}}\gamma ^\mu \left(1\gamma ^5\right)\left\{\stackrel{~}{g}^\text{V}\rho _e(X)\delta _{\mu 0}\stackrel{~}{g}^\text{A}\rho _e^0(X)\delta _{\mu z}\right\}`$ (50)
$`+`$ $`\delta ^{a\nu _{\mu ,\tau }}\delta ^{b\nu _{\mu ,\tau }}{\displaystyle \frac{G_F}{\sqrt{2}}}\gamma ^\mu \left(1\gamma ^5\right)\left\{g^\text{V}\rho _e(X)\delta _{\mu 0}g^\text{A}\rho _e^0(X)\delta _{\mu z}\right\},`$ (51)
where the magnetic field is parallel to the Z-axis. The electron number density in the lowest Landau level is represented as $`\rho _e^0`$. It is again true that the other components of the self-energy are zero except for $`\mathrm{\Sigma }_{F\text{RR}}=\mathrm{\Sigma }_{F\text{LL}}`$ with $`(1\gamma ^5)(1+\gamma ^5)`$ for the Majorana neutrino. It is easily understood that neutrino-neutrino scatterings can be treated just in the same way.
Now that we obtain the specific form of the neutrino self-energy, we can apply it to the left hand side of Eq. (37). Suppressing the flavor and spinor indices and writing only the chirality components in matrix form, we obtain $`D`$ in Eq. (37) using Eq. (27) as
$$D=\left(\begin{array}{cc}D_{\text{LL}}& D_{\text{LR}}\\ D_{\text{RL}}& D_{\text{RR}}\end{array}\right)=\left(\begin{array}{cc}k_\mu \gamma ^\mu \mathrm{\Phi }\gamma ^0\mathrm{\Phi }_\text{B}\gamma ^z& M^{}\\ M& k_\mu \gamma ^\mu +\mathrm{\Phi }^{}\gamma ^0+\mathrm{\Phi }_\text{B}^{}\gamma ^z\end{array}\right).$$
(52)
Here the potentials are defined as
$`\mathrm{\Phi }`$ $`=`$ $`\delta ^{a\nu _e}\delta ^{b\nu _e}\sqrt{2}G_F\stackrel{~}{g}^\text{V}\rho _e+\delta ^{a\nu _{\mu ,\tau }}\delta ^{b\nu _{\mu ,\tau }}\sqrt{2}G_Fg^\text{V}\rho _e,`$ (53)
$`\mathrm{\Phi }_\text{B}`$ $`=`$ $`\delta ^{a\nu _e}\delta ^{b\nu _e}\sqrt{2}G_F\stackrel{~}{g}^\text{A}\rho _e^0\delta ^{a\nu _{\mu ,\tau }}\delta ^{b\nu _{\mu ,\tau }}\sqrt{2}G_Fg^\text{A}\rho _e^0,`$ (54)
with $`\mathrm{\Phi }^{}=\mathrm{\Phi }`$ and $`\mathrm{\Phi }_\text{B}^{}=\mathrm{\Phi }_\text{B}`$ for the Majorana neutrino, and $`\mathrm{\Phi }^{}=0`$ and $`\mathrm{\Phi }_\text{B}^{}=0`$ for the Dirac neutrino. It is understood in the above equations that $`\mathrm{\Phi }_\text{B}=0`$ in the case of no magnetic field. The dispersion relations for quasi-particles are obtained from the eigen values of $`D`$.
We first make an order estimate of each term in the advection part. Defining $`R`$ as a typical length scale of the matter distribution, the density scale height, for example, $`E_\nu `$ as a typical energy of neutrino, and $`\mathrm{\Delta }m_\nu ^2`$ as a square mass difference, we find
$`{\displaystyle \frac{D}{k}}{\displaystyle \frac{n}{X}}`$ $``$ $`{\displaystyle \frac{n}{R}},`$ (55)
$`{\displaystyle \frac{D}{X}}{\displaystyle \frac{n}{k}}`$ $``$ $`{\displaystyle \frac{\mathrm{\Phi }}{R}}{\displaystyle \frac{n}{E_\nu }}{\displaystyle \frac{\mathrm{\Delta }m_\nu ^2}{E_\nu ^2}}{\displaystyle \frac{n}{R}},`$ (56)
$`Dn`$ $``$ $`{\displaystyle \frac{\mathrm{\Delta }m_\nu ^2}{E_\nu }}n{\displaystyle \frac{R}{\lambda }}{\displaystyle \frac{n}{R}}.`$ (57)
Here $`\lambda `$ is a wave length corresponding to $`\mathrm{\Delta }m_\nu ^2/E_\nu `$ : $`\lambda 0.1\mathrm{cm}[\mathrm{\Delta }m_\nu ^2/1\mathrm{e}\mathrm{V}^2]^1[E_\nu /1\mathrm{M}\mathrm{e}\mathrm{V}]`$. For the typical mass difference and energy of neutrino, $`\mathrm{\Delta }m_\nu ^2/E_\nu ^210^{12}`$. Hence the second term in the left hand side of Eq. (37), which represents the potential force exerted on neutrino by surrounding matter, is much smaller than the first term, which corresponds to the ordinary advection term in the Boltzmann equation. We ignore the former in the following discussion.
Next we show that the $`n_{\text{LL}}`$ can be decoupled from the other components assuming that $`\mathrm{\Delta }m_\nu ^2/E_\nu ^2`$ is neglected. We perform two matrix manipulations for $`DnnD`$ : (1) multiply the first row with $`k_\mu \gamma ^\mu +\mathrm{\Phi }^{}\gamma ^0+\mathrm{\Phi }_\text{B}^{}\gamma ^z`$ from the left and add to it the second row multiplied with $`M^{}`$ from the left. (2) then multiply the first column with $`k_\mu \gamma ^\mu +\mathrm{\Phi }^{}\gamma ^0+\mathrm{\Phi }_\text{B}^{}\gamma ^z`$ from the right and add to it the second column multiplied with $`M`$ from the right. Taking into account that $`n_{\text{LR}}n_{\text{RL}}(\mathrm{\Delta }m_\nu /E_\nu )n_{\text{LL}}`$ and $`\mathrm{\Phi }/E_\nu \mathrm{\Phi }_\text{B}/E_\nu \mathrm{\Delta }m_\nu ^2/E_\nu ^2`$, we obtain the equation for $`n_{\text{LL}}`$ as
$`[k_\mu k^\mu M^2`$ $``$ $`k_\mu \gamma ^\mu (\mathrm{\Phi }\gamma ^0+\mathrm{\Phi }_\text{B}\gamma ^z)+(\varphi ^{}\gamma ^0+\mathrm{\Phi }_\text{B}^{}\gamma ^z)k_\mu \gamma ^\mu ]n_{\text{LL}}k_\mu \gamma ^\mu `$ (58)
$``$ $`k_\mu \gamma ^\mu n_{\text{LL}}\left[k_\mu k^\mu M^2+k_\mu \gamma ^\mu \left(\mathrm{\Phi }^{}\gamma ^0+\mathrm{\Phi }_\text{B}^{}\gamma ^z\right)\left(\mathrm{\Phi }\gamma ^0+\mathrm{\Phi }_\text{B}\gamma ^z\right)k_\mu \gamma ^\mu \right].`$ (59)
The same manipulations are done for $`D/k_\mu n/X^\mu +n/X^\mu D/k_\mu `$ to obtain
$$k_\mu \gamma ^\mu \gamma ^\nu \frac{n_{\text{LL}}}{X^\nu }k_\sigma \gamma ^\sigma +k_\sigma \gamma ^\sigma \frac{n_{\text{LL}}}{X^\nu }\gamma ^\nu k_\mu \gamma ^\mu .$$
(60)
In the next subsection, it is shown that $`n_{\text{LL}}`$ can be also separated from the other components in the collision terms by applying the same procedures.
To the leading order of $`\mathrm{\Delta }m_\nu /E_\nu `$, $`n_{\text{LL}}`$ is a scalar with respect to the spinor index. The familiar form of advection terms in the Boltzmann equation is obtained by taking the trace with respect to the spinor indices after multiplying Eqs. (58) and (60) with $`\gamma ^0(1\gamma ^5)/2`$ from the left :
$$k^\mu \frac{n_{\text{LL}}^{ab}}{X^\mu }+ik^0\left\{\left(\frac{M^{2ac}}{2k^0}+\mathrm{\Phi }^{ac}+\frac{k^z}{k^0}\mathrm{\Phi }_\text{B}^{ac}\right)n_{\text{LL}}^{cb}n_{\text{LL}}^{ac}\left(\frac{M^{2cb}}{2k^0}+\mathrm{\Phi }^{cb}+\frac{k^z}{k^0}\mathrm{\Phi }_\text{B}^{cb}\right)\right\}.$$
(61)
Here the indices of flavor are explicitly included. From this equation we see that the resultant equation is identical for the Dirac and Majorana neutrinos up to the leading order of $`\mathrm{\Delta }m_\nu /E_\nu `$ and that $`D`$ is effectively replaced by $`D_{\text{eff}}`$ in the flavor space of the left handed neutrinos :
$$D_{\text{eff}}=\frac{1}{2}\left(k_\mu k^\mu M^22k^0\mathrm{\Phi }2k^z\mathrm{\Phi }_\text{B}\right).$$
(62)
The positive and negative zeros of $`D_{\text{eff}}`$ correspond to the energies of the neutrino and the anti-neutrino, respectively. The transport equation for the anti-neutrino is obtained with the replacements : $`k^\mu k^\mu `$, $`n_{\text{LL}}n_{\text{LL}}`$ in Eq. (61). In the following, we consider only the transport of on-shell neutrinos. Since we are not interested in the small difference $`\mathrm{\Delta }m_\nu ^2/E_\nu `$ of the on-shell energies among different flavors except in the terms responsible for the flavor mixing, we take $`k^0=|𝒌|`$ in Eq. (61).
In order to illuminate the structure of the advection part of the transport equation obtained above, we discuss only the two-flavor case of electron- and muon-neutrinos. Then Eq. (61) multiplied with $`i`$ becomes on the flavor basis
$$ik^\mu \frac{n^{ab}}{X^\mu }+k^0[H,n]^{ab},$$
(63)
with
$$H=\frac{\mathrm{\Delta }M^2}{2k^0}\mathrm{\Phi }\frac{k^z}{k^0}\mathrm{\Phi }_\text{B}=\frac{1}{2k^0}\left(\begin{array}{cc}\frac{\mathrm{\Delta }_0}{2}\mathrm{cos}2\theta _0\hfill & \frac{\mathrm{\Delta }_0}{2}\mathrm{sin}2\theta _0\hfill \\ \frac{\mathrm{\Delta }_0}{2}\mathrm{sin}2\theta _0\hfill & \frac{\mathrm{\Delta }_0}{2}\mathrm{cos}2\theta _0\text{}\hfill \end{array}\right)\left(\begin{array}{cc}\sqrt{2}G_F\rho _e\frac{k^z}{k^0}\sqrt{2}G_F\rho _e^0& 0\\ 0& 0\text{}\end{array}\right),$$
(64)
where $`\mathrm{\Delta }M^2`$ is a mass matrix with a diagonal matrix $`(m_1^2+m_2^2)/2\mathrm{𝟏}`$ subtracted. The eigen values of $`M^2`$ are $`m_1^2`$ and $`m_2^2`$ with the latter larger than the former, and their difference is defined to be $`\mathrm{\Delta }_0=m_2^2m_1^2`$. The mixing angle in vacuum is denoted as $`\theta _0`$ and the following relation holds :
$`M^2`$ $`=`$ $`U\left(\begin{array}{cc}m_1^2& 0\\ 0& m_2^2\end{array}\right)U^{},`$ (67)
$`U`$ $`=`$ $`\left(\begin{array}{cc}\mathrm{cos}\theta _0\hfill & \mathrm{sin}\theta _0\hfill \\ \mathrm{sin}\theta _0\hfill & \mathrm{cos}\theta _0\hfill \end{array}\right).`$ (70)
Eqs. (63) and (64) are the same equations as obtained by Sirera and Pérez.
Taking the bases on which $`H`$ in Eq. (64) is diagonalized at each point in space,
$$U_\text{M}^{}(x)H(x)U_\text{M}(x)=\frac{\stackrel{~}{M}^2(x)}{2k^0}\frac{1}{2k^0}\left(\begin{array}{cc}\stackrel{~}{M}_1^2(x)& 0\\ 0& \stackrel{~}{M}_2^2(x)\end{array}\right),$$
(71)
with
$$U_\text{M}(x)=\left(\begin{array}{cc}\mathrm{cos}\theta _\text{M}(x)\hfill & \mathrm{sin}\theta _\text{M}(x)\hfill \\ \mathrm{sin}\theta _\text{M}(x)\hfill & \mathrm{cos}\theta _\text{M}(x)\hfill \end{array}\right),$$
(72)
we obtain
$`U_\text{M}^{}(x)\left[ik^\mu {\displaystyle \frac{n}{X^\mu }}\right]U_\text{M}(x)`$ $`=`$ $`ik^\mu {\displaystyle \frac{\stackrel{~}{n}}{X^\mu }}[\stackrel{~}{n},U_\text{M}^{}(x)ik^\mu {\displaystyle \frac{U_\text{M}(x)}{X^\mu }}],`$ (73)
$`=`$ $`ik^\mu {\displaystyle \frac{\stackrel{~}{n}}{X^\mu }}[\stackrel{~}{n},\left(\begin{array}{cc}\hfill 0& \hfill 1\\ \hfill 1& \hfill 0\end{array}\right)ik^\mu {\displaystyle \frac{\theta _\text{M}(x)}{X^\mu }}].`$ (76)
Here the mass eigen values and mixing angle in matter are denoted as $`\stackrel{~}{M}_1`$, $`\stackrel{~}{M}_2`$ and $`\theta _\text{M}`$, respectively, and the distribution function in this representation is defined as $`\stackrel{~}{n}`$. In order to see the oscillation among different flavors, we write down each component of the above equation :
$`ik^\mu {\displaystyle \frac{\stackrel{~}{n}^{11}}{X^\mu }}`$ $``$ $`ik^\mu {\displaystyle \frac{\theta _\text{M}(x)}{X^\mu }}\left(\stackrel{~}{n}^{12}+\stackrel{~}{n}^{21}\right),`$ (77)
$`ik^\mu {\displaystyle \frac{\stackrel{~}{n}^{12}}{X^\mu }}`$ $`+`$ $`ik^\mu {\displaystyle \frac{\theta _\text{M}(x)}{X^\mu }}\left(\stackrel{~}{n}^{11}\stackrel{~}{n}^{22}\right)+{\displaystyle \frac{\mathrm{\Delta }_\text{M}}{2}}\stackrel{~}{n}^{12},`$ (78)
$`ik^\mu {\displaystyle \frac{\stackrel{~}{n}^{21}}{X^\mu }}`$ $``$ $`ik^\mu {\displaystyle \frac{\theta _\text{M}(x)}{X^\mu }}\left(\stackrel{~}{n}^{11}\stackrel{~}{n}^{22}\right){\displaystyle \frac{\mathrm{\Delta }_\text{M}}{2}}\stackrel{~}{n}^{21},`$ (79)
$`ik^\mu {\displaystyle \frac{\stackrel{~}{n}^{22}}{X^\mu }}`$ $`+`$ $`ik^\mu {\displaystyle \frac{\theta _\text{M}(x)}{X^\mu }}\left(\stackrel{~}{n}^{12}+\stackrel{~}{n}^{21}\right),`$ (80)
where the mass square difference in matter is defined as $`\mathrm{\Delta }_\text{M}=\stackrel{~}{M}_2^2\stackrel{~}{M}_1^2`$. Ignoring the collision terms for a moment and adding Eqs. (77) and (80), we obtain the relation
$$ik^\mu \frac{}{X^\mu }\left(\stackrel{~}{n}^{11}+\stackrel{~}{n}^{22}\right)=0,$$
(81)
which expresses the number conservation of neutrinos. From Eqs. (78) and (79) the following equations are obtained,
$`ik^\mu {\displaystyle \frac{}{X^\mu }}\left({\displaystyle \frac{\stackrel{~}{n}^{12}+\stackrel{~}{n}^{21}}{2}}\right)`$ $`+`$ $`{\displaystyle \frac{\mathrm{\Delta }_\text{M}}{2}}\left({\displaystyle \frac{\stackrel{~}{n}^{12}\stackrel{~}{n}^{21}}{2}}\right)+ik^\mu {\displaystyle \frac{\theta _\text{M}(x)}{X^\mu }}\left(\stackrel{~}{n}^{11}\stackrel{~}{n}^{22}\right)=0,`$ (82)
$`ik^\mu {\displaystyle \frac{}{X^\mu }}\left({\displaystyle \frac{\stackrel{~}{n}^{12}\stackrel{~}{n}^{21}}{2}}\right)`$ $`+`$ $`{\displaystyle \frac{\mathrm{\Delta }_\text{M}}{2}}\left({\displaystyle \frac{\stackrel{~}{n}^{12}+\stackrel{~}{n}^{21}}{2}}\right)=0.`$ (83)
Although we can infer the oscillating nature of the solution, it is better seen by eliminating $`\stackrel{~}{n}^{12}\stackrel{~}{n}^{21}`$ and taking only the leading terms of $`\mathrm{\Delta }m_\nu /E_\nu `$. The resultant equation roughly becomes
$$\frac{d^2}{d\mathrm{}^2}\left(\stackrel{~}{n}^{12}+\stackrel{~}{n}^{21}\right)+\left(\frac{\mathrm{\Delta }_\text{M}}{2E_\nu }\right)^2\left(\stackrel{~}{n}^{12}+\stackrel{~}{n}^{21}\right)+\frac{1}{R^2}\left(\stackrel{~}{n}^{11}\stackrel{~}{n}^{22}\right)0,$$
(84)
where $`\mathrm{}`$ is the path length and $`R`$ is the typical scale length of matter distribution. It is evident that $`\stackrel{~}{n}^{12}+\stackrel{~}{n}^{21}`$ have an oscillating part with an oscillation length of $`\lambda ^1=\mathrm{\Delta }_\text{M}/\mathrm{\hspace{0.17em}2}E_\nu `$ and a non-oscillating part which is negligible when the adiabatic condition $`\lambda /R1`$ is fulfilled. Hence we can ignore the non-diagonal components $`\stackrel{~}{n}^{12}`$ and $`\stackrel{~}{n}^{21}`$ of the matrix distribution function if we are interested only in the variation of the neutrino population on the length scale much longer than $`\lambda `$ and consider only the mass difference and energy of neutrino which satisfy the above adiabatic condition, as is usually true for the supernova cores and proto neutron stars. The non-diagonal components $`\stackrel{~}{n}^{12}`$ and $`\stackrel{~}{n}^{21}`$ can be ignored in the collision terms after taking the average of the rapidly oscillating terms over the length scale much larger than $`\lambda `$.
In the following, we set $`\stackrel{~}{n}^{12}=\stackrel{~}{n}^{21}=0`$ and consider the equations governing the diagonal components of the matrix distribution function for neutrinos. Following Raffelt et al., we represent $`\stackrel{~}{n}^{11}`$ and $`\stackrel{~}{n}^{22}`$ in terms of $`n^{\nu _e}`$ and $`n^{\nu _\mu }`$, the diagonal components on the flavor basis. From the relation
$$U_\text{M}^{}nU_\text{M}=\stackrel{~}{n}=\left(\begin{array}{cc}\stackrel{~}{n}^{11}& 0\\ 0& \stackrel{~}{n}^{22}\end{array}\right),$$
(85)
we obtain the distribution functions on the flavor basis as
$$n=\left(\begin{array}{cc}n^{\nu _e}& \frac{1}{2}\mathrm{tan}2\theta _\text{M}\left(n^{\nu _e}n^{\nu _\mu }\right)\\ \frac{1}{2}\mathrm{tan}2\theta _\text{M}\left(n^{\nu _e}n^{\nu _\mu }\right)& n^{\nu _\mu }\end{array}\right).$$
(86)
Inversely transforming the equation for $`\stackrel{~}{n}`$ in this approximation,
$$ik^\mu \frac{\stackrel{~}{n}}{X^\mu }=(\text{collision terms}),$$
(87)
we finally obtain the equation for $`n`$ in the limit of the adiabatic oscillation between two flavors as
$`ik^\mu {\displaystyle \frac{n}{X^\mu }}+\{\mathrm{tan}2\theta _\text{M}(n^{\nu _e}n^{\nu _\mu })\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right)`$ $``$ $`(n^{\nu _e}n^{\nu _\mu })\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)\}ik^\mu {\displaystyle \frac{\theta _\text{M}(X)}{X^\mu }}`$ (92)
$`=`$ $`(\text{collision terms})\text{}.`$ (93)
Here the mixing angle in matter is given as $`\mathrm{tan}2\theta _\text{M}=\mathrm{sin}2\theta _0/\{\mathrm{cos}2\theta _0[\mathrm{\hspace{0.17em}2}\sqrt{2}G_F(\rho _ek^z/k^0\rho _e^0)E_\nu ]/\mathrm{\Delta }_0\}`$.
We briefly discuss here the general relativistic correction terms in the advection part. We take an arbitrary point in space-time and consider a small patch of space-time around it of the size of $`\mathrm{}`$ which satisfies the condition $`R\mathrm{}1/E_\nu `$. Then we can take a local coordinates $`X_\text{R}`$ in this small region, which has a Minkowskian metric up to the second order of $`\mathrm{}/R`$. We also define an orthonormal tetrad $`𝒆_\text{R}^i`$ aligned to this coordinate and use it to project the four momentum of neutrino on it. On this coordinate in the small patch of space time, the above derivation for the advection terms in the Boltzmann equation is still valid, that is, we obtain Eq. (61) or Eq. (92) for the adiabatic two-flavor mixing with trivial replacements of $`X^\mu `$ with $`X_\text{R}^\mu `$ and $`k^\mu `$ with $`k_\text{R}^\mu `$. Thus, there is no additional mixing due to the general relativistic gravity under the current assumption that the tiny mass difference of neutrinos is ignored except for the mixing term, that is, the second terms of Eq. (61) or Eq. (92). All we have to do now is to make a coordinate transformation $`X_\text{R}^\mu X^\mu `$ and an associated momentum transformation $`k_\text{R}^\mu k^\mu `$. The latter is, in fact, induced by the transformation of the tetrads, $`𝒆_\text{R}^i𝒆^i`$, the latter of which is given globally. Employing the orthogonal transformation between two tetrads $`𝒆^i(X)=T_j^i(X)𝒆_\text{R}^j`$ and the transformation matrix
$$\left(\begin{array}{cc}\frac{X^\mu }{X_\text{R}^\nu }& \frac{k^j}{X_\text{R}^\nu }\\ \frac{X^\mu }{k_\text{R}^i}& \frac{k^j}{k_\text{R}^i}\text{}\end{array}\right)=\left(\begin{array}{cc}\frac{X^\mu }{X_\text{R}^\nu }& \frac{X^\rho }{X_\text{R}^\nu }\frac{k_\text{R}^m}{X^\rho }\left(T^1\right)_m^j\\ 0& \left(T^1\right)_i^j\text{}\end{array}\right),$$
(94)
we can perform the transformation for the advection term as follows:
$`k_\text{R}^i𝒆_{\text{R}i}^\mu {\displaystyle \frac{n}{X_\text{R}^\mu }}`$ $`=`$ $`k^i𝒆_i^\mu {\displaystyle \frac{n}{X^\mu }}k^i𝒆_i^\mu k^m{\displaystyle \frac{𝒆_m𝒆_\text{R}^j}{X^\mu }}𝒆_{\text{R}j}𝒆_n{\displaystyle \frac{n}{k^n}},`$ (95)
$`=`$ $`k^i𝒆_i^\mu {\displaystyle \frac{n}{X^\mu }}\omega _{im}^nk^ik^m{\displaystyle \frac{n}{k^n}}.`$ (96)
In the above equation, the inner product of two vectors is denoted as $`𝒗_1𝒗_2`$ and the component of the connection 1-form is designated as $`\omega _{ij}^k=_{𝒆_i}𝒆_j𝒆^k`$ with $`_{𝒆_i}`$ the covariant derivative in the direction of $`𝒆_i`$. The second term of Eq. (96) is a familiar correction term due to the general relativity, which accounts for the red shift and ray bending of neutrino in the gravitational field. Since the mixing term and the collision terms (see below) do not contain spatial derivatives, they are unaffected by the above transformation. However, $`\theta _\text{M}/X_\text{R}^\mu `$ in Eq. (92) is affected just in the same way as $`n/X_\text{R}^\mu `$ shown above. It is noted that this term actually originates from the advection term $`n/X_\text{R}^\mu `$ due to our pointwise choosing of the local mass eigen state basis. What remains to be done is, thus, to calculate the connection 1-form.
#### 2 collision part
In this section, we derive collision terms in the Born approximation for the neutrino self-energy. It is well known that the approximation of the self-energy for the advection terms is different from that for the collision terms. The Born approximation is conveniently represented by the Feynman diagram shown in Fig. 2. Only the first term of the right hand side of Eq. (30) contributes to $`\mathrm{\Sigma }_\pm `$ in the collision part. As done in the previous section, we evaluate the self-energy coming from various processes separately in the following.
For the nucleon scattering, the self-energies are obtained as
$`i\mathrm{\Sigma }_{\text{LL}}^{ab}(x,y)`$ $`=`$ $`{\displaystyle \frac{G_F^2}{2}}\left[\gamma ^\mu \left(1\gamma ^5\right)iG_{\text{LL}}^{ab}(x,y)\gamma ^\nu \left(1\gamma ^5\right)\right]S_{\text{N}\mu \nu }(x,y),`$ (97)
$`i\mathrm{\Sigma }_{+\text{LL}}^{ab}(x,y)`$ $`=`$ $`{\displaystyle \frac{G_F^2}{2}}\left[\gamma ^\mu \left(1\gamma ^5\right)iG_{+\text{LL}}^{ab}(x,y)\gamma ^\nu \left(1\gamma ^5\right)\right]S_{\text{N}\nu \mu }(y,x).`$ (98)
In the above equations, the dynamical structure function for nucleon is defined with the weak neutral current of nucleon $`J_\text{N}^\mu (x)`$ as $`S_\text{N}^{\mu \nu }(x,y)=J_\text{N}^\mu (x)J_\text{N}^\nu (y)`$. The weak neutral current for nucleon is given by $`J_\text{N}^\mu =\overline{\psi }_\text{N}\gamma _\mu \left(h_\text{N}^\text{V}h_\text{N}^\text{A}\gamma ^5\right)\psi _\text{N}`$. The other components of the matrix Green function are zero for the Dirac neutrinos. For the Majorana neutrinos, $`\mathrm{\Sigma }_{\text{LR}}`$ is obtained, for example, by replacing $`G_{\text{LL}}`$ with $`G_{\text{LR}}`$ and $`\gamma ^\nu \left(1\gamma ^5\right)`$ with $`\gamma ^\nu \left(1+\gamma ^5\right)`$. Recalling the relations $`n_{\text{LR}}n_{\text{RL}}(\mathrm{\Delta }m_\nu /E_\nu )n_{\text{LL}}(\mathrm{\Delta }m_\nu /E_\nu )n_{\text{RR}}`$ and $`\mathrm{\Sigma }_{\text{LR}}\mathrm{\Sigma }_{\text{RL}}(\mathrm{\Delta }m_\nu /E_\nu )\mathrm{\Sigma }_{\text{LL}}(\mathrm{\Delta }m_\nu /E_\nu )\mathrm{\Sigma }_{\text{RR}}`$, we find that the $`LL`$-component can be decoupled from the other components after the same matrix manipulations as done for the advection part and that the resulting collision terms are identical for the Dirac and Majorana neutrinos, if we take only the leading terms of $`\mathrm{\Delta }m_\nu /E_\nu `$. Note that the exchanged terms are added in the collision part unlike in the advection part. From the term $`ni\mathrm{\Sigma }_{}`$, for example, we obtain $`\left[k_\mu \gamma ^\mu n_{\text{LL}}\right]\left[i\mathrm{\Sigma }_{\text{LL}}k_\mu \gamma ^\mu \right]`$ as the $`LL`$-component.
Following the procedures taken for the advection part, we multiply the collision terms with $`\gamma ^0\left(1\gamma ^5\right)/2`$ from the left, take the trace with respect to the spinor indices and divide by $`4k^0`$. We then obtain from $`ni\mathrm{\Sigma }_{}`$, for example, the following:
$$\frac{d^4q}{(2\pi )^4}\frac{1}{4k^0}Tr\left\{\gamma ^0\frac{1\gamma ^5}{2}k_\rho \gamma ^\rho n_{\text{LL}}\frac{G_F^2}{2}\gamma ^\mu \left(1\gamma ^5\right)iG_{\text{LL}}(kq,X)\gamma ^\nu \left(1\gamma ^5\right)k_\sigma \gamma ^\sigma \right\}S_{\text{N}\mu \nu }(q,X).$$
(99)
Ignoring again the tiny masses of neutrinos and $`A`$ in deriving $`G_\pm `$ from $`D`$, $`A`$ and $`n`$, we obtain
$$iG_{\pm \text{LL}}(k)\frac{1\gamma ^5}{2}k_\mu \gamma ^\mu 2\pi \delta (k^2)\left[\mathrm{\Theta }(k^0)\left\{\begin{array}{c}n_\nu (𝒌)\\ 1n_\nu (𝒌)\end{array}\right\}+\mathrm{\Theta }(k^0)\left\{\begin{array}{c}1n_{\overline{\nu }}(𝒌)\\ n_{\overline{\nu }}(𝒌)\end{array}\right\}\right].$$
(100)
In the above equation, it is explicitly indicated that the negative energy contribution to the number density of the neutrino corresponds to the number density of the anti-neutrino. It is noted that the number density is a function of $`𝒌`$ after we ignored $`A`$ and imposed an on-shell condition $`k^2=0`$. The upper (lower) components in the columns correspond to $`G_+`$ ($`G_{}`$). Inserting this relation to Eq. (99) and recalling that $`n`$ is a scalar with respect to the spinor indices, we obtain the following collision term for the neutrino distribution function:
$$\frac{d^3k^{}}{(2\pi )^3}\frac{1}{2k^{}_{}{}^{}0}\frac{1}{2}n_\nu (k)\left[\mathrm{\hspace{0.17em}1}n_\nu (k^{})\right]\frac{G_F^2}{2}L^{\mu \nu }(k,k^{})S_{\text{N}\mu \nu }(q,X),$$
(101)
where $`k^{}=kq`$ is the four momentum of the scattered neutrino, and the tensor $`L^{\mu \nu }`$ is given as
$$L^{\mu \nu }(k,k^{})=Tr\left\{k_\rho \gamma ^\rho \gamma ^\mu (\mathrm{\hspace{0.17em}1}\gamma ^5)k_\sigma ^{}\gamma ^\sigma \gamma ^\nu (\mathrm{\hspace{0.17em}1}\gamma ^5)\right\}=8\left\{k^\mu k^{}_{}{}^{}\nu +k^\nu k^{}_{}{}^{}\mu g^{\mu \nu }k^\rho k_\rho ^{}i\epsilon ^{\mu \nu \rho \sigma }k_\rho k_\sigma ^{}\right\}.$$
(102)
Here the metric tensor is denoted as $`g^{\mu \nu }`$ and the anti-symmetric tensor as $`\epsilon ^{\mu \nu \rho \sigma }`$ with $`\epsilon ^{0123}=1`$. From the term $`i\mathrm{\Sigma }_{}n`$ we obtain the collision term which is obtained from Eq. (101) by replacing $`n_\nu (k)\left[\mathrm{\hspace{0.17em}1}n_\nu (k^{})\right]`$ with $`\left[\mathrm{\hspace{0.17em}1}n_\nu (k^{})\right]n_\nu (k)`$. Just in the same way we obtain from the term $`[1n][i\mathrm{\Sigma }_+]`$ the following collision term:
$$\frac{d^3k^{}}{(2\pi )^3}\frac{1}{2k^{}_{}{}^{}0}\frac{1}{2}\left[\mathrm{\hspace{0.17em}1}n_\nu (k)\right]n_\nu (k^{})\frac{G_F^2}{2}L^{\mu \nu }(k,k^{})S_{\text{N}\nu \mu }(q,X).$$
(103)
For the term $`[i\mathrm{\Sigma }_+][1n]`$ we replace $`\left[\mathrm{\hspace{0.17em}1}n_\nu (k)\right]n_\nu (k^{})`$ with $`n_\nu (k^{})\left[\mathrm{\hspace{0.17em}1}n_\nu (k)\right]`$ in the above equation. Using the relation $`S_{\text{N}\nu \mu }(q)=e^{\beta q^0}S_{\text{N}\mu \nu }(q)`$ for the matter in equilibrium, which stands for the detailed balance, we finally obtain the collision terms as
$`{\displaystyle \frac{d^3k^{}}{(2\pi )^3}\frac{1}{2k^{}_{}{}^{}0}}`$ $`{\displaystyle \frac{1}{2}}`$ $`\{e^{\beta q^0}{\displaystyle \frac{n_\nu (k^{})\left[\mathrm{\hspace{0.17em}1}n_\nu (k)\right]+\left[\mathrm{\hspace{0.17em}1}n_\nu (k)\right]n_\nu (k^{})}{2}}`$ (105)
$`{\displaystyle \frac{n_\nu (k)\left[\mathrm{\hspace{0.17em}1}n_\nu (k^{})\right]+\left[\mathrm{\hspace{0.17em}1}n_\nu (k^{})\right]n_\nu (k)}{2}}\}{\displaystyle \frac{G_F^2}{2}}L^{\mu \nu }(k,k^{})S_{\text{N}\mu \nu }(q,X).`$
If the matrix distribution function is diagonal, the above equation reduces to the ordinary collision term.
If the mixing occurs adiabatically, the above collision term is further simplified. For the two-flavor case ($`\nu _e`$ and $`\nu _\mu `$, for example), we insert the matrix distribution function given by Eq. (86) into Eq. (105). Then we obtain the collision term for the $`\nu _e`$ distribution function as
$`{\displaystyle \frac{d^3k^{}}{(2\pi )^3}\frac{1}{2k^{}_{}{}^{}0}}`$ $`{\displaystyle \frac{1}{2}}`$ $`\{e^{\beta q^0}[n^{\nu _e}(k^{})[\mathrm{\hspace{0.17em}1}n^{\nu _e}(k)]{\displaystyle \frac{1}{2}}\mathrm{tan}2\theta _\text{M}(k^{})[n^{\nu _e}(k^{})n^{\nu _\mu }(k^{})]{\displaystyle \frac{1}{2}}\mathrm{tan}2\theta _\text{M}(k)[n^{\nu _e}(k)n^{\nu _\mu }(k)]]`$ (107)
$`[n^{\nu _e}(k)[\mathrm{\hspace{0.17em}1}n^{\nu _e}(k^{})]{\displaystyle \frac{1}{2}}\mathrm{tan}2\theta _\text{M}(k)[n^{\nu _e}(k)n^{\nu _\mu }(k)]{\displaystyle \frac{1}{2}}\mathrm{tan}2\theta _\text{M}(k^{})[n^{\nu _e}(k^{})n^{\nu _\mu }(k^{})]]\}`$
$`\times `$ $`{\displaystyle \frac{G_F^2}{2}}L^{\mu \nu }(k,k^{})S_{\text{N}\mu \nu }(q,X).`$ (108)
The collision term for the $`\nu _\mu `$ distribution function is obtained by replacing $`n^{\nu _e}`$ with $`n^{\nu _\mu }`$ in the above equation. It is noted that the correction terms due to mixing cancel each other for iso-energetic scatterings, which we commonly assume for the neutrino-nucleon scattering in the supernova cores and proto neutron stars.
The collision terms for the neutrino-electron scattering are essentially the same as those obtained for the nucleon scattering. The main difference originates from the fact that the electron weak current has flavor dependence, which gives rise to non-trivial contractions of flavor indices between the electron structure function and the neutrino distribution function such as $`n_\nu ^{ac}(k)\left[1n_\nu ^{cb}\right]S_e^{cb}`$, where the superscripts $`a,b,c`$ represent flavors. The structure function $`S_e`$ is an electron counter part of the nucleon structure function $`S_\text{N}`$ and is defined as $`S_{e\mu \nu }^{ab}(x,y)=J_{e\mu }^a(x)J_{e\nu }^b(y)`$. Here the weak current for electron is given as $`J_e^\mu =\overline{\psi }_e\gamma ^\mu \left(\stackrel{~}{g}^\text{V}\stackrel{~}{g}^\text{A}\gamma ^5\right)\psi _e`$ for the $`\nu _e`$ scattering and $`J_e^\mu =\overline{\psi }_e\gamma ^\mu \left(g^\text{V}g^\text{A}\gamma ^5\right)\psi _e`$ for the $`\nu _\mu `$ and $`\nu _\tau `$ scatterings, respectively. As a result, the collision term for the neutrino-electron scattering becomes for the electron type neutrino in the case of the adiabatic two-flavor mixing as
$`{\displaystyle \frac{d^3k^{}}{(2\pi )^3}\frac{1}{2k^{}_{}{}^{}0}}`$ $`{\displaystyle \frac{1}{2}}`$ $`\{\text{}[e^{\beta q^0}n^{\nu _e}(k^{})[\mathrm{\hspace{0.17em}1}n^{\nu _e}(k)]n^{\nu _e}(k)[\mathrm{\hspace{0.17em}1}n^{\nu _e}(k^{})]]S_{e\mu \nu }^{\nu _e\nu _e}(q,X)`$ (109)
$``$ $`\left[e^{\beta q^0}\mathrm{\hspace{0.17em}1}\right]{\displaystyle \frac{1}{2}}\mathrm{tan}2\theta _\text{M}(k^{})\left[n^{\nu _e}(k^{})n^{\nu _\mu }(k^{})\right]{\displaystyle \frac{1}{2}}\mathrm{tan}2\theta _\text{M}(k)\left[n^{\nu _e}(k)n^{\nu _\mu }(k)\right]`$ (110)
$`\times `$ $`{\displaystyle \frac{S_{e\mu \nu }^{\nu _e\nu _\mu }(q,X)+S_{e\mu \nu }^{\nu _\mu \nu _e}(q,X)}{2}}\}`$ (111)
$`\times `$ $`{\displaystyle \frac{G_F^2}{2}}L^{\mu \nu }(k,k^{}).`$ (112)
The $`\nu _\mu `$ counterpart is obtained by the replacement of $`\nu _e\nu _\mu `$ in the flavor indices in the above equation.
Although we assumed in the above derivation that the four momentum transfer is space-like to describe the scatterings, it is obvious that the same Feynman diagram represents the annihilation and creation of neutrino pairs if the transferred four momentum is time-like. As stated above, the transport equations for the anti-neutrinos are obtained from the negative energy part of the distribution function. in that case, it is noted that the mixing angle $`\theta _\text{M}`$ should be also calculated for the negative energy. As is obvious from Eq. (64) the sign of the potentials is changed for the anti-neutrinos and the resonance conversion does not occur in this case as is well known. It is noted that neutrino-neutrino scatterings are treated just in the same way by substituting the neutrino structure function, which in turn should be evaluated with the neutrino Green functions, Eq. (100).
Next we consider the neutrino emission and absorption on nucleons. For the temperature and neutrino energy of current interest, the muon is not abundant and only the electron-type neutrino is involved in this process. The interaction Lagrangian density is
$$_{int}^{abs}=\frac{G_F}{\sqrt{2}}\left[\overline{\psi }_\text{e}\gamma ^\mu \left(1\gamma ^5\right)\psi _\text{L}^{\nu _e}\right]\left[\overline{\psi }_\text{p}\gamma _\mu \left(g^\text{V}g^\text{A}\gamma ^5\right)\psi _\text{n}\right]+H.C.,$$
(113)
where the coupling constants $`g^\text{V}=1.0`$ and $`g^\text{A}=1.23`$, and $`H.C.`$ stands for the Hermite conjugate. In the Born approximation, the self-energy is given by
$`i\mathrm{\Sigma }_{\text{LL}}^{\nu _e\nu _e}(x,y)`$ $`=`$ $`{\displaystyle \frac{G_F^2}{2}}\left[\gamma ^\mu \left(1\gamma ^5\right)iG_{}^e(x,y)\gamma ^\nu \left(1\gamma ^5\right)\right]S_{\text{pn}\mu \nu }(x,y),`$ (114)
$`i\mathrm{\Sigma }_{+\text{LL}}^{\nu _e\nu _e}(x,y)`$ $`=`$ $`{\displaystyle \frac{G_F^2}{2}}\left[\gamma ^\mu \left(1\gamma ^5\right)iG_+^e(x,y)\gamma ^\nu \left(1\gamma ^5\right)\right]S_{\text{np}\nu \mu }(y,x),`$ (115)
where the Green functions for electrons are denoted as $`G_\pm ^e`$ and the structure functions for the charged weak currents of nucleons are defined, for example, as $`S_{pn}^{\mu \nu }(x,y)=J_{pn}^\mu (x)J_{np}^\nu (y)`$ with the charged current given by $`J_{np}^\mu =\overline{\psi }_\text{p}\gamma _\mu \left(g^\text{V}g^\text{A}\gamma ^5\right)\psi _\text{n}`$. The Green functions $`G_\pm ^e`$ are essentially the same as $`G_{\pm \text{LL}}`$ in Eq.(100) with the self-evident substitution of $`n_\nu `$ with the electron distribution function $`f_e`$. Following the same procedure as shown above and employing the detailed balance relation satisfied by nucleons in thermal equilibrium,
$$S_{np}^{\nu \mu }(q)=e^{\beta (q^0+\mathrm{\Delta }\mu _{np})}S_{pn}^{\mu \nu }(q),$$
(116)
with the difference of the chemical potentials, $`\mathrm{\Delta }\mu _{np}=\mu _n\mu _p`$, we obtain the collision term for the emission and absorption of neutrinos on nucleons as
$$\frac{d^3p_e}{(2\pi )^3}\frac{1}{2E_e}\frac{1}{2}\left\{e^{\beta (q^0+\mathrm{\Delta }\mu _{np})}f_e(p_e)\left[\mathrm{\hspace{0.17em}1}n^{\nu _e}(k)\right]n^{\nu _e}(k)\left[\mathrm{\hspace{0.17em}1}f_e(p_e)\right]\right\}\frac{G_F^2}{2}L_{\mu \nu }(k,p_e)S_{pn}^{\mu \nu }(q),$$
(117)
in the adiabatic mixing case. Here $`p_e`$ and $`E_e`$ are the momentum and energy of electrons, respectively, and the transfer four momentum is $`q=kp_e`$. Since there is no other components of self-energy in flavor space than the $`\nu _e\nu _e`$ component, the resulting term is identical to those with no neutrino mixing.
## III summary
With a view of application to the simulations of supernova explosion and proto neutron star cooling, we have derived a Boltzmann equation with the neutrino flavor mixing being taken into account. The derivation is based on the nonequilibrium field theory, and the ordinary gradient expansion has been performed. We assumed that the typical neutrino wave length is much shorter than the scale height of the background matter distribution, which is true for the supernova cores and proto neutron stars. The neutrino distribution matrix which is non-diagonal in the neutrino flavor space is introduced. Following the common practice, the advection part has been obtained in the mean field approximation, where the self-energy of neutrino is non-diagonal in the flavor space. This self-energy gives rise to the term in the advection part, which is responsible for the neutrino mixing and does not appear in the ordinary transport equation. The collision terms, on the other hand, have been calculated in the Born approximation. The collision terms also have corrections due to the mixing. In these derivations, the relativistic kinematics is taken into consideration. We have further simplified the Boltzmann equation for the adiabatic flavor mixing, which is a good approximation in the supernova cores and proto neutron stars. The advection terms thus derived are essentially the same as those derived by Sirera and Pérez, although they employed the Wigner function formalism in the mean field approximation and did not give collision terms. The collision terms derived here, on the other hand, have the same structure as those found by Raffelt et al. in the non-relativistic density matrix method. We have also shown the general relativistic correction term which accounts for the red shift and ray bending in the gravitational field and is commonly taken into account in the supernova and proto neutron star simulations.
The applications of the Boltzmann equation found here remain to be done. Since the corrections due to the flavor mixing are rather minor, particularly in the case of the adiabatic mixing, it will be simple to implement them in the neutrino transport code we have now at our disposal. This is already underway. Since the mixing angle in matter is dependent on the neutrino energy and the direction of momentum with respect to the magnetic field if it exists. In the analyses of the neutrino flavor mixing in the supernova core, it is usually assumed that the neutrinos are flowing out radially. However, they have an angular distribution near the neutrino sphere. Different positions of the resonant conversion due to different directions of flight of neutrinos will lead to the reduction of the neutrino flavor conversion. This will also be true in the absence of the magnetic field if the energy distribution of neutrinos and the coupling between neutrinos with different energies are taken into account. These possibilities and their implications to the mechanism of the supernova explosion, kick velocity of pulsars, and nucleosynthesis of heavy elements will be studied in the forthcoming papers.
###### Acknowledgements.
I gratefully acknowledge some comments by A. Pérez. This work is partially supported by the Grants-in-Aid for the Center-of-Excellence (COE) Research of the Ministry of Education, Science, Sports and Culture of Japan to RESCEU (No.07CE2002). |
warning/0002/hep-th0002252.html | ar5iv | text | # Black Holes and U-duality in Diverse Dimensions
## 1 Preliminaries
The purpose of this work is to investigate some properties of black hole solutions of different supergravity theories in terms of the corresponding U-duality group.
Black holes are solutions of supergravity theories with point-like sources. They are rotationally invariant and are, in general, charged with respect to the vector fields of the theory. We will consider extremal black hole solutions, that is, solutions on which some of the supersymmetry charges $`Q`$ are null. In this case, the solutions are parametrized by the set of charges $`q^\mathrm{\Lambda }`$ and the values of the asymptotic values of the scalar fields $`\varphi ^i`$ of the theory at $`\mathrm{}`$ (moduli).
We consider the $`N`$-extended super Poincaré algebra in a $`D=4,\mathrm{}9`$ dimensional space-time (we follow the notation of Ref.). For $`D=2n+1`$ the irreducible representation of the Clifford algebra $`𝒞(1,D1)`$, of dimension $`2^n`$, is also irreducible with respect to so(1,$`D1`$). This is the only spinor representation and consequently it is self-conjugated. For $`D=9`$ it is real, and the anticommutator of two supercharges is given by
$$\{Q_{\alpha i},Q_{\beta j}\}=\left(\mathrm{\Gamma }_\mu Cp^\mu \right)_{\alpha \beta }\delta _{ij}+C_{\alpha \beta }Z_{ij}$$
(1)
where $`i,j=1,\mathrm{}N`$, $`\alpha ,\beta =1,\mathrm{}2^4`$. $`C`$ is the charge conjugation matrix that in this case is symmetric, so the central charge extension of the super Poincaré algebra, $`Z_{ij}`$, is also symmetric.
For $`D=5,7`$ the spinor representation is pseudoreal, so the anticommutator is instead
$$\{Q_{\alpha i},Q_{\beta j}\}=\left(\mathrm{\Gamma }_\mu Cp^\mu \right)_{\alpha \beta }\mathrm{\Omega }_{ij}+C_{\alpha \beta }Z_{ij},$$
(2)
where $`\mathrm{\Omega }`$ is the symplectic bilinear form. For $`D=5`$, $`C=C^T`$, so $`Z`$ is antisymmetric and for $`D=7`$ $`C=C^T`$ so $`Z`$ is symmetric.
For $`D=2n`$ the irreducible representation of the Clifford algebra $`𝒞(1,D1)`$ splits into two irreducible pieces under so(1,$`D1`$), (2$`{}_{}{}^{n1})_\pm `$. The projectors on each piece are
$$𝒫_\pm =\frac{1\pm \mathrm{\Gamma }_{D+1}}{2},Q_\pm =𝒫_\pm Q.$$
For $`D=4,8`$ these representations are complex and pairwise conjugate. The anticommutator is
$`\{Q_{\alpha i+},Q_j^{\dot{\beta }}\}`$ $`=`$ $`\left(𝒫_+\mathrm{\Gamma }_\mu Cp^\mu \right)_\alpha ^{\dot{\beta }}\delta _{ij}`$
$`\{Q_{\alpha i+},Q_{\beta j+}\}`$ $`=`$ $`(𝒫_+C)_{\alpha \beta }Z_{ij}.`$ (3)
When $`D=4`$, $`C=C^T`$, so $`Z`$ is antisymmetric and when $`D=8`$, $`C=C^T`$ and $`Z`$ is symmetric.
Finally, for $`D=6`$ the two spinor representations are pseudoreal and independent. This means that one can choose independently the number of supersymmetry charges with chirality $`+`$ or $``$ ($`N_+`$ and $`N_{}`$).
$`\{Q_{\alpha i+},Q_{\beta j+}\}`$ $`=`$ $`\left(𝒫_+\mathrm{\Gamma }_\mu Cp^\mu \right)_{\alpha \beta }\mathrm{\Omega }_{ij}`$
$`\{Q_{\alpha i+},Q_{\beta j}\}`$ $`=`$ $`(𝒫_+CZ_{ij})_{\alpha \beta }`$
$`\{Q_{\alpha i},Q_{\beta j}\}`$ $`=`$ $`\left(𝒫_{}\mathrm{\Gamma }_\mu Cp^\mu \right)_{\alpha \beta }\mathrm{\Omega }_{ij}.`$ (4)
There are some rotations of the charge vectors $`Q_{\alpha i}`$ that leave invariant the momentum term in the anticommutator of supercharges. These transformations are automorphisms of the super Poincaré algebra, and they form the $`R`$-symmetry group of the algebra. It is a compact group that we will denote by $`H`$. The nature of this group obviously depends on the reality properties the spinors. For $`D=4,\mathrm{}9`$ these are
$`R`$-symmetry group $`H`$
$`D=9`$ SO(N)
$`D=8`$ $`\text{SU(N)}\times \text{U(1)}`$
$`D=7`$ USp(N)
$`D=6`$ $`\text{USp}(N_+)\times \text{USp}(N_{})`$
$`D=5`$ USp(N)
$`D=4`$ $`\text{SU(N)}\times \text{U(1)}`$
We consider theories with a maximal number of supesymmetries, $`n=32`$. From the symmetry properties of $`Z`$ stated above, it follows that the central charge transforms in the following representation of the $`R`$-symmetry group,
Central charge representation of the $`R`$-symmetry
$`D=9`$ 3 of SO(2), ( real symmetric tensor). (5)
$`D=8`$ 3(+) of SU(2)$`\times `$U(1), (complex triplet).
$`D=7`$ 10 of USp(4), (real symmetric tensor).
$`D=6`$ $`\text{16}\text{ of USp(4)}\times \text{USp(4)},`$
(bispinor (4,4) of O(5)$`\times `$O(5)).
$`D=5`$ 27 of USp(8), ($`\mathrm{\Omega }`$-traceless
symplectic antisymmetric tensor).
$`D=4`$ 28 of SU(8), (complex antisymmetric
tensor).
We will see that it is always possible to put the central charge matrix $`Z_{ij}`$ in normal form (diagonal or skew diagonal) using a transformation of the $`R`$-symmetry group (we denote the eignevalues by $`Z_A`$). Then, using the theory of induced representations, one can go to the rest frame $`(p=(m,0,\mathrm{},0)`$ and the supercharge anticommutator becomes
$`\{S_A,S_B^{}\}(m+Z_A)\delta _{AB},`$
$`\{\widehat{S}_A,\widehat{S}_B^{}\}(mZ_A)\delta _{AB},`$
so we obtain a Clifford algebra. It is clear that if $`m=|Z_A|`$ for some $`A`$ (in general one has $`m|Z_A|`$), there is one pair of oscillators in the Clifford algebra that decouple. Representations with such value of the central charge are representations in which certain supersymmetry charges become trivial . They are called BPS states. If all eigenvalues are equal, half of the oscillators decouple and we say that we have a representation with $`\frac{1}{2}`$BPS. For less restrictive conditions one obtains representations with $`\frac{1}{4}`$BPS or $`\frac{1}{8}`$BPS.
The corresponding supergravity theories can be obtained by compactifying $`D=11`$, $`N=1`$ supergravity on a torus $`T^{d+1}`$ or compactifying $`D=10`$, type IIA and IIB supergravity on a torus $`T^d`$. The scalar fields of the theory parametrize a coset manifold of the form $`G/H`$, where $`G`$ is a non compact group whose maximally compact subgroup is $`H`$. For any $`d`$, $`G=E_{d+1(d+1)}`$ (the factor U(1) for $`D=4`$ must be supressed to have $`H`$ as a subgroup of $`G`$). $`G`$ is called the $`U`$-duality group. $`G`$ acts on the non linear manifold of the scalars, and the vector fields also transform in a representation of $`G`$. This representation naturally extends (5) and is given in the following table for each case ,
Central charge representation of the $`U`$-duality group
$`D=9`$ 2+1 of E<sub>2</sub>=SL(2)$`\times `$ O(1,1) (6)
$`D=8`$ (3,2) of E<sub>3</sub>=Sl(3)$`\times `$Sl(2)
$`D=7`$ 10 of E<sub>4</sub>=Sl(5)
$`D=6`$ 16 of E<sub>5</sub>=O(5,5)
$`D=5`$ 27 of E<sub>6(6)</sub>
$`D=6`$ 56 of E<sub>7(7)</sub>
Let $`\mathrm{\Phi }^i`$ be the scalar fields which parametrize $`G/H`$ and consider a local section on the principal bundle $`G`$ over the base $`G/H`$, $`L(\varphi )G`$ (coset representative). The action of $`G`$ can be expressed as
$$L(\varphi _g)=gL(\varphi )h(\varphi _g)$$
(7)
where $`h(\varphi _g)H`$ is a transformation in the fiber over $`\varphi _g`$. Choosing a representation of $`G`$ the coset representative becomes a matrix $`L_a^\mathrm{\Lambda }(\mathrm{\Phi })`$. The indices $`a,\mathrm{\Lambda }`$ run over the same represenation of $`G`$ (and of $`H`$), the different names used to remind the transformation rule (7).
The charges of a black hole solution can be computed by integrating the Hodge dual of the field strengths of the vector fields present in the theory on a $`D2`$ spacial surface enclosing the source. We denote these charges by $`q^\mathrm{\Lambda }`$. Then, the central charges of the black hole solutions are given by
$$Z_a(q,\mathrm{\Phi })=q_\mathrm{\Lambda }^TL_a^\mathrm{\Lambda }(\mathrm{\Phi })$$
(8)
Notice that the vectors fields are in a certain representation of $`G`$, which necesarily coincides with the representation in (1).
When the central charges correspond to BPS states, some supersymmetry generators are null when acting on the black hole solutions, so the solutions are ”supersymmetric”. The condition to be a BPS black hole is a condition on $`Z`$, that is, on $`q`$ and $`\mathrm{\Phi }`$. Due to the form of (8), it is clear that any condition $`E_\alpha (Z)=0`$ (where $`\alpha `$ runs over some representation $`T`$ of $`G`$) that is covariant under $`G`$, that is
$$E_\alpha (Zg)=E_\beta (Z)T(g)_\alpha ^\beta ,$$
(9)
will become simply a condition on $`q`$, $`E_\alpha (q)=0`$. It was shown in that this condition is actually moduli independent or $`U`$-duality invariant in all theories in dimensions $`D=4,\mathrm{},9`$ with $`n=32,16`$ supersymmetries. Here we will review these results, in particular $`D=4`$, $`N=8`$ and $`D=5`$, $`N=4`$.
## 2 Maximal supergravity in dimensions $`D=4,\mathrm{},9`$
We study the diagonalization of the matrix of central charges in the different cases (Table 1).
A matrix is quaternionic (or symplectic) if
$$Z^{}=\mathrm{\Omega }Z\mathrm{\Omega }.$$
A matrix satisfying this conditions can be understood as a matrix whose entries are quaternions, and these are in the representation
$$\text{Id}_{2\times 2},i\sigma _1,i\sigma _2,i\sigma _3.$$
One has the following results:
1. Any matrix can be brought to a diagonal form by making a transformation
$$Z_D=U_1ZU_2^{}$$
where $`U_i`$ are orthogonal matrices if $`Z`$ is real, unitary if $`Z`$ is complex and unitary symplectic if $`Z`$ is quaternionic (In the representation of $`n\times n`$ quaternionic matrices by complex, $`2n\times 2n`$ matrices, it is $`\mathrm{\Omega }Z`$ which is brought to diagonal form).
2. In the case that the original matrices are symmetric, hermitian or symplectic hermitian respectively, $`U_1=U_2`$ and the eigenvalues are real.
3. An antisymmetric matrix can be brought to skew diagonal form by a transformation
$$Z_{SD}=UZU^T$$
where $`U`$ belongs to the appropriate group as before and the eigenvalues are real.
For $`D=9`$ we can trivially diagonalize $`Z`$ with an $`R`$-symmetry transformation by using the result 2. For $`D=8`$, an $`R`$-transformation $`U`$ brings the matrix $`Z`$ to $`UZU^T`$. It is easy to see that to diagonalize the $`2\times 2`$ matrix it is enough to use the result 1 with $`U_2^{}=U_1^T`$, being the eigenvalues real. For $`D=7`$, $`Z`$ is four dimensional and since the spinors are pseudoreal it is quaternionic. Since it is also symmetric, we can take advantage of the isomorphism Sp(4)$``$O(5) and decompose $`Z`$ as
$$Z_{ij}=Z_{IJ}(\gamma ^{IJ})_{ab},IJ=1,\mathrm{}5$$
where $`Z^{IJ}`$ is real and antisymmetric and
$$\gamma ^{IJ}=\frac{1}{2}[\gamma ^I,\gamma ^J],$$
and $`\gamma ^I`$ are the gamma matrices of O(5). We can skew diagonalize $`Z^{IJ}`$ using the result 3, with only two independent eigenvalues. For $`D=6`$ one can use the result 1 and diagonalize the matrix with an element $`(U_1,U_2)`$USp(4)$`\times `$USp(4). The quaternionic property is preserved by a transformation $`Z_D=U_1ZU_2^{}`$, but since the eigenvalues are not real there are four independent real quantities. Finally, for $`D=5,4`$ one can use the result 3.
For $`D=7,8,9`$ there are only two independent eigenvalues. In the generic situation, when the two eigenvalues are different, we have that the bound $`m=|Z_A|`$ can be reached only by the highest eigenvalue. Then, $`\frac{1}{4}`$ of the oscillators of the supersymmetry alegra are null on the solution. We say that we have a $`\frac{1}{4}`$BPS state. If all the eigenvalues are equal, the half of the oscillators decouple and we have a $`\frac{1}{2}`$BPS state. For $`D=4,5,6`$ we have four eignevalues and the generic case, when all of them are different preserves $`\frac{1}{8}`$ of the supersymmetry. If the eigenvalues are equal by pairs we have $`\frac{1}{4}`$BPS states and if they are all equal we have $`\frac{1}{2}`$BPS states. Exotic cases like having only one pair of equal eigenvalues and the rest different or having three equal eigenvalues and a different fourth one are forbidden on physical grounds as we will see in the next example.
### 2.1 BPS states in $`D=4`$, $`N=8`$ Supergravity.
As we have seen in this case $`Z`$ is a $`8\times 8`$ complex antisymmetric matrix, so it can be skew diagonalized
$$Z_{SD}=UZU^T$$
with 4 independent eigenvalues. Since we can use only transformations $`U`$ SU(8), the eigenvalues are not real. Instead there is an overall phase that cannot be removed . In fact, this phase is an extra parameter of the solution which usually is set to zero and all the eigenvalues real. Their absolute value, can be computed as the square root of the eigenvalues of the matrix $`ZZ^{}`$.
#### $`\frac{1}{2}`$BPS
Since all the eigenvalues are real, the matrix $`ZZ^{}`$ is a multiple of the identity,
$$ZZ^{}=\text{Tr}(ZZ^{})\frac{1}{8}\text{Id}.$$
(10)
This is an SU(8) covariant constraint. The idea is to find an E<sub>7,7</sub> covariant constraint that necessarily implies (10). Then we will be in the situation (9), where the $`\frac{1}{2}`$BPS condition is moduli independent. Nevertheless, the eigenvalues will depend on the moduli.
We consider the irreducible representation of E<sub>7,7</sub> 56 (the number indicates the dimensionality of the representation). Under SU(8) it decomposes as $`\mathrm{𝟓𝟔}=\mathrm{𝟐𝟖}+\overline{\mathrm{𝟐𝟖}}`$. $`\mathrm{𝟐𝟖}`$ is the twofold antisymmetric representation of SU(8) in which $`Z`$ sits. We will write a vector of 56 symbolically as $`\stackrel{~}{Z}=(Z,\overline{Z})`$. We consider now the quartic invariant of this representation of E<sub>7,7</sub> ,
$$I=4\text{Tr}(Z\overline{Z})^2(\text{Tr}Z\overline{Z})^2+2^4(\mathrm{𝑃𝑓}Z+\mathrm{𝑃𝑓}\overline{Z}),$$
(11)
where
$$\mathrm{𝑃𝑓}Z=\frac{1}{2^44!}ϵ^{ABCDRPGH}Z_{AB}Z_{CD}Z_{RP}Z_{GH}$$
is the Pfaffian of $`Z`$. Let us take the second derivative of this invariant
$$\frac{^2I}{\stackrel{~}{Z}\stackrel{~}{Z}}.$$
It is a quadratic polynomial which is a symmetric tensor, so it sits in the $`(\mathrm{𝟓𝟔}\times \mathrm{𝟓𝟔})_S=\mathrm{𝟏𝟓𝟗𝟔}`$ representation of E<sub>7</sub> which is reducible and decomposes as 1463 \+ 133. 133 is the adjoint of E<sub>7</sub>, so we can take the projection
$$\frac{^2I}{\stackrel{~}{Z}\stackrel{~}{Z}}|_{Adj_{E_7}}.$$
(12)
Since 133 decomposes as 63+70 under SU(8) (63 is the adjoint representation of SU(8)), the expression (12) splits into two SU(8) covariant polynomials
$`V_A^C={\displaystyle \frac{^2I}{Z_{AB}\overline{}Z^{CB}}}|_{Adj_{SU(8)}}`$
$`(Z_{AB}\overline{Z}^{CB}{\displaystyle \frac{1}{8}}\delta _A^CZ_{PQ}\overline{Z}^{PQ}),`$ (13)
$`V_{[ABCD]}^+={\displaystyle \frac{^2I}{Z_{AB}\overline{}Z^{CB}}}|_{\mathrm{𝟕𝟎}}{\displaystyle \frac{^2I}{Z_{[AB}Z_{CD]}}}`$
$`{\displaystyle \frac{1}{4!}}ϵ^{ABCDPQRS}{\displaystyle \frac{^2I}{\overline{Z}^{[AB}\overline{Z}^{CD]}}}.`$ (14)
The 63 piece, (13) is just (10), which is the constraint we want. But the 70 part $`V_{[ABCD]}^+`$ is not, in principle, zero. Here is where the $`G/H`$ structure of the theory enters. Consider $`L(\mathrm{\Phi }):G/HG`$, the local section or coset representative of $`G/H`$. Using this map we can make a pull-back of the Maurer-Cartan left invariant forms on $`G`$ to an open set of $`G/H`$,
$$\mathrm{\Omega }(\varphi )=L^1(\varphi )dL(\varphi )=\omega _iT_i+P_\alpha T_\alpha .$$
(15)
$`T_i`$ are the generators in the Lie algebra of $`H`$ and $`T_\alpha 𝒦`$, where $`𝒢=+𝒦`$ is a Cartan decomposition of $`𝒢`$, the Lie algebra of $`G`$. $`\omega _i`$ are the components of the spin connection of the bundle $`GG/H`$ and $`P_\alpha `$ is the vielbein of the invariant metric on $`G/H`$.
In the case of E<sub>7,7</sub>/SU(8), the equation (15) takes the form
$$_{SU(8)}Z_{AB}=\frac{1}{2}ϵ_{ABCD}\overline{Z}^{CD}.$$
Taking the covariant derivative $`_{SU(8)}`$ of the equation $`V_A^C=0`$ and using the Maurer-Cartan equations above, one obtains
$$V_+^{ABCD}=0,$$
as we wanted to show.
Notice that the reality of the eigenvalues follows in fact from (14), so it follows from the E<sub>7,7</sub> invariance.
#### $`\frac{1}{4}`$BPS.
In this case we have that the eignevalues are equal by pairs. It can be shown that the E<sub>7,7</sub> covariant condition for this to happen is
$$\frac{I}{Z_{AB}}=0(\frac{I}{\overline{Z}^{AB}}=0),$$
where $`I`$ is again the quartic invariant. The reality also follows from the E<sub>7,7</sub> invariance.
$`\frac{1}{8}`$BPS is the generic case when all the eigenvalues are different. One can always make the highest eigenvalue (which is the mass of the BPS state) real, so no equation is needed to assure reality. Having three or two eigenvalues equal and the rest different makes the quartic invariant negative . Since this invariant has an interpretation as the entropy of the black hole squared, $`IS^2`$ it cannot be negative, so these cases are excluded.
For all the other theories in Table 1 a similar analysis can be carried out
## 3 Supergravities with 16 supersymmetries. BPS states of $`D=5`$, $`N=4`$ supergravity
Theories with sixteen supersymmetries are obtained as compactifications of the heterotic string on tori T<sup>d</sup> ($`1d6`$). They can also be obtained by compactifying $`D=11`$ supergravity or $`D=10`$ Type IIA and IIB supergravity on manifolds preserving less supersymmetries, as K<sub>3</sub>. In these theories there are matter fields and the $`U`$-duality group $`G`$ depends on the matter content as well as on the dimension of space-time. The maximal compact subgroup is a direct product $`H_R\times H_M`$ where $`H_R`$ is the R-symmetry of the corresponding supersymmetry algebra and $`H_M`$ is the group acting on the matter multiplets. If $`n`$ is the number of matter multiplets this group is $`H_M=\text{O}(n)`$. $`G`$ is of the form O(10-$`D,n`$)$`\times `$O(1,1) for $`5D9`$ while for $`D=4`$ it is SL(2)$`\times `$O(6,$`n`$). The R-symmetry groups are O(10-$`D`$) for $`5D9`$ and O(6)$`\times `$O(2)$``$SU(4)$`\times `$U(1) for $`D`$=4.
The $`G`$ and $`H_R`$ representations of the central charges are given in the following tables.
Central charge representation of $`H_R`$.
$`D=9`$ $`\mathrm{𝟏}`$ $`\text{O}(1)=\text{Id}`$
$`D=8`$ $`\mathrm{𝟏}^𝐜\text{complex}`$ $`\text{U}(1)\text{O}(2)`$
$`D=7`$ $`\mathrm{𝟑}\text{real}`$ $`\text{SU}(2)\text{USp}(2)`$
$`D=6`$ $`\mathrm{𝟒}\text{real}`$ $`\text{O}(4)\text{USp}(2)\times \text{USp}(2)`$
$`D=5`$ $`\mathrm{𝟏}+\mathrm{𝟓}\text{real}`$ $`\text{O}(5)\text{USp}(4)`$
$`D=4`$ $`\mathrm{𝟔}^𝐜\text{complex}`$ $`\text{O}(6)\times \text{O}(2)`$
$`\text{SU}(4)\times \text{U}(1)`$
From the above table, and according to our previous analysis, it follows that the matrix of central charges, $`Z`$, has only one independent real eigenvalue for $`D=6,\mathrm{}9`$ and two independent eigenvalues for $`d=5,6`$. Therefore, for $`D=6,\mathrm{}9`$ only 1/2 BPS states can occur while for $`D=4,5`$ both, 1/2 and 1/4 BPS states can occur.
Central charge representation of $`G`$
$`D=6,\mathrm{}9`$ $`𝐝+𝐧\text{real vector}`$
$`\text{O}(d,n)\times \text{O}(1,1)`$
$`D=5`$ $`\mathrm{𝟏}+(\mathrm{𝟓}+𝐧)\text{(singlet+vector)}r`$
$`\text{O}(5,n)\times \text{O}(1,1)`$
$`D=4`$ $`(\mathrm{𝟐},\mathrm{𝟔}+𝐧)`$
$`\text{Sl}(2)\times \text{SO}(6,n)`$
We will briefly outline here the case of $`D=5`$. It corresponds to heterotic string on $`T^5`$ or $`D=11`$ supergravity on $`K^3\times T^2`$. The number of matter multiplets is $`n=21`$, although our analysis is independent of such number.
The central charge $`\widehat{Z}`$ is an antisymmetric quaternionic matrix. This implies that the matrix $`Z=\widehat{Z}\mathrm{\Omega }`$ is hermitian and quatenionic. The $`4\times 4`$ matrix depends on 6 real parameters. We want to exploit the isomorphism O(5)$``$USp(4), so we decompose $`Z`$ as
$$Z=Z^a\gamma _a+Z^0\text{Id},$$
where $`\gamma _a`$, $`a=1,\mathrm{}5`$ are the O(5) $`\gamma `$-matrices and $`Z^a,Z^0`$ are real numbers. $`Z^0`$ is a singlet under O(5) and $`Z^a`$ is a vector. It follows that
$`\text{Tr}Z=4Z^0`$
$`(\text{detZ)}^{1/2}={\displaystyle \frac{1}{8}}(\text{Tr}Z)^2{\displaystyle \frac{1}{4}}\text{Tr}Z^2=`$
$`Z_{}^{0}{}_{}{}^{2}\stackrel{}{Z}^2`$ (16)
The characteristic equation of $`Z`$ (or better, its square root) is
$`\sqrt{\text{det}Z\lambda \text{Id}}=`$
$`\lambda ^2{\displaystyle \frac{1}{2}}\text{Tr}Z\lambda +(\text{det}Z)^{1/2}=0,`$
implying that $`Z`$ has two coinciding eigenvalues (in absolute value) either if
$$\text{Tr}Z=0\text{or}\frac{1}{4}(\text{Tr}Z)^2=4(\text{det}Z)^{1/2}$$
Using (16), the above equation directly implies
$$Z_0Z_a=0.$$
(17)
The eigenvalues are given by
$$\lambda _{1,2}=\frac{1}{2}\left(\frac{1}{2}\text{Tr}Z\pm \sqrt{\text{Tr}Z^2\frac{1}{4}(\text{Tr}Z)^2}\right),$$
being the plus sign the mass squared of the BPS state.
The central charge representation of $`G`$ (and $`H`$), singlet + vector, is reducible and the coset representative $`L_a^\mathrm{\Lambda }`$ splits into blocks
$$\left(\begin{array}{cc}e^{2\sigma }& \\ & e^\sigma \left(\begin{array}{c}M\end{array}\right)\end{array}\right),$$
where $`\sigma `$ parametrizes O(1,1) and $`M`$ is in the fundamental representation of O(5,n). Since $`Z_a=q_\mathrm{\Lambda }L_a^\mathrm{\Lambda }`$, then $`Z_0=e^{2\sigma }q_0`$. The condition $`Z_0=0`$ implies $`q_0=0`$, which is a singlet of O(5+n), and then, $`U`$-duality invariant.
If $`Z_I,I=1,\mathrm{}n`$ are the matter charges associated to the $`n`$ matter multiplets, we have that, because of (15),
$$_{O(5)}Z_a=\text{Tr}(\gamma _aP_I)Z^I+Z_ad\sigma ,$$
therefore $`Z_a=0`$ implies $`Z_I=0`$. This is also an O(5,n) invariant statement since, it comes by differentiating the quadratic invariant polynomial
$$I=\underset{a=1}{\overset{5}{}}Z_aZ^a\underset{I=1}{\overset{M}{}}Z_IZ^I.$$
Therefore, $`Z_a=Z_I=0`$ implies $`q^\mathrm{\Lambda }=0`$ where $`q^\mathrm{\Lambda }`$, $`\mathrm{\Lambda }=1,\mathrm{}5+n`$, is a fixed charge vector of O(5,M), as found in .
Acknowledgments
I want to thank to my collaborators R. D’Auria and S. Ferrara for reading the manuscript an suggesting many improvements. |
warning/0002/astro-ph0002436.html | ar5iv | text | # Constraining the Radio-Submillimetre Redshift Indicator using data from the SCUBA Local Universe Galaxy Survey.
## 1 Introduction
Ever since the discovery of Ultra-luminous Infrared Galaxies (ULIRGs) (Joseph et al. 1985; Sanders et al. 1988) which emit up to 99 per cent of their bolometric luminosity in the far-IR, there have been suspicions that our optical/UV view of the early universe might be biased by large amounts of dust obscuration. These were supported by the discovery of the Cosmic infrared Background (CIB) (Puget et al. 1996; Fixsen et al. 1998), which contains up to twice the luminosity of the integrated optical/UV background. It thus seems likely that much of the star formation activity in the early universe takes place in obscured environments. Many deep surveys have been conducted with the SCUBA bolometer array on the James Clerk Maxwell Telescope (JCMT)<sup>1</sup><sup>1</sup>1The JCMT is operated by the Joint Astronomy Center on behalf of the UK Particle Physics and Astronomy Research Council, the Netherlands Organization for scientific Research and the Canadian National Research Council. in order to uncover the sources responsible for this background. Significant numbers of sources, far in excess of no evolution models have been found at both submm (Smail, Ivison & Blain 1997; Barger et al. 1998; Hughes et al. 1998; Eales et al. 1999) and far-IR (Puget et al. 1999; Kawara et al. 1998) wavelengths.
The next step after finding these sources is identifying them and determining their redshifts, thus allowing the true star formation history of the universe to be traced. This has proved to be a difficult task, due to the large SCUBA beam size at 850$`\mu `$m (15 arcsec) and the several arcsec positional uncertainty inherent in all of the SCUBA maps (Ivison et al. 1998,2000; Downes et al. 2000; Barger, Cowie & Richards 2000; Eales et al. 2000). The SCUBA sources can sometimes be identified with several plausible optical candidates, or in other cases they remain optical blank fields, even to the magnitudes probed by HST and Keck. A possible solution to the problem of poor positional accuracy lies in the tight correlation between thermal FIR emission and synchrotron radio emission seen in the local universe (Helou et al. 1985; Condon 1992). The mechanism for this relationship is believed to be the massive stars from recent star-formation, which both heat the dust to produce the FIR flux and, when they explode as supernovae, provide relativistic electrons which constitute the synchrotron radio emission. This means that radio observations with the VLA are also sensitive probes of star formation, albeit not as sensitive as the submm at higher redshifts given the unfavourable K-correction at radio wavelengths. For example, Barger, Cowie & Richards (2000) have recently targeted 1.4 GHz sources in the deep VLA images (Richards 1999) of the Hubble Flanking Fields (HFF) with SCUBA and found that the majority of bright ($`>6`$ mJy) submm sources have radio counterparts. By observing the submm sources from the deep surveys at radio frequencies (particularly 1.4 GHz), more accurate positions can be obtained, allowing more secure identifications of the submm sources with optical or near IR counterparts. However, even with the accurate VLA positions, many SCUBA sources still have no detectable optical/IR counterparts, or they are too faint for spectroscopy. This creates a great problem in determining the redshifts for the sources, necessary to trace their evolution.
One possible route for obtaining redshift estimates for these objects has been suggested by Carilli & Yun (1999). This is based on the FIR-radio correlation plus the break in the spectral slope at around 3 mm, where the thermal dust emission takes over from the declining synchrotron tail. This puts the radio on one side of the spectral break and the FIR (or submm in this case) on the other, creating a redshift sensitive ratio. Carilli & Yun (1999) used the submm-radio spectral index, defined as:
$$\alpha _{1.4}^{850}=0.42\times \mathrm{log}\left(\frac{S_{850}}{S_{1.4}}\right)$$
where $`S_{850}`$ and $`S_{1.4}`$ are the fluxes at 850$`\mu `$m and 1.4 GHz.
The exact dependence of $`\alpha _{1.4}^{850}`$ on redshift is sensitive to a few parameters:
1. The slope of the thermal Rayleigh-Jeans tail in the FIR and submm. This depends on frequency as $`\nu ^{2+\beta }`$ where the dust emissivity index ($`\beta `$) is thought to lie between 1 and 2 (Hildebrand 1983). It is the steepness of this slope which potentially makes the indicator sensitive to redshift
2. The slope of the radio synchrotron emission. For star forming regions this is assumed to be $`0.70.8`$ (Condon 1992). At higher frequencies this will flatten due to radio free-free emission but this is not simple to model. Free-free absorption at lower frequencies can also lead to a flattening of the radio spectra, the likely effect of this on $`\alpha _{1.4}^{850}`$ is discussed in Carilli & Yun (2000) and they conclude that it is unlikely to be a dominant cause of scatter in the relationship. In any case the slope at radio frequencies is much shallower than in the submm and any uncertainties in it have less effect.
3. The temperature of the dust. This determines the redshift at which the thermal spectrum turns over. When this happens the ratio only depends weakly on redshift, and so becomes useless as a redshift indicator. Accurate dust temperatures are notoriously difficult to estimate and require measurements at many different wavelengths throughout the submm and FIR.
4. Contribution by AGN. Galaxies containing a radio-loud AGN will have higher radio emission than a purely star-forming galaxy and so lower values of $`\alpha _{1.4}^{850}`$. This could lead to ambiguity between high redshift objects with AGN and lower redshift starbursts.
5. The FIR-radio correlation and its variability with redshift. In using this technique to estimate redshifts, an inherent assumption is made that the FIR-radio correlation does not vary with redshift. Possible causes of such a change could be a variation in magnetic field strengths in galaxies in the past, or changes in the dust mass opacity coefficient which relates the FIR/submm emission to the mass of dust. The latter would depend on the dust composition and grain sizes so there is scope for it to have been different at the epoch of galaxy formation, when metallicities would have been lower.
6. Linked to (v) is the possibility that the FIR-radio, or more importantly, the submm-radio relationship itself depends on some other galaxy property such as dust temperature, luminosity etc.
The possible impact of points i–vi above will be discussed later.
Several authors have already used this technique to determine redshift estimates or limits for sources discovered in deep submm surveys with SCUBA (Hughes et al. 1998; Lilly et al. 1999; Smail et al. 1999; Barger et al. 2000; Eales et al. 2000). This paper uses data from the SCUBA Local Universe Galaxy Survey (SLUGS) to examine the usefulness of the submm-radio ratio as a redshift indicator, to try to normalise it with respect to the local universe and to estimate the uncertainties in it.
We will use $`H_0=75`$ km s<sup>-1</sup> Mpc<sup>-1</sup> and $`q_0=0.5`$ throughout.
## 2 Constraining the Redshift Indicator
The low redshift data is taken from the SLUGS survey (Dunne et al. 2000) which is a complete sample of 104 galaxies selected from the IRAS Bright Galaxy Sample. The sample includes mainly infra-red bright objects ($`L_{\mathrm{fir}}=7\times 10^92\times 10^{12}\mathrm{L}_{}`$), some ULIRGS as well as many ‘more normal’ galaxies. The objects were mapped at 850$`\mu `$m with SCUBA and so should provide reliable submillimetre fluxes. We fitted a single-temperature dust spectral energy distribution (SED) to the measured fluxes for each galaxy, allowing the dust emissivity index ($`\beta `$) as well as the temperature to vary to find the best fit. While the assumption of a single temperature is probably wrong, and indeed the low values of $`\beta `$ we derive suggest that there is a cold dust component, the SEDs are still a good empirical fit to the data, which is all that we require for the present analysis.
The original Carilli-Yun models for the behaviour of the radio-submm spectral index with redshift were simple two power-law models, along with some empirical SED data for 2 local galaxies. They have since modified their original estimator by using data from nearby galaxies (17 objects) to try to provide a normalisation and to estimate the scatter in the submm-radio ratio with redshift (Carilli & Yun 2000). We have independently undertaken a similar analysis but using the full SLUGS data set (104 objects), thereby providing a more statistically valid measure of the relationship and its scatter.
We have taken the real spectral energy distributions of the 104 SLUGS galaxies as determined from the IRAS 60 and 100$`\mu `$m fluxes and the SCUBA 850$`\mu `$m data, and the 1.4 GHz radio flux (Condon et al. 1990). The measured radio spectral indices were used where available (for about half of the sample) and the median $`\alpha _{radio}`$ was found to be $`0.7`$. This was used in the SEDs of those galaxies with no measured $`\alpha _{radio}`$. If the change in each of these SEDs with redshift is plotted, we get 104 separate curves from which the median can be selected. The uncertainty in this method (at any given redshift) can be estimated from the spread in the 104 curves. In Figure 1, the thick solid line shows the median curve with the upper and lower solid curves being those ranked 16 per cent from the top and bottom respectively, to give the $`\pm 1\sigma `$ errors (defined as the range of $`\alpha _{1.4}^{850}`$ in which 68 per cent of the lines lie). So at any given redshift, the horizontal distance between the thin lines provides an estimate of the $`\pm 1\sigma `$ uncertainty on the redshift provided by the median curve. The rms scatter in $`\alpha _{1.4}^{850}`$ is almost constant with redshift, decreasing very slightly from 0.097 at $`z=0`$ to 0.080 at $`z=6`$. A rough guide to the redshift uncertainties at various redshifts are given in Table 1. The median curve can be approximated by the expression
$$z=0.5516.652\alpha +25.57\alpha ^230.56\alpha ^3+13.75\alpha ^4$$
Our estimator starts to loose its effectiveness at redshifts greater than about 4–5, where it starts to turn over due to the peak in the thermal dust spectrum. This can be seen from the increase in $`\sigma _z`$ with redshift (Table 1). Also, as mentioned in Carilli & Yun (1999) at $`z>6`$ the radio emission begins to become quenched by inverse Compton losses off the microwave background. The dashed line in Fig. 1 shows the new Carilli-Yun estimator (with $`\pm 1\sigma `$ uncertainties) using the data from 17 local galaxies (Carilli & Yun 2000). It can be seen that their mean estimator lies outside the 1$`\sigma `$ uncertainties from our analysis, however the upper $`1\sigma `$ curve from theirs is consistent with our lower $`1\sigma `$ limit, with agreement being better at lower redshifts ($`<2`$). The differences in redshift estimates from the two indicators are listed in Table 2 for various values of $`\alpha _{1.4}^{850}`$.
There are two main differences between the present results and those in Carilli & Yun (2000). The first is in the zero redshift normalisation in the two samples. The mean value of $`\alpha _{1.4}^{850}`$ at $`z=0`$ in Carilli & Yun (2000) is $`0`$ with an rms scatter of 0.14, while for our sample the mean $`\alpha _{1.4}^{850}=0.18`$ at $`z=0`$, with a lower rms scatter of 0.097. The shape of the Carilli-Yun estimator is also different due to the SED shapes of the galaxies in their sample (which in general had higher values of $`\beta `$ than ours). These differences are partly due to our much larger sample size, but there are also discrepancies in the submillimetre fluxes of the individual galaxies which are present in both the samples used by ourselves and Carilli & Yun (2000). Most of their submm data comes from a sample of 19 bright IRAS galaxies (Lisenfeld, Isaak & Hills 2000). Seven of the galaxies in this sample were also observed by us, and in several cases the fluxes we derived were as much as a factor of two higher. This was attributed to Lisenfeld et al. missing flux from extended objects (Lisenfeld, private communication), and has been corrected in the revised version of the Lisenfeld et al. paper. The consequences of using under-estimated submm fluxes would be to i) lower the mean $`\alpha _{1.4}^{850}`$ and ii) increase the value of $`\beta `$ determined in single temperature fits to the 60, 100 and 850$`\mu `$m fluxes, which explains the different shape of the Carilli-Yun estimator. We also note a possible error in the 1.4 GHz fluxes used for Zw049.057 and NGC 5936 by Carilli & Yun (2000). An erratum has since been published to accompany Carilli & Yun (2000) and this is discussed further in Section 5.
If we now place the handful of SCUBA sources from the deep surveys which have both spectroscopic redshifts and radio fluxes (Smail et al. 1999 and references therein; Eales et al. 1999; Lilly et al. 1999; Barger et al. 2000; Eales et al. 2000; Ivison et al. 2000) on the redshift estimator we have created from our observed SEDs (Fig. 2), we see that the agreement is satisfactory given the uncertainties. The object at $`z=2.8`$ is SMM02399-0136 from Ivison et al. (1998) and is known to harbour an AGN, which explains why it lies below the predicted lines. Given that sources may have AGN activity but cannot be confirmed as such, this redshift estimator can give only reliable lower limits. Additionally, if taken with another redshift estimation, such as photometric redshift, it could help to identify the presence of AGN (a point first noted by Carilli & Yun 1999).
## 3 The FIR and Submm Vs. Radio Properties
The SLUGS sample provides an ideal test for the variation of the FIR-radio and submm-radio relationships with galaxy properties, such as dust temperatures, luminosities etc. which may affect the reliability of $`\alpha _{1.4}^{850}`$ as a redshift indicator, and which is also interesting in its own right. Figure 3 shows the FIR-radio correlation for the SLUGS galaxies along with the 850$`\mu `$m–1.4 GHz relationship. Both are highly significant (Table 3 lists all the correlation coefficients in this discussion) but the scatter in the FIR-radio correlation is smaller (0.06 (15%)) than that which exists for the submm-radio (0.097 (25%)). This suggests that it is actually the FIR which is physically related to the radio rather than the submm. There is also a difference in slope between the two, which we will return to later. The FIR-radio correlation in the local universe does not depend on galaxy luminosity (except at very low $`L_{\mathrm{fir}}`$) or dust temperature, something which is confirmed by this sample. However if the submm-radio relationship is tested in this way by plotting $`\alpha _{1.4}^{850}`$ against $`L_{\mathrm{fir}}`$, $`L_{1.4}`$ and $`L_{850}`$ (Figure 4) then it is clear that there is a strong dependence of $`\alpha _{1.4}^{850}`$ on both $`L_{\mathrm{fir}}`$ and $`L_{1.4}`$ but not on $`L_{850}`$. The link with $`L_{\mathrm{fir}}`$ was noted by Carilli & Yun (2000) but the trend is even more evident with radio luminosity. The most likely explanation for these correlations is that both $`L_{\mathrm{fir}}`$ and $`L_{1.4}`$ are very sensitive to recent star formation in a galaxy but $`L_{850}`$ is not as sensitive, as a significant fraction of the 850$`\mu `$m flux may be produced by colder dust heated by the general interstellar radiation field (ISRF). Therefore $`L_{850}`$ will not change by as great a fraction as $`L_{\mathrm{fir}}`$ with increasing star formation. This explains the larger scatter and steeper slope in the 850$`\mu `$m–1.4 GHz correlation.
One of the two crucial relations for determining whether $`\alpha _{1.4}^{850}`$ is a biased redshift indicator is that between $`\alpha _{1.4}^{850}`$ and $`L_{850}`$. Since SLUGS was selected at 60$`\mu `$m, it is quite possible that the 850$`\mu `$m luminosities of SCUBA-selected galaxies are distinctly different from those in SLUGS. If there were also a significant relationship between $`\alpha _{1.4}^{850}`$ and $`L_{850}`$, the calibration of the Carilli-Yun method using SLUGS galaxies would lead to biased redshift estimates. Fig. 4(c) shows, however, that there is no significant correlation. To examine this, the dot-dash line on Fig. 4(c) shows the lower end of the rest-frame 850$`\mu `$m luminosity range for a galaxy with $`S_{850}=4`$mJy at redshifts of $`>2`$ ($`\beta =2`$ and $`T_\mathrm{d}=50`$ K, lower $`\beta `$ and $`T_\mathrm{d}`$ give slightly higher rest-frame $`L_{850}`$ for the same flux and redshift). If we look at the values of $`\alpha _{1.4}^{850}`$ for the SLUGS galaxies, which correspond to the luminosities of the high redshift objects (i.e. the crosses to the right of the dot-dash line in Fig. 4(c)), they do not display a lower mean value than that for the whole sample. Given how little we know about the properties of the high-z SCUBA galaxies, and whether they are related to galaxies of the same luminosity at low redshift (density evolution) or to galaxies of lower luminosity (luminosity evolution), we do not feel that there is a need for a luminosity dependent indicator. This is further supported by Fig. 2, showing a relatively good agreement of the indicator with the high-$`z`$ data, and in particular, there is no systematic bias from the line in one direction, which may be indicative of changes in $`\alpha _{1.4}^{850}`$ with redshift/luminosity.
The second important relationship is that between $`\alpha _{1.4}^{850}`$ and dust temperature, since an 850$`\mu `$m selected sample may well have a different mean dust temperature than one selected at 60$`\mu `$m. It has also been suggested by Blain (1999) that there should be a dependence of $`\alpha _{1.4}^{850}`$ on dust temperature, and that this produces a degeneracy in the redshift indicator whereby a hot galaxy at high redshift may be indistinguishable from a colder galaxy at lower redshift. We have tested for a temperature dependence by plotting $`\alpha _{1.4}^{850}`$ against dust temperature, and we find no significant correlation (Figure 5(a)). The slight correlation seen between $`\alpha _{1.4}^{850}`$ and dust emissivity index $`\beta `$ (Figure 5(b)) is probably linked to the relationship of $`\alpha _{1.4}^{850}`$ with $`L_{1.4}`$ and $`L_{\mathrm{fir}}`$. Despite the significant correlation given by the statistic to this relationship, it appears to only hold for the lower values of $`\beta `$ (0.9–1.4). This could be because a low $`\beta `$ (when produced by fits to 60,100 and 850$`\mu `$m points only) is another possible indicator of the fraction of 850$`\mu `$m emission produced by the ISRF, rather than directly by star forming regions. Evidence for this comes from observations of a sample of optically selected galaxies presently under way with SCUBA and also from studies of NGC 891 and the Milky Way (Alton et al. 1998; Masi et al. 1995), where the SEDs using only 60, 100 and 850$`\mu `$m fluxes give a very low $`\beta `$ of $`0.7`$. These galaxies are known to have large fractions of cold dust at $`T<20`$ K, however this does not lead to a dependence of $`\alpha _{1.4}^{850}`$ on the fitted dust temperature because that is determined by the FIR fluxes which are dominated by the warmer dust.
## 4 Discussion
The submm-radio spectral index appears to provide a moderately satisfactory redshift indicator out to redshifts of $`45`$ where it ceases to be as sensitive to redshift. The uncertainties in the estimated redshifts are obtainable using the spread of the local SEDs, and range from $`<\pm 0.3`$ at $`z<1`$ to $`>\pm 1`$ at $`z>4`$. Sources which harbour radio loud AGN will not follow this pattern as they will lie below the indicators, leading to AGN sources which are not recognised as such being given misleadingly low redshifts.
Let us now return to the assumptions and limitations discussed in section 1, and apply them to our data sets.
1. The submillimetre emissivity index has a mean value of 1.3 for the SLUGS sample, as determined by single temperature fits to the 60, 100 and 850$`\mu `$m fluxes. As discussed at length in Dunne et al. (2000), if a colder dust component is present this may not represent the true value of $`\beta `$, the true value being higher than that determined using a single component model. If the true value of $`\beta `$ were nearer to 2 this would increase the rate of change of the ratio with redshift as the submm flux in the Rayleigh-Jeans tail $`\nu ^{2+\beta }`$ (see (iii)).
2. The radio spectral index is also somewhat uncertain, but being flatter it has much less impact than the submm slope. However, if high-redshift dust sources have systematically different spectral indices than those at low redshift, this would lead to a bias. The range of measured indices for the SLUGS sample is $`0.2`$ to $`1.0`$. The situation at high redshifts is harder to determine as the value of $`\alpha _{radio}`$ for $`\mu `$Jy sources may depend on the selection frequency, with sources selected at $`\nu >5`$ GHz likely to have flatter radio spectra than those selected at 1.4 GHz (Richards 1999). For example, the sources from the Canada-UK SCUBA survey have 5 GHz radio fluxes from Fomalont (1991) and an average $`\alpha _{radio}=0.38`$, while SCUBA observations of 1.4 GHz sources in the Hubble Flanking Fields were associated with radio sources with steeper spectra, $`\alpha _{rad}0.7`$ (Barger et al. 2000). This is most probably only a selection effect, as only relatively flat spectrum sources would be detected at 5 GHz. In general, it is preferable to survey at 1.4 GHz as this gives the best sensitivity to high redshift objects.
3. The mean dust temperature in the local galaxies could be over-estimated. If a colder component is present then it will produce a flattening of the thermal dust spectrum at longer wavelengths when compared to warm dust, thus lowering the redshift at which the estimator loses its sensitivity to redshift. We can model the implications of points (i) and (iii) as they are inter-related. If there is a cold dust component in the local galaxies, the true value of $`\beta `$ is likely to be nearer 2 and the mean dust temperature lower. If we make assumptions about the temperatures of the two components and the relative masses in each (and that $`\beta =2`$), we can then produce a new indicator. Figure 6 shows these ‘grey’ indicators for a few values of the various parameters which reflect current observational evidence (Frayer et al. 1999; Alton et al. 1998; Calzetti et al. 2000). The parameters used are given in Table 4, and all ‘grey’ models were normalised at zero redshift to have the mean value for the sample (0.18). There is generally no great difference until $`z2`$, and at this point the different parameters assumed start to make a larger impact. In general, we might expect that the galaxies detected by both SCUBA and the VLA would have $`T_w>30`$ K, since it is the recent star formation which the VLA is sensitive to (and which also produces higher dust temperatures). This makes the upper ($`T_w=40`$ K) curves more plausible although their shape is still sensitive to the mass fraction and temperature of the cold dust (colder dust or more cold dust will cause the indicator to turn over earlier and flatten more), and is also dependent on the value of $`\beta `$. If in reality, $`\beta `$ lies somewhere between $`1.32`$, then this will make the ‘grey’ models more like the median one at low redshifts. Figure 6 suggests that the uncertainties in whether there is actually a single dust temperature, or in the true value of $`\beta `$, do not add much to the uncertainty we have derived from the single-temperature fits. Currently, it is not possible to be more specific about the nature of any cold dust components as there are only a handful of local galaxies with enough FIR and submm fluxes to make a decomposition of the SED feasible. Our knowledge of dust properties in local galaxies should improve in the near future (Dunne et al. 2000).
4. Contributions to the radio flux by AGN will produce misleading redshift estimates if the object is not recognised as such. Since the fraction of deep SCUBA sources harbouring radio-loud AGN is still not well determined, the estimator must be treated as a statistical tool rather than a redshift indicator for individual objects. The redshifts given by the upper $`1\sigma `$ curve should be treated as a robust lower limit.
5. A variation in the FIR-radio correlation with redshift is quite possible if either magnetic fields or dust properties were different in the past, although currently this is difficult to test. Condon (1992) does however point out that the FIR-radio relation in the local universe holds over a large range of magnetic field strengths. The tendency of the galaxies observed so far in both the radio and submillimetre at high redshifts, to lie in agreement with the estimator (Fig. 2) suggests that there has been no dramatic evolution in the underlying FIR-radio relation, although more data is needed to fully investigate this.
## 5 Conclusions
We have used the submm data from a large, complete sample of local IRAS galaxies, along with complementary radio data, to define a median redshift estimator using the change in spectral slope between the submm and the radio at $`3`$mm. The scatter in the data has been used to provide $`1\sigma `$ uncertainties on the relationship. The limited number of submm sources from the deep surveys with both radio fluxes and spectroscopic redshifts generally agree with the estimator (within the uncertainties), except for one object which is known to be an AGN (Fig 2). The estimator is useful in a statistical sense rather than for predicting the redshifts of individual objects as there are many uncertainties, especially the possible contribution to the radio flux by AGN which would lower the data relative to the models thus leading to an under-estimated redshift. Our redshift indicator differs from the recent one of Carilli & Yun (2000) in terms of shape and normalisation. The difference is significant, particularly at redshifts $`>2`$ (Fig 1). The discrepancy can be attributed primarily to the larger sample size used in this work and also, in part, to the under-estimation of submm fluxes used in CY 2000. An erratum to CY 2000 has since been published which produces a revised redshift estimator using the corrected Lisenfeld et al. (2000) submm fluxes. This removes $`30`$ per cent of the discrepancy between the estimators, and we postulate that the remaining difference is a result of sample selection (the Lisenfeld sample is not a complete sample and contains a higher proportion of radio bright objects compared to SLUGS), and the smaller sample size (17 objects compared to 104 in SLUGS).
The submm-radio spectral index decreases as radio and FIR luminosity increase, but shows no strong trend with 850$`\mu `$m luminosity (Fig. 4). The correlations are believed to be due to the sensitivity of the non-thermal radio and thermal FIR emission to recent star formation, while a significant fraction of the 850$`\mu `$m flux may be due to colder dust, and not linked to the star formation rate. The absence of a correlation between $`\alpha _{1.4}^{850}`$ and 850$`\mu `$m luminosity means that any differences in $`L_{850}`$ between SLUGS and the SCUBA-selected galaxies will not lead to biased redshift estimates. Importantly, we also find no significant evidence for a systematic variation of $`\alpha _{1.4}^{850}`$ with dust temperature (Fig. 5(a)) implying that the redshift-temperature degeneracy hypothesised by Blain (1999) does not play a dominant part in the scatter of the $`\alpha _{1.4}^{850}`$ – redshift relation.
We have investigated the impact of neglecting a possible colder dust component in the SEDs of the local galaxies (Fig 6). While sensitive to the parameters assumed, in general the differences are not very significant for most plausible models and the agreement is still within the $`1\sigma `$ uncertainties at lower redshifts ($`<2`$). If cold components are present but unaccounted for, using the median line in Fig 1 to estimate redshifts will most likely lead to an over-estimate of $`z`$.
We would like to thank Chris Carilli and Ute Lisenfeld for useful discussions. The support of PPARC is gratefully aknowledged by L. Dunne, S. Eales and D. Clements. |
warning/0002/hep-th0002032.html | ar5iv | text | # 1 Introduction
## 1 Introduction
The problem of string quantization with Ramond-Ramond flux has received more attention due to the AdS/CFT correspondence. The construction of the worldsheet superconformal field theory of $`typeIIB`$ on $`AdS_3\times S^3`$ with Ramond-Ramond flux makes use of the $`N=4`$ topological string formulated by Berkovits and Vafa. The Berkovits-Vafa quantization is a hybrid of Ramond-Neveu-Schwarz and Green-Schwarz formalisms. The spacetime part is constructed from GS-like variables and the compactification variables are in RNS formalism. The supersymmetry is made manifest using GS-like variables and the supercharges do not have a singular OPE with the compactification variables.
The purpose of was to address string quantization with R-R flux, but $`AdS_3\times S^3`$ background with NS-NS flux is also considered in the paper. In this paper we compute the constraint equations for the vertex operator for this simpler theory to make contact with earlier literature concerning NS-NS flux (e.g.). Also, see and for brief history of string theory in $`AdS_3`$.
$`AdS_3\times S^3`$ background with NS-NS flux is an exact background whose conformal field theory is a WZW model with target space $`SU^{}(22)`$. The WZW model on $`SU^{}(22)`$ is a special limit of a two parameter conformal field theory. The two parameters are the curvature of the target space and the level of the affine Kac-Moody algebra associated with the supergroup. The generators of $`SU^{}(22)`$ are then incorporated in a straightforward way in constructing the $`N=4`$ topological superconformal algebra. The $`N=4`$ algebra is then used to impose physical and gauge conditions on vertex operators to obtain supersymmetric string constraints leading to string equations of motion. As in , we will show the eqiuvalence of the string equations and the supergravity field equations in the $`AdS_3\times S^3`$ background, but with NS-NS flux and different field identifications. This differs from the RNS formalism, when one has to calculate various amplitudes and then derive the form of the Lagrangian density to obtain the supergravity field equations.
In section two and three, we calculate supersymmetric string constraint equations for the spacetime vertex operators. The constraint equations are used to solve for the vertex operator. We then compute the string equations of motion and gauge conditions for the physical degrees of freedom carried by the vertex operator. In section four, we linearize $`N=(0,2),D=6`$ supergravity with one tensor multiplet around the $`AdS_3\times S^3`$ metric and self-dual three form field strength. With suitable field definitions we show that the string and supergravity equations of motion match. In the final section, we discuss the degrees of freedom remaining massless under the Kaluza-Klein compactification on $`S^3`$ and how they can be represented by vertex operators introduced by Giveon, Kutasov and Seiberg.
## 2 $`N=4`$ Algebra and Supersymmetric String Constraint Equations
The topological $`N=4`$ superconformal algebra(SCA) is the key ingredient for computing the supersymmetric string constraints for the compactification independent massless vertex operator for the $`typeIIB`$ on $`AdS_3\times S^3`$ background with NS-NS flux. $`N=4`$ algebra can be constructed in the following way. One takes the critical $`N=1`$ RNS superconformal field theory and constructs $`N=2,c=6`$ SCA by combining the matter and the ghosts. For example,
$$T=T_{N=1}+\frac{1}{2}(bc+\xi \eta )=T_{N=1}\frac{1}{2}J,$$
$$G^+=J_{BRST}+^2c+(c\xi \eta ),$$
$$G^{}=b,$$
$$J=cb+\eta \xi ,$$
where $`T_{N=1}`$ is the combined energy-momentum tensor of the original $`N=1`$ matter and ghost fields and the $`(\beta ,\gamma )`$ are the super-reparameterization ghosts bosonized as ($`\beta =e^\varphi \xi ,\gamma =e^\varphi \eta ).`$ For $`N=2`$ SCA, $`c=6`$ is the critical dimension allowing the $`N=2`$ string to be reformulated as an $`N=4`$ string. For instance, if $`J`$ is the $`U(1)`$ current of $`N=2,c=6`$ SCA, then the three $`SU(2)`$ currents are $`J,e^{\scriptscriptstyle J}`$ and $`e^{+{\scriptscriptstyle J}}`$. The four supercurrents are $`G^\pm ,\stackrel{~}{G}^{}=[e^{\scriptscriptstyle J},G^+]`$ and $`\stackrel{~}{G}^+=[e^{+{\scriptscriptstyle J}},G^{}].`$ Normally, the four supercurrents $`G^\pm `$ and $`\stackrel{~}{G}^\pm `$ in $`N=4`$ SCA have conformal dimensions of $`\frac{3}{2}`$ and are charged under the $`U(1)`$ current. In the topological string the the $`U(1)`$ current $`J`$ is constructed out of ghosts and the stress-energy tensor is modified or twisted by $`\frac{1}{2}J`$. Then the conformal dimensions of the generators are shifted by $`\frac{1}{2}`$ of their $`U(1)`$ charge. The conformal dimensions of $`G^+`$ and $`\stackrel{~}{G}^+`$ are then one and their zero modes are used as “BRST” charges. The $`N=4`$ topological SCA for the $`SU^{}(22)`$ target space is given by,
$$T=T_{SU^{}(22)}\frac{1}{2}\rho \rho \frac{1}{2}\sigma \sigma +\frac{3}{2}^2(\rho +i\sigma )+T_C,$$
$$G^+=e^{2\rho i\sigma }(p)^4+\frac{i}{2}e^\rho (p_ap_bK^{ab}+\frac{p_ap^a}{2k})$$
$$e^{i\sigma }[T_{SU^{}(22)}+\frac{1}{2}(\rho +i\sigma )(\rho +i\sigma )\frac{1}{2}^2(\rho +i\sigma )]+G_C^+,$$
$$G^{}=e^{i\sigma }+G_C^{},$$
$$\stackrel{~}{G}^+=e^{iH_C+\rho }+e^{\rho +i\sigma }\stackrel{~}{G}_C^+,$$
$$\stackrel{~}{G}^{}=e^{iH_C}\{e^{3\rho i2\sigma }(p)^4+\frac{i}{2}e^{2\rho i\sigma }(p_ap_bK^{ab}+\frac{p_ap^a}{2k})$$
$$e^\rho [T_{SU^{}(22)}+\frac{1}{2}(\rho +i\sigma )(\rho +i\sigma )\frac{1}{2}^2(\rho +i\sigma )]\}+e^{\rho i\sigma }\stackrel{~}{G}^+_C,$$
$$J=(\rho +i\sigma )+J_C,$$
$$J^{++}=e^{\rho +i\sigma }J_C^{++},$$
$$J^{}=e^{\rho i\sigma }J_C^{}.$$
$`T_{SU^{}(22)}`$ is the stress tensor for the $`SU^{}(22)`$ WZW model,
$$T_{SU^{}(22)}=\frac{1}{8k}ϵ_{abcd}K^{ab}K^{cd}\frac{1}{2k}ϵ_{\alpha \beta }S^{a\alpha }S_a^\beta $$
$$=\frac{1}{8k}ϵ_{abcd}j^{ab}j^{cd}+p^a\theta _a.$$
The cohomology of $`\stackrel{~}{G}_0^+`$ is taken to be trivial so that the physical vertex operator $`\mathrm{\Phi }^+`$ is written in terms of $`U(1)`$ neutral vertex operator $`V`$, $`\mathrm{\Phi }^+=\stackrel{~}{G}_0^+V`$ ($`O_n\mathrm{\Phi }`$ denotes $`(h+n)`$th pole in the OPE between the operator $`O`$ of conformal dimension $`h`$ and $`\mathrm{\Phi }`$). $`\mathrm{\Phi }^+`$ is BRST closed under $`G_0^+`$ because the OPE $`G^+\stackrel{~}{G}^+`$ is nonsingular. The gauge invariance of $`V`$ follows from the fact that $`V`$ is defined up to the equivalence class of the “BRST” charges. The gauge is fixed by choosing
$$G_0^{}V=\stackrel{~}{G}_0^{}V=0.$$
Furthermore, the gauge-fixing condition and $`G_0^+\mathrm{\Phi }^+=0`$ imply $`T_0V=0.`$ Similar conditions also arise in the right-moving sector. In summary, we have the following physical and gauge conditions which $`V`$ must satisfy (the right movers are barred):
$`T_0V=\overline{T}_0V=0,`$
$`G_0^{}V=\stackrel{~}{G}_0^{}V=\overline{G}_0^{}V=\overline{\stackrel{~}{G}}_0^{}V=0`$
$`G_0^+\stackrel{~}{G}_0^+V=\overline{G}_0^+\overline{\stackrel{~}{G}}_0^+V=0.`$ (1)
$`V`$ is used to describe the most general massless vertex operators that are independent of compactification fields. It is
$$V=\underset{m,n=\mathrm{}}{\overset{+\mathrm{}}{}}e^{m(i\sigma +\rho )+n(i\overline{\sigma }+\overline{\rho })}V_{m,n}(x,\theta ,\overline{\theta }).$$
(2)
$`G_0^{}V=\overline{G}_0^{}V=0`$ implies that $`V_{m,n}=0`$ for $`m>1`$ or $`n>1`$. $`\stackrel{~}{G}_0^{}V=\overline{\stackrel{~}{G}}_0^{}V=0`$ tells us that $`V_{m,n}=0`$ for $`m<1`$ or $`n<1.`$ This leaves
$$V_{1,1},V_{1,0},V_{1,1},$$
$$V_{0,1},V_{0,0},V_{0,1},$$
$$V_{1,1},V_{1,0},V_{1,1}$$
to consider.
$`\stackrel{~}{G}_0^{}V=\overline{\stackrel{~}{G}}_0^{}V=0`$ gives the following conditions on $`V_{m,n}(m,n=1,0,1)`$:
$$^4V_{1,n}=0,$$
(3)
$$\frac{1}{6}ϵ^{abcd}_b_c_dV_{1,n}=i(j_{0}^{}{}_{}{}^{ab}\frac{1}{4k}\delta ^{ab})_bV_{0,n},$$
(4)
$$\frac{1}{2}ϵ^{abcd}_c_dV_{o,n}=ij_0^{ab}V_{1,n},$$
(5)
$$j_{0}^{}{}_{}{}^{ab}_a_bV_{1,n}=0,$$
(6)
$$_aV_{1,n}=0,$$
(7)
and
$$\overline{}^4V_{m,1}=0,$$
(8)
$$\frac{1}{6}ϵ^{\overline{a}\overline{b}\overline{c}\overline{d}}\overline{}_{\overline{b}}\overline{}_{\overline{c}}\overline{}_{\overline{d}}V_{m,1}=i(j_{0}^{}{}_{}{}^{\overline{a}\overline{b}}\frac{1}{4k}\delta ^{\overline{a}\overline{b}})\overline{}_{\overline{b}}V_{m,0},$$
(9)
$$\frac{1}{2}ϵ^{\overline{a}\overline{b}\overline{c}\overline{d}}\overline{}_{\overline{c}}\overline{}_{\overline{d}}V_{m,0}=ij_0^{\overline{a}\overline{b}}V_{m,1},$$
(10)
$$j_{0}^{}{}_{}{}^{\overline{a}\overline{b}}\overline{}_{\overline{a}}\overline{}_{\overline{b}}V_{m,1}=0,$$
(11)
$$\overline{}_{\overline{a}}V_{m,1}=0.$$
(12)
Ansatz for $`V_{1,1}`$ is provided by ,
$`V_{1,1}=\theta ^a\overline{\theta }^{\overline{a}}V_{a\overline{a}}^{}+\theta ^a\theta ^b\overline{\theta }^{\overline{a}}\sigma _{ab}^m\overline{\xi }_{m\overline{a}}^{}+\theta ^a\overline{\theta }^{\overline{a}}\overline{\theta }^{\overline{b}}\sigma _{\overline{a}\overline{b}}^m\xi _{ma}^{}`$
$`+\theta ^a\theta ^b\overline{\theta }^{\overline{a}}\overline{\theta }^{\overline{b}}\sigma _{ab}^m\sigma _{\overline{a}\overline{b}}^n(g_{mn}+b_{mn}+\varphi \overline{g}_{mn})+\theta ^a(\overline{\theta }^3)_{\overline{a}}A_a^{+\overline{a}}+(\theta ^3)_a\overline{\theta }^{\overline{a}}A_{\overline{a}}^{+a}`$
$`+\theta ^a\theta ^b(\overline{\theta }^3)_{\overline{a}}\sigma _{ab}^m\overline{\chi }_m^{+\overline{a}}+(\theta ^3)^a\overline{\theta }^{\overline{a}}\overline{\theta }^{\overline{b}}\sigma _{\overline{a}\overline{b}}^m\chi _m^{+a}+(\theta ^3)_a(\overline{\theta }^3)_{\overline{a}}F^{++a\overline{a}}.`$ (13)
The remaining eight superfields can be computed up to integration constants using the constraints (3) - (13). These will contain new component fields, but these can be gauged away by the following gauge transformation:
$$\delta V=G_0^+\mathrm{\Lambda }+\overline{G}_0^+\overline{\mathrm{\Lambda }}+\stackrel{~}{G}_0^+\stackrel{~}{\mathrm{\Lambda }}+\overline{\stackrel{~}{G}}_0^+\overline{\stackrel{~}{\mathrm{\Lambda }}}.$$
The gauge parameters $`\mathrm{\Lambda },\stackrel{~}{\mathrm{\Lambda }},\overline{\mathrm{\Lambda }}`$, and $`\overline{\stackrel{~}{\mathrm{\Lambda }}}`$ are not arbitrary. $`\mathrm{\Lambda }`$ and $`\stackrel{~}{\mathrm{\Lambda }}`$ must be annhilated by $`\overline{G}_0^+\overline{\stackrel{~}{G}}_0^+`$, $`\overline{\mathrm{\Lambda }}`$ and $`\overline{\stackrel{~}{\mathrm{\Lambda }}}`$ annihilated by $`G_0^+\stackrel{~}{G}_0^+`$, and all the gauge parameters are annihilated by $`T_0`$, $`\overline{T}_0`$, $`G_0^{}`$, $`\stackrel{~}{G}_0^{}`$, $`\overline{G}_0^{}`$, and $`\overline{\stackrel{~}{G}}_0^{}`$. The parameters are chosen as follows,
$$\mathrm{\Lambda }=e^{2\rho +i\sigma +n(\overline{\rho }+i\overline{\sigma })}\lambda _n(x,\theta ,\overline{\theta }),\overline{\mathrm{\Lambda }}=e^{2\overline{\rho }+i\overline{\sigma }+n(\rho +i\sigma )}\overline{\lambda }_n(x,\theta ,\overline{\theta }),$$
$$\stackrel{~}{\mathrm{\Lambda }}=e^{\rho iH_C^{GS}+n(\overline{\rho }+i\overline{\sigma })}\stackrel{~}{\lambda }_n(x,\theta ,\overline{\theta })+\stackrel{~}{G}_0^{}\overline{\stackrel{~}{G}}_0^+\overline{\stackrel{~}{G}}_0^{}\widehat{\mathrm{\Lambda }},$$
$$\overline{\stackrel{~}{\mathrm{\Lambda }}}=e^{\overline{\rho }i\overline{H}_C^{GS}+n(\rho +i\sigma )}\overline{\stackrel{~}{\lambda }}_n(x,\theta ,\overline{\theta })+\stackrel{~}{G}_0^{}\overline{\stackrel{~}{G}}_0^+\overline{\stackrel{~}{G}}_0^{}\widehat{\mathrm{\Lambda }},$$
and $`\widehat{\mathrm{\Lambda }}`$ is defined as $`\frac{1}{2}e^{(\rho +i\sigma )+(\overline{\rho }+i\overline{\sigma })}\widehat{\lambda }.`$ In short, $`V_{1,1}`$ carries all the degrees of freedom and we will show that $`V_{1,1}`$ describes $`D=6,N=(0,2)`$ supergravity with one tensor multiplet.
Next we apply (6) and (11) on (13) to obtain gauge conditions on the fields. They are
$$j_0^{ab}(\sigma _{ab}^m\overline{\xi }_{m\overline{a}}^{})=0,$$
(14)
$$j_0^{ab}(\sigma _{ab}^m\sigma _{\overline{a}\overline{b}}^nG_{mn})=0,$$
(15)
$$j_0^{ab}(\sigma _{ab}^m\overline{\chi }_m^{+\overline{a}})=0,$$
(16)
$$ϵ_{abcd}j_0^{cd}A_{\overline{a}}^{+a}=0,$$
(17)
$$ϵ_{abcd}j_0^{cd}(\sigma _{\overline{a}\overline{b}}^m\chi _m^{+a})=0,$$
(18)
$$ϵ_{abcd}j_0^{cd}F^{++a\overline{a}}=0.$$
(19)
$$j_0^{\overline{a}\overline{b}}(\sigma _{\overline{a}\overline{b}}^m\xi _{ma}^{})=0,$$
(20)
$$j_0^{\overline{a}\overline{b}}(\sigma _{ab}^m\sigma _{\overline{a}\overline{b}}^nG_{mn})=0,$$
(21)
$$j_0^{\overline{a}\overline{b}}(\sigma _{\overline{a}\overline{b}}^m\chi _m^{+a})=0,$$
(22)
$$ϵ_{\overline{a}\overline{b}\overline{c}\overline{d}}j_0^{\overline{c}\overline{d}}A_a^{+\overline{a}}=0,$$
(23)
$$ϵ_{\overline{a}\overline{b}\overline{c}\overline{d}}j_0^{\overline{c}\overline{d}}(\sigma _{ab}^m\overline{\chi }_m^{+\overline{a}})=0,$$
(24)
$$ϵ_{\overline{a}\overline{b}\overline{c}\overline{d}}j_0^{\overline{c}\overline{d}}F^{++a\overline{a}}=0,$$
(25)
where $`G_{mn}=g_{mn}+b_{mn}+\overline{g}_{mn}\varphi `$, and $`\overline{g}_{mn}`$ is the metric of $`AdS_3\times S^3`$. Now we use the definitions of invariant derivatives on the $`SO(4)`$ group manifold to obtain the gauge conditions for the bosonic fluctuations.
$$D^m(g_{mn}+\overline{g}_{mn}\varphi )\frac{1}{2}Z_{npq}b^{pq}=0,$$
(26)
$$D^mb_{mn}=0,$$
(27)
$$\sigma _{ab}^pD_pA_d^{+a}\frac{1}{2}ϵ_{abcd}A^{+ca}=0,$$
(28)
$$\sigma _{ab}^pD_pA_d^{+b}+\frac{1}{2}ϵ_{abcd}A^{+cb}=0,$$
(29)
$$\sigma _{ab}^pD_pF^{++ae}\frac{1}{2}ϵ_{abcd}\delta ^{ed}F^{++ac}=0,$$
(30)
$$\sigma _{ab}^pD_pF^{++ea}\frac{1}{2}ϵ_{abcd}\delta ^{ed}F^{++ac}=0.$$
(31)
Here $`Z_{mnp}=(\sigma _m\sigma _n\sigma _p)_{ab}\delta ^{ab}`$ is the selfdual combination of sigma matrices.
## 3 String Equations of Motion
The physical condition that gives rise to string equations is $`T_0V=\overline{T}_0V=0,`$ or $`\frac{1}{8}ϵ_{abcd}j^{ab}j^{cd}V_{1,1}=\frac{1}{8}ϵ_{\overline{a}\overline{b}\overline{c}\overline{d}}j^{\overline{a}\overline{b}}j^{\overline{c}\overline{d}}V_{1,1}=0.`$ The string equations for the bosonic fields are
$`D^pD_p(g_{mn}+b_{mn}+\overline{g}_{mn}\varphi )=`$
$`Y_{(m}^qg_{n)q}Y_{pmnq}g^{pq}+2Z_{pq(n}D^pb_{m)}^q+2Z_{mnp}D^p\varphi `$
$`+2Y_{mn}\varphi {\displaystyle \frac{3}{2}}Y_{[m}^pb_{n]q}+2Z_{pq[n}D^pg_{m]}^q,`$ (32)
$$D^pD_pF^{++ab}=0$$
(33)
$$D^pD_pA_b^{+a}\sigma _{cb}^pD_pA^{+ca}=0,$$
(34)
$$D^pD_pA_b^{+a}+\sigma _{cb}^pD_pA^{+ca}=0,$$
(35)
$$D^pD_pV_{ab}^{}+\delta ^{cd}\sigma _{bc}^pD_pV_{ad}^{}+\delta ^{cd}\sigma _{ca}^pD_pV_{db}^{}\frac{1}{2}ϵ_{ab}^{cd}V_{cd}^{}=0,$$
(36)
where $`Y^{mn}\sigma _{ab}^m\sigma _{cd}^n\delta ^{ac}\delta ^{bd}`$, and
$$Y_{mnpq}\frac{1}{4}(\overline{g}_{mp}Y_{nq}+\overline{g}_{nq}Y_{mp}\overline{g}_{np}Y_{mq}\overline{g}_{mq}Y_{np}).$$
$`\overline{g}_{mn}=\frac{1}{2}\sigma _m^{ab}\sigma _{nab}`$ follows from the algebra of sigma matrices, $`\sigma _m^{ab}\sigma _{nad}+\sigma _n^{ab}\sigma _{mad}=\overline{g}_{mn}\delta _d^b.`$ Furthermore, we have used the following identity,
$`Y^{mnpq}={\displaystyle \frac{1}{4}}\delta ^{ah}\delta ^{kg}[\sigma _{ka}^m\sigma _{ge}^n\sigma _{hf}^p+\sigma _{ak}^m\sigma _{he}^n\sigma _{gf}^p+\sigma _{kf}^m\sigma _{ga}^n\sigma _{he}^p+\sigma _{af}^m\sigma _{hk}^n\sigma _{ge}^p]\sigma ^{qef}.`$ (37)
## 4 Linearized Supergravity Field Equations in $`AdS_3\times S^3`$
In the bosonic sector, the $`N=(0,2),D=6`$ supergravity field equations are
$$R_{mn}=H_{mpq}^iH_n^{ipq}+K_{mpq}K_n^{pq}+D_m\varphi ^iD_n\varphi ^i,$$
(38)
$$D^mD_m\varphi ^i=\frac{2}{3}H^{imnp}K_{mnp},$$
(39)
$$H_{mnp}^i=\frac{1}{3!}e_{mnpqrs}H^{iqrs},K_{mnp}=\frac{1}{3!}e_{mnpqrs}K^{qrs},$$
(40)
where $`e_{mnpqrs}\frac{1}{\sqrt{g}}g_{mm^{}}\mathrm{}g_{ss^{}}ϵ^{m^{}n^{}p^{}q^{}r^{}s^{}}`$, and $`ϵ^{012345}=1,1i5.`$ The Ricci tensor is defined as $`R_{mn}=g^{pq}R_{pmqn}`$. In $`AdS_3\times S^3`$ background the following relation holds, $`\overline{R}_{pmqn}=\frac{1}{4}(\overline{g}_{pq}\overline{R}_{mn}+\overline{g}_{mn}\overline{R}_{pq}\overline{g}_{mq}\overline{R}_{pn}\overline{g}_{pn}\overline{R}_{mq}).`$ Fields in (38)-(40) will be linearized around a background with non-vanishing $`AdS_3\times S^3`$ metric $`\overline{g}_{mn}`$ and selfdual three form field strength $`\overline{G}_{mnp}^5`$, i.e. only one of the components of the selfdual tensor is non-zero. The fluctuations about this background are expressed as
$`g_{mn}=\overline{g}_{mn}+h_{mn},H_{mnp}^i=\overline{G}_{mnp}^i+g_{mnp}^i,`$ (41)
$`K_{mnp}=\overline{K}_{mnp}+g_{mnp}^6+\overline{G}_{mnp}^j\varphi ^j.`$ (42)
From the scalar equation, the zeroth order is given by $`\overline{G}_{mnp}^i\overline{K}^{mnp}=0`$. To satisfy this we choose $`\overline{K}_{mnp}=0`$. From the Einstein equation, we get the identity $`\overline{R}_{mn}=\overline{G}_{mpq}\overline{G}_n^{pq}`$, where $`\overline{G}_{mnp}\overline{G}_{mnp}^5`$.
The linearized bosonic field equations that obey the zeroth order are
$$\frac{1}{3!}\overline{e}_{mnp}^{qrs}g_{qrs}^i=g_{mnp}^i3h_{q[m}\overline{G}_{np]}^{iq}+\frac{1}{2}h_q^q\overline{G}_{mnp}^i,$$
(43)
$$\frac{1}{3!}\overline{e}_{mnp}^{qrs}g_{qrs}^6=g_{mnp}^62\varphi ^i\overline{G}_{mnp}^i,$$
(44)
$$D^pD_p\varphi ^i=\frac{2}{3}\overline{G}^{imnp}g_{mnp}^6,$$
(45)
and
$`D^pD_ph_{mn}=2\overline{R}_{(m}^qh_{n)q}+2\overline{R}_{pmnq}h^{pq}+2D_{(m}D^ph_{n)p}`$
$`D_mD_nh_q^q4\overline{G}_{pq(m}g_{n)}^{5pq}+4\overline{G}_{mp}^q\overline{G}_n^{pr}h_{qr}.`$ (46)
$`g_{mnp}^i=3D_{[m}b_{np]}^i`$ and $`g_{mnp}^6=3D_{[m}b_{np]}^6`$ are no longer selfdual and anti-selfdual, respectively, but they are exact.
To obtain second order field equations for the tensors we take a covariant divergence of (43) and (44). We write the NS-NS two form as $`b_{mn}=Ab_{mn}^6+Bb_{mn}^5`$ rather than the most general form $`Ab_{mn}^6+B^ib_{mn}^i`$ because we will have to set $`B^x=0`$ for $`x=1,2,3,4`$ to satisfy (32). Since the covariant divergence of the dual of an exact form vanishes, we get
$`D^pD_pb_{mn}=\overline{R}_{[m}^sb_{n]s}2AD^p\varphi ^5\overline{G}_{pmn}`$
$`+{\displaystyle \frac{B}{C}}\overline{G}_{mnp}D_qg^{pq}+{\displaystyle \frac{2B}{C}}D^pg_{q[m}\overline{G}_{n]p}^q2D_{[m}D^pb_{n]p},`$ (47)
and
$$D^pD_pb_{mn}^x=\overline{R}_{[m}^pb_{n]p}^x.$$
(48)
We have expanded the fluctation $`h_{mn}`$ in terms of the graviton $`g_{mn}`$, $`h_{mn}=\frac{1}{C}g_{mn}+\frac{1}{6}\overline{g}_{mn}h_q^q`$, for some constant C to be determined. Next we use the gauge conditions (26) and (27). (47) then becomes
$`D^pD_pb_{mn}=\overline{R}_{[m}^pb_{n]p}2AD^p\varphi ^5\overline{G}_{mnp}`$
$`+{\displaystyle \frac{B}{C}}(D_q\varphi +{\displaystyle \frac{1}{2}}Z_{qpr})\overline{G}_{mn}^q+{\displaystyle \frac{2B}{C}}D^pg_{q[m}\overline{G}_{n]p}^q.`$ (49)
Comparing (49) with the antisymmetric part of (32) yields the following identifications: $`\overline{G}_{mnp}=\frac{C}{B}Z_{mnp}`$, $`\varphi =\frac{2AC}{3B}\varphi ^5`$ and $`\overline{R}_{mn}=\frac{1}{2}Y_{mn}`$, using $`Z_r^{mn}Z^{rpq}b_{pq}=2Y^{q[m}b_q^{n]}.`$
The scalar equation (39) reduces
$$D^pD_p\varphi ^5=\frac{2}{3A}\overline{G}^{pqr}H_{pqr}\frac{B}{AC}\overline{R}^{pq}g_{pq},$$
(50)
and
$$D^pD_p\varphi ^x=0.$$
(51)
We expand (46) by substituting $`h_{mn}=\frac{1}{C}g_{mn}+\frac{1}{6}\overline{g}_{mn}h_q^q`$. We get
$`D^pD_pg_{mn}+{\displaystyle \frac{C}{6}}\overline{g}_{mn}D^pD_ph_q^q=2\overline{R}_{(m}^qg_{n)q}+2\overline{R}_{pmnq}g^{pq}+2D_{(m}D^pg_{n)p}`$
$`{\displaystyle \frac{2C}{3}}D_mD_nh_q^q{\displaystyle \frac{4C}{B}}\overline{G}_{(m}^{pq}H_{n)pq}+4\overline{G}_{mp}^q\overline{G}_n^{pr}g_{qr}`$
$`+{\displaystyle \frac{2C}{3}}\overline{R}_{mn}h_q^q+{\displaystyle \frac{4CA}{B}}({\displaystyle \frac{1}{6A}}\overline{g}_{mn}\overline{G}^{pqr}H_{pqr}{\displaystyle \frac{B}{4AC}}\overline{g}_{mn}\overline{R}^{pq}g_{pq}\overline{R}_{mn}\varphi ^5),`$ (52)
where $`H_{mnp}=Ag_{mnp}^6+Bg_{mnp}^5.`$ To reduce (52) further, we trace it to obtain an expression for $`D^pD_ph_q^q`$, use the gauge condition (26) and use the identity $`\overline{G}_{mp}^q\overline{G}_n^{pr}g_{qr}=\overline{R}_{pmnq}g^{rq}`$.
Equation (52) will reduce to
$`D^pD_pg_{mn}{\displaystyle \frac{1}{5}}\overline{g}_{mn}D^pD_p\varphi ={\displaystyle \frac{1}{30}}Z_{pqr}H^{pqr}\overline{g}_{mn}{\displaystyle \frac{4}{5}}\overline{g}_{mn}\overline{R}^{pq}g_{pq}`$
$`+2\overline{R}_{(m}^qg_{n)q}2\overline{R}_{pmnq}g^{pq}(2D_{(m}D_{n)}\varphi +{\displaystyle \frac{2C}{3}}D_{(m}D_{n)}h_q^q)`$
$`+D_{(m}b^{pq}Z_{n)pq}{\displaystyle \frac{4C}{B}}D_{(m}b^{pq}\overline{G}_{n)pq}+{\displaystyle \frac{8C}{B}}\overline{G}_{pq(m}D^pb_{n)}^q+{\displaystyle \frac{2C}{3}}\overline{R}_{mn}h_q^q`$
$`+{\displaystyle \frac{2C}{3B}}\overline{g}_{mn}\overline{G}^{pqr}H_{pqr}{\displaystyle \frac{4CA}{B}}\overline{R}_{mn}\varphi ^5.`$ (53)
Next trace (53) to solve for $`D^pD_p\varphi `$ and use it to put (52) in the same form as the symmetric part of (32). The string equation does not have terms of the form $`D_{(m}D_{n)}\varphi `$. The simplest choice is to set $`\varphi =\frac{C}{3}h_q^q.`$ To satisfy (32) we must also have $`\overline{R}_{mn}=\frac{1}{2}Y_{mn}`$ and $`\frac{2C}{3}h_q^q\frac{4CA}{B}\varphi ^5=4\varphi .`$ With $`\varphi =\frac{C}{3}h_q^q`$, the latter expression gives $`\varphi =\frac{2CA}{3B}\varphi ^5.`$ This is the same expression we obtained from the tensor equation of motion. Thus (53) will reduce to
$`D^pD_pg_{mn}+D^pD_p\varphi =Y_{(m}^qg_{n)q}g^{pq}Y_{pmnq}`$
$`+2Y_{mn}\varphi +Z_{pq(m}D_{n)}b^{pq}(1{\displaystyle \frac{4C^2}{B^2}})`$
$`{\displaystyle \frac{8C^2}{B^2}}D_pb_{q(m}Z_{n)}^{pq}+\overline{g}_{mn}Z^{pqr}H_{pqr}({\displaystyle \frac{2C^2}{3B^2}}{\displaystyle \frac{1}{6}}).`$ (54)
We conclude that $`\frac{C^2}{B^2}=\frac{1}{4}`$ or $`\frac{C}{B}=\pm \frac{1}{2}.`$ For consistency we take $`\frac{C}{B}=\frac{1}{2}.`$
In summary, the following identifications
$`\overline{g}_{mn}={\displaystyle \frac{1}{2}}\sigma _m^{ab}\sigma _{nab},`$
$`\overline{R}_{mn}={\displaystyle \frac{1}{2}}\sigma _m^{ab}\sigma _n^{cd}\delta _{ac}\delta _{bd},`$
$`\overline{G}^{mnp}={\displaystyle \frac{1}{2}}\sigma _{ab}^m\sigma _{cd}^n\sigma ^{pbd}\delta ^{ac},`$
$`B=2C,\varphi ={\displaystyle \frac{1}{3}}\varphi ^{}={\displaystyle \frac{C}{3}}h_q^q.`$ (55)
will allow the supergravity equation to be recovered from the string equation (32).
The supergravity equations of tensor and scalar fields, (48) and (51), are identified with the string equations (33)-(36) in the following way. We know that the bispinor $`F^{++ab}`$ contains $`H_{mnp}^{++}`$ and a scalar $`\varphi ^{++}`$. In fact, we can use (30) and (31) to deduce
$$F^{++ab}=\sigma ^{pab}D_p\varphi ^{++}+(\sigma ^m\sigma ^n\sigma ^p)^{(ab)}H_{mnp}^{++}+\delta ^{ab}\varphi ^{++},$$
(56)
and
$$D^pD_p\varphi ^{++}=0.$$
(57)
(52) with (30) or (31) implies
$$D^pH_{pmn}^{++}=0,\mathrm{or}$$
$$D^pD_pb_{mn}^{++}=\overline{R}_{[m}^pb_{n]p}^{++},$$
(58)
We can do the same for $`b_{mn}^+`$ and $`b_{mn}^+`$. Define $`F^{+a\overline{a}}=(j^{\overline{a}\overline{b}}+\delta ^{\overline{a}\overline{b}})A_{\overline{b}}^{+a}`$ and $`F^{+a\overline{a}}=(j^{ab}+\delta ^{ab})A_b^{+\overline{a}}`$, then the gauge conditions for $`A_{\overline{a}}^{+a}`$ and $`A_a^{+\overline{a}}`$ are satisfied if we write
$$ϵ_{abcd}j^{cd}F^{+a\overline{a}}=0,ϵ_{\overline{a}\overline{b}\overline{c}\overline{d}}j^{\overline{c}\overline{d}}F^{+a\overline{a}}=0,$$
$$ϵ_{abcd}j^{cd}F^{+a\overline{a}}=0,ϵ_{\overline{a}\overline{b}\overline{c}\overline{d}}j^{\overline{c}\overline{d}}F^{+a\overline{a}}=0.$$
These definitions provide us with equations for $`b_{mn}^+`$, $`\varphi ^+`$ and $`b_{mn}^+`$, $`\varphi ^+`$ as in (56) and (57). Now we are left with (36) to consider. We define $`F^{a\overline{a}}=(j^{ab}+\delta ^{ab})(j^{\overline{a}\overline{b}}+\delta ^{\overline{a}\overline{b}})V_{b\overline{b}}^{}`$, then we can satisfy $`ϵ_{abcd}j^{cd}F^{a\overline{a}}=ϵ_{\overline{a}\overline{b}\overline{c}\overline{d}}j^{\overline{c}\overline{d}}F^{a\overline{a}}=0`$ so that $`ϵ_{abcd}j^{ab}j^{cd}F^{a\overline{a}}=0.`$ Therefore, we can express the string equations (33)-(36) as
$`D^pD_pb_{mn}^x=\overline{R}_{[m}^pb_{n]p}^x,`$
$`D^pD_p\varphi ^x=0,`$ (59)
where $`x=++,+,+,\mathrm{and}`$. (59) is identical to (48) and (51).
## 5 The Massless Multiplets and the GKS Construction
Following , numerous interesting papers on $`typeIIB`$ on $`AdS_3`$ appeared . We summarize the basic things we need to discuss the massless multiplet on $`AdS_3`$.
$`V_{1,1}`$ contains degrees of freedom to describe $`D=6,N=(0,2)`$ supergravity and one tensor multiplet, with representations
$$[(3,3)+5(3,1)+4(3,2)]+[(1,3)+5(1,1)+4(1,2)]$$
in $`Spin(4)SU(2)\times SU(2)`$, the lightcone little group.
We denote the anti-selfdual tensor as $`(1,3)`$. The Kaluza-Klein compactification to $`AdS_3\times S^3`$ will leave us with
$$(2,2)+4(1,1)+2(1,2)+2(2,1)$$
(60)
as massless propagating degrees of freedom, where the representation in (60) are under the $`SO(4)SU(2)\times SU(2)`$ gauge group that comes from $`S^3`$. The four massless degrees of freedom, $`(2,2)`$, comes from the linear combination of the anti-selfdual part of the NS-NS tensor and the dilaton. The four moduli $`4(1,1)`$ are inherited from the four R-R scalars $`\varphi ^x`$. The superpartners of the 8 scalars transform as $`2(1,2)+2(2,1)`$ under the gauge group. These states come from the Kaluza-Klein reduction of the six dimensional fermions, $`4(1,2).`$ So intuitively, we can see how the Kaluza-Klein compactification of $`typeIIB`$ on $`AdS_3\times S^3\times M^4(M^4=T^4`$ or $`K3)`$ gives $`n`$ massless mulitplet of (60);$`n=5`$ for $`T^4`$ and $`n=21`$ for $`K3`$. As an example, $`(2,2)`$ come from $`b_{mn}^6`$ and $`\varphi `$ of the $`\theta ^a\theta ^b\overline{\theta }^{\overline{a}}\overline{\theta }^{\overline{b}}`$ component of $`V_{1,1}`$ and $`(2,2)`$ is represented by $`\varphi _+^{(l+1,0)(l+1,0)}`$ with $`l=0`$, (we use the results of )
$`\theta ^a\theta ^b\overline{\theta }^{\overline{a}}\overline{\theta }^{\overline{b}}[\sigma _{ab}^\mu \sigma _{\overline{a}\overline{b}}^\nu (\overline{e}_{\mu \nu \rho }{\displaystyle \underset{l=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{l+2}{2(l+1)}}^\rho \varphi _+^{(l+1,0)(l+1,0)}(x)Y^{(l+1,0)}(y)`$
$`+\overline{g}_{\mu \nu }{\displaystyle \underset{l=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{l}{2(l+2)}}\varphi _+^{(l+1,0)(l+1,0)}(x)Y^{(l+1),0}(y))+\mathrm{}]`$ (61)
The greek indices label the coordinates of $`AdS_3`$. The spherical harmonics of $`S^3`$ are represented by $`Y^{(l_1,l_2)}`$ and $`(l_1,l_2)`$ label the highest weight of the gauge group $`SO(4)`$. They are related to $`SU(2)\times SU(2)`$ isospin $`(j,\overline{j})`$ as $`j=\frac{1}{2}(l_1+l_2)`$ and $`\overline{j}=\frac{1}{2}(l_1l_2)`$, and $`l_1l_2`$ since $`j,\overline{j}=0,\frac{1}{2},1,\mathrm{}`$
$`\varphi _+(x)`$ is labeled by the $`AdS_3`$ energy $`E_o`$, helicity $`s_o`$ and $`(l_1,l_2)`$, i.e. $`\varphi _+^{(E_o,s_o)(l_1,l_2)}`$. The $`AdS_3`$ energy $`E_o`$ and helicity $`s_o`$ are related to the spacetime conformal dimensions $`h_{st}=\frac{1}{2}(E_o+s_o),\overline{h}_{st}=\frac{1}{2}(E_os_o).`$ $`\varphi _+`$ is a chiral primary in the sense that the $`SU(2)`$ isospin is equal to the spacetime conformal dimension, $`(j,\overline{j})=(h_{st},\overline{h}_{st})`$.
The massless multiplet (60) can be described in terms of vertex operators used by . To be specific we take $`TypeIIB`$ on $`AdS_3\times S^3\times T^4`$. The holomorphic NS and its superpartner R vertex operators are (see for conventions)
$$𝒲_j=e^\varphi (\psi V_{jm})_{j1}V_{jm^{}}^{}$$
(62)
and
$$𝒴_j=e^{\frac{\varphi }{2}}S^{\dot{A}}(SV_{jm}V_{jm^{}}^{})_{j\frac{1}{2},j\frac{1}{2}}.$$
(63)
To describe (60) we take $`j=\frac{1}{2}`$(see below). In (62) and (63) we have used the following definitions:
$$(\psi V_{jm})_{j1}=(\psi ^3V_{jm}\frac{1}{2}\psi ^+V_{jm1}\frac{1}{2}\psi ^{}V_{jm+1})_{j1}$$
and
$$(SV_{jm}V_{jm^{}}^{})_{j\frac{1}{2},j\frac{1}{2}}=$$
$$S^{\dot{1}}V_{jm\frac{1}{2}}V_{jm\frac{1}{2}}^{}S^{\dot{2}}V_{jm\frac{1}{2}}V_{jm^{}+\frac{1}{2}}^{}$$
$$+S^{\dot{3}}V_{jm+\frac{1}{2}}V_{jm^{}\frac{1}{2}}^{}S^{\dot{4}}V_{jm+\frac{1}{2}}V_{jm^{}+\frac{1}{2}}^{}.$$
$`S^{\dot{A}}`$, $`\dot{A}=1,2`$, is the spin field of $`T^4`$ and $`S^{\dot{\alpha }}`$, $`\alpha =1,\mathrm{},4`$, is the spin field of $`AdS_3\times S^3`$ given (up to a cocycle) by $`S^{\dot{\alpha }}=e^{\frac{i}{2}(ϵ_1H_1+ϵ_2H_2+ϵ_3H_3)}.`$ $`ϵ_i=\pm `$ and $`ϵ_1ϵ_2ϵ_3=1.`$ $`S^\alpha `$ is given by $`ϵ_1ϵ_2ϵ_3=+1.`$ We have chosen
$$S^{\dot{1}}=(ϵ_1,ϵ_2,ϵ_3)=(,,),$$
$$S^{\dot{2}}=(,+,+),$$
$$S^{\dot{3}}=(+,,+),$$
$$S^{\dot{4}}=(+,+,).$$
Setting $`j=j^{}`$ gives worldsheet conformal dimension of 1 for (62) and (63). The subscript below the closed parenthesis of $`(\psi V_{jm})_{j1}`$ means that it has $`j1`$ representation under the $`SU(1,1)`$ Kac-Moody currents. Similarly, $`(SV_{jm}V_{jm^{}}^{})_{j\frac{1}{2},j\frac{1}{2}}`$ has $`(j\frac{1}{2},j\frac{1}{2})`$ representation under $`SU(1,1)\times SU(2)`$. (62) and (63) are BRST invariant and describe chiral primaries in spacetime because the spacetime conformal dimension is equal to the $`SU(2)`$ quantum number; $`h_{st}=j_{SU(2)}=j_{SU(1,1)}+1`$ for the NS vertex operator and $`h_{st}=j_{SU(2)}=j_{SU(1,1)}+\frac{1}{2}`$ for the R vertex operator. $`j_{SU(2)}`$ and $`j_{SU(1,1)}`$ refers to the total representation of the vertex operator under the Kac-Moody currents, i.e. $`j_{SU(1,1)}=j1,j_{SU(2)}=j`$ for (62), and $`j_{SU(1,1)}=j_{SU(2)}=j\frac{1}{2}`$ for (63). The $`AdS_3`$ part of NS vertex operators belong to the discrete unitary series $`D_{\stackrel{~}{j}}^+`$ with $`\stackrel{~}{j}=j_{SU(1,1)}=\frac{1}{2},0,\frac{1}{2},1,\frac{3}{2},\mathrm{}`$. The NS highest weight state is $`j_{SU(1,1)}=\frac{1}{2}`$ or $`j=\frac{1}{2}`$ and the R highest weight state has $`j_{SU(1,1)}=0`$. Then the four NS-NS scalars are described by
$$(j_{SU(2)},\overline{j}_{SU(2)})=(\frac{1}{2},\frac{1}{2})=(h_{st},\overline{h}_{st}),$$
and the four R-R scalars are labelled with
$$(j_{SU(2)},\overline{j}_{SU(2)})=(0,0)=(h_{st},\overline{h}_{st}).$$
Acknowledgement: The author would like to thank Louise Dolan for many helpful discussions. |
warning/0002/gr-qc0002038.html | ar5iv | text | # 1 Introduction
## 1 Introduction
There is no energy-momentum tensor of gravitational field in general relativity. Instead there are infinitely many pseudotensors . The concept of nonlocalizability of gravitational energy density was introduced to explain this unusual situation. Yet this concept seems unnecessary, if the gravity theory is build by field-theoretical means without requiring general covariance or if one assumes the existence of privileged frame. In general relativity the harmonic coordinates provide natural candidacy for privileged coordinates for an isolated system .
Weinberg’s pseudotensor (the suffix pseudo is dropped in the following) is singled out by the fact that it is the source of gravitational field . For this reason it is interesting to find it in harmonic coordinates, which goes over to Minkovskiian ones far away from gravitating body. Due to unusual properties of space-time beyond horizon and interchanged role of time and radial coordinates there , we expect that it will be impossible to get the collapsar total energy by integrating energy density over system’s volume. Calculations confirm this: the energy-momentum tensor have unintegrable singularities on the horizon.
## 2 Calculation of energy-momentum tensor
We use on the whole Weinberg notation, but denote harmonic coordinates $`x_i`$, $`|\stackrel{}{x}|=r`$ with small letters. Indices of $`h_{\mu \nu }`$, $`R_{\mu \nu }^{(1)}`$, $`\frac{}{x^\mu }`$ are raised and lowered with the help of $`\eta `$, indices of generally covariant tensors such as $`R_{\mu \nu }`$ are raised and lowered with the help of $`g`$. Latin indices run from 1 to 3.
$$g_{\mu \nu }=\eta _{\mu \nu }+h_{\mu \nu },\eta _{\mu \nu }=diag(1,1,1,1),\varphi =\frac{GM}{r},$$
$`(1)`$
$$d\tau ^2=g_{\mu \nu }dx^\mu dx^\nu =\frac{1+\varphi }{1\varphi }dt^2(1\varphi )^2d\stackrel{}{x}^2\frac{1\varphi }{1+\varphi }\varphi ^2\frac{(\stackrel{}{x}d\stackrel{}{x})^2}{r^2}.$$
$`(2)`$
Now we reproduce eqs. (7.6.3) (7.6.4) from . The exact Einstein equations are written there in the form
$$R_{\mu \kappa }^{(1)}\frac{1}{2}\eta _{\mu \kappa }R^{(1)}=8\pi G[T_{\mu \kappa }+t_{\mu \kappa }],$$
$`(3)`$
where
$$t_{\mu \kappa }=\frac{1}{8\pi G}[R_{\mu \kappa }\frac{1}{2}g_{\mu \kappa }RR_{\mu \kappa }^{(1)}+\frac{1}{2}\eta _{\mu \kappa }R^{(1)}],$$
$`(4)`$
and $`R_{\mu \kappa }^{(1)}`$ is linear in $`h`$ part of $`R_{\mu \kappa }:`$
$$R_{\mu \kappa }^{(1)}=\frac{1}{2}[h_{,\mu \kappa }h^\lambda {}_{\mu ,\lambda \kappa }{}^{}h^\lambda {}_{\kappa ,\lambda \mu }{}^{}+h_{\mu \kappa ,\lambda }{}_{}{}^{\lambda }],h_{,\mu }=\frac{}{x^\mu }h.$$
$`(5)`$
Eq.(3) has the form of wave equation for spin-2 field. Its source is $`T_{\mu \kappa }+t_{\mu \kappa }`$. Hence, $`t_{\mu \kappa }`$ (i.e. energy-momentum tensor of gravitational field) is also a source of gravitational field. Eq. (3) is suggested by solution of Einstein equations by iteration. In linear approximation $`h_{\mu \nu }=h^{(1)_{\mu \nu }}`$ is generated by material tensor. Inserting this solution of linearized equation in (4) and retaining quadratic in $`h^{(1)}`$ terms, we get $`t_{\mu \kappa }^{(2)}`$. \[See eq.(7.6.14) in , in which it is not indicated explicitly that figuring there $`h_{\mu \nu }`$ are $`h_{\mu \nu }^{(1)}`$. The expression for $`t_{\mu \nu }^{(2)}`$, obtained from that eq., for Newtonian center is given in . It coincides with eqs. (11), (13) below, which were found from exact expressions.\] Further, $`t_{\mu \nu }^{(2)}`$ according to wave equation generates $`h_{\mu \nu }^{(2)}`$ and so on. The sum over all approximations gives $`h_{\mu \kappa }`$, which exactly satisfies Einstein equations.
Now we assume up to eq.(14) that the matter is absent in the considered region. Then the energy-momentum tensor has the form
$$t_{\mu \nu }=\frac{1}{8\pi G}[\frac{1}{2}\eta _{\mu \nu }R^{(1)}R_{\mu \nu }^{(1)}],R^{(1)}=R_\lambda ^{(1)}{}_{}{}^{\lambda }=h_{,\lambda }{}_{}{}^{\lambda }h^{\mu \nu }{}_{,\mu \nu }{}^{},h=h_\lambda {}_{}{}^{\lambda }.$$
$`(6)`$
From (1-2) and (5-6) we find
$$h=2\varphi ^24\varphi 4+\frac{2}{1\varphi }+\frac{2}{1+\varphi },h_{,\lambda }{}_{}{}^{\lambda }=\frac{4}{r^2}\varphi ^2[\frac{1}{(1\varphi )^3}+\frac{1}{(1+\varphi )^3}+1],$$
$$h_{ij}=(1\varphi )^2\delta _{ij}+\frac{\varphi ^2\varphi ^3}{1+\varphi }\frac{x_ix_j}{r^2}\delta _{ij},$$
$`(7)`$
$$h_{ij,kl}=\frac{2x_ix_jx_kx_l}{r^6}\left(12\varphi ^2+15\varphi 8+\frac{3}{1+\varphi }+\frac{3}{(1+\varphi )^2}+\frac{2}{(1+\varphi )^3}\right)+$$
$$\frac{2\delta _{ij}x_kx_l}{r^4}(4\varphi ^23\varphi )+\frac{2\delta _{ij}\delta _{kl}}{r^2}(\varphi \varphi ^2)+\frac{\delta _{ik}\delta _{jl}+\delta _{il}\delta _{jk}}{r^2}\left(\varphi ^2+2\varphi 2+\frac{2}{1+\varphi }\right)$$
$$+\frac{2}{r^4}[x_ix_j\delta _{kl}+x_kx_j\delta _{il}+x_ix_k\delta _{jl}+x_jx_l\delta _{ik}+x_ix_l\delta _{jk}]\times $$
$$\times \left(2\varphi ^23\varphi +2\frac{1}{1+\varphi }\frac{1}{(1+\varphi )^2}\right).$$
$`(8)`$
With the help of this expressions we obtain
$$R_{kl}^{(1)}=\frac{x_kx_l}{r^4}\left(2\varphi ^22+\frac{1}{1+\varphi }+\frac{1}{(1+\varphi )^2}\frac{1}{1\varphi }\frac{1}{(1\varphi )^2}+\frac{2}{(1\varphi )^3}\right)+$$
$$\frac{\delta _{kl}}{r^2}\left(2\frac{3}{1+\varphi }+\frac{1}{(1+\varphi )^2}+\frac{1}{1\varphi }\frac{1}{(1\varphi )^2}\right),$$
$$R^{(1)}=\frac{1}{r^2}\left(2\varphi ^2+4+\frac{4}{1\varphi }\frac{8}{(1\varphi )^2}+\frac{4}{(1\varphi )^3}\frac{8}{1+\varphi }+\frac{4}{(1+\varphi )^2}\right),$$
$`(9)`$
$$t_{kl}=\frac{1}{8\pi G}[\frac{x_kx_l}{r^4}(2\varphi ^2+2\frac{1}{1+\varphi }\frac{1}{(1+\varphi )^2}+\frac{1}{1\varphi }+\frac{1}{(1\varphi )^2}\frac{2}{(1\varphi )^3})+$$
$$\frac{\delta _{kl}}{r^2}(\varphi ^2\frac{1}{1+\varphi }+\frac{1}{(1+\varphi )^2}+\frac{1}{1\varphi }\frac{3}{(1\varphi )^2}+\frac{2}{(1\varphi )^3})].$$
$`(10)`$
For $`\varphi 1`$ we have
$$t_{ik}|_{\varphi 1}=\frac{\varphi ^2}{8\pi G}[\frac{7\delta _{ik}}{r^2}\frac{14x_ix_k}{r^4}].$$
$`(11)`$
Similarly, we find
$$t_{00}=\frac{1}{8\pi Gr^2}\left(\frac{4}{1+\varphi }\frac{2}{(1+\varphi )^2}\varphi ^22\right),$$
$`(12)`$
$$t_{00}|_{\varphi 1}=\frac{3}{8\pi G}(\varphi )^2=\frac{3GM^2}{8\pi r^4}.$$
$`(13)`$
Now we check that conservation laws $`t^{\mu \nu }{}_{,\nu }{}^{}=0`$ are fulfilled. As $`t_{i0}=0`$, we need to verify that $`t_{ni,n}=0`$. Straightforward calculation gives
$$R_{ni,n}^{(1)}=\frac{1}{2}R_{,i}^{(1)}=$$
$$\frac{x_i}{r^4}[44\varphi ^2+\frac{4}{1+\varphi }+\frac{4}{(1+\varphi )^2}\frac{4}{(1+\varphi )^3}\frac{2}{1\varphi }\frac{2}{(1\varphi )^2}+\frac{10}{(1\varphi )^3}\frac{6}{(1\varphi )^4}],$$
Q.E.D.
Introducing tensor
$$Q^{\rho \nu \lambda }=\frac{1}{2}[h^{,\nu }\eta ^{\rho \lambda }h^{,\rho }\eta ^{\nu \lambda }h^{\mu \nu }{}_{,\mu }{}^{}\eta _{}^{\rho \lambda }+h^{\mu \rho }{}_{,\mu }{}^{}\eta _{}^{\nu \lambda }+h^{\nu \lambda ,\rho }h^{\rho \lambda ,\nu }],$$
$`(14)`$
with property $`Q^{\rho \nu \lambda }=Q^{\nu \rho \lambda }`$, we have, see. Ch.7, §6 :
$$Q^{\rho \nu \lambda }{}_{,\rho }{}^{}=R^{(1)\nu \lambda }\frac{1}{2}\eta ^{\nu \lambda }R^{(1)}=8\pi G\tau ^{\nu \lambda },$$
$`(15)`$
$$\tau ^{\nu \lambda }=\eta ^{\mu \nu }\eta ^{\lambda \kappa }[T_{\mu \kappa }+t_{\mu \kappa }].$$
Due to this relation for smooth tensor $`Q^{\rho \nu \lambda }`$ the integral of total (gravitational and material) energy density over the volume of system may be written as an integral over remote surface (see )
$$P^0=\frac{1}{8\pi G}Q^{i00}{}_{,i}{}^{}d_{}^{3}x=\frac{1}{8\pi G}Q^{i00}n_ir^2𝑑\mathrm{\Omega }=M,$$
$`(16)`$
$$Q^{i00}=\frac{1}{2}(h_{jj,i}h_{ij,j})=\frac{x_i}{r^2}\left(2\varphi ^2\frac{2}{1+\varphi }\right),$$
$`n_i`$ are components of external normal to the surface. Yet in case of horizon we see from (10) and (12) that at $`\varphi =1`$ (i.e. at $`r=GM`$ in harmonic system) tensor $`t_{\mu \kappa }`$ has unintegrable singularity. This prevent us from using Gauss theorem (i.e. from going to the second equation in (16)) in the whole volume of system. But it is easy to find the energy outside a sphere of radius $`r_1=GM(1+\delta )`$:
$$\mathrm{\Delta }P^0=\frac{r}{2G}\left(2\varphi ^2\frac{2}{1+\varphi }\right)|_{r_1}^{\mathrm{}}$$
$`(17)`$
For $`0<\delta 1`$ we get
$$\mathrm{\Delta }P^0=M\left(\frac{1}{\delta }+\frac{1}{2}\right).$$
$`(19)`$
It is interesting to compare (19) with Dehnen result . Using standard
Schwarzschild coordinates ($`r^{}=r+GM`$) Dehnen find for his tensor
$$\mathrm{\Delta }P^0=GM^2_{r_1^{}}^{\mathrm{}}\frac{dr^{}}{r^2\left(1\frac{r_g}{r^{}}\right)^{\frac{3}{2}}}=M\left(1\frac{1}{\sqrt{1\frac{r_g}{r_1^{}}}}\right).$$
$`(20)`$
This expression has square root divergence for $`r_1^{}r_g`$.
It is interesting that in Box 23.1 in arguments are given in favor of localazability of gravitational energy density in case of spherical symmetry. According to this arguments the gravitational energy outside the matter ball is zero. In Newtonian limit this way of accountig for gravitational energy corresponds to accounting for gravitational attractions between different parts of matter. It is not clear how to reconcile this with (19) or (20) and with the concept that nonlinear corrections to gravitational field are generated by gravitational energy-momentum tensor. As to the nonuniqueness of energy-momentum tensors, one may note that if some tensor correctly describes the interaction with gravitons then it is natural to consider this tensor as the right one.
Although Schwarzschild singularity is considered fictitious (from the time of Lemaître, see box 31.1 in ), it is difficult to be reconciled with this. Physically the singularity manifest itself in impossibility for cosmonaut, crossing it, to return back, in unlimited growth of acceleration (in static frame) of freely falling particle nearing horizon, in absence of static frame beyond horizon and in many other unusual things.
One can consider these singularities as a hint that in the regime of strong field the theory will be modified in the future and no drastic changes in space-time topology will be needed.
The author is grateful to V.I.Ritus for useful discussions and constructive remarks.
## 3 ץ
1. Misner C.W., Thorne K.S., Wheeler J.A.,Gravitation, San Francisco (1973).
2. Goldberg J.N., Phys. Rev. 111, 315 (1958).
3. Fock V., The of Theory of Space-Time and Gravitation (2nd revised edition, Pergamon Press, New York, 1964).
4. Weinberg S., Gravitation and Cosmology, New York (1972).
5. Novikov I.D., Frolov V.P. Physics of black holes, Nauka, Moscow (1986).
6. A.I.Nikishov, gr-qc/9912034.
7. Dehnen H. Zeitschr. für Phys., 179, 76 (1964). |
warning/0002/gr-qc0002056.html | ar5iv | text | # Large Quantum Gravity Effects: Cylindrical Waves in Four Dimensions
## 1 Introduction
The quantization of spacetimes with two commuting spacelike Killing fields has deserved intensive study during recent years \[1-11\]. The main reason for this interest is the ability of this type of spacetimes to provide a suitable framework in which one can discuss conceptual problems and develop mathematical methods for the quantization of general relativity. The families of solutions with two Killing fields that have been quantized in the literature, though simple enough as to be tractable, still possess an infinite number of degrees of freedom, so that they are expected to retain the field complexity that should be present in the elusive theory of quantum gravity.
Additional motivation for the quantum analysis of spacetimes with two commuting Killing fields comes from their possible application to cosmology and astrophysics. Most of the spacetimes of this kind that have been subject to quantization can in fact be interpreted as gravitational waves propagating in a source-free background. Particular examples are the Gowdy cosmologies with the spatial topology of the three-torus , the family of purely gravitational plane waves , and the set of cylindrical waves in vacuum gravity . Among all of them, it is the family of cylindrical waves whose quantization has received more attention and is probably best understood \[1-7,14,15\].
The pioneer of this midisuperspace approach to quantum gravity was Kuchař . He made a preliminary discussion of the quantum mechanics for Einstein-Rosen waves , i.e., cylindrical gravitational waves in four dimensions with linear polarization. The classical description of these waves is equivalent to that corresponding to a rotationally symmetric massless scalar field coupled to three-dimensional gravity and, therefore, to the description of a rotationally symmetric Einstein-Maxwell model in three dimensions . Recently, a rigorous quantization of this three-dimensional counterpart of the Einstein-Rosen model has been carried out by Ashtekar and Pierri , superseding previous work on the subject by Allen . Some technical details concerning the self-adjointness of the metric operators in this three-dimensional system have been revised by Varadarajan . On the other hand, a quantum theory for the most general family of cylindrical waves in four-dimensional gravity, which exploits the group-theoretical properties of the system, has been presented by Korotkin and Samtleben, although no explicit construction has been provided for the metric operators .
An issue that has been investigated with special interest in this context is the existence of quantum gravitational states that can be approximated by a classical solution, or semiclassical one if quantum matter fields are included. By analyzing the three-dimensional theory that is obtained from the Einstein-Rosen waves via dimensional reduction, Ashtekar has proved that, at least in a certain sector of quantum gravity, the semiclassical approximation may become meaningless owing to the appearance of huge quantum gravity effects . Namely, in the rotationally symmetric Einstein-Maxwell model in three dimensions, Ashtekar has considered all coherent states of the Maxwell field and computed their expectation values and quantum fluctuations in the three-metric at large distances from the center (i.e., the axis of symmetry). For coherent states that are sharply peaked around a characteristic wave number $`k_0`$, the asymptotic expectation value of the three-metric is peaked around a classical solution if and only if $`N(e^{k_0}1)^21`$, where $`N`$ is the number of photons contained in the state. In addition, if one requires the quantum uncertainties in the Maxwell field to be relatively small, the semiclassical description is accurate only when $`N1`$ . Here, and in the rest of the paper, we have used units in which $`c=\mathrm{}=8G_3=1`$, $`G_3`$ being the gravitational constant in three dimensions or, equivalently, the effective Newton constant per unit length in the direction of the symmetry axis.
The possibility of finding states in the Einstein-Maxwell system with improved coherence in the three-metric at the expense of increasing the dispersion in the Maxwell field was proved by Gambini and Pullin . Large quantum gravity effects similar to those detected by Ashtekar were also found in the rotationally symmetric Einstein-Maxwell model by employing non-local variables , and in a three-dimensional model with toroidal symmetry . Only one purely four-dimensional gravitational system has been discussed in which the quantum fluctuations invalidate the classical description of the geometry: a midisuperspace model for linearly polarized plane waves in vacuum gravity . In this case, the huge fluctuations appear in a region where null geodesics are focused, and not in the asymptotic region.
The aim of the present work is to revisit Ashtekar’s results about the existence of large quantum effects in cylindrical gravity from a strictly four-dimensional point of view. The classical equivalence of the Einstein-Rosen and the three-dimensional Einstein-Maxwell systems does not necessarily imply their quantum equivalence. On the other hand, since the Einstein-Rosen model and its three-dimensional counterpart have different metrics, all questions about the existence of quantum states peaked around classical geometries in general relativity should be addressed from a four-dimensional perspective. In fact, as we will see, the Einstein-Rosen metric can be expressed as a function of the Maxwell field and the metric in three dimensions which is highly non-linear in the matter field. As a consequence, coherence in the four-metric does not generally follow from coherence in the three-metric and the field.
The plan of the paper is the following. In Sec. 2, we construct a midisuperspace model for cylindrical waves starting from the Hamiltonian formulation of general relativity for spacetimes that possess two commuting spacelike Killing fields. We adopt a gauge-fixing procedure that is similar to that introduced by Ashtekar and Pierri in three dimensions , and calculate the reduced Hamiltonian of the model by a careful analysis of surface terms in the gravitational action. This framework is then particularized to the case of linearly polarized waves via symmetry reduction. Employing the quantum theory put forward in Ref. for the rotationally symmetric three-dimensional system, we present a complete quantum theory for the Einstein-Rosen waves in Sec. 3. In particular, we obtain regularized, positive operators to describe all components of the four-metric. The behavior of these operators on the quantum states that are coherent in the Maxwell field is discussed in Sec. 4. We first analyze the quantum gravitational effects at large distances of the symmetry axis, showing that the conclusions reached by Ashtekar for the three-dimensional metric are valid as well for the metric of the Einstein-Rosen waves, not only qualitatively, but also quantitatively. From the point of view of the four-metric, however, the requirement of a classical behavior for the Maxwell field is now spurious, so that the condition $`N1`$ is no longer necessary to reach an acceptable classical approximation in the asymptotic region. Using our four-dimensional formalism, we are also able to study the quantum fluctuations in the metric on the symmetry axis. We argue that these fluctuations cannot be neglected for any of the coherent states. We summarize our results and conclude in Sec. 5. Finally, two appendices are added. In Appendix A we prove some useful operator identities, while Appendix B contains some calculations employed in the discussion of the metric fluctuations.
## 2 The Midisuperspace Model
Let us first construct a gauge-fixed midisuperspace model to describe cylindrical waves in vacuum gravity. Since this family of waves can be regarded as a particular class of spacetimes that possess two commuting Killing vectors, we can start our analysis with the Hamiltonian formulation for spacetimes of this kind, which is discussed in Sec. 3 of Ref. . For convenience, we adopt the notation $`\{x^i\}\{Z,\theta ,R\}`$ ($`i=1,2,3`$) for the spatial coordinates and assume that the two commuting Killing vector fields are $`_{x^a}`$ ($`a=1,2`$), so that the metric is independent of $`\theta `$ and $`Z`$. In addition, we impose that $`R0`$ and $`\theta S^1`$ (with $`S^1`$ being the unit circle). With this terminology, $`Z`$ denotes the coordinate of the symmetry axis, whereas $`R`$ and $`\theta `$ are the radial and angular coordinates on each surface of constant $`Z`$ and time $`t`$.
The momentum constraints corresponding to the coordinates $`x^a`$ can be eliminated by requiring that the induced three-metric $`h_{ij}`$ is block-diagonal, namely, $`h_{aR}=0`$ (where we have adopted the alternative notation $`h_{aR}`$ instead of $`h_{a3}`$). The gauge fixing is almost identical to that explained in Ref. for the case of plane waves, and we will not repeat details here. Apart from the different domains of definition for the spatial coordinates, the only modification that must be introduced concerns the system of units. In the cited paper, the authors set $`c=\mathrm{}=4G_3=1`$, where $`G_3=G/(𝑑Z)`$ is the effective Newton constant per unit length. In the present work, however, we have fixed $`8G_3=1`$ (to facilitate comparison of our results with those of Ashtekar and Pierri). We can nevertheless take account of this discrepancy by simply multiplying all gravitational constraints in Ref. by a factor of two and dividing the canonical momenta of the metric functions by the same factor.
As shown in Ref. , the dynamical stability of the gauge-fixing conditions $`h_{aR}=0`$ requires that $`4N\sqrt{h_{RR}}f_a=\sqrt{\mathrm{det}h_{cd}}h_{ab}(N^b)^{}`$, where $`f_a`$ are two constants (independent of $`R`$) that determine the momenta of $`h_{aR}`$. Here, $`N`$ is the lapse function, $`N^i`$ is the shift vector, and the prime stands for the derivative with respect to $`R`$. Since the two-metric $`h_{ab}`$ becomes degenerate on the symmetry axis (that we suppose located at $`R=0`$), the regularity of the four-metric on this axis implies that the constants $`f_a`$ must vanish. As a consequence, we conclude that the components $`N^a`$ of the shift vector are independent of the spatial coordinates, and can be absorbed by a redefinition of $`x^a`$. It hence turns out that the condition of regularity on the axis suffices to ensure that the orbits spanned by the two Killing vectors admit orthogonal surfaces.
The remaining momentum constraint can be eliminated in a very similar way to that discussed at the end of Sec. 3 and the beginning of Sec. 4 in Ref. . One only needs to change the choice of the strictly increasing function $`z_0`$ that determines the coordinate $`R`$. We now select $`z_0=\mathrm{ln}R`$. In this way, the radial coordinate $`R`$ is set to coincide with the square root of the determinant of the metric on Killing orbits. We notice that our gauge fixing for the momentum constraint associated with $`R`$ is analogous to that performed by Ashtekar and Pierri in three dimensions .
The resulting reduced system has a configuration space with three degrees of freedom which, with the conventions of Ref. , can be chosen as the three metric functions $`v`$, $`y`$, and $`w`$. In order to adopt a notation similar to that employed in the three-dimensional Einstein-Maxwell model , it is convenient to introduce the definitions $`\psi =\mathrm{ln}Ry/2`$ and $`\gamma =2w`$. The system has still one constraint, namely, the Hamiltonian constraint, which can now be written
$``$ $`=`$ $`{\displaystyle \frac{e^{(\psi \gamma )/2}}{2R}}\left[R^2(\psi ^{})^22R\gamma ^{}+e^{2\psi }(v^{})^2+p_\psi ^2+R^2e^{2\psi }p_v^2\right]`$ (2.1)
$`+`$ $`e^{(\psi \gamma )/2}p_\gamma (p_vv^{}+p_\psi \psi ^{}+p_\gamma \gamma ^{}2p_\gamma ^{}).`$
The corresponding gauge freedom can be eliminated by imposing the vanishing of the momentum canonically conjugate to $`\gamma `$: $`p_\gamma =0`$. This condition is inspired by the gauge fixing carried out in the three-dimensional counterpart of our model . It is straightforward to check that the gauge fixing is well posed. In addition, the gauge condition is preserved by the dynamical evolution provided that $`\{p_\gamma ,𝑑RN\}(e^{(\psi \gamma )/2}N)^{}=0`$, where the symbols $`\{,\}`$ and $``$ denote Poisson brackets and weak identity, respectively. Hence, the lapse function must be of the form $`N=f(t)e^{(\gamma \psi )/2}`$, with $`f(t)`$ being a function of time (that can generally be absorbed by a redefinition of $`t`$). We will choose this function equal to $`e^{\gamma _{\mathrm{}}/2}`$, where $`\gamma _{\mathrm{}}`$ is the value of the metric function $`\gamma `$ when $`R\mathrm{}`$. As we will see below, this choice guarantees that $`_t`$ is a unit asymptotic time translation.
On the other hand, the solution to the Hamiltonian constraint with $`p_\gamma =0`$ is
$$\gamma =\frac{1}{2}_0^R𝑑\overline{R}\overline{R}\left[(\psi ^{})^2+\frac{p_\psi ^2}{\overline{R}^2}+e^{2\psi }\frac{(v^{})^2}{\overline{R}^2}+e^{2\psi }p_v^2\right],$$
(2.2)
where we have imposed that $`\gamma `$ vanish at $`R=0`$ in order to obtain (with suitable boundary conditions on $`\psi `$ and $`v`$) a regular metric on the axis of symmetry. After our gauge fixing, the line element has the expression
$$ds^2=e^\psi \left[e^\gamma (e^\gamma _{\mathrm{}}dt^2+dR^2)+R^2d\theta ^2\right]+e^\psi (dZvd\theta )^2.$$
(2.3)
Assuming as a boundary condition (see Ref. for a detailed discussion in the linearly polarized case) that the metric functions $`\psi `$ and $`v`$ fall off sufficiently fast as $`R\mathrm{}`$ (so that, in particular, $`\gamma _{\mathrm{}}`$ is finite), we get that the above metric describes an asymptotically flat spacetime with a generally non-zero deficit angle. In this asymptotic region, as we anticipated, $`_t`$ is a unit timelike vector.
The reduced model obtained in this way is free of constraints and has only two metric degrees of freedom, described by the variables $`\psi `$ and $`v`$. Its reduced symplectic structure is $`\mathrm{\Omega }=𝑑R(𝐝p_\psi 𝐝\psi +𝐝p_v𝐝v)`$, where $`𝐝`$ and $``$ denote, respectively, the exterior derivative and product. The Hamiltonian that generates the dynamics of the model, on the other hand, can be obtained by reducing the gravitational Hilbert-Einstein action supplemented with appropriate boundary terms . Let us explain this point in more detail. In our gauge-fixing procedure, we have removed some of the original degrees of freedom by expressing them in terms of the remaining canonical variables and, possibly, of the coordinates. All the expressions employed are in fact local, except in the very last step of the procedure, where relation (2.2) has been introduced. It is not difficult to realize then that, in our discussion of the gauge fixing, the dynamical equations that we have computed via Poisson brackets are actually valid in the interior of our manifold, even though we have not explicitly included surface terms in the Hamiltonian. This fact ensures that our gauge fixing has been carried out consistently. Furthermore, it then follows that the Hamiltonian of the reduced model is actually given by the reduction of the total Hamiltonian (including surface terms) of our original system. Since, as we have pointed out, relation (2.2) is not local, this reduced Hamiltonian may be non-trivial.
The boundary terms for the gravitational Hamiltonian have been recently analyzed by Hawking and Hunter . To apply their results to our reduced model, let us first consider a manifold that, on each section $`\mathrm{\Sigma }_t`$ of constant time, has a two-dimensional boundary $`B_t`$ which is a cylinder of radius $`R_f`$ . In addition, we assume that the spacetime metric has the form (2.3). Then, in the limit $`R_f\mathrm{}`$ we clearly reach the family of cylindrical waves that we want to study. Since all the constraints have been eliminated in the process of gauge fixing and the shift vector vanishes in Eq. (2.3), it is not difficult to conclude from the discussion in Ref. that the Hamiltonian of our reduced model comes exclusively from boundary terms on $`B_t`$, and is given by
$$H=\underset{R_f\mathrm{}}{lim}2N\sqrt{\sigma }(\kappa \kappa _0).$$
(2.4)
Here, we have made $`8G_3=1`$, $`\sigma `$ is the determinant of the two-metric induced on $`B_t`$, and $`\kappa `$ and $`\kappa _0`$ are the trace of the extrinsic curvature of this metric embedded, respectively, in $`\mathrm{\Sigma }_t`$ and in a three-dimensional Minkowski background. It is straightforward to check that $`\kappa =e^{\gamma _{\mathrm{}}/2}/(NR_f)`$ and $`\kappa _0=1/R_f`$, while $`\sigma =R_f^2`$. Therefore, we obtain that the reduced Hamiltonian that generates time evolution in the coordinate $`t`$ is $`H=2(1e^{\gamma _{\mathrm{}}/2})`$. In particular, this implies that $`\gamma _{\mathrm{}}`$ is a constant of motion, because it commutes with the Hamiltonian. So, given any classical solution, it is possible to absorb the factor $`e^\gamma _{\mathrm{}}`$ in the line element by a mere rescaling of the time coordinate: $`T=e^{\gamma _{\mathrm{}}/2}t`$ (off-shell, one would have $`T=_0^t𝑑\overline{t}e^{\gamma _{\mathrm{}}/2}`$).
Let us now particularize our considerations to the simpler case of linearly polarized cylindrical waves. For the Einstein-Rosen waves, we have $`v=0`$. We can impose this restriction as a symmetry condition in our model. Its compatibility with the Hamiltonian evolution leads to the secondary constraint $`p_v=0`$. One can also check that there are not tertiary constraints. The symmetry conditions $`v=p_v=0`$ form a pair of second-class constraints that allow the reduction of the model by removing two canonical degrees of freedom. The resulting system has the following metric and symplectic structure:
$`ds^2`$ $`=`$ $`e^\psi [e^\gamma (dT^2+dR^2)+R^2d\theta ^2]+e^\psi dZ^2,`$ (2.5)
$`\gamma `$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle _0^R}𝑑\overline{R}\overline{R}\left[(\psi ^{})^2+{\displaystyle \frac{p_\psi ^2}{\overline{R}^2}}\right],`$ (2.6)
$`\mathrm{\Omega }`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}𝑑R𝐝p_\psi 𝐝\psi .`$ (2.7)
Note that the term between square brackets in Eq.(2.5) is precisely the gauge-fixed metric of the dimensionally reduced Einstein-Maxwell model discussed in Ref. . In addition, the reduced Hamiltonian coincides also with that found by Ashtekar and Pierri in three-dimensions. Finally, it is straightforward to see that, in terms of the time coordinate $`T`$, the dynamical equations for the field $`\psi `$ are exactly those satisfied by a rotationally symmetric massless scalar field in three-dimensional Minkowski spacetime . This scalar field can be interpreted as the dual of a Maxwell field . In this way, one recovers the Einstein-Maxwell analog in three dimensions of the Einstein-Rosen waves.
## 3 Quantum Theory
Since the field $`\psi `$ is a rotationally symmetric solution to the massless Klein-Gordon equation in three dimensions that is regular at the origin $`R=0`$ , all classical solutions admit the mode expansion
$$\psi (R,T)=\frac{1}{\sqrt{2}}_0^{\mathrm{}}𝑑kJ_0(kR)\left[A(k)e^{ikT}+A^{}(k)e^{ikT}\right].$$
(3.1)
The constants of motion $`A(k)`$ and $`A^{}(k)`$ are complex conjugate to each other, because $`\psi `$ and $`J_0`$ (i.e., the zeroth-order Bessel function of the first kind) are real. Employing the identity $`2\pi J_0(kR)=𝑑\theta e^{ikR\mathrm{cos}\theta }`$, we can write the above expression in the alternative form
$$\psi (R,T)=\frac{1}{2\pi }_{IR^2}\frac{d^2k}{\sqrt{2}|\stackrel{}{k}|}\left[A(|\stackrel{}{k}|)e^{i(\stackrel{}{k}\stackrel{}{x}|\stackrel{}{k}|T)}+A^{}(|\stackrel{}{k}|)e^{i(\stackrel{}{k}\stackrel{}{x}|\stackrel{}{k}|T)}\right],$$
(3.2)
with $`R=|\stackrel{}{x}|`$. Taking then into account that, from the Hamiltonian equations of motion, $`p_\psi =R\dot{\psi }`$, where the overdot stands for the derivative with respect to $`T`$, substitution of Eq. (3.2) in the symplectic form leads to $`\mathrm{\Omega }=i_0^{\mathrm{}}𝐝A^{}(k)𝐝A(k)`$. Therefore, $`A(k)`$ and $`A^{}(k)`$ can be understood as annihilation and creation like variables. In addition, a trivial calculation using Eqs. (2.6) and (3.2) shows that $`\gamma _{\mathrm{}}`$ equals the Hamiltonian of the massless scalar field : $`\gamma _{\mathrm{}}=_0^{\mathrm{}}𝑑kkA^{}(k)A(k)`$.
Essentially, the quantization of our basic field $`\psi `$ can then be carried out by introducing a Fock space in which $`\psi (R,T)`$ goes over to an operator-valued distribution $`\widehat{\psi }(R,T)`$, obtained by representing $`A(k)`$ and $`A^{}(k)`$ as standard annihilation and creation operators . The Fock space in this representation is that over the Hilbert space of square integrable functions on the positive real axis, $`L^2(IR^+,dk)`$. Using such a representation, a complete quantization of the Einstein-Maxwell counterpart of our system has been recently proposed . Our aim in this section is to show how the quantization put forward by Ashtekar and Pierri in three dimensions can be employed to construct a consistent quantum theory which fully describes the four-dimensional metric of the Einstein-Rosen model.
As a first step towards the introduction of meaningful metric operators, let us regularize the basic field $`\widehat{\psi }(R,T)`$, which is defined only as an operator-valued distribution \[the reason being that $`J_0(kR)`$ does not belong to $`L^2(IR^+,dk)`$ for any $`R0`$\]. Given that $`J_0`$ is bounded in $`IR^+`$, the regularization can be achieved by simply multiplying the factor $`J_0(kR)`$ in the quantum version of Eq. (3.1) by a square integrable real function, $`gL^2(IR^+,dk)`$,
$$\widehat{\psi }(R,T|g)=\frac{1}{\sqrt{2}}_0^{\mathrm{}}𝑑kJ_0(kR)g(k)\left[\widehat{A}(k)e^{ikT}+\widehat{A}^{}(k)e^{ikT}\right].$$
(3.3)
This regularization can be justified from a physical point of view, e.g., by admitting the existence of a cut-off $`k_c`$ in momentum space . The corresponding function $`g(k)`$ equals then the unity on the compact interval $`[0,k_c]`$ and vanishes outside. In this sense, it is worth pointing out that the model itself provides an energy scale, namely, $`c^4/G_3`$ (adopting a general system of units). Thus, a natural candidate for $`k_c`$ could be $`c^3/(\mathrm{}G_3)`$, which has dimensions of an inverse length.
As an alternative motivation for the regularization, one can just smear the operator-valued distribution $`\widehat{\psi }(\stackrel{}{x},T)\widehat{\psi }(R=|\stackrel{}{x}|,T)`$, defined via the quantum analog of Eq. (3.2), with a test function in two dimensions $`f(\stackrel{}{x})`$ that is also rotationally symmetric. We assume that $`f(\stackrel{}{x})`$ belongs to the Schwartz space $`𝒮(IR^2)`$ of smooth functions on the plane with rapid decrease at infinity. In order to interpret the smearing as an average, we further accept that $`f(\stackrel{}{x})`$ is real and has a unit integral over $`IR^2`$. Then, a simple calculation proves that the smeared operator $`d^2x_0f(\stackrel{}{x}_0)\widehat{\psi }(\stackrel{}{x}\stackrel{}{x}_0)`$ is rotationally symmetric and can be expressed in the form (3.3), with $`g(k)`$ given by
$$2\pi \stackrel{~}{f}(k)=_{IR^2}d^2xf(\stackrel{}{x})e^{i\stackrel{}{k}\stackrel{}{x}}=2\pi _0^{\mathrm{}}𝑑RRJ_0(kR)f(R).$$
(3.4)
Here, $`\stackrel{~}{f}(\stackrel{}{k})`$ denotes the Fourier transform of $`f(\stackrel{}{x})`$ in two dimensions, and the notation $`\stackrel{~}{f}(k)`$ and $`f(R)`$ indicates that these functions depend only on $`k=|\stackrel{}{k}|`$ and $`R=|\stackrel{}{x}|`$, respectively. The last identity in the above equation shows that $`\stackrel{~}{f}(k)`$ is real; hence, so is $`g(k)`$. The first identity, together with the properties of the Fourier transform and the fact that $`f(\stackrel{}{x})`$ belongs to the Schwartz space, implies that $`\stackrel{~}{f}(\stackrel{}{k})𝒮(IR^2)`$. Then, we have that $`g(k)`$ belongs to the Hilbert space $`L^2(IR^+,dk)`$. In addition, since $`f(\stackrel{}{x})`$ has unit integral, it follows that $`g(0)=2\pi \stackrel{~}{f}(0)=1`$.
For any choice of the real function $`gL^2(IR^+,dk)`$, the operator (3.3), with domain given by the dense subspace of the Fock space consisting of all finite particle vectors , is symmetric and admits a self-adjoint extension , which we will denote again with the symbol $`\widehat{\psi }(R,T|g)`$. The standard spectral theorems ensure then that the exponential operators $`e^{\pm \widehat{\psi }(R,T|g)}`$ are well-defined and positive. Besides, recalling the definition of normal ordering and the Campbell-Baker-Hausdorff (CBH) formula $`e^{\widehat{a}}e^{\widehat{b}}=e^{[\widehat{a},\widehat{b}]/2}e^{\widehat{a}+\widehat{b}}`$, which is valid for operators whose commutator is a $`c`$-number , we conclude
$$:e^{\pm \widehat{\psi }(R,T|g)}:=e^{g_R^2/2}e^{\pm \widehat{\psi }(R,T|g)},$$
(3.5)
where $`||.||`$ denotes the norm in $`L^2(IR^+,dk)`$ and
$$g_R(k)=\frac{1}{\sqrt{2}}J_0(kR)g(k).$$
(3.6)
The diagonal $`\theta `$ and $`Z`$ components of the four-metric can then be represented by the positive operators
$$\widehat{h}_{\theta \theta }(R,T|g)=R^2:e^{\widehat{\psi }(R,T|g)}:,\widehat{h}_{ZZ}(R,T|g)=:e^{\widehat{\psi }(R,T|g)}:$$
(3.7)
Note that the normal ordering in these definitions guarantees that the vacuum expectation values reproduce the classical values of $`h_{\theta \theta }`$ and $`h_{ZZ}`$ in Minkowski spacetime.
On the other hand, the representation of the metric function (2.6) by a regularized operator $`:\widehat{\gamma }(f_R,T):`$ was discussed in Refs. . The symbol $`f_R`$ denotes a smearing function employed in the regularization, namely, $`f_R(r)`$ is a function on the positive real axis that equals the unity for all $`rR`$, decreases smoothly in $`[R,R+ϵ]`$ and vanishes for $`rR+ϵ`$, with $`ϵ>0`$ being a certain parameter with dimensions of length . It has been recently shown that this regularized operator has a well-defined action on a dense subspace of the Fock space which is contained in the set of finite particle vectors. In that domain of definition, the operator $`:\widehat{\gamma }(f_R,T):`$ is symmetric . As a straightforward consequence, so is $`:\widehat{\gamma }(f_R,T):\widehat{\psi }(R,T|g)`$ provided that the real function $`g`$ belongs to $`L^2(IR^+,dk)`$. In addition, Varadarajan has argued that $`:\widehat{\gamma }(f_R,T):`$ admits a self-adjoint extension, because it is (formally) possible to find a conjugation that leaves invariant the domain of definition of this operator and commutes with it . In fact, the same argument supports the existence of a self-adjoint extension of $`:\widehat{\gamma }(f_R,T):\widehat{\psi }(R,T|g)`$, because it is easy to check that the considered conjugation commutes as well with $`\widehat{\psi }(R,T|g)`$ when $`g`$ is real. Using the spectral theorem, we would then conclude that the exponential of this self-adjoint extension,
$$\widehat{\mathrm{\Gamma }}(R,T|g,f_R)e^{:\widehat{\gamma }(f_R,T):\widehat{\psi }(R,T|g)},$$
(3.8)
is a well-defined, positive operator.
We can then represent the remaining non-trivial components of the four-metric (i.e., the diagonal $`R`$ and $`T`$ components) by the operator
$$\widehat{h}_{RR}(R,T|g,f_R)=e^{\overline{g}_R^2}\widehat{\mathrm{\Gamma }}(R,T|g_1,f_R),$$
(3.9)
where we have adopted the notation
$$\overline{g}_R(k)=\frac{\sqrt{e^k1k}}{e^k1}g_R(k),g_1(k)=\frac{k}{e^k1}g(k),$$
(3.10)
and used definition (3.6). It is readily seen that the functions $`\overline{g}_R`$ and $`g_1`$ belong to $`L^2(IR^+,dk)`$ if so does the function $`g`$. According to our discussion above, the introduced operator should then be positive if the function $`g`$ is real and square integrable on the positive axis. In our definition (3.9), the numerical factor $`e^{\overline{g}_R^2}`$, as well as the replacement of $`g`$ with $`g_1`$ as the regularization function used in $`\widehat{\mathrm{\Gamma }}`$, can be understood as a convenient choice of factor ordering. Indeed, after restoring the dimensional constants $`c`$, $`\mathrm{}`$, and $`G_3`$ in our calculations, it is possible to check that, when $`\mathrm{}0`$, the factor $`e^{\overline{g}_R^2}`$ tends to the unity, whereas $`g_1g`$. The selected factor ordering is motivated by the following considerations.
In the limit $`R\mathrm{}`$, the smearing function $`f_R`$ tends to the unit function, and the operator $`:\widehat{\gamma }(f_R,T):`$ becomes
$$:\widehat{\gamma }_{\mathrm{}}:=_0^{\mathrm{}}𝑑kk\widehat{A}^{}(k)\widehat{A}(k),$$
(3.11)
which is the normal ordered Hamiltonian of a rotationally symmetric, massless scalar field in three dimensions . We then obtain that, in the asymptotic region $`R\mathrm{}`$, the purely radial component of the quantum metric is given by $`lim_{\overline{R}\mathrm{}}\widehat{h}_{RR}(\overline{R},T|g,1)`$. On the other hand, it is shown in Appendix A that
$$\widehat{h}_{RR}(\overline{R},T|g,1)=e^{_0^{\mathrm{}}𝑑kg_{\overline{R}}(k)e^{ikT}\widehat{A}^{}(k)}e^{:\widehat{\gamma }_{\mathrm{}}:}e^{_0^{\mathrm{}}𝑑kg_{\overline{R}}(k)e^{ikT}\widehat{A}(k)}.$$
(3.12)
Therefore, our factor ordering ensures that, at least in the asymptotic region, the vacuum expectation value of the metric operator (3.9) coincides with the classical value of $`h_{RR}`$ in Minkowski spacetime, a value which is equal to the unity. In addition, the factor ordering adopted is also very convenient from a practical point of view, because, for polynomials of the operator (3.12), all matrix elements between coherent states of the basic field $`\psi `$ are explicitly computable. For such coherent states, one can then complete the calculation of the asymptotic fluctuations in $`\widehat{h}_{RR}`$. Moreover, the operator $`e^{:\widehat{\gamma }_{\mathrm{}}:}`$ that appears in Eq. (3.12) is precisely the operator employed by Ashtekar and Pierri to represent the purely radial component of the metric in the three-dimensional counterpart of our system . As we will see in the next section, this fact leads to a simple relation between the coherent expectation values and fluctuations obtained for the radial component of the asymptotic metric in the four and three-dimensional models.
## 4 Metric fluctuations
We are now in an adequate position to study the quantum geometry of the model and discuss whether the conclusions obtained by Ashtekar in three dimensions about the existence of large quantum gravity effects in the asymptotic region generalize to the four-dimensional model describing Einstein-Rosen waves. Like in the analysis of Ref. , we will only consider quantum states that are coherent in the basic field $`\psi `$. These states show the most classical behavior that is allowed for the fundamental field of the theory . As such, they are natural candidates in the search for states that admit an approximate classical description of the geometry.
Given any complex function $`CL^2(IR^+,dk)`$, there exists an associated coherent state $`|C`$ of unit norm, which has the form
$$|C=e^{C^2/2}e^{_0^{\mathrm{}}𝑑kC(k)\widehat{A}^{}(k)}|0.$$
(4.1)
Here, $`|0`$ is the unique vacuum of the Fock space. For any coherent state, the expectation value of the (regularized) field $`\widehat{\psi }(R,T|g)`$ coincides, at all values of $`R`$ and $`T`$, with the classical field solution obtained by replacing the annihilation and creation operators with the functions $`C(k)`$ and its complex conjugate:
$$\widehat{\psi }(R,T|g)_C=2_0^{\mathrm{}}𝑑kg_R(k)\mathrm{Re}[C(k)e^{ikT}],$$
(4.2)
with $`\mathrm{Re}[.]`$ denoting the real part. In addition, using definitions (3.7) and the CBH formula, one can check that the coherent expectation values of the diagonal $`\theta `$ and $`Z`$ components of the metric are also equal to the corresponding classical expressions:
$$\widehat{h}_{\theta \theta }(R,T|g)_C=R^2\left(\widehat{h}_{ZZ}(R,T|g)_C\right)^1=R^2e^{\widehat{\psi }(R,T|g)_C}.$$
(4.3)
The calculation of the expectation value of the purely radial component of the metric is much more involved, and we will only analyze the asymptotic case $`R\mathrm{}`$. According to our discussion at the end of Sec. 3, this asymptotic expectation value is equal to the limit of $`\widehat{h}_{RR}(\overline{R},T|g,1)`$ when $`\overline{R}\mathrm{}`$. Employing Eq. (3.12) and the operator identities (A.5), one arrives at
$$\widehat{h}_{RR}(\overline{R},T|g,1)_C=e^{:\widehat{\gamma }_{\mathrm{}}:}_Ce^{\widehat{\psi }(\overline{R},T|g)_C}.$$
(4.4)
Here, $`e^{:\widehat{\gamma }_{\mathrm{}}:}_C`$ is precisely the coherent expectation value obtained in three dimensions for the diagonal $`R`$ component of the asymptotic metric :
$$e^{:\widehat{\gamma }_{\mathrm{}}:}_C=e^{_0^{\mathrm{}}𝑑k(e^k1)|C(k)|^2}.$$
(4.5)
Notice that Eq. (4.4) can be understood as the quantum counterpart of the classical relation $`h_{RR}=e^{\gamma \psi }`$ when $`\gamma `$ is set equal to its asymptotic value. In fact, this non-trivial result is due to the factor ordering adopted in Eq. (3.9). Furthermore, assuming that there exist strictly positive constants $`k_1`$ and $`\alpha `$ such that the function $`g(k)k^{1/2\alpha }`$ is bounded in the interval $`[0,k_1]`$, we prove in Appendix B that the limit $`R\mathrm{}`$ of the right-hand side of Eq. (4.2) vanishes. Therefore, the asymptotic expectation value of $`\widehat{h}_{RR}`$ in a coherent state turns out to coincide then with $`e^{:\widehat{\gamma }_{\mathrm{}}:}_C`$. Once again, this coincidence can be interpreted as the analog of the classical identity $`h_{RR}(R=\mathrm{})=e^\gamma _{\mathrm{}}`$, which incorporates the boundary condition that $`\psi `$ vanish at infinity. Taking into account that the expectation value $`e^{:\widehat{\gamma }_{\mathrm{}}:}_C`$ equals the classical value of $`e^\gamma _{\mathrm{}}`$ (at least) if the wave profile $`C(k)`$ has negligible high-energy contributions , we conclude that all coherent states in the low-energy sector would admit an approximate classical description of the four-dimensional geometry in the asymptotic region provided that they have small relative fluctuations in the metric when $`R\mathrm{}`$.
Before continuing our analysis, let us briefly comment on the assumption introduced above about the real function $`gL^2(IR^+,dk)`$ employed in the regularization. The existence of a bound in an interval starting at the origin is clearly satisfied for the function $`g`$ itself (i.e., with $`\alpha =1/2`$) if $`g(k)`$ is a cut-off in momentum space; in that case, $`g(k)1`$ on the positive axis. In addition, if the adopted regularization can be interpreted as a smooth spatial smearing, the function $`g(k)`$ is bounded again on the whole semiaxis $`k0`$, because $`g(\stackrel{}{k})g(k=|\stackrel{}{k}|)`$ given by Eq. (3.4) is a Schwartz test function in $`IR^2`$. These facts strongly support our hypothesis and show its compatibility with a wide class of feasible regularizations.
As a first step in the calculation of the metric fluctuations in the asymptotic region, one can check that
$`\mathrm{\Xi }_C\widehat{h}_{\theta \theta }(\overline{R},T|g)`$ $`=`$ $`\mathrm{\Xi }_C\widehat{h}_{ZZ}(\overline{R},T|g)=e^{g_{\overline{R}}^2}1,`$ (4.6)
$`\mathrm{\Xi }_C\widehat{h}_{RR}(\overline{R},T|g,1)`$ $`=`$ $`e^{\stackrel{˘}{C}^2}e^{\widehat{\psi }(\overline{R},T|g)_{\stackrel{˘}{C}}}e^{g_{\overline{R}}^2}1,`$ (4.7)
where $`\stackrel{˘}{C}(k)=C(k)(e^k1)`$ and $`\mathrm{\Xi }_C\widehat{a}=(\widehat{a}^2_C/\widehat{a}_C^2)1`$ is the square of the relative uncertainty in the operator $`\widehat{a}`$ for the coherent state $`|C`$. It is worth noticing that, when $`g=0`$, Eq. (4.7) reproduces the asymptotic fluctuations in the radial component of the three-metric studied by Ashtekar . In order to deduce the value of the asymptotic fluctuations, one only needs to take the limit $`\overline{R}\mathrm{}`$ in the above expressions. Actually, with our assumption about the existence of a segment to the right of the origin where the function $`g(k)k^{1/2\alpha }`$ is bounded for some choice of $`\alpha >0`$, it is shown in Appendix B that the asymptotic limits of $`g_{\overline{R}}`$ and $`\widehat{\psi }(\overline{R},T|g)_{\stackrel{˘}{C}}`$ vanish. Hence, for any of the coherent states, all metric operators display a classical behavior in the asymptotic region, except the operator that describes the purely radial component. Moreover, the square of the relative uncertainty in this last operator is given by the quantity $`e^{\stackrel{˘}{C}^2}1`$, which is precisely the value of the corresponding uncertainty in the three-dimensional Einstein-Maxwell analog of our cylindrical system .
As a straightforward consequence, it turns out that all the results reached by Ashtekar in three dimensions about the appearance of large quantum gravity effects apply as well to the four-dimensional model constructed here for the Einstein-Rosen waves. Indeed, the conclusions reached by Ashtekar are not only qualitatively valid from a four-dimensional point of view, but also quantitatively accurate. The only existing difference is that, as far as the four-metric is concerned, one does not need to demand that the relative fluctuations in the basic field $`\psi `$ (and in the physical quantities associated with it, like, e.g., the Hamiltonian) be negligible. Then, it is not necessary that the coherent states contain a large number of elementary excitations, a condition that is imposed in the three-dimensional system . From this perspective, there exist more coherent states that admit a classical description of the four-metric in the asymptotic region than those that provide a meaningful semiclassical solution to the Einstein-Maxwell model obtained by dimensional reduction.
Summarizing, in order for the classical approximation to be acceptable in the asymptotic region only two conditions must be verified :
$$_0^{\mathrm{}}𝑑k|C(k)|^2(e^k1k)1,_0^{\mathrm{}}𝑑k|C(k)|^2(e^k1)^21.$$
(4.8)
The first condition ensures that the coherent expectation value of $`e^{:\widehat{\gamma }_{\mathrm{}}:}`$ coincides with the classical value of $`e^\gamma _{\mathrm{}}`$. The second condition guarantees that the asymptotic fluctuations in the radial component of the metric are sufficiently small. In fact, the latter of these inequalities turns out to imply the former. In particular, for a wave profile $`C(k)`$ peaked around a certain wave number $`k_0`$ and with expected number of “particles” equal to $`N=𝑑k|C(k)|^2`$ , the above conditions reduce to $`N(e^{k_0}1)^21`$.
Finally, let us notice that Eq. (4.6) determines the metric fluctuations in the $`\theta `$ and $`Z`$ components at all points of the spacetime, and not just in the asymptotic region. It is then possible to obtain a useful estimate of those fluctuations also on the symmetry axis $`R=0`$, at least for a physically reasonable class of regularization functions $`g`$. Taking into account the definition of $`g_R`$ given in Eq. (3.6) and that the Bessel function $`J_0`$ equals the unity at the origin, one can check that the square norm $`g_R^2`$ becomes equal to $`g^2/2`$ when one approaches the axis. Suppose then that we further demand that the real regularization function $`gL^2(IR^+,dk)`$ take on the constant unit value in an interval starting at $`k=0`$. This interval will have the generic form $`[0,k_c]`$, where $`k_c`$ is a positive but otherwise arbitrary parameter. Notice that, in this case, one can make $`k_1k_c`$ and $`\alpha =1/2`$, because the function $`g`$ is bounded in an interval containing $`[0,k_c]`$. More importantly, according to our discussion in Sec. 3, all cut-off functions satisfy our new condition, with the parameter $`k_c`$ being the cut-off introduced in momentum space. One can then interpret every function $`g`$ in the considered family as a kind of generalized cut-off. Besides, it is clear that the regularization can still be viewed as a smooth spatial smearing if, in addition, $`g(\stackrel{}{k})g(k=|\stackrel{}{k}|)`$ belongs to $`𝒮(IR^2)`$. For this class of regularization functions, one readily obtains that $`g^2k_c`$, so that, on the axis,
$$\mathrm{\Xi }_C\widehat{h}_{\theta \theta }=\mathrm{\Xi }_C\widehat{h}_{ZZ}e^{k_c/2}1.$$
(4.9)
The relative uncertainties in the diagonal $`\theta `$ and $`Z`$ components of the metric will thus become relevant on the symmetry axis unless $`k_c1`$. However, one would expect that, in our model, a physically reasonable (generalized) cut-off parameter $`k_c`$ should be at least of the order of the inverse of the natural length scale provided by the system, i.e., $`c^3/(\mathrm{}G_3)k_P`$ (in a general system of units). With our conventions, $`c=\mathrm{}=8G_3=1`$, and thus $`k_P=8`$. But, for $`k_ck_P=8`$, we get from Eq. (4.9) that $`\mathrm{\Xi }_C\widehat{h}_{\theta \theta }=\mathrm{\Xi }_C\widehat{h}_{ZZ}>50`$. So, quantum gravity effects are huge on the symmetry axis for all of the considered regularizations and, therefore, also in the limit in which the cut-off is removed. In particular, this fact seems to indicate that the requirement of regularity on the axis of rotational symmetry is meaningless from a quantum mechanical point of view.
## 5 Conclusions
We have constructed a complete quantum theory that describes the metric of the family of Einstein-Rosen waves. This theory is based on the quantization carried out in Ref. for the Einstein-Maxwell model obtained by the dimensional reduction of linearly polarized cylindrical gravity.
We have started with the Hamiltonian formulation of general relativity for spacetimes that admit two commuting spacelike Killing vectors. Introducing suitable gauge-fixing conditions adapted to cylindrical symmetry, we have been able to remove all the gravitational constraints. In this way, we have arrived at a reduced model for the most general family of cylindrical waves in vacuum gravity. We have also calculated the symplectic structure induced from general relativity and the Hamiltonian that generates the time evolution. This Hamiltonian has been computed by reducing the gravitational Einstein-Hilbert action supplemented with appropriate surface terms. Such terms include the contribution of the timelike boundary located at $`R\mathrm{}`$ (where $`R`$ is the radial coordinate), and have been normalized to vanish in Minkowski spacetime.
We have then imposed the requirement of linear polarization as a symmetry condition. This has led to a reduced midisuperspace model whose classical solutions are precisely the Einstein-Rosen waves. The model has only one degree of freedom in configuration space, given by a cylindrically symmetric field $`\psi `$, and is indeed classically equivalent to a rotationally symmetric, massless scalar field (dual to a Maxwell field) coupled to three-dimensional gravity. The non-zero components of the four-metric in our reduced model are exponentials of the basic field $`\psi `$ multiplied either by trivial functions or by the purely radial component of the three-metric in the Einstein-Maxwell system. Employing the quantum theory proposed in Ref. for this three-dimensional model, we have then achieved a full quantization of the metric for Einstein-Rosen waves. Since, using a Fock space representation in which the field $`\psi `$ is represented as an operator-valued distribution, Ashtekar and Pierri had already succeeded in constructing a (presumably ) positive operator for the diagonal radial component of the metric in three dimensions, our quantization process has been reduced, basically, to the following two steps. Firstly, we have regularized the field $`\psi `$ to reach a well-defined operator and avoid ultraviolet divergences. Secondly, owing to the non-linearity of the metric in $`\psi `$, we have introduced a reasonable choice of factor ordering for the metric operators.
We have also analyzed whether there exist large quantum gravity effects in the system, as happens to be the case in the Einstein-Maxwell counterpart of the model. We have shown that, with the chosen factor ordering, the expectation values of the diagonal $`\theta `$ and $`Z`$ components of the four-metric correspond in fact to classical trajectories in all of the coherent states of the field $`\psi `$. In addition, we have seen that, like in the three-dimensional model, the asymptotic expectation value of the radial component is that predicted by the classical theory (semiclassical theory in three dimensions), provided that the coherent state has negligible contributions from the high-energy sector. In the derivation of this result, we have introduced the very weak assumption that the function $`gL^2(IR^+,dk)`$, employed in the regularization of the field $`\psi `$, is bounded in a certain interval around $`k=0`$ when multiplied by a factor of the form $`k^{1/2\alpha }`$, with $`\alpha `$ being a positive constant. Such an assumption is satisfied by all cut-off regularizations, as well as by those regularizations that can be interpreted as a smooth spatial smearing of the field, and therefore implies no restriction in physically relevant situations.
For such regularizations we have also computed the value of the quantum uncertainties in the metric when one approaches the asymptotic region. The fluctuations in the $`\theta `$ and $`Z`$ components turn out to vanish when $`R\mathrm{}`$. Therefore, these metric operators display a classical asymptotic behavior. As far as they are concerned, the boundary condition that the basic field vanish asymptotically is respected quantum mechanically in all coherent states. The asymptotic fluctuations in the diagonal radial component, on the other hand, are exactly the same as in the three-dimensional model . These results prove the validity of the analysis carried out by Ashtekar in three dimensions, not only qualitatively, but also quantitatively. There is only one caveat: in order for the four-dimensional metric to admit a classical description, it is not needed that the physical quantities associated with the field $`\psi `$ possess small relative uncertainties. The requirement that the number of fundamental excitations contained in the coherent state be large, necessary for a meaningful semiclassical approximation in the three-dimensional model , is no longer present. In this sense, the set of coherent states that are peaked around a classical four-metric in the asymptotic region is bigger than that corresponding to acceptable semiclassical solutions in three dimensions. In particular, the vacuum is contained only in the former of these sets. Of course, apart from the quantum metric, one could be interested in considering other operators related with the four-dimensional geometry (like those which describe the spacetime curvature). Demanding that such operators have negligible asymptotic fluctuations might well impose further conditions on the family of coherent states that display a classical behavior and, perhaps, restrict again their number of “particles”.
We have also analyzed the quantum fluctuations in the diagonal $`\theta `$ and $`Z`$ components of the metric when one approaches the symmetry axis. We have shown that, for all regularizations that do not modify the mode decomposition of the field $`\psi `$ up to wave numbers of the order of the natural scale $`k_P=c^3/(\mathrm{}G_3)`$, the relative uncertainties in the metric at $`R=0`$ are large. In particular, they explode in the limit in which the regularization is removed. One should then expect significant quantum effects on the axis. It is thus unclear up to what extent the condition of regularity of the four-geometry on the symmetry axis is sensible from a quantum mechanical perspective.
Although there exist other possible quantizations of our model, the quantum theory constructed presents clear advantages. In fact, it has been constructed in such a way that the relation between the metric operators in three and four dimensions are as simple and natural as possible. This fact has allowed us to compare the physical results for the Einstein-Rosen waves with those obtained by Ashtekar in the three-dimensional Einstein-Maxwell model, and prove that the latter are indeed relevant in four dimensions. Regarding the factor ordering, we have checked that the existence of large quantum fluctuations in the metric is rather insensitive to the operator ordering. However, it generally affects the expectation value of the metric in coherent states, so that, for factor orderings other than the one selected, such value would only reproduce a classical solution in the limit $`\mathrm{}0`$. Obviously, this is one of the reasons that motivated our choice of ordering.
On the other hand, it is worth noticing that our discussion about the expectation value of the metric and its uncertainty in the asymptotic region is in fact regularization independent, apart from the more than reasonable hypothesis that the regularization function (possibly multiplied by a factor $`k^{1/2\alpha }`$, with $`\alpha >0`$) be bounded in a neighborhood of the origin of wave numbers, an assumption that, as we have commented, involves no physical limitation in practice. Our analysis of the metric uncertainties on the symmetry axis, nevertheless, has been restricted to a particular (though quite general) family of regularizations, which can be interpreted as a generalized cut-off. In this sense, our results about the fluctuations on the axis depend on the regularization adopted. However, since those fluctuations are always significant when the regularization is removed at scales below the inverse-length parameter $`k_P`$, naturally provided by the system, one would expect the existence of important quantum gravity effects on the axis of cylindrical symmetry in all physically plausible situations. Finally, since coherence in the basic field $`\psi `$ is not a requisite for the validity of the classical approximation from a purely four-dimensional viewpoint, it would be interesting to investigate whether the quantum fluctuations in the four-metric can be diminished by considering other families of quantum states, like, e.g., those analyzed by Gambini and Pullin .
### Acknowledgments
M. E. A. acknowledges the financial support provided by C.S.I.C. during the completion of this work. G. A. M. M. was supported by DGESIC under the Research Projects No. PB97-1218 and HP1988-0040.
## Appendix A
In this appendix, we will prove relation (3.12). We will make use of the operator expansion theorem
$$e^{x\widehat{a}}\widehat{b}e^{x\widehat{a}}=\underset{n=0}{\overset{\mathrm{}}{}}\frac{x^n}{n!}[\widehat{a},\widehat{b}]_{(n)}$$
(A.1)
and the identity
$$e^{\widehat{a}}e^{\widehat{b}}e^{\widehat{a}}=\mathrm{exp}\left(e^{\widehat{a}}\widehat{b}e^{\widehat{a}}\right).$$
(A.2)
In these expressions, $`\widehat{a}`$ and $`\widehat{b}`$ denote two generic operators and $`[\widehat{a},.]_{(n)}`$ is the $`n`$-th application of the commutator with $`\widehat{a}`$. Particularizing these equations to the case in which $`\widehat{b}=:\widehat{\gamma }_{\mathrm{}}:`$, $`x=1`$, and
$$\widehat{a}=\frac{1}{\sqrt{2}}_0^{\mathrm{}}\frac{dk}{k}J_0(k\overline{R})g_1(k)\left[\widehat{A}^{}(k)e^{ikT}\widehat{A}(k)e^{ikT}\right]\widehat{D}(\overline{R},T|g),$$
(A.3)
we obtain
$$\widehat{h}_{RR}(\overline{R},T|g,1)=e^{\stackrel{ˇ}{g}_{\overline{R}}^2}e^{\widehat{D}(\overline{R},T|g)}e^{:\widehat{\gamma }_{\mathrm{}}:}e^{\widehat{D}(\overline{R},T|g)}.$$
(A.4)
Here, $`\stackrel{ˇ}{g}_R(k)=g_R(k)/\sqrt{e^k1}`$ and we have employed definitions (3.6) and (3.10). On the other hand, a repeated application of the operator expansion theorem to calculate the commutator of $`e^{:\widehat{\gamma }_{\mathrm{}}:}`$, firstly with the smeared version of the creation and annihilation operators, and then with their exponentials, leads to
$`e^{:\widehat{\gamma }_{\mathrm{}}:}e^{_0^{\mathrm{}}𝑑kf(k)\widehat{A}^{}(k)}`$ $`=`$ $`e^{_0^{\mathrm{}}𝑑kf(k)e^k\widehat{A}^{}(k)}e^{:\widehat{\gamma }_{\mathrm{}}:},`$
$`e^{_0^{\mathrm{}}𝑑kf(k)\widehat{A}(k)}e^{:\widehat{\gamma }_{\mathrm{}}:}`$ $`=`$ $`e^{:\widehat{\gamma }_{\mathrm{}}:}e^{_0^{\mathrm{}}𝑑kf(k)e^k\widehat{A}(k)}.`$ (A.5)
Using these relations, together with the CBH formula, one can readily check that the right-hand sides of Eqs. (3.12) and (A.4) coincide.
## Appendix B
We want to prove that the expectation value $`\widehat{\psi }(R,T|g)_C`$ and the norm $`g_R`$ vanish in the asymptotic limit $`R\mathrm{}`$ if the functions $`C`$ and $`g`$ belong to the Hilbert space $`L^2(IR^+,dk)`$ and, for some choice of positive constant $`\alpha `$, the function $`g(k)k^{1/2\alpha }`$ is bounded in an interval of the form $`[0,k_1]`$. Here, $`k_1`$ is a strictly positive number and $`g_R(k)=J_0(kR)g(k)/\sqrt{2}`$. In fact, we only need to show that $`g_R`$ vanishes in the asymptotic region, because, using Eq. (4.2), the triangle inequality for complex numbers, and the Schwarz inequality on $`L^2(IR^+,dk)`$, one gets
$$\left|\widehat{\psi }(R,T|g)_C\right|2||C||||g_R||.$$
(B.1)
Obviously, the same arguments apply to the value of $`\widehat{\psi }(R,T|g)_{\stackrel{˘}{C}}`$ appearing in Eq. (4.7). Let us then write
$$g_R^2=\frac{1}{2}_0^{k_1}𝑑kJ_0^2(kR)|g(k)|^2+\frac{1}{2}_{k_1}^{\mathrm{}}𝑑kJ_0^2(kR)|g(k)|^2.$$
(B.2)
The second term on the right-hand side vanishes when $`R\mathrm{}`$ because, in that limit, $`J_0^2(kR)`$ tends to zero uniformly in $`k[k_1,\mathrm{})`$, with $`k_1>0`$. As for the first term, let $`G`$ be the upper bound of $`|g(k)k^{1/2\alpha }|`$ in $`[0,k_1]`$. Then
$$\frac{1}{2}_0^{k_1}𝑑kJ_0^2(kR)|g(k)|^2\frac{G^2}{2R^{2\alpha }}_0^{k_1R}\frac{dk}{k^{12\alpha }}J_0^2(k).$$
(B.3)
Recalling that $`\alpha >0`$ and $`J_0(k)\mathrm{cos}(k\pi /4)\sqrt{2/(\pi k)}`$ (up to subdominant terms) for $`k1`$, one can finally show that the limit of the above expression when $`R\mathrm{}`$ is zero. |
warning/0002/hep-th0002069.html | ar5iv | text | # 1 Introduction
## 1 Introduction
Branes or the spatially extended solitons play the central role in the study of string theory or M-theory. They are known to form a variety of bound states, and the corresponding classical solutions should exist in the effective supergravity theories. However, a difficulty arises in constructing classical solutions corresponding to the bound states of intersecting branes. The solutions for two intersecting branes have been found only when one of the two branes are smeared along the directions longitudinal to the other. (There are a few exceptions such as and ) Recently a fully localized solutions have been found for M5-branes intersecting on three-branes and for M2-brane junctions. There has also been a perturbative analysis of Born-Infeld theory coupled to supergravity. In this paper we try to obtain another solutions for localized branes intersecting one another.
In this paper we shall work in the eleven-dimensional supergravity. It contains a three-form potential besides the graviton and the gravitino, and M2-branes and M5-branes are the electric and magnetic sources of the potential. The presence of a single BPS M2-brane or M5-brane preserves a half of supersymmetry. An M2-brane and an M5-brane are known to form a $`1/4`$-supersymmetric bound state when they intersect on a string. It is also believed that an M2-brane can end on an M5-brane. In this case the M2-brane is “open” and has a $`(1+1)`$-dimensional boundary on the M5-brane. Interestingly, such open M2-branes can be created between two M5-branes with $`(1+1)`$ commonly longitudinal directions when the two M5-branes pass through each other. Such creation of branes are familiar in the D-brane worldvolume gauge field theory, and a supergravity analysis has been made for a specific example.
In the following we shall find the generic BPS configuration corresponding to the bound states of M2-branes$`(012)`$, M5-branes$`(013456)`$ and M5-branes$`(01789\mathrm{})`$. The numbers specify the directions the branes are lying along, and we denoted the eleventh direction by $`\mathrm{}`$. The analysis of the BPS condition $`\delta _{\mathrm{SUSY}}\psi _m=0`$ is performed in section 2 for spherical symmetric configurations, and the result is fully generalized in section 3. The resultant expression depends on three arbitrary functions of $`x_{2,3,\mathrm{},9,\mathrm{}}`$. These functions have to satisfy the equations of motion given in section 4 in order to describe a bound state of localized(delta-functional) branes. They are not solved, but we show that if the M5-branes are localized, the boundaries of M2-branes stretching between them are automatically localized. We also propose in section 5 an idea for understanding the brane creation in supergravity.
## 2 Analysis of the BPS Condition
Here we analyze the BPS condition under the assumption of spherical symmetry. The solution is summarized in (23) and (24), and the following arguments help the reader to understand to what extent the result is general.
We parameterize 01-directions by $`t,\sigma `$, 2-direction by $`x_2`$, 3456-directions by $`x_1,\theta _\mathrm{I},\varphi _\mathrm{I},\psi _\mathrm{I}`$ and 789$`\mathrm{}`$-directions by $`x_3,\theta _{\mathrm{I}\mathrm{I}},\varphi _{\mathrm{I}\mathrm{I}},\psi _{\mathrm{I}\mathrm{I}}`$. We assume the following form for the metric:
$$ds^2=g_0^2(dt^2+d\sigma ^2+\varpi ^1\varpi ^1+\varpi ^2\varpi ^2+\varpi ^3\varpi ^3)+g_\mathrm{I}^2d\mathrm{\Omega }_\mathrm{I}^2+g_{\mathrm{I}\mathrm{I}}^2d\mathrm{\Omega }_{\mathrm{I}\mathrm{I}}^2$$
(1)
$$d\mathrm{\Omega }_i^2=d\theta _i^2+\mathrm{cos}^2\theta _id\varphi _i^2+\mathrm{sin}^2\theta _id\psi _i^2,i=\mathrm{I}\mathrm{or}\mathrm{I}\mathrm{I}$$
where $`\varpi ^{1,2,3}`$ are linear combinations of $`dx^{1,2,3}`$ with coordinate-dependent coefficients. We define the vielbein as follows:
$$\begin{array}{ccccccccccccccccc}\hfill e^{\underset{¯}{t}}& =& g_0dt\hfill & ;& \hfill e^{\underset{¯}{1}}& =& g_0\varpi ^1\hfill & ;& e^{\underset{¯}{\theta _\mathrm{I}}}\hfill & =& \hfill g_\mathrm{I}d\theta _\mathrm{I}& ;& e^{\underset{¯}{\theta _{\mathrm{I}\mathrm{I}}}}\hfill & =& \hfill g_{\mathrm{I}\mathrm{I}}d\theta _{\mathrm{I}\mathrm{I}}& & \\ \hfill e^{\underset{¯}{\sigma }}& =& g_0d\sigma \hfill & ;& \hfill e^{\underset{¯}{2}}& =& g_0\varpi ^2\hfill & ;& e^{\underset{¯}{\varphi _\mathrm{I}}}\hfill & =& \hfill g_\mathrm{I}\mathrm{cos}\theta _\mathrm{I}d\varphi _\mathrm{I}& ;& e^{\underset{¯}{\varphi _{\mathrm{I}\mathrm{I}}}}\hfill & =& \hfill g_{\mathrm{I}\mathrm{I}}\mathrm{cos}\theta _{\mathrm{I}\mathrm{I}}d\varphi _{\mathrm{I}\mathrm{I}}& & \\ & & & & \hfill e^{\underset{¯}{3}}& =& g_0\varpi ^3\hfill & ;& e^{\underset{¯}{\psi _\mathrm{I}}}\hfill & =& \hfill g_\mathrm{I}\mathrm{sin}\theta _\mathrm{I}d\psi _\mathrm{I}& ;& e^{\underset{¯}{\psi _{\mathrm{I}\mathrm{I}}}}\hfill & =& \hfill g_{\mathrm{I}\mathrm{I}}\mathrm{sin}\theta _{\mathrm{I}\mathrm{I}}d\psi _{\mathrm{I}\mathrm{I}}& & \end{array}$$
Hereafter we denote local Lorentz indices with underbar. Under the spherical symmetry, the most generic form for the four-form field strength is given by:
$`F_{(4)}`$ $`=`$ $`F_{(4)}^{[M2]}+F_{(4)}^{[M5]}+F_{(4)}^{[M5^{}]}`$ (2)
$`=`$ $`g_0^1\left({\displaystyle \frac{1}{2}}e^{\underset{¯}{t\sigma ij}}E_kϵ^{ijk}+e^{\underset{¯}{\theta _{\mathrm{I}\mathrm{I}}\varphi _{\mathrm{I}\mathrm{I}}\psi _{\mathrm{I}\mathrm{I}}i}}B_i+e^{\underset{¯}{\theta _\mathrm{I}\varphi _\mathrm{I}\psi _\mathrm{I}i}}H_i\right)`$
$`=`$ $`dtd\sigma {\displaystyle \frac{1}{2}}g_0^3ϵ^{ijk}E_i\varpi ^{jk}+d^3\mathrm{\Omega }_{\mathrm{I}\mathrm{I}}g_{\mathrm{I}\mathrm{I}}^3B_i\varpi ^i+d^3\mathrm{\Omega }_\mathrm{I}g_\mathrm{I}^3H_i\varpi ^i.`$
Here the overall factor $`g_0^1`$ in the second line is simply for later convenience.
To obtain supersymmetric solutions we have to focus on the supersymmetry transformation law of gravitino
$`e_{\underset{¯}{m}}^m\delta _ϵ\psi _m`$ $`=`$ $`e_{\underset{¯}{m}}^m𝒟_mϵ{\displaystyle \frac{1}{288}}(3\mathrm{\Gamma }^{\underset{¯}{pqrs}}\mathrm{\Gamma }_{\underset{¯}{m}}\mathrm{\Gamma }_{\underset{¯}{m}}\mathrm{\Gamma }^{\underset{¯}{pqrs}})ϵF_{\underset{¯}{pqrs}},`$ (3)
and find field configurations that admit $`\delta _ϵ\psi _m=0`$ for some nonzero $`ϵ`$. Here and throughout this section we use “mostly Hermitian” Gamma matrices satisfying
$$\{\mathrm{\Gamma }^{\underset{¯}{a}},\mathrm{\Gamma }^{\underset{¯}{b}}\}=2\eta ^{\underset{¯}{ab}}=2\mathrm{d}\mathrm{i}\mathrm{a}\mathrm{g}(++\mathrm{}+),\mathrm{\Gamma }^{\underset{¯}{a_1\mathrm{}a_{11}}}=ϵ^{a_1\mathrm{}a_{11}},ϵ^{t\sigma 123\theta _\mathrm{I}\varphi _\mathrm{I}\psi _\mathrm{I}\theta _{\mathrm{I}\mathrm{I}}\varphi _{\mathrm{I}\mathrm{I}}\psi _{\mathrm{I}\mathrm{I}}}=1.$$
Let us evaluate (3) term by term. To begin with, the covariant derivative of a spinor is defined by
$$e_{\underset{¯}{m}}^m𝒟_mϵ=e_{\underset{¯}{m}}^m(_m+\frac{1}{4}\mathrm{\Omega }_{m\underset{¯}{pq}}\mathrm{\Gamma }^{\underset{¯}{pq}})ϵ(_{\underset{¯}{m}}+\frac{1}{4}\mathrm{\Omega }_{\underset{¯}{mpq}}\mathrm{\Gamma }^{\underset{¯}{pq}})ϵ;_{\underset{¯}{m}}e_{\underset{¯}{m}}^m_m$$
(4)
where $`\mathrm{\Omega }_{\underset{¯}{pq}}=dx^m\mathrm{\Omega }_{m\underset{¯}{pq}}`$ is the spin connection. Under the assumption of spherical symmetry we can obtain some components of the spin connection from the torsion-free condition alone:
$$𝒟e^{\underset{¯}{p}}=de^{\underset{¯}{p}}+\mathrm{\Omega }_{\underset{¯}{q}}^{\underset{¯}{p}}e^{\underset{¯}{q}}=0.$$
(5)
The only nonzero components of $`\mathrm{\Omega }_{\underset{¯}{mnp}}`$ except for those with $`(m,n,p=1,2,3)`$ are
$`\mathrm{\Omega }_{\underset{¯}{tti}}=\mathrm{\Omega }_{\underset{¯}{\sigma \sigma i}}`$ $`=`$ $`_{\underset{¯}{i}}\mathrm{ln}g_0`$
$`\mathrm{\Omega }_{\underset{¯}{\theta _\mathrm{I}\theta _\mathrm{I}i}}=\mathrm{\Omega }_{\underset{¯}{\varphi _\mathrm{I}\varphi _\mathrm{I}i}}=\mathrm{\Omega }_{\underset{¯}{\psi _\mathrm{I}\psi _\mathrm{I}i}}`$ $`=`$ $`_{\underset{¯}{i}}\mathrm{ln}g_\mathrm{I}`$
$`\mathrm{\Omega }_{\underset{¯}{\theta _{\mathrm{I}\mathrm{I}}\theta _{\mathrm{I}\mathrm{I}}i}}=\mathrm{\Omega }_{\underset{¯}{\varphi _{\mathrm{I}\mathrm{I}}\varphi _{\mathrm{I}\mathrm{I}}i}}=\mathrm{\Omega }_{\underset{¯}{\psi _{\mathrm{I}\mathrm{I}}\psi _{\mathrm{I}\mathrm{I}}i}}`$ $`=`$ $`_{\underset{¯}{i}}\mathrm{ln}g_{\mathrm{I}\mathrm{I}}`$
$`\mathrm{\Omega }_{\underset{¯}{\varphi _\mathrm{I}\varphi _\mathrm{I}\theta _\mathrm{I}}}={\displaystyle \frac{1}{g_\mathrm{I}}}{\displaystyle \frac{\mathrm{sin}\theta _\mathrm{I}}{\mathrm{cos}\theta _\mathrm{I}}}`$ ; $`\mathrm{\Omega }_{\underset{¯}{\psi _\mathrm{I}\psi _\mathrm{I}\theta _\mathrm{I}}}={\displaystyle \frac{1}{g_\mathrm{I}}}{\displaystyle \frac{\mathrm{cos}\theta _\mathrm{I}}{\mathrm{sin}\theta _\mathrm{I}}}`$
$`\mathrm{\Omega }_{\underset{¯}{\varphi _{\mathrm{I}\mathrm{I}}\varphi _{\mathrm{I}\mathrm{I}}\theta _{\mathrm{I}\mathrm{I}}}}={\displaystyle \frac{1}{g_{\mathrm{I}\mathrm{I}}}}{\displaystyle \frac{\mathrm{sin}\theta _{\mathrm{I}\mathrm{I}}}{\mathrm{cos}\theta _{\mathrm{I}\mathrm{I}}}}`$ ; $`\mathrm{\Omega }_{\underset{¯}{\psi _{\mathrm{I}\mathrm{I}}\psi _{\mathrm{I}\mathrm{I}}\theta _{\mathrm{I}\mathrm{I}}}}={\displaystyle \frac{1}{g_{\mathrm{I}\mathrm{I}}}}{\displaystyle \frac{\mathrm{cos}\theta _{\mathrm{I}\mathrm{I}}}{\mathrm{sin}\theta _{\mathrm{I}\mathrm{I}}}}`$
Using these expressions we can rewrite the BPS condition $`\delta _ϵ\psi _m=0`$ in the following way:
$`2\mathrm{\Gamma }^{\underset{¯}{\theta _i}}_{\underset{¯}{\theta _i}}ϵ`$ $`=`$ $`2\mathrm{\Gamma }^{\underset{¯}{\varphi _i}}_{\underset{¯}{\varphi _i}}ϵ+\mathrm{\Gamma }^{\underset{¯}{\theta _i}}ϵ_{\underset{¯}{\theta _i}}\mathrm{ln}\mathrm{cos}\theta _i`$ (6)
$`=`$ $`2\mathrm{\Gamma }^{\underset{¯}{\psi _i}}_{\underset{¯}{\psi _i}}ϵ+\mathrm{\Gamma }^{\underset{¯}{\theta _i}}ϵ_{\underset{¯}{\theta _i}}\mathrm{ln}\mathrm{sin}\theta _i,i=\mathrm{I}\mathrm{or}\mathrm{I}\mathrm{I}`$
$`g_0(4_{\underset{¯}{i}}+\mathrm{\Omega }_{\underset{¯}{ijk}}\gamma ^{jk}2_{\underset{¯}{j}}\mathrm{ln}g_0\gamma ^i\gamma ^j)ϵ`$ $`=`$ $`2(E_i\gamma ^{123}+B_i\gamma ^0H_i\gamma ^5)ϵ`$
$`3g_0/\mathrm{ln}g_0ϵ`$ $`=`$ $`(2/E\gamma ^{123}/B\gamma ^0+/H\gamma ^5)ϵ`$
$`6g_0\mathrm{\Gamma }^{\underset{¯}{t\sigma \theta _\mathrm{I}}}_{\underset{¯}{\theta _\mathrm{I}}}ϵ+3g_0/\mathrm{ln}g_\mathrm{I}ϵ`$ $`=`$ $`(/E\gamma ^{123}/B\gamma ^02/H\gamma ^5)ϵ`$
$`6g_0\mathrm{\Gamma }^{\underset{¯}{t\sigma \theta _{\mathrm{I}\mathrm{I}}}}_{\underset{¯}{\theta _{\mathrm{I}\mathrm{I}}}}ϵ+3g_0/\mathrm{ln}g_{\mathrm{I}\mathrm{I}}ϵ`$ $`=`$ $`(/E\gamma ^{123}+\mathrm{\hspace{0.17em}2}/B\gamma ^0+/H\gamma ^5)ϵ.`$ (7)
Here we have introduced a set of new gamma-matrices
$$(\gamma ^0,\gamma ^1,\gamma ^2,\gamma ^3)=(\mathrm{\Gamma }^{\underset{¯}{\theta _{\mathrm{I}\mathrm{I}}\varphi _{\mathrm{I}\mathrm{I}}\psi _{\mathrm{I}\mathrm{I}}}},\mathrm{\Gamma }^{\underset{¯}{t\sigma 1}},\mathrm{\Gamma }^{\underset{¯}{t\sigma 2}},\mathrm{\Gamma }^{\underset{¯}{t\sigma 3}}),\gamma ^5=\gamma ^0\gamma ^1\gamma ^2\gamma ^3=\mathrm{\Gamma }^{\underset{¯}{\theta _\mathrm{I}\varphi _\mathrm{I}\psi _\mathrm{I}}}$$
(8)
and used some short-hand notations
$$/\gamma ^1_{\underset{¯}{1}}+\gamma ^2_{\underset{¯}{2}}+\gamma ^3_{\underset{¯}{3}},/E\gamma ^1E_1+\gamma ^2E_2+\gamma ^3E_3,\mathrm{etc}.$$
(9)
$`\gamma ^0`$ and $`\gamma ^5`$ are anti-hermitian while $`\gamma ^{1,2,3}`$ are hermitian. We also assumed in the above that $`ϵ`$ does not depend on $`t`$ and $`\sigma `$.
We would like to obtain the bosonic field configuration that satisfies the BPS condition $`\delta _ϵ\psi _m=0`$ with the following $`ϵ`$:
$$ϵ=f(x_i)U_\mathrm{I}(\theta _\mathrm{I},\varphi _\mathrm{I},\psi _\mathrm{I})U_{\mathrm{I}\mathrm{I}}(\theta _{\mathrm{I}\mathrm{I}},\varphi _{\mathrm{I}\mathrm{I}},\psi _{\mathrm{I}\mathrm{I}})ϵ_0,$$
(10)
where $`f(x_i)`$ is a scale factor, $`U_{\mathrm{I},\mathrm{I}\mathrm{I}}`$ are two mutually commuting local Lorentz transformations and $`ϵ_0`$ is a constant spinor. We also assume that the residual supersymmetry is characterized by
$$\mathrm{\Gamma }^{\underset{¯}{t\sigma 2}}ϵ=\mathrm{\Gamma }^{\underset{¯}{t\sigma 3\theta _{\mathrm{I}\mathrm{I}}\varphi _{\mathrm{I}\mathrm{I}}\psi _{\mathrm{I}\mathrm{I}}}}ϵ=\mathrm{\Gamma }^{\underset{¯}{t\sigma 1\theta _\mathrm{I}\varphi _\mathrm{I}\psi _\mathrm{I}}}ϵ=ϵ\text{or}\gamma ^2ϵ=\gamma ^{30}ϵ=\gamma ^{51}ϵ=ϵ.$$
(11)
We can rather easily find the solution of the angular equation (6) satisfying also (11). It is given by the following $`U_\mathrm{I}`$ and $`U_{\mathrm{I}\mathrm{I}}`$:
$`U_\mathrm{I}`$ $`=`$ $`\mathrm{exp}({\displaystyle \frac{\theta _\mathrm{I}\mathrm{\Gamma }^{\underset{¯}{\theta _\mathrm{I}1}}}{2}})\mathrm{exp}({\displaystyle \frac{\varphi _\mathrm{I}\mathrm{\Gamma }^{\underset{¯}{\varphi _\mathrm{I}1}}}{2}})\mathrm{exp}({\displaystyle \frac{\psi _\mathrm{I}\mathrm{\Gamma }^{\underset{¯}{\psi _\mathrm{I}\theta _\mathrm{I}}}}{2}})`$
$`U_{\mathrm{I}\mathrm{I}}`$ $`=`$ $`\mathrm{exp}({\displaystyle \frac{\theta _{\mathrm{I}\mathrm{I}}\mathrm{\Gamma }^{\underset{¯}{\theta _{\mathrm{I}\mathrm{I}}3}}}{2}})\mathrm{exp}({\displaystyle \frac{\varphi _{\mathrm{I}\mathrm{I}}\mathrm{\Gamma }^{\underset{¯}{\varphi _{\mathrm{I}\mathrm{I}}3}}}{2}})\mathrm{exp}({\displaystyle \frac{\psi _{\mathrm{I}\mathrm{I}}\mathrm{\Gamma }^{\underset{¯}{\psi _{\mathrm{I}\mathrm{I}}\theta _{\mathrm{I}\mathrm{I}}}}}{2}})`$ (12)
These local Lorentz transformations are understood as relating the polar frame (where $`\mathrm{\Gamma }^{\underset{¯}{t},\underset{¯}{\sigma },\underset{¯}{1},\mathrm{},\underset{¯}{\psi _{\mathrm{I}\mathrm{I}}}}`$ are coordinate independent) to the orthonormal frame (where $`\mathrm{\Gamma }^{\underset{¯}{0},\underset{¯}{1},\mathrm{},\underset{¯}{\mathrm{}}}`$ are coordinate independent). The remaining equations (7) determine the dependence on the coordinates $`x_{1,2,3}`$. Using (10) and (11) we can rewrite them as follows:
$$2_{\underset{¯}{i}}ϵ=_{\underset{¯}{i}}\mathrm{ln}g_0ϵ$$
(13)
$$g_0(\frac{1}{2}\mathrm{\Omega }_{\underset{¯}{ipq}}\gamma ^{pq}_{\underset{¯}{p}}\mathrm{ln}g_0\gamma ^{ip})=E_i\gamma ^{31}B_i\gamma ^{23}+H_i\gamma ^{12}$$
(14)
$$E_2=B_1=H_3$$
(15)
$$\begin{array}{ccccccc}\hfill 3g_0/\mathrm{ln}g_0& =& \gamma ^1(2E_3+H_2)\hfill & +& \gamma ^2(B_3H_1)\hfill & +& \gamma ^3(B_2+2E_1)\hfill \\ \hfill 3g_0g_\mathrm{I}^1\gamma ^1+3g_0/\mathrm{ln}g_\mathrm{I}& =& \gamma ^1(E_32H_2)\hfill & +& \gamma ^2(B_3+2H_1)\hfill & +& \gamma ^3(B_2E_1)\hfill \\ \hfill 3g_0g_{\mathrm{I}\mathrm{I}}^1\gamma ^3+3g_0/\mathrm{ln}g_{\mathrm{I}\mathrm{I}}& =& \gamma ^1(E_3+H_2)\hfill & +& \gamma ^2(\mathrm{\hspace{0.33em}2}B_3H_1)\hfill & +& \gamma ^3(2B_2E_1)\hfill \end{array}$$
(16)
The first equation (13) determines the scale factor of $`ϵ`$ as follows:
$$ϵ=g_0^{1/2}U_\mathrm{I}U_{\mathrm{I}\mathrm{I}}ϵ_0.$$
(17)
The next equations (14) and (15) relate the components of the spin connection and the gauge field strength:
$$\begin{array}{ccccccccccc}\hfill \mathrm{\Omega }_{\underset{¯}{112}}& =& g_0^1H_1+_{\underset{¯}{2}}\mathrm{ln}g_0\hfill & ;& \hfill \mathrm{\Omega }_{\underset{¯}{212}}& =& g_0^1H_2_{\underset{¯}{1}}\mathrm{ln}g_0\hfill & ;& \hfill \mathrm{\Omega }_{\underset{¯}{312}}& =& g_0^1H_3\hfill \\ \hfill \mathrm{\Omega }_{\underset{¯}{123}}& =& g_0^1B_1\hfill & ;& \hfill \mathrm{\Omega }_{\underset{¯}{223}}& =& g_0^1B_2+_{\underset{¯}{3}}\mathrm{ln}g_0\hfill & ;& \hfill \mathrm{\Omega }_{\underset{¯}{323}}& =& g_0^1B_3_{\underset{¯}{2}}\mathrm{ln}g_0\hfill \\ \hfill \mathrm{\Omega }_{\underset{¯}{131}}& =& g_0^1E_1_{\underset{¯}{3}}\mathrm{ln}g_0\hfill & ;& \hfill \mathrm{\Omega }_{\underset{¯}{231}}& =& g_0^1E_2\hfill & ;& \hfill \mathrm{\Omega }_{\underset{¯}{331}}& =& g_0^1E_3+_{\underset{¯}{1}}\mathrm{ln}g_0\hfill \end{array}$$
Since the torsion-free condition (5) relates the spin connection to the vielbein, the above relations allow us to express the components of the gauge field strength in terms of the vielbein. The most elegant way is:
$`d\varpi ^1`$ $`=`$ $`H_1\varpi ^{21}+E_1\varpi ^{13}`$
$`d\varpi ^2`$ $`=`$ $`H_2\varpi ^{21}+2E_2\varpi ^{13}+B_2\varpi ^{23}`$
$`d\varpi ^3`$ $`=`$ $`B_3\varpi ^{23}+E_3\varpi ^{13}`$ (18)
The most convenient choice of coordinates that is compatible with the above would be
$`\varpi ^1`$ $`=`$ $`h_1dx^1`$
$`\varpi ^2`$ $`=`$ $`h_2(dx^2+A_1dx^1+A_3dx^3)`$
$`\varpi ^3`$ $`=`$ $`h_3dx^3`$ (19)
The next equations (16) enable us to express $`g_{0,\mathrm{I},\mathrm{I}\mathrm{I}}`$ in terms of the components of $`\varpi ^i`$. A careful analysis of them with the help of (18) and (19) yields that $`\frac{g_\mathrm{I}}{h_1g_0}`$ depends only on $`x_1`$, and similarly $`\frac{g_{\mathrm{I}\mathrm{I}}}{h_3g_0}`$ depends only on $`x_3`$. Using the diffeomorphism degrees of freedom we can therefore set
$$g_\mathrm{I}=x_1h_1g_0,g_{\mathrm{I}\mathrm{I}}=x_3h_3g_0.$$
Using the above relation we then find $`g_0g_\mathrm{I}g_{\mathrm{I}\mathrm{I}}=x_1x_3`$. Hence
$$g_0=(h_1h_3)^{\frac{1}{3}},g_\mathrm{I}=x_1h_1^{\frac{2}{3}}h_3^{\frac{1}{3}},g_{\mathrm{I}\mathrm{I}}=x_3h_1^{\frac{1}{3}}h_3^{\frac{2}{3}}$$
(20)
The equations (16) also yield the following equations
$$_1\left(\frac{h_2h_3}{h_1}\right)=_2\left(\frac{A_1h_2h_3}{h_1}\right),_3\left(\frac{h_2h_1}{h_3}\right)=_2\left(\frac{A_3h_2h_1}{h_3}\right).$$
(21)
Hence we put
$$h_2^2=_2X_2Y,h_1^2=\frac{_2Y}{_2Z},h_3^2=\frac{_2X}{_2Z},A_1=\frac{_1X}{_2X},A_3=\frac{_3Y}{_2Y}.$$
(22)
Thus the most generic solution of the BPS condition with spherical symmetry is summarized as follows:
$`ds^2`$ $`=`$ $`\left(_2X_2Y_2Z\right)^{1/3}[_2Z(dt^2+d\sigma ^2)+_2Y(dx_1^2+x_1^2d\mathrm{\Omega }_\mathrm{I}^2)+_2X(dx_3^2+x_3^2d\mathrm{\Omega }_{\mathrm{I}\mathrm{I}}^2)`$ (23)
$`+_2X_2Y_2Z(dx_2+{\displaystyle \frac{_1X}{_2X}}dx_1+{\displaystyle \frac{_3Y}{_2Y}}dx_3)^2],`$
$`F_{(4)}`$ $`=`$ $`dtd\sigma {\displaystyle \frac{1}{2}}d\left[dx_1D_1Z+dx_3D_3Z\right]`$ (24)
$`+d^3\mathrm{\Omega }_{\mathrm{I}\mathrm{I}}{\displaystyle \frac{1}{2}}\left[d(x_3^3D_3X)+x_3^3dx_3\left({\displaystyle \frac{_2}{_2Y}}{\displaystyle \frac{_2X}{_2Z}}+x_3^3D_3x_3^3D_3X\right)\right]`$
$`+d^3\mathrm{\Omega }_\mathrm{I}{\displaystyle \frac{1}{2}}\left[d(x_1^3D_1Y)+x_1^3dx_1\left({\displaystyle \frac{_2}{_2X}}{\displaystyle \frac{_2Y}{_2Z}}+x_1^3D_1x_1^3D_1Y\right)\right].`$
Here $`D_1`$ and $`D_3`$ are the “coordinate covariant derivatives” defined as follows:
$$D_1_1\frac{_1X}{_2X}_2_1|_{X\mathrm{fixed}},D_3_3\frac{_3Y}{_2Y}_2_3|_{Y\mathrm{fixed}}.$$
These are obviously invariant under the change of the coordinate $`x_2x_2^{}=f(x_1,x_2,x_3)`$. This is a residual diffeomorphism symmetry, and owing to this symmetry we may parameterize the 2-direction by any of $`X,Y,Z`$.
## 3 Generalization
From the previous result (23), (24) we can guess the expression for more general solutions without spherical symmetry. It is expressed by three arbitrary functions $`X,Y,Z`$ of $`x_{2,3,\mathrm{},\mathrm{}}`$ as follows:
$`ds^2`$ $`=`$ $`\left(_2X_2Y_2Z\right)^{1/3}[_2Z(dt^2+d\sigma ^2)+_2Ydx_i^2+_2Xdx_p^2`$ (25)
$`+_2X_2Y_2Z(dx_2+{\displaystyle \frac{_iX}{_2X}}dx_i+{\displaystyle \frac{_pY}{_2Y}}dx_p)^2]`$
$`2F_{(4)}`$ $`=`$ $`dtd\sigma d\left[dx_iD_iZ+dx_pD_pZ\right]`$ (26)
$`+\left[{\displaystyle \frac{1}{6}}ϵ^{pqrs}d(D_pX)dx_qdx_rdx_sdx_7dx_8dx_9dx_{10}\left({\displaystyle \frac{_2}{_2Y}}{\displaystyle \frac{_2X}{_2Z}}+D_pD_pX\right)\right]`$
$`+\left[{\displaystyle \frac{1}{6}}ϵ^{ijkl}d(D_iY)dx_jdx_kdx_ldx_3dx_4dx_5dx_6\left({\displaystyle \frac{_2}{_2X}}{\displaystyle \frac{_2Y}{_2Z}}+D_iD_iY\right)\right]`$
$$i,j,k,l=(3,4,5,6),p,q,r,s=(7,8,9,10).$$
The coordinate covariant derivatives are defined as follows:
$$D_i_i\frac{_iX}{_2X}_2_i|_{X\mathrm{fixed}},D_p_p\frac{_pY}{_2Y}_2_p|_{Y\mathrm{fixed}}.$$
We can safely say that the above expression is the most generic BPS configuration, because it is the unique generalization of the most generic spherical symmetric configuration obtained in the previous section. We would like to note here again that one of $`X,Y,Z`$ is a residual diffeomorphism degree of freedom.
## 4 Equation of Motion
The equation of motion in the absence of the source is given by
$$dF_{(4)}=dF_{(7)}F_{(4)}F_{(4)}=0$$
(27)
$$\frac{1}{4}\left[R_{mn}\frac{1}{2}g_{mn}R\right]=\frac{1}{12}\left[F_{mpqr}F_n^{pqr}\frac{1}{8}g_{mn}F_{pqrs}F^{pqrs}\right]$$
(28)
A careful analysis of these equations shows that, under the assumption of the BPS condition some of the above equations turn out equivalent. The result is that the solution of (27) automatically satisfies (28). Therefore we concentrate on (27) in the following.
In the presence of the source the equation of motion is modified as
$$dF_{(4)}=j_5,dF_{(7)}F_{(4)}F_{(4)}=j_8.$$
(29)
BPS condition now relates the components of the stress tensor to the components of $`j_5`$ and $`j_8`$. Our generic BPS configuration (25), (26) has the following currents:
$`2j_5`$ $`=`$ $`d^4x_{789\mathrm{}}df^{[M5]}d^4x_{3456}df^{[M5^{}]}`$ (30)
$`2j_8`$ $`=`$ $`d^8x_{3456789\mathrm{}}f^{[M2]}`$ (31)
$`d^4x_{789\mathrm{}}Dx_2_2Y{\displaystyle \frac{1}{6}}ϵ^{ijkl}D_if^{[M5]}d^3x_{jkl}`$
$`d^4x_{3456}Dx_2_2X{\displaystyle \frac{1}{6}}ϵ^{pqrs}D_pf^{[M5^{}]}d^3x_{qrs}`$
$`dtd\sigma Dx_2d^4x_{3456}_2Zdf^{[M5^{}]}dtd\sigma Dx_2d^4x_{789\mathrm{}}_2Zdf^{[M5]}`$
$$Dx_2dx_2+dx_i\frac{_iX}{_2X}+dx_p\frac{_pY}{_2Y},d^nx_{j_1\mathrm{}j_n}dx_{j_1}\mathrm{}dx_{j_n},$$
where the three functions $`f^{[M2,M5,M5^{}]}`$ are defined as follows:
$`f^{[M5]}`$ $`=`$ $`{\displaystyle \frac{_2}{_2Y}}{\displaystyle \frac{_2X}{_2Z}}+D_pD_pX`$
$`f^{[M5^{}]}`$ $`=`$ $`{\displaystyle \frac{_2}{_2X}}{\displaystyle \frac{_2Y}{_2Z}}+D_iD_iY`$ (32)
$`f^{[M2]}`$ $`=`$ $`D_iD_i{\displaystyle \frac{_2X}{_2Z}}+D_pD_p{\displaystyle \frac{_2Y}{_2Z}}2D_iD_pXD_pD_iY+{\displaystyle \frac{_2}{_2Y}}{\displaystyle \frac{_2X}{_2Z}}{\displaystyle \frac{_2}{_2X}}{\displaystyle \frac{_2Y}{_2Z}}`$
These encode the position of the sources. The source-free equations of motion are hence given by $`f^{[M2]}=f^{[M5]}=f^{[M5^{}]}=0`$.
If both M5-branes and M5’-branes are present, they possibly bend each other. However, bending of branes is a notion that depends on the choice of coordinates. We may say that there is no bending effects if we can find in a natural way a coordinate frame in which both M5 and M5’-branes are flat. But the following consideration leads us to conclude that this is not the case.
Looking at the expressions for currents carefully, one finds that the fourth and the fifth terms in $`j_8`$ of (31) correspond to M2-brane charges with Euclidean worldvolume. Hence it is reasonable to require them to vanish even in the presence of the source. We therefore impose the following condition:
$$D_if^{[M5]}D_pf^{[M5^{}]}0.$$
(33)
This is equivalent to requiring that $`f^{[M5]}`$ is a function of $`(x_p,X)`$ and $`f^{[M5^{}]}`$ is a function of $`(x_i,Y)`$. Under the above condition the currents take the following simple form:
$`2j_5`$ $`=`$ $`d^4x_{789\mathrm{}}Dx_2_2f^{[M5]}d^4x_{3456}Dx_2_2f^{[M5^{}]}`$
$`2j_8`$ $`=`$ $`d^8x_{3456789\mathrm{}}f^{[M2]}.`$
The classical solution for some isolated M5 and M5’-branes is thus obtained by solving
$`f^{[M5]}={\displaystyle \frac{_2}{_2Y}}{\displaystyle \frac{_2X}{_2Z}}+D_pD_pX`$ $`=`$ $`{\displaystyle \underset{j}{}}Q_j\delta ^4(x_pa_p^{(j)})\theta (Xa_2^{(j)})`$
$`f^{[M5^{}]}={\displaystyle \frac{_2}{_2X}}{\displaystyle \frac{_2Y}{_2Z}}+D_iD_iY`$ $`=`$ $`{\displaystyle \underset{j}{}}Q_j^{}\delta ^4(x_ib_i^{(j)})\theta (Yb_2^{(j)}).`$ (34)
The solution corresponds to the system of M5-branes of charge $`Q_j`$ at $`(X,x_p)=(a_2^{(j)},a_p^{(j)})`$ and M5’-branes of charge $`Q_j^{}`$ at $`(Y,x_i)=(b_2^{(j)},b_i^{(j)})`$. We find that M5-branes are flat in $`x_2=X`$ frame while M5’-branes are flat in $`x_2=Y`$ frame. Hence we conclude that the M5-branes and M5’-branes in general bend each other.
Choosing one of $`X,Y,Z`$ as the $`x_2`$-coordinate we can regard (34) as two equations for two unknown functions. They are nonlinear and highly complicated equations, $`(X,Y,Z)`$ appearing as coordinates as well as functions. Moreover the solution of (34) must not be unique because there is a freedom to put an arbitrary number of M2-branes. At present the generic solution for them is not known. It is known, however, that under the assumption
$$_2X=H_5(x_p),_2Y=H_5^{}(x_i),(_2Z)^1=H_2(x_i,x_p).$$
the equations of motion are reduced to the following linear differential equations:
$$_p_pH_5=_i_iH_5^{}=(H_5_i_i+H_5^{}_p_p)H_2=0.$$
This type of equations has been analyzed in in different contexts. The above equations describe the system of M2-branes together with some M5 and M5’-branes smeared along the $`x_2`$-direction. Since all the fields are $`x_2`$-independent the solutions cannot represent M2-branes ending on M5-branes.
The third equation $`f^{[M2]}=0`$ remains to be analyzed. In analyzing this, recall that one of $`X,Y,Z`$ is the gauge degree of freedom. Therefore if $`f^{[M2]}=f^{[M5]}=f^{[M5^{}]}=0`$ were three independent equations, the system would be over-determined. This is not the case. The point is that the $`x_2`$-derivative of $`f^{[M2]}`$ is zero where $`f^{[M5]}=f^{[M5^{}]}=0`$. Indeed, using (33) we find
$$_2f^{[M2]}=_2f^{[M5]}\left(f^{[M5^{}]}D_iD_iY\right)+_2f^{[M5^{}]}\left(f^{[M5]}D_pD_pX\right).$$
(35)
Since $`_2f^{[M2]}`$ represents the boundaries of M2-branes, the above equality means that M2-branes can have boundaries only on M5-branes.
## 5 Brane Creation
We would like to give an idea for how the brane creation can be seen in supergravity. Let us consider the system of an M5-brane and an M5’-brane. Then the functions $`f^{[M5]}`$ and $`f^{[M5^{}]}`$ have support on semi-infinite six-planes that are bounded by M5 and M5’-branes, respectively. Assume that one of the two six-planes is on the left of the M5-brane, and the other is on the right of the M5’-brane, as depicted in the Figure 1.
Note that one can change whether a six-plane appears on the left or on the right of an M5-brane by the shift $`XX+f(x_p)`$ or $`YY+f(x_i)`$. Then, according to the relative position of two M5-branes the two six-planes may or may not have an intersection. Since the component $`f^{[M2]}`$ of the M2-brane current $`j_8`$ satisfies
$`_2f^{[M2]}`$ $`=`$ $`_2\left(f^{[M5]}f^{[M5^{}]}\right)_2f^{[M5]}D_iD_iY_2f^{[M5^{}]}D_pD_pX`$
$`\text{or}f^{[M2]}=f^{[M5]}f^{[M5^{}]}+\mathrm{},`$
there is an M2-brane precisely on the intersection of two six-planes, and its charge is proportional to the product of the charges of the two M5-branes. This explains the brane creation in supergravity, namely, when an M5-brane pass through an M5’-brane, an M2-brane is created between them.
It is expected that all the other types of M2-branes, namely those with semi-infinite worldvolume or those stretching between M5-M5 or M5’-M5’ are described by the second and third terms in (5).
We give here a simple example. Let us solve the equation of motion under the assumption that $`X,Y,Z`$ depend only on $`x_2`$. The solution representing the system of an M5-brane at $`Z=a`$ and an M5’-brane at $`Z=b`$ is obtained by solving
$$\frac{_Z^2X}{_ZY}=\theta (Za),\frac{_Z^2Y}{_ZX}=\theta (bZ).$$
(37)
Assuming $`a<b`$, the solution is given in terms of a function $`f(Z)`$ satisfying $`f^{\prime \prime }(Z)=f(Z)`$ as follows:
$$\begin{array}{c}(Za)\hfill \\ \{\begin{array}{c}_ZX=f(a)\hfill \\ _ZY=f^{\prime \prime }(a)(Za)+f^{}(a)\hfill \end{array}\hfill \end{array}\begin{array}{c}(aZb)\hfill \\ \{\begin{array}{c}_ZX=f(Z)\hfill \\ _ZY=f^{}(Z)\hfill \end{array}\hfill \end{array}\begin{array}{c}(bZ)\hfill \\ \{\begin{array}{c}_ZX=f^{}(b)(Zb)+f(b)\hfill \\ _ZY=f^{}(b)\hfill \end{array}\hfill \end{array}$$
Then $`f^{[M2]}`$ takes the following form as expected:
$$f^{[M2]}=\frac{_Z^2X}{_ZY}\frac{_Z^2Y}{_ZX}=\theta (Za)\theta (bZ).$$
(38)
This represents the M2-brane stretching between the two M5-branes, completely de-localized in the $`x_{3,4,\mathrm{},9,\mathrm{}}`$-directions. If the right-hand sides of the equations (37) are shifted by constants, the solutions will contain some M2-branes with semi-infinite worldvolume. It is straightforward to find such solutions.
## 6 Conclusion
In this article we have found the most generic BPS configuration for M5-branes(013456), M5’-branes(01789$`\mathrm{}`$) and M2-branes(012). We have also given and studied the equation of motion for localized sources. The equations are highly nonlinear, and it seems very difficult to obtain the generic solution. In fact, it is not clear whether or not the solution for localized M5 and M5’-branes indeed exists. But the analysis of the equations of motion themselves has lead to some interesting results.
By focusing on a specific term in the M2-brane current we have given an explanation for the brane creation in supergravity. Strictly speaking, however, this is no more than a conjecture because we have no justification for picking up a specific term in the current. Constructing a solution for the equations (34) will help us in great deal in understanding the mechanism of brane creation in supergravity and checking if the above conjecture indeed holds.
Acknowledgment
The author thanks J. Hashiba for collaboration at the early stage of this work. The author is also thankful to T. Eguchi, Y. Sugawara and S. Terashima for discussions and comments. The work of the author was supported in part by JSPS Research Fellowships for Young Scientists. |
warning/0002/astro-ph0002395.html | ar5iv | text | # Properties of spherical galaxies and clusters with an NFW density profile
## 1 Introduction
A universal profile of dark matter haloes was introduced as a result of high-resolution $`N`$-body simulations performed by Navarro, Frenk & White (1995, 1996, 1997, hereafter NFW) for power-law as well as CDM initial power spectra of density fluctuations. NFW found that in a large range of masses the density profiles of dark haloes can be fitted with a simple formula with only one fitting parameter. The density profile steepens from $`r^1`$ near the centre of the halo to $`r^3`$ at large distances. The NFW profile has been confirmed in cosmological simulations by Cole & Lacey (1996), Tormen, Bouchet & White (1997), Huss, Jain & Steinmetz (1999a), Jing (2000), Bullock et al. (1999), while Huss, Jain & Steinmetz (1999b) have shown that the NFW profile also arises from non-cosmological initial conditions. It is worthwhile noting that some (but not all) recent very high resolution cosmological simulations produce steeper density profiles, with inner slopes $`1.5`$ (Fukushige & Makino 1997, Moore et al. 1998, Ghigna et al. 1999, see also Jing & Suto 2000). The density profiles in the cosmological simulations also display considerable scatter (Avila-Reese et al. 1999, Bullock et al. 1999), and Avila-Reese et al. find that the outer slopes of galaxy size haloes are steeper than the NFW slope of $`3`$ when selected within clusters ($`4`$) and slightly shallower within groups ($`2.7`$). Although the exact properties of dark matter haloes are still under debate, the NFW profile is presently considered to provide the reference frame for any further numerical research on density profiles of dark haloes.
Simple cosmological derivations of the density profiles of bound objects are difficult, essentially because one needs to work in the non-linear regime of the growth of gravitational instabilities. Nevertheless, using the spherical top-hat model of Gunn & Gott (1972), density profiles typically varying as $`r^{9/4}`$ were derived by Gott (1975), Gunn (1977), Fillmore & Goldreich (1984) and Bertschinger (1985). Hoffman & Shaham (1985) applied the spherical infall model to the hierarchical clustering scenario and predicted that the density profiles of haloes should depend on $`\mathrm{\Omega }`$ as well as the initial power spectrum of density fluctuations. However, for $`\mathrm{\Omega }=1`$ they obtained power-law profiles in contradiction with the steepening slopes found in the current $`N`$-body simulations described above. In a recent study, Łokas (2000) has improved the model of Hoffman & Shaham (1985) by a generalization of the initial density distribution, the introduction of a cut-off in this distribution at half the inter-peak separation and by a proper calculation of the collapse factor. The improved model reproduces the changing slope of the density profile and its dependence on halo mass and the type of cosmological power spectrum found by NFW. The NFW profile is also reproduced in studies taking into account the merging mechanism (see Lacey & Cole 1993) in the halo formation scenario (e.g. Salvador-Solé, Solanes & Manrique 1998, Avila-Reese, Firmani & Hernandez 1998). Therefore the numerical and analytical considerations seem to converge on the statement that the density profiles of dark matter haloes are indeed well described by the universal formula proposed by NFW.
The ultimate test of both the analytical and numerical results must come from the observations of density profiles of galaxies and galaxy clusters. Three recent studies of clusters (Carlberg et al. 1997, Adami et al. 1998, van der Marel et al. 2000) claim good agreement between cluster observations and the NFW mass density profile. But for galaxies, the situation is less satisfying. Flores & Primack (1994) show that the NFW profile is incompatible with the rotation curves of spiral galaxies, while Kravtsov et al. (1998) estimate that the inner slope of the density profile of dwarf irregular and LSB galaxies is $`0.3`$ instead of $`1`$. However, these conclusions were obtained with a number of assumptions and approximations concerning the very unclear issues of biasing, non-sphericity of objects and so on. Besides, as pointed out by van den Bosch et al. (2000), Swaters, Madore & Trewhella (2000) and van den Bosch & Swaters (2000), the observed rotation curves of these galaxies are too uncertain to discriminate between cores and cusps.
The main motivation for this research is to explore analytically the physical properties of objects with NFW density profiles. The aim is to check whether these properties are acceptable from the physical point of view and thus to test the validity of density profiles obtained in cosmological $`N`$-body simulations. Additionally, this paper presents formulae for observable quantities that can be used for comparisons between the theoretical predictions (such as the NFW profile) and observations.
The paper is organized as follows: after a short presentation of the universal formula for the density profile proposed by NFW, in Section 2 we describe physical properties of spherical systems following from this density profile. Section 3 is devoted to a simple comparison between the projected NFW density profile and the surface brightness of elliptical galaxies. A more thorough comparison is beyond the scope of the present paper and will be given elsewhere (Mamon & Łokas, in preparation). The discussion follows in Section 4.
## 2 Properties of the NFW model
### 2.1 Basic properties
NFW established that the density profiles of dark matter haloes in high resolution cosmological simulations for a wide range of masses and for different initial power spectra of density fluctuations are well fitted by the formula
$$\frac{\rho (r)}{\rho _c^0}=\frac{\delta _{\mathrm{char}}}{(r/r_\mathrm{s})(1+r/r_\mathrm{s})^2}$$
(1)
with a single fitting parameter $`\delta _{\mathrm{char}}`$, the characteristic density. The so-called scale radius $`r_\mathrm{s}`$ is defined by
$$r_\mathrm{s}=\frac{r_v}{c},$$
(2)
where $`r_v`$ is the virial radius usually defined as the distance from the centre of the halo within which the mean density is $`v`$ times the present critical density, $`\rho _c^0`$. The value of the virial overdensity $`v`$ is often assumed to be $`v=178`$, a number predicted by the simplest version of the spherical model for $`\mathrm{\Omega }=1`$. For other cosmological models it can be lower by a factor of 2 or more (Lacey & Cole 1993, Eke, Cole & Frenk 1996). However, according to the improved spherical infall model (Łokas 2000) $`v`$ can be as low as 30 even for $`\mathrm{\Omega }=1`$. In the following, $`v`$ is kept as a free parameter.
The quantity $`c`$ introduced in equation (2) is the concentration parameter, which is related to the characteristic density by
$$\delta _{\mathrm{char}}=\frac{vc^3g(c)}{3},$$
(3)
where
$$g(c)=\frac{1}{\mathrm{ln}(1+c)c/(1+c)}.$$
(4)
The concentration parameter will be used hereafter as the only parameter describing the shape of density profile. From cosmological $`N`$-body simulations (Navarro et al. 1997, Jing 2000, Bullock et al. 1999, Jing & Suto 2000), extended Press-Schechter theory (Navarro et al. 1997, Salvador-Solé, Solanes & Manrique 1998), and the spherical infall model (Łokas 2000), we know that $`c`$ depends on the mass of object and the form of the initial power spectrum of density fluctuations. For all initial power spectra, the observed trend is for lower concentration parameter in higher mass objects, with $`4<c<22`$ in cosmological simulations with CDM initial power spectra and $`c`$ up to 90 for the less realistic scale-free power spectra. More precisely, in the $`\mathrm{\Lambda }`$CDM cosmology, $`c=5`$ corresponds to the masses of clusters of galaxies, while $`c=10`$ corresponds to the masses of bright galaxies.
It is convenient to express the distance from the centre of the object in units of the virial radius $`r_v`$:
$$s=\frac{r}{r_v}$$
(5)
and the density profile of equation (1) then becomes
$$\frac{\rho (s)}{\rho _c^0}=\frac{vc^2g(c)}{3s(1+cs)^2}.$$
(6)
The mass of the halo is usually defined as the mass within the virial radius:
$$M_v=\frac{4}{3}\pi r_v^3v\rho _c^0.$$
(7)
The distribution of mass in units of the virial mass follows from equation (6):
$$\frac{M(s)}{M_v}=g(c)\left[\mathrm{ln}(1+cs)\frac{cs}{1+cs}\right]$$
(8)
and we see that it diverges at large $`s`$, which is a disadvantage of the model from a physical point of view.
The gravitational potential associated with the density distribution (6) is
$$\frac{\mathrm{\Phi }(s)}{V_v^2}=g(c)\frac{\mathrm{ln}(1+cs)}{s},$$
(9)
where $`V_v`$ is the circular velocity at $`r=r_v`$:
$$V_v^2=V^2(r_v)=\frac{GM_v}{r_v}=\frac{4}{3}\pi Gr_v^2v\rho _c^0.$$
(10)
Hence, from equation (9) we see that the gravitational potential at the centre, $`\mathrm{\Phi }(0)=cg(c)V_v^2`$, is finite.
Equations (8) and (10) lead to a circular velocity that obeys
$$\frac{V^2(s)}{V_v^2}=\frac{g(c)}{s}\left[\mathrm{ln}(1+cs)\frac{cs}{1+cs}\right].$$
(11)
Equations (8), (9) and (11) were first derived by Cole & Lacey (1996).
The radial velocity dispersion $`\sigma _\mathrm{r}(r)`$ can be obtained by solving the Jeans equation
$$\frac{1}{\rho }\frac{\mathrm{d}}{\mathrm{d}r}(\rho \sigma _\mathrm{r}^2)+2\beta \frac{\sigma _\mathrm{r}^2}{r}=\frac{\mathrm{d}\mathrm{\Phi }}{\mathrm{d}r},$$
(12)
where $`\beta =1\sigma _\theta ^2(r)/\sigma _\mathrm{r}^2(r)`$ is a measure of the anisotropy in the velocity distribution. In the simplest case of isotropic orbits, $`\sigma _\theta (r)=\sigma _\mathrm{r}(r)`$ and $`\beta =0`$. This value of $`\beta `$ is also close to the results of $`N`$-body simulations: Cole & Lacey (1996) and Thomas et al. (1998) show that, in a variety of cosmological models, the ratio $`\sigma _\theta /\sigma _\mathrm{r}`$ is not far from unity and decreases slowly with distance from the centre to reach $`0.8`$ at the virial radius. However, Huss, Jain & Steinmetz (1999a) find $`\sigma _\theta /\sigma _\mathrm{r}0.6`$ at $`r_v`$.
First we consider the case of $`\beta `$=const. Then the solution to the equation (12) with the condition of $`\sigma _\mathrm{r}0`$ at $`s\mathrm{}`$ is
$`{\displaystyle \frac{\sigma _\mathrm{r}^2}{V_v^2}}(s,\beta =\mathrm{const})`$ $`=`$ $`g(c)(1+cs)^2s^{12\beta }`$ (13)
$`\times `$ $`{\displaystyle _s^{\mathrm{}}}\left[{\displaystyle \frac{s^{2\beta 3}\mathrm{ln}(1+cs)}{(1+cs)^2}}{\displaystyle \frac{cs^{2\beta 2}}{(1+cs)^3}}\right]ds.`$
For $`\beta =0`$, 0.5 and 1, reasonably simple analytical solutions to this equation can be found:
$`{\displaystyle \frac{\sigma _\mathrm{r}^2}{V_v^2}}(s,\beta =0)`$ $`=`$ $`{\displaystyle \frac{1}{2}}c^2g(c)s(1+cs)^2[\pi ^2\mathrm{ln}(cs){\displaystyle \frac{1}{cs}}`$ (14)
$``$ $`{\displaystyle \frac{1}{(1+cs)^2}}{\displaystyle \frac{6}{1+cs}}+\left(1+{\displaystyle \frac{1}{c^2s^2}}{\displaystyle \frac{4}{cs}}{\displaystyle \frac{2}{1+cs}}\right)`$
$`\times `$ $`\mathrm{ln}(1+cs)+3\mathrm{ln}^2(1+cs)+6\mathrm{Li}_2(cs)],`$
$`{\displaystyle \frac{\sigma _\mathrm{r}^2}{V_v^2}}(s,\beta =0.5)`$ $`=`$ $`cg(c)(1+cs)^2[{\displaystyle \frac{\pi ^2}{3}}+{\displaystyle \frac{1}{2(1+cs)^2}}`$ (15)
$`+`$ $`{\displaystyle \frac{2}{1+cs}}+{\displaystyle \frac{\mathrm{ln}(1+cs)}{cs}}+{\displaystyle \frac{\mathrm{ln}(1+cs)}{1+cs}}`$
$``$ $`\mathrm{ln}^2(1+cs)2\mathrm{L}\mathrm{i}_2(cs)],`$
$`{\displaystyle \frac{\sigma _\mathrm{r}^2}{V_v^2}}(s,\beta =1)`$ $`=`$ $`g(c)(1+cs)^2{\displaystyle \frac{1}{s}}[{\displaystyle \frac{\pi ^2}{6}}{\displaystyle \frac{1}{2(1+cs)^2}}`$ (16)
$``$ $`{\displaystyle \frac{1}{1+cs}}{\displaystyle \frac{\mathrm{ln}(1+cs)}{1+cs}}+{\displaystyle \frac{\mathrm{ln}^2(1+cs)}{2}}+\mathrm{Li}_2(cs)].`$
In the above expressions $`\mathrm{Li}_2(x)`$ is the dilogarithm, a special function which can be conveniently dealt with using *Mathematica* packages. Otherwise, it can be approximated by
$$\mathrm{Li}_2(x)=_x^0\frac{\mathrm{ln}(1t)\mathrm{d}t}{t}x\left[1+10^{0.5}(x)^{0.62/0.7}\right]^{0.7}.$$
(17)
The fit is accurate to better than 1.5% in the range $`100<x<0`$.
We included the predictions for $`\beta =1`$ just as a limiting case. In fact such a model with purely radial orbits and NFW density profile is not physical since its distribution function is not everywhere non-negative. As pointed out by e.g. Richstone & Tremaine (1984, see also Łokas & Hoffman 2000), such velocity anisotropy requires the inner density profile to be $`r^2`$ or steeper for the model to be physical.
A more realistic description of velocity anisotropy is provided by a model proposed by Osipkov (1979) and Merritt (1985) with $`\beta `$ dependent on distance from the centre of the object
$$\beta _{\mathrm{OM}}=\frac{s^2}{s^2+s_\mathrm{a}^2}$$
(18)
where $`s_\mathrm{a}`$ is the anisotropy radius determining the transition from isotropic orbits inside to radial orbits outside. As mentioned above, the results of $`N`$-body simulations suggest $`\sigma _\theta /\sigma _\mathrm{r}0.8`$ and therefore $`\beta 0.36`$ at $`s=1`$, which gives $`s_\mathrm{a}4/3`$, a value that we adopt here for all numerical calculations.
For the Osipkov-Merritt model the solution to the Jeans equation (with the same boundary condition as before) reads
$`{\displaystyle \frac{\sigma _\mathrm{r}^2}{V_v^2}}(s,\beta _{\mathrm{OM}})`$ $`=`$ $`{\displaystyle \frac{g(c)s(1+cs)^2}{s^2+s_\mathrm{a}^2}}`$ (19)
$`\times `$ $`{\displaystyle _s^{\mathrm{}}}\left[{\displaystyle \frac{(s^2+s_\mathrm{a}^2)\mathrm{ln}(1+cs)}{s^3(1+cs)^2}}{\displaystyle \frac{c(s^2+s_\mathrm{a}^2)}{s^2(1+cs)^3}}\right]ds`$
and the integration gives
$`{\displaystyle \frac{\sigma _\mathrm{r}^2}{V_v^2}}(s,\beta _{\mathrm{OM}})`$ $`=`$ $`{\displaystyle \frac{g(c)s(1+cs)^2}{2(s^2+s_\mathrm{a}^2)}}`$
$`\times `$ $`\{{\displaystyle \frac{cs_\mathrm{a}^2}{s}}c^2s_\mathrm{a}^2\mathrm{ln}(cs)+c^2s_\mathrm{a}^2\mathrm{ln}(1+cs)(1+{\displaystyle \frac{1}{c^2s^2}}{\displaystyle \frac{4}{cs}})`$
$``$ $`(1+c^2s_\mathrm{a}^2)\left[{\displaystyle \frac{1}{(1+cs)^2}}+{\displaystyle \frac{2\mathrm{ln}(1+cs)}{1+cs}}\right]`$
$`+`$ $`(1+3c^2s_\mathrm{a}^2)[{\displaystyle \frac{\pi ^2}{3}}{\displaystyle \frac{2}{1+cs}}+\mathrm{ln}^2(1+cs)+2\mathrm{L}\mathrm{i}_2(cs)]\}.`$
Figure 1 shows the radial dependence of the radial velocity dispersion. The upper panel of the Figure presents how the results depend on the concentration parameter in the isotropic case, while the lower panel compares predictions for different anisotropy models with $`c=10`$.
### 2.2 The energy distributions
The potential energy associated with the mass distribution of equation (8) is
$`W(s)`$ $`=`$ $`{\displaystyle \frac{1}{r_v}}{\displaystyle _0^s}{\displaystyle \frac{GM(s)}{s}}{\displaystyle \frac{\mathrm{d}M(s)}{\mathrm{d}s}}ds`$ (21)
$`=`$ $`W_{\mathrm{}}\left[1{\displaystyle \frac{1}{(1+cs)^2}}{\displaystyle \frac{2\mathrm{ln}(1+cs)}{1+cs}}\right],`$
where
$$W_{\mathrm{}}=\underset{s\mathrm{}}{lim}W(s)=\frac{cg^2(c)GM_v^2}{2r_v}.$$
(22)
The kinetic energy for arbitrary $`\beta `$ is given by
$$T(s,\beta )=2\pi r_v^3_0^s(32\beta )\rho (s)\sigma _\mathrm{r}^2(s,\beta )s^2ds.$$
(23)
For the three cases of $`\beta `$=0, 0.5 and 1, we obtain respectively
$`T(s,\beta =0)`$ $`=`$ $`{\displaystyle \frac{1}{2}}W_{\mathrm{}}\{3+{\displaystyle \frac{3}{1+cs}}2\mathrm{ln}(1+cs)`$ (24)
$`+`$ $`cs[5+3\mathrm{ln}(1+cs)]c^2s^2[7+6\mathrm{ln}(1+cs)]`$
$`+`$ $`c^3s^3[\pi ^2\mathrm{ln}c\mathrm{ln}s+\mathrm{ln}(1+cs)`$
$`+`$ $`3\mathrm{ln}^2(1+cs)+6\mathrm{L}\mathrm{i}_2(cs)]\},`$
$`T(s,\beta =0.5)`$ $`=`$ $`{\displaystyle \frac{1}{3}}W_{\mathrm{}}\{3+{\displaystyle \frac{3}{1+cs}}3\mathrm{ln}(1+cs)`$ (25)
$`+`$ $`6cs[1+\mathrm{ln}(1+cs)]c^2s^2[\pi ^2`$
$`+`$ $`3\mathrm{ln}^2(1+cs)+6\mathrm{L}\mathrm{i}_2(cs)]\},`$
$`T(s,\beta =1)`$ $`=`$ $`{\displaystyle \frac{1}{2}}W_{\mathrm{}}\{2\mathrm{ln}(1+cs)+cs[{\displaystyle \frac{\pi ^2}{3}}{\displaystyle \frac{1}{1+cs}}`$ (26)
$`+`$ $`\mathrm{ln}^2(1+cs)+2\mathrm{L}\mathrm{i}_2(cs)]\},`$
where we have used in each case the corresponding expression for $`\sigma _\mathrm{r}^2(s,\beta )`$ from equations (14)-(16). For the Osipkov-Merritt model the calculation has to be done numerically.
The results for the potential and kinetic energy (21)-(26) lead to a virial ratio $`lim_s\mathrm{}2T/|W|=1`$ for any value of $`c`$, in agreement with the virial theorem. Figure 2 shows how the virial ratio depends on distance for three different values of the concentration parameter in the isotropic case (upper panel) and compares the ratios obtained for different $`\beta `$ with $`c=10`$ (lower panel). At low radii, the virial ratio is large, especially for low concentration parameters and models with much anisotropy. However, as demonstrated by Figure 3, at the virial radius $`r_v(s=1)`$, $`2T/|W|`$ is still greater than unity and grows with the amount of anisotropy in the model. We see that the virial theorem is better satisfied at $`s=1`$ for objects with larger concentration parameters, as $`lim_c\mathrm{}2T/|W|(s=1)=1`$. Since objects of smaller mass have larger concentration parameters, they are closer to dynamical equilibrium.
The scalar virial theorem we referred to above is expected to be satisfied for self-gravitating systems in steady state. In more realistic situations, the system is never isolated and experiences an external gravitational field; there is also continuous infall of matter. We may conclude from the results above that objects with NFW density profiles and different velocity distributions are close to dynamical equilibrium. However, the virial ratio cannot be used to define the boundary of the virialized object.
### 2.3 Structural parameters
A useful quantity is the half-mass radius. Unfortunately, the divergence of the mass of the NFW profile forces one to define the half-mass radius within a cutoff radius $`r_{\mathrm{cut}}`$. The most natural choice is $`r_{\mathrm{cut}}=r_v`$, since the density distribution is only reliable out to the virial radius. With $`r_{\mathrm{cut}}=r_v`$, the half-mass radius $`r_\mathrm{h}`$ satisfies the following relation for the mass of dimensionless radius:
$$M\left(\frac{r_\mathrm{h}}{r_v}\right)=\frac{M(1)}{2}.$$
(27)
Numerical values of $`r_\mathrm{h}/r_v`$ are easily obtained using equation (8) and over the range $`1<c<100`$ they can be approximated to better than 2% accuracy by
$`{\displaystyle \frac{r_\mathrm{h}}{r_v}}`$ $`=`$ $`0.60820.1843\mathrm{log}c`$ (28)
$`0.1011\mathrm{log}^2c+0.03918\mathrm{log}^3c.`$
The lowest thick solid line in Figure 4 shows how $`r_\mathrm{h}/r_v`$ decreases with increasing concentration parameter.
It is useful to estimate the concentration $`\gamma `$ of a dynamical system, such that
$$\sigma ^2=\gamma \frac{GM}{r_\mathrm{h}},$$
(29)
where $`\sigma ^2=\sigma _\mathrm{r}^2+\sigma _\theta ^2+\sigma _\varphi ^2`$ is the mass weighted mean-square velocity dispersion. As first noted by Spitzer (1969) for polytropes, many realistic density profiles have $`\gamma =0.4`$. For example, it is easy to show that for the Hernquist (1990) model with $`\beta =0`$, $`\gamma =(1+\sqrt{2})/60.403`$ (Mamon 2000).
Using equation (29) and limiting again the mass to $`r_{\mathrm{cut}}=r_v`$, we define $`\gamma `$ with
$$\gamma =\frac{r_\mathrm{h}\sigma ^2_{rr_v}}{GM(1)}=2\frac{r_\mathrm{h}T(1,\beta )}{GM^2(1)},$$
(30)
where we made use of
$$T(x,\beta )=\frac{1}{2}M(x)\sigma ^2_{rxr_v}.$$
(31)
The values of $`\gamma `$ for different velocity anisotropy models, derived from equations (7), (8), (22), (22)-(26), (28), and (30) are shown in Figure 4 and in the case of $`\beta =0`$ yield numbers closest to 0.4: $`\gamma =0.56`$ for $`c=5`$ and $`\gamma =0.51`$ for $`c=10`$. Thus the NFW model produces $`\gamma `$s that are higher than the canonical value of 0.4, especially if more velocity anisotropy is assumed. This may be caused by the ill-defined cutoff radius.
In models with homogeneous cores, the central density, the core radius $`r_\mathrm{c}`$ and the central 3-D velocity dispersion $`\sigma ^2(0)`$ are related through
$$4\pi G\rho (0)r_\mathrm{c}^2=\frac{1}{3}\eta \sigma ^2(0).$$
(32)
King (1966) models have $`\eta =9`$. In models with cuspy cores, we propose the scaling relation
$$4\pi G\rho (r_\mathrm{s})r_\mathrm{s}^2=\frac{1}{3}\eta \sigma ^2_{r<r_\mathrm{s}}.$$
(33)
Using equations (2), (6) and (7), one has $`4\pi G\rho (r_\mathrm{s})r_\mathrm{s}^2=cg(c)V_v^2/4`$ and from equation (31) for $`x=1/c`$ one obtains
$$\eta =\frac{3cg(c)V_v^2M(1/c)}{8T(1/c,\beta )}.$$
(34)
For different velocity anisotropy models we then have
$$\eta (\beta =0)=\frac{3(2\mathrm{ln}21)}{2(\pi ^278\mathrm{ln}2+6\mathrm{ln}^22)}2.797,$$
(35)
$$\eta (\beta =0.5)=\frac{9(12\mathrm{ln}2)}{4(\pi ^296\mathrm{ln}2+6\mathrm{ln}^22)}2.138,$$
(36)
$$\eta (\beta =1)=\frac{9(2\mathrm{ln}21)}{2(\pi ^2312\mathrm{ln}2+6\mathrm{ln}^22)}1.212,$$
(37)
where we have used equations (8) and (24)-(26), and the fact that $`\mathrm{Li}_2(1)=\pi ^2/12`$. Note that $`\eta `$ is independent of $`c`$ in all cases with $`\beta =`$const. For the Osipkov-Merritt model $`\eta `$ is no longer a constant but we find $`1.902<\eta <2.797`$ in the range $`1<c<100`$ with the limiting cases of $`\eta \eta (\beta =1)`$ for $`c0`$ and $`\eta \eta (\beta =0)`$ for $`c\mathrm{}`$. Such limiting behaviour is due to the fact that for large $`c`$ the integration of $`T(1/c,\beta )`$, equation (23), probes only the range of $`s`$ where $`\beta `$ is close to zero, while for small $`c`$ the integral is dominated by contribution from large $`s`$ where $`\beta `$ is close to unity.
Finally, we consider the structural parameter
$$\text{WUM}=\frac{W(s)}{M(s)\mathrm{\Phi }(0)}$$
(38)
brought forward by Seidov & Skvirsky (2000) with the motivation of WUM being constant for different self-gravitating objects of simple geometry. Using equations (8), (9) and (21) we find that for the NFW model
$$\text{WUM}=\frac{cs(2+cs)2(1+cs)\mathrm{ln}(1+cs)}{2(1+cs)[cs+(1+cs)\mathrm{ln}(1+cs)]}$$
(39)
so the parameter turns out to be a function of $`cs=r/r_\mathrm{s}`$ only. It grows with $`s`$ from zero at $`s0`$ reaching a maximum value of $`0.196`$ at $`r/r_\mathrm{s}=4.62`$ and decreases to zero again as $`s\mathrm{}`$. The values of this parameter at the virial radius are $`0.196`$, $`0.187`$ and $`0.125`$ respectively for $`c=5,10`$ and $`100`$.
### 2.4 The distribution function
A quantity of great dynamical importance is the distribution function. For a spherical system with an isotropic velocity tensor, the distribution function depends on the phase-space coordinates only through the energy (e.g. Binney & Tremaine 1987), and can be derived through the Eddington (1916) formula (e.g. Binney & Tremaine 1987):
$$f()=\frac{1}{\sqrt{8}\pi ^2}\left[_0^{}\frac{\mathrm{d}^2\rho }{\mathrm{d}\mathrm{\Psi }^2}\frac{\mathrm{d}\mathrm{\Psi }}{\sqrt{\mathrm{\Psi }}}+\frac{1}{^{1/2}}\left(\frac{\mathrm{d}\rho }{\mathrm{d}\mathrm{\Psi }}\right)_{\mathrm{\Psi }=0}\right],$$
(40)
where $``$ and $`\mathrm{\Psi }`$ are the conventionally defined relative energy and potential; here $`=E`$, where $`E`$ is the total energy per unit mass and $`\mathrm{\Psi }=\mathrm{\Phi }`$, where $`\mathrm{\Phi }`$ is given by equation (9).
It is easy to show that, given equations (6) and (9), the second term in brackets in equation (40) is zero. The simplest way to perform the integration of the first term is to introduce dimensionless variables $`\stackrel{~}{\mathrm{\Psi }}=\mathrm{\Psi }/C_1`$ and $`\stackrel{~}{\rho }=\rho /C_2`$, where $`C_1=g(c)V_v^2`$ and $`C_2=c^2g(c)M_v/(4\pi r_v^3)`$. Then the integration variable should be changed to $`s`$ and the limit of integration corresponding to $``$ found numerically for each $``$ by solving equation $`\mathrm{\Psi }(s)=`$. Otherwise, with a few percent accuracy, the integration in (40) can be done directly with an approximation $`s_{\mathrm{apx}}=1.75\mathrm{ln}(\stackrel{~}{\mathrm{\Psi }}/c)/\stackrel{~}{\mathrm{\Psi }}`$.
The calculations of the distribution function are usually performed in units such that $`G=M=R_\mathrm{e}=1`$ (Binney & Tremaine 1987), where $`M`$ is the total mass of the system and $`R_\mathrm{e}`$ is its effective radius. Since in the case of NFW profile the total mass is infinite a reasonable choice seems to be to put $`M_v=1`$. The effective radius is not well defined either but can be approximated as $`r_v/2`$ (see the next section). Therefore we choose the units so that $`G=M_v=r_v/2=1`$ and arrive at the numerical results shown in Figure 5. This choice of normalization is equivalent to measuring $`f`$ in units of $`\sqrt{8}M_v/(r_vV_v)^3`$ and $`E`$ in units of $`V_v^2`$.
Figure 5 proves that the distribution function turns out to be similar to the distribution functions obtained from other density profiles (see e.g. Figure 4-12 in Binney & Tremaine 1987), except that the NFW distribution functions do not display the cutoff at nearly unbound energies characteristic of King (1966) models. The results shown in Figure 5 indicate a proper behaviour of the distribution function (it is nowhere negative). Quantitative comparisons with other models should, however, be made with caution because of the aforementioned problem with normalization. Distribution functions for more realistic velocity dispersion models, like the Osipkov-Merritt model, were recently considered in detail by Widrow (2000).
### 2.5 Projected distributions
Of primary importance for comparisons with observations are the projected distributions. The surface mass density of an object is obtained by integrating the density along the line of sight:
$`\mathrm{\Sigma }_M(R)`$ $`=`$ $`2{\displaystyle _R^{\mathrm{}}}{\displaystyle \frac{r\rho (r)}{(r^2R^2)^{1/2}}}dr`$ (41)
$`=`$ $`{\displaystyle \frac{c^2g(c)}{2\pi }}{\displaystyle \frac{M_v}{r_v^2}}{\displaystyle \frac{1|c^2\stackrel{~}{R}^21|^{1/2}C^1[1/(c\stackrel{~}{R})]}{(c^2\stackrel{~}{R}^21)^2}},`$
where
$$C^1(x)=\{\begin{array}{cc}\mathrm{cos}^1(x)\hfill & \text{if }R>r_\mathrm{s}\hfill \\ \mathrm{cosh}^1(x)\hfill & \text{if }R<r_\mathrm{s}\text{ .}\hfill \end{array}$$
(42)
In the above expressions $`R`$ is the projected radius and $`\stackrel{~}{R}=R/r_v`$. For the singular case $`R=r_\mathrm{s}`$ the $`\stackrel{~}{R}`$-dependent expression in equation (41) equals $`1/3`$ and we have $`\mathrm{\Sigma }_M(R)=c^2g(c)M_v/(6\pi r_v^2)`$ . An analytical formula equivalent to equation (41) was derived independently by Bartelmann (1996).
The projected mass is then given by
$`M_\mathrm{p}(R)`$ $`=`$ $`2\pi {\displaystyle _0^R}R\mathrm{\Sigma }_M(R)dR`$ (43)
$`=`$ $`g(c)M_v\left[{\displaystyle \frac{C^1[1/(c\stackrel{~}{R})]}{|c^2\stackrel{~}{R}^21|^{1/2}}}+\mathrm{ln}\left({\displaystyle \frac{c\stackrel{~}{R}}{2}}\right)\right],`$
which is logarithmically divergent at large $`\stackrel{~}{R}`$. $`C^1(x)`$ is again given by equation (42).
Another important projected quantity is the line-of-sight velocity dispersion which for a spherical non-rotating system is (Binney & Mamon 1982)
$$\sigma _{\mathrm{los}}^2(R)=\frac{2}{\mathrm{\Sigma }_M(R)}_R^{\mathrm{}}\left(1\beta \frac{R^2}{r^2}\right)\frac{\rho \sigma _\mathrm{r}^2(r,\beta )r}{\sqrt{r^2R^2}}dr,$$
(44)
where $`\mathrm{\Sigma }_M(R)`$ is given by equation (41) and the radial velocity dispersions $`\sigma _\mathrm{r}(r,\beta )`$ for our four models are given by equations (14)-(16) and (2.1). For circular orbits, $`\sigma _\mathrm{r}=0`$, and one has
$$\sigma _{\mathrm{los}}^2(R)=\frac{1}{\mathrm{\Sigma }_M(R)}_R^{\mathrm{}}\left(\frac{R}{r}\right)^2\frac{\rho V^2r}{\sqrt{r^2R^2}}dr,$$
(45)
where $`V`$ is the circular velocity given by equations (10) and (11). The upper panel of Figure 6 shows the profiles of line-of-sight velocity dispersion (with isotropic orbits), obtained through numerical integration of equation (44) for different concentration parameters. The lower panel of Figure 6 compares the radial profiles of line-of-sight velocity dispersions obtained for $`c=10`$ for different velocity anisotropy models.
For more distant or intrinsically small galaxies, as well as for groups and clusters, spectroscopic observations are often limited to a single large aperture centred on the object. The mean velocity dispersion within an aperture (hereafter, aperture velocity dispersion) is
$$\sigma _{\mathrm{ap}}^2(R)=\frac{S^2(R)}{M_\mathrm{p}(R)}$$
(46)
where
$$S^2(R)=2\pi _0^R\mathrm{\Sigma }_M(P)\sigma _{\mathrm{los}}^2(P)PdP.$$
(47)
In the above expressions $`R`$ is the radius of the aperture, $`\mathrm{\Sigma }_M(P)`$ is the surface mass distribution, equation (41), and $`M_\mathrm{p}(R)`$ is the projected mass given by equation (43).
Inserting the expression for $`\sigma _{\mathrm{los}}`$ (eq. ) into equation (47), we obtain a double integral, which after inversion of the order of integration is reduced to an easily computable single integral:
$`S^2(R)`$ $`=`$ $`c^2g(c)M_v\{{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{\sigma _\mathrm{r}^2(s,\beta )s}{(1+cs)^2}}(1{\displaystyle \frac{2\beta }{3}})\mathrm{d}s`$ (48)
$`+`$ $`{\displaystyle _{\stackrel{~}{R}}^{\mathrm{}}}{\displaystyle \frac{\sigma _\mathrm{r}^2(s,\beta )(s^2\stackrel{~}{R}^2)^{1/2}}{(1+cs)^2}}[{\displaystyle \frac{\beta (\stackrel{~}{R}^2+2s^2)}{3s^2}}1]\mathrm{d}s\},`$
where as before, $`\stackrel{~}{R}=R/r_v`$, $`s=r/r_v`$ and $`\sigma _\mathrm{r}^2(s,\beta )`$ for different $`\beta `$ are given by equations (14)-(16) and (2.1). Analogous expression for circular orbits can be obtained from (48) by replacing $`\sigma _\mathrm{r}^2`$ by $`V^2`$, keeping only the terms proportional to $`\beta `$ and dividing by $`(2\beta )`$.
Figure 7 displays the radial profiles of aperture velocity dispersion, computed numerically from equation (48). From the upper panel of the Figure we see that in the isotropic case the dependence of the results on the concentration parameter is rather strong and monotonic for a given $`R`$. The lower panel of the Figure compares the predictions for different velocity anisotropy models.
## 3 Comparison with observations
Comparisons of the surface mass density to surface brightness observations are usually performed with the assumption of constant mass-to-light ratio $`\mathrm{{\rm Y}}=\mathrm{const}`$. This assumption is not likely to be physical, because of the different physics involved in the assemblies of the dark matter and baryonic components of galaxies. In particular, the baryons in elliptical galaxies may well settle at an early epoch, within a radius that is the lower of the radius with virial overdensity $`v200`$ and the radius at which gas can cool to form molecular clouds and later stars. The baryons in ellipticals will then sit today in a region of overdensity $`v200`$, and one then expects $`\mathrm{{\rm Y}}`$ to rise with $`r`$, at least at large radii.
Nevertheless, for simplicity, we check whether the observations of elliptical galaxies are consistent with the idea that stars are distributed within elliptical galaxies according to the NFW density profile, characterized by a virial radius where the mean overdensity is 200. Such a situation may arise if the dark matter were negligible within elliptical galaxies or distributed precisely like the luminous matter. In a forthcoming paper (Mamon & Łokas, in preparation), we will check in more detail whether the observations of elliptical galaxies are compatible or not with NFW density profiles for the mass distribution.
For constant mass-to-light ratio we have $`\mathrm{\Sigma }_M(R)=\mathrm{{\rm Y}}I(R)`$, where $`I`$ is the surface brightness. The radial profiles of $`I=\mathrm{\Sigma }_M/\mathrm{{\rm Y}}`$ and $`M_\mathrm{p}`$ are shown in Figure 8. Both quantities are normalized to their values at the virial radius. Figure 8 shows that the surface mass density depends weakly on the concentration parameter, especially at larger distances from the centre.
Since the surface mass density (eq. ) behaves as $`1/R^2`$ at large distances, one may therefore compare it with the Hubble-Reynolds formula (Reynolds 1913), which was the first model used to describe the surface brightness profiles of elliptical galaxies:
$$I_{\mathrm{HR}}(R)=\frac{I_0}{(1+R/R_{\mathrm{HR}})^2}.$$
(49)
$`R_{\mathrm{HR}}`$ is the characteristic radius of the distribution, where the surface brightness falls to one-quarter of its central value. The thin curves of Figure 8 show that the surface mass density of the NFW model (eq. ) is very well fitted by equation (49) and the best-fit values of $`\stackrel{~}{R}_{\mathrm{HR}}=R_{\mathrm{HR}}/r_v`$ are $`0.119`$, $`0.0640`$ and $`0.00743`$ respectively for $`c=5,10`$ and $`100`$.
The surface brightness profiles of astrophysical objects are often scaled with the effective radius, which we denote $`R_\mathrm{e}`$, where the projected luminosity is half the total luminosity. Given the divergence of the projected mass, we are forced again to introduce a cut-off at some scale $`R_{\mathrm{cut}}=\stackrel{~}{R}_{\mathrm{cut}}r_v`$. We then have
$$M_\mathrm{p}(R_\mathrm{e})=M_\mathrm{p}(R_{\mathrm{cut}})/2.$$
(50)
Figure 9 shows the effective radius, calculated numerically from equations (43) and (50). For $`\stackrel{~}{R}_{\mathrm{cut}}=1`$, a useful approximation, good to better than 2% relative accuracy, is:
$`R_\mathrm{e}/r_v`$ $`=`$ $`0.55650.1941\mathrm{log}c`$ (51)
$`0.0756\mathrm{log}^2c+0.0331\mathrm{log}^3c.`$
The prediction for the surface brightness $`I=\mathrm{\Sigma }_M/\mathrm{{\rm Y}}`$ with $`\mathrm{\Sigma }_M`$ given by equation (41) expressed in terms of the effective radius and the corresponding effective brightness $`I_\mathrm{e}=I(R_\mathrm{e})`$ is shown in the upper panel of Figure 10 for different values of the concentration parameter $`c`$. For comparison, we also show the de Vaucouleurs (1948) $`R^{1/4}`$ law describing the observed surface brightness distribution in giant elliptical galaxies:
$$I(R)=I_\mathrm{e}\mathrm{exp}\{b[(R/R_\mathrm{e})^{1/4}1]\},$$
(52)
where $`b=7.67`$. Clearly, the NFW surface brightness profiles are poorly fitted by the $`R^{1/4}`$ law, when using $`R_{\mathrm{cut}}=r_v`$ to define the effective radius of the NFW profile.
The lower panel of Figure 10 shows how the results depend on the choice of cut-off for $`c=10`$ and $`\stackrel{~}{R}_{\mathrm{cut}}=3`$, 3.5, 4, 4.5, and 5. At first glance, it seems that the NFW profile is well fitted by the $`R^{1/4}`$ law, especially for $`\stackrel{~}{R}_{\mathrm{cut}}4`$. However, the range of surface mass densities where the fit is excellent is roughly $`10^2`$, and the fit is adequate for a range smaller than $`10^3`$. In contrast, the surface brightness profile of the nearby giant elliptical galaxy NGC 3379 (M 105) follows the $`R^{1/4}`$ law in a range of 10 magnitudes (de Vaucouleurs & Capaccioli 1979), i.e. a factor $`10^4`$ in intensity.
In order to see how good is the de Vaucouleurs’s fit in this case in both panels of Figure 10 we plotted a number of data points equally spaced in $`R^{1/4}`$. Since de Vaucouleurs & Capaccioli (1979) do not provide the error bars for their data, the error bars shown in the Figure were taken from Goudfrooij et al. (1994). The excess of the data above the $`R^{1/4}`$ law for small $`R`$ was already noted by de Vaucouleurs & Capaccioli (1979). The error bars are negligible for $`R<R_\mathrm{e}`$ and smaller than 15% out to $`2.5R_\mathrm{e}`$, the maximum distance from the centre reached in the data of Goudfrooij et al. (1994).
According to de Vaucouleurs & Capaccioli (1979), in this galaxy the $`R^{1/4}`$ surface brightness profile extends to $`R_{\mathrm{lim}}=7.5R_\mathrm{e}=26.4\mathrm{kpc}`$, given a distance of 12.4 Mpc to NGC 3379 (Salaris & Cassisi 1998). Within $`R_{\mathrm{lim}}`$, de Vaucouleurs & Capaccioli (1979) report a blue magnitude, corrected for galactic extinction of $`B=10.10`$, yielding a total blue luminosity of $`2.2\times 10^{10}L_{}`$, hence a blue luminosity density of $`2.8\times 10^5L_{}\mathrm{kpc}^3`$. Since the mass within $`R_{\mathrm{lim}}`$ must be greater than the mass in stars, we infer that within this radius, $`\mathrm{{\rm Y}}_B>8`$ (the typical mass to blue luminosity ratio for old stellar populations), yielding an overdensity of the galaxy, relative to the critical density $`\rho _c`$ of $`v>1.6\times 10^4/(H_0/70\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1)^2`$. Therefore, since $`v100`$ (the value at $`r_v`$), we conclude that $`R_{\mathrm{lim}}r_v`$, hence $`R_\mathrm{e}r_v/7.5`$. In contrast, with $`\stackrel{~}{R}_{\mathrm{cut}}=1`$, the effective radius of the NFW model ($`c=10`$) is $`0.3r_v`$ (Figure 9). This discrepancy in $`R_\mathrm{e}/r_v`$ between NFW and $`R^{1/4}`$ law gets even worse if one adopts $`\stackrel{~}{R}_{\mathrm{cut}}=4`$, which provides the best fits of the NFW surface mass density to the $`R^{1/4}`$ law: indeed, Figure 9 indicates $`R_\mathrm{e}0.8r_v`$ for the NFW model.
In summary, the NFW surface mass density profile resembles an $`R^{1/4}`$ law in a fairly wide range of radii, but 1) one has to resort to an abnormally large effective radius, very close to the virial radius, and assume that the effective radius measures half the projected light (or mass) within 4 times the virial radius, and 2) the fit is good in a considerably smaller range of radii than is observed in the nearby giant elliptical NGC 3379.
The generalization of the $`R^{1/4}`$ law into an $`R^{1/m}`$ law, first proposed by Sérsic (1968), is known to fit the surface brightness profiles of elliptical galaxies within a much larger mass range than the de Vaucouleurs law (Caon, Capaccioli & D’Onofrio 1993). The surface brightness of the Sérsic profile is
$$I(R)=I_\mathrm{e}\mathrm{exp}\{b(m)[(R/R_\mathrm{e})^{1/m}1]\},$$
(53)
where $`b(m)`$ is tabulated by Ciotti (1991), who gives the empirical relation $`b(m)2m0.324`$, good to 0.1% relative accuracy. The de Vaucouleurs law is reproduced for $`m=4`$, while $`m=1`$ corresponds to an exponential law as in spiral disks.
In Figure 11, we plot the NFW surface brightness $`I=\mathrm{\Sigma }_M/\mathrm{{\rm Y}}`$, with $`\mathrm{\Sigma }_M`$ given by equation (41) and $`c=10`$, as a function of $`(R/R_\mathrm{e})^{1/m}`$ for various values of the Sérsic parameter $`m`$. We compare them to the Sérsic profiles given by the straight dashed lines. The agreement is good for all values of $`m`$, within ranges of $`I/I_\mathrm{e}`$ that increase with increasing $`m`$. Comparison of the plots for different $`m`$ shows that the Sérsic models with lower $`m`$ generally agree better with the NFW surface brightness for smaller radii, while those with larger $`m`$ are in better agreement at larger radii, closer to the virial radius. Overall, the NFW profile matches best the $`m=3`$ Sérsic law, over a factor of $`10^3`$ in intensity (7.5 magnitudes).
For a more quantitative comparison, we performed two-parameter fits of the Sérsic models (53) to the projected NFW formula (41). The NFW profile was sampled in the range of $`0.01<\stackrel{~}{R}<1`$ with a given $`c`$. The fitted Sérsic parameters $`1/m`$ and $`R_\mathrm{e}/r_v`$ obtained for different $`c`$ are shown in Figure 12. Figure 13 compares the two projected profiles for $`c=10`$. The best-fit parameters of the Sérsic model in this case are $`m=3.07`$ and $`R_\mathrm{e}/r_v=0.55`$.
While Caon et al. (1993) find similar ranges of agreement between observed profiles and Sérsic laws, this range in intensity is still smaller than the range of $`10^4`$ found for NGC 3379 by de Vaucouleurs & Capaccioli (1979). Moreover, while Caon et al. (1993) find that the best fitting Sérsic models for elliptical galaxies have indices spanning a wide range, from $`m=2`$ for faint ellipticals to $`m=10`$ for bright ellipticals, the Sérsic laws that match the NFW models span a much smaller range, roughly $`m=3\pm 0.5`$ ($`2.71<m<3.41`$ for $`5<c<15`$, see Figure 12). Moreover, the problem of very high values of $`R_\mathrm{e}/r_v`$ ($`0.46<R_\mathrm{e}/r_v<0.81`$ for $`5<c<15`$, see Figure 12), remains in the fits of Sérsic profiles to projected NFW models.
## 4 Discussion
The main disadvantage of the NFW model is the logarithmic divergence of its mass (and luminosity for constant mass-to-light ratio). In contrast, the Jaffe (1983) and Hernquist (1990) models converge in mass, and their properties can be expressed in units of their asymptotic mass. For the NFW model, one is restricted to a mass at a physical radius such as the virial radius. This mass divergence also complicates the analysis of surface brightness profiles, which involve the effective radius where the aperture luminosity is half its asymptotic value. However, independently of the radial cut-off introduced to define the effective radius, the projected NFW density profile is consistent with constant mass-to-light ratio, given the observed Sérsic profiles of elliptical galaxies, but only in a limited range of radii, with unusually high values of $`R_\mathrm{e}`$ and in a smaller interval of Sérsic shape parameters than observed. On the other hand, the Hernquist (1990) model, whose density profile scales as $`r^4`$ at large radii, produces better fits to the $`R^{1/4}`$ law.
The upper panel of Figure 10 suggests that, for reasonable effective radii, if indeed dark matter follows the NFW profile, the mass-to-light ratio, $`\mathrm{{\rm Y}}`$, is not constant but increases with radius, not only in the outer regions, as is inferred from the commonly accepted picture of galaxies embedded in more spatially extended dark haloes, but also in the inner regions. This is at odds with the observed kinematics of ellipticals that Bertola et al. (1993) inferred from observations of ionised and neutral gas around specific ellipticals. Moreover, increasing $`\mathrm{{\rm Y}}`$ throughout the galaxy implies radial velocity anisotropy throughout elliptical galaxies, whereas violent relaxation should cause isotropic cores.<sup>1</sup><sup>1</sup>1Note that recent, state of the art observations and modelling by Saglia et al. (2000) and Gebhardt et al. (2000) do not strongly constrain the gravitational potentials of elliptical galaxies, although NFW potentials may turn out to be inconsistent with the current data. On the other hand, Kronawitter et al. (2000) are able to rule out constant $`\mathrm{{\rm Y}}`$ for some elliptical galaxies. Thus it appears difficult to reconcile the photometry and kinematics of elliptical galaxies with NFW models. In a forthcoming paper (Mamon & Łokas, in preparation), we will omit the assumption of mass follows light in a more detailed assessment of the compatibility of the observations of elliptical galaxies with the NFW model.
The results presented in this paper can be directly applied to the analysis of the mass and light distribution in clusters of galaxies. A standard procedure to do it is to measure the surface brightness and the light-of-sight velocity dispersion and assuming some form of velocity distribution or mass-to-light ratio calculate the luminosity density and the velocity dispersion by solving the Abel integral equations (41) and (44) and the Jeans equation (Binney & Mamon 1982, Tonry 1983, Solanes & Salvador-Solé 1990, Dejonghe & Merritt 1992). The results of this procedure are uncertain because it involves derivatives of observed quantities which are usually noisy. One also experiences a degeneracy because different models fit the data equally well (Merritt 1987). Instead of solving the Abel equations one can also model the luminosity density and velocity dispersion with simple functions and fit their parameters so that they reproduce their projected counterparts (Carlberg et al. 1997).
Our results are useful for the simpler approach of assuming realistic forms of the density distribution, velocity distribution and mass-to-light ratio. Here we provide the tools for modelling the NFW density profile with different velocity distributions and constant mass-to-light ratio ($`\mathrm{{\rm Y}}=\mathrm{const}`$), and obtain exact predictions for the surface brightness and the line-of-sight as well as aperture velocity dispersion that can be directly compared to observations.
## Acknowledgements
We thank Daniel Gerbal and Bernard Fort for useful conversations, and an anonymous referee for helpful comments. ELŁ acknowledges hospitality of Institut d’Astrophysique de Paris, where part of this work was done. This research was partially supported by the Polish State Committee for Scientific Research grant No. 2P03D00813 and 2P03D02319 as well as the Jumelage program Astronomie France Pologne of CNRS/PAN. |
warning/0002/physics0002002.html | ar5iv | text | # The Multifractal Time and Irreversibility in Dynamic Systems
## I Introduction
The dynamic equations of the physical theories are reversible, it is well known. The kinetic equations of the statistical theory are irreversible. The irreversibility the Boltzmann’s statistical theory was in due time the main reason of non-recognition by Poincare of Boltzmann’s statistical theory. For the irreversibility introducing in the physical dynamical equations (for example, in the equation of the Liouville) it is necessary to introduce the dissipation terms , or the functionals of a microscopic entropy and time ensuring realization of the second law of thermodynamics. In the first case irreversibility in the dynamic equations arises as a sequence of mathematical approaches. Prigogin’s the point of view is consists in recognizing the primacy of irreversible processes and it seems intuitively more reasonable. Are more general, than mentioned, a methods of introducing of the irreversibility in the dynamical equations of physics exist? Is it possible the reversibility of the equations of the dynamical theories to introduce as result of approximate transition from more rigorous the dynamical irreversible equations (these dynamic equations may be obtained as a result of generalization of the known equations) to the idealized and reversible, but approximate equations? The purpose of a note is to introduce one of a possible generalizations of the dynamical theories of physics realized by replacement of time with topological dimension equal unity on ”multifractal” time (for the first time it was introduced in ). In the mentioned theory the time is characterized in each time and space points by fractional (fractal) dimensions (FD)$`d_t(𝐫(t),t)`$. The marked replacement dimension of time by fractional dimensions gives in the origin of irreversibility in the dynamical equations of physics and the existence of irreversibility in our world may be interpreted by new reason. It is not contradicts (for FD is small differs from unity for weak physical fields on the Earth) an experimental data, and allows to receive a new interesting physical results. The equations of the classical mechanics (the equations of the Hamilton - Liouville) are chosen for research as an typical example of dynamic systems.
## II Equation of a mechanics with time defined on multifractal sets
The method of generalization of the classical mechanics equations is founded on the new model of an approach to the problem of a nature of time .This model is consists in replacement of usual time by the time defined on a multifractal subsets $`s_t`$ of continuously set $`M_t`$ (the measure carrier). The multifractal set $`S_t`$ consists of subsets $`s_t`$, i.e. very small time intervals (in further named ”points”), which are also multifractal, and each of them is characterized in turn by its global fractal dimension FD) $`d_t(𝐫(t),t)`$ ( defined as box-dimension, , and so on) that depends of a nature of sets $`s_t`$, and depends at coordinates and time (see , ). So each time subset is characterized by its global fractal dimension $`d_t(𝐫(t),t)`$ which characterize the scaling characteristics for this subset. The continuity $`d_t(𝐫(t),t)`$ is supposed. The new approach to a nature of time is consists in the replacing the usual time points of time axe by selection for describing of the time’s intervals (the ”points” on a time axis consisting of sets $`s_t`$ that defined on the measure carrier $`M_t`$) only the ”points” that characterized by sets $`s_t`$. The time axe (or, in other selections time plane or time volume $`R_n`$) is the carrier of measure of all the multifractal time subsets $`s_t`$ defined on it. The researching of the dynamic equations and physical quantities (in particular, the entropy) with time ”points” with fractal characteristics defined on multifractal sets $`s_t`$ with FD $`d_t(𝐫(t),t)`$, lead to irreversibility of all dynamic theories used in physics. It is stipulated by an openness of dynamic systems with multifractal time (role of a thermostat plays set $`R_n`$) that is appears in a time dependencies of all physical and mathematical (except for zero) objects. For describing of small changes of functions defined on multifractal time sets, it is impossible to apply ordinary or fractional (in sense of the Riemann - Liouville) derivatives and integrals, since to different time points there corresponds to different fractal dimensions. For describing of changes of such functions need’s introduction of generalized fractional derivatives and integrals . In this note the multifractal properties of the space sets $`s_𝐫`$ are not considered, since the irreversibility of the dynamic physical equations arises already at the using only of multifractal time (see also )and global FD of small time sets intervals $`s_t`$.
## III Generalized fractional derivatives and integrals on multifractal set $`S_t`$ of time points
It is necessary if we want to describe the dynamics of time-dependent functions determined on multifractal set $`S_t`$ to enter the functionals that extends the fractional derivatives and integrals of the Riemann - Liouville on the set $`S_t`$ with FD $`d_t(𝐫(t),t)`$) that is different in each subsets $`s_t`$ )
$$D_{+,t}^{d_t}f(t)=\left(\frac{d}{dt}\right)^n\underset{a}{\overset{t}{}}𝑑t^{}\frac{f(t^{})}{\mathrm{\Gamma }(nd_t(t^{}))(tt^{})^{d_t(t^{})n+1}}$$
(1)
$`D_{,t}^{d_t}f(t)=(1)^n\times `$ (2)
$`\times \left({\displaystyle \frac{d}{dt}}\right)^n{\displaystyle \underset{t}{\overset{b}{}}}𝑑t^{}{\displaystyle \frac{f(t^{})}{\mathrm{\Gamma }(nd_t(t^{}))(t^{}t)^{d_t(t^{})n+1}}}`$ (3)
where $`\mathrm{\Gamma }`$-is Euler gamma -function,$`a<b`$, $`a`$ and $`b`$ is stationary values selected on an axis (from $`\mathrm{}`$ to $`\mathrm{}`$), $`n1d_t<n`$, $`n=\{d_t\}+1`$, $`\{d_t\}`$ is an integer part of $`d_t0`$, $`n=0`$ for $`d_t<0`$, $`d_t=d_t(𝐫(t),t)`$-is the fractal dimensions (FD). The dependencies FD from time and space coordinates are defined by the Lagrangians’s densities of a viewed problem , , . The generalized fractional derivative (GFD) (1)-(2) coincide with fractional derivatives or fractional integrals of the Riemann - Liouville in the case $`d_t=const`$. At $`d_t=n+\epsilon (t)`$, $`\epsilon 0`$ GFD are represented by usual derivatives and integrals . The functions and integrals in (1)-(2) are considered as generalized functions given on the set of finitary functions . The definitions the GFD (1)-(2) allow to describe the dynamics of functions defined on multifractal sets and GFD substitute (for such functions) the usual or fractional differentiation and integration (GFD partially conserve the memory about of the last time events).
## IV Hamilton equations
The Hamilton equations for system from $`N`$ of classical particles with identical masses m on the set with multifractal time (i.e. time defined on multifractal set $`S_t`$) reads:
$$D_{+,t}^{d_t}𝐫_i=\frac{H}{𝐩_i},D_{,t}^{d_t}𝐩_i=\frac{H}{𝐫_i},𝐩_i=D_{+,t}^{d_t}𝐫_i$$
(4)
$$H=\underset{i=1..N}{}\frac{𝐩_i^2}{2m}+\frac{1}{2}\underset{ij=1..N}{}V(\left|𝐫_i𝐫_j\right|)$$
(5)
The equations (refeq3)-(ref4) differ from the classical Hamilton equations by replacement the derivatives with respect to time by GFD (refeq1)-(ref2) and coincide with the classical equations of a mechanics at $`d_t=1`$. The equations for arbitrary function $`B`$ dynamic variable $`𝐩,𝐫`$ will look like
$$D_{+,t}^{d_t}B=\stackrel{~}{D}_{+,t}^{d_t}B+\frac{H}{𝐩_i}\frac{B}{𝐫_i}\frac{H}{𝐫_i}\frac{B}{𝐩_i}$$
(6)
The figure $`\stackrel{~}{D}_{+,t}^{d_t}`$ in (6) differs from $`D_{+,t}^{d_t}`$ in (1) by replacement the complete derivative with respect to time $`t`$ on a partial differential with respect to time $`t`$. Let’s show, that the Hamiltonian function $`H`$ in space with multifractal time is not integral of a motion of the equation (6), i.e. does not convert a right part of (6) in zero. Substitution $`H`$ in (6) gives in
$$D_{+,t}^{d_t}H=\stackrel{~}{D}_{+,t}^{d_t}H$$
(7)
From the (1) follows, that equation (7) is of the form (for $`d_t(𝐫(t),t)=1+\epsilon (𝐫(t),t)`$, $`\epsilon 0`$, when a simplifying assumption about lack at $`d_t`$, $`\epsilon `$ of explicit dependence from $`t`$ is valid)
$$D_{+,t}^{d_t}H=\stackrel{~}{D}_{+,t}^{d_t}H=\frac{\epsilon H}{\mathrm{\Gamma }(1+\epsilon )t^{d_t}}$$
(8)
and is equal to zero when $`d_t(𝐫(t),t)=1`$. So, in space with multifractal time, at classical system with a Hamiltonian $`H`$ that not depends explicitly at time (conservative systems) the GFD with respect to a total energy depends on time and decreases with the increases of time, i.e. in the model of multifractal time the rigorously conservative classical systems does not exist. For differs $`d_t(𝐫(t),t)`$ from unity by a little bit (that it is valid about it represents experimental data about time and results of ) the changing of energy of system will be very small. Let’s consider the problem of change $`H`$ with change of time in more general case ($`\epsilon =\epsilon (𝐫(t),t)`$). For this purpose we shall be restricted to a case, when FD of time $`d_t(𝐫(t),t)=1`$ is not considerably differs from unity: $`d_t=1+\epsilon (𝐫(t),t)=1`$, $`|\epsilon |1`$. In this case GFD is represented as (integral is calculated as the total of a principal value and residue in a singular point)
$`\stackrel{~}{D}_{+,t}^{1+\epsilon }H`$ $``$ $`{\displaystyle \frac{}{t}}H{\displaystyle \frac{1}{2}}{\displaystyle \frac{}{t}}[{\displaystyle \frac{\epsilon (𝐫(t),t)H}{\mathrm{\Gamma }(1+\epsilon (𝐫(t),t))}}]+`$ (9)
$`+`$ $`{\displaystyle \frac{\epsilon H}{\mathrm{\Gamma }(1+\epsilon )t^{d_t}}},\epsilon >0,d_t<1`$ (10)
$`\stackrel{~}{D}_{+,t}^{1+\epsilon }H`$ $``$ $`{\displaystyle \frac{}{t}}[{\displaystyle \frac{1}{\mathrm{\Gamma }(1\epsilon )}}H]\pm {\displaystyle \frac{1}{2}}{\displaystyle \frac{}{t}}[{\displaystyle \frac{\epsilon (𝐫(t),t)H}{\mathrm{\Gamma }(1\epsilon )}}]+`$ (11)
$`+`$ $`{\displaystyle \frac{\epsilon H}{\mathrm{\Gamma }(1+\epsilon )t^{d_t}}},\epsilon >0,d_t>1`$ (12)
The selection of signs (plus or minus) in (9)-(11) is determined by sign of $`\epsilon `$ and requirements of a regularization of integrals and selection of FD $`d_t`$ (greater or smaller unity). Let $`d_t(𝐫(t),t)<1`$. In this case from (9) follows (for $`H`$ do not containing explicit time dependence)
$$D_{+,t}^{d_t}H=\stackrel{~}{D}_{+,t}^{d_t}H\pm \frac{1}{2}H\frac{\epsilon (𝐫(t),t)}{t}+\frac{\epsilon H}{\mathrm{\Gamma }(1+\epsilon )t^{d_t}}$$
(13)
For $`t\mathrm{}`$, the basic contribution in the (13) imports corrections proportional to velocity of change FD $`d_t(𝐫(t),t)`$. The total energy conservative (in sense, that the Hamiltonian has not an explicit dependence at time ) systems now in space with multifractal time is not conservative systems. It changes can be at $`t\mathrm{}`$ of any sign, and depend on a sign of derivatives with respect to time from the fractional correction to dimension of time $`\epsilon `$. For $`\epsilon =0`$ the total energy of system is conserves and all relations coincide with known relations following from the dynamic equations of classical systems mechanics.
## V Liouville equation
The equation of the Liouville for a N-partial distribution function $`\rho (X,t)`$ ($`X`$\- are coordinate and impulses of particles of a $`6N`$-dimension phase space) is equivalent to the Hamilton equations for system from $`N`$ of classical particles and is invariant in relation to transformations
$$𝐫_i𝐫_i,𝐩_i𝐩_i,tt,i=1,2,3,..N$$
(14)
(equation is reversible). On set $`S_t`$ of multifractal time complete derivative $`\rho (X,t)`$ will reads:
$`D_{+,t}^{d_t}\rho (X,t)=\stackrel{~}{D}_{+,t}^{d_t}\rho (X,t)+{\displaystyle \frac{H}{𝐩_i}}{\displaystyle \frac{\rho (X,t)}{𝐫_i}}`$ (15)
$`{\displaystyle \frac{H}{𝐫_i}}{\displaystyle \frac{\rho (X,t)}{𝐩_i}}=\stackrel{~}{D}_{+,t}^{d_t}\rho (X,t)L\rho (X,t)`$ (16)
$`𝐩_i=D_{+,t}^{d_t}𝐫_i`$ (17)
where $`D_{+,t}^{d_t}\rho `$ is defined by (1), $`L`$-functional of the Liouville and differs from the functional of the known equation of the Liouville by replacement of ordinary derivatives with respect to time on GFD. For a demonstration of the irreversibility of expression (15) to transformations (14) we shall mark the following: the distribution functions $`\rho (X,t)`$ and $`\rho (X_0,t_0)`$ viewed in different moments $`t_0`$ and $`t`$ are connected by the relation
$$\rho (X,t)dX=\rho (X_0,t_0)dX_0$$
(18)
in which because of change of fractal dimension with a change time $`dXdX_0`$. Therefore $`\rho (X,t)\rho (X_0,t_0)`$ and $`\rho (X,t)`$ evolves with a time. Complete derivative $`D_{+,t}^{d_t}\rho (X,t)`$, in particular and in that connection, is not equal to zero. Intuitively it is clear, that derivative $`D_{+,t}^{d_t}\rho (X,t)`$ is determined by a functional from function $`\rho (X,t)`$ equal to zero at $`d_t=1+\epsilon =1`$. Let’s designate this functional describing change $`\rho (X,t)`$ owing to interaction with ”thermostat”, by $`\phi =\phi _0(\rho ,d_t,t)\epsilon (𝐫(t),t)`$. The role of the thermostat plays the set $`M_t`$ (being the carrier of a measure of a subsets of time points $`s_t`$ and belonging to one of spaces $`R^n`$), as was already marked. The appealing of the functional $`\phi `$ is caused not to interior processes happening with change of energy inside system, but is determined by different properties of time sets $`s_t`$ in different instant of time (change of dimension of $`s_t`$ with a time changes). Complete derivative in this case will be equal to a functional $`\phi `$ and (15) will reads as the equation
$`\stackrel{~}{D}_{+,t}^{d_t}\rho (X,t)`$ $`+`$ $`{\displaystyle \frac{H}{𝐩_i}}{\displaystyle \frac{\rho (X,t)}{𝐫_i}}{\displaystyle \frac{H}{𝐫_i}}{\displaystyle \frac{\rho (X,t)}{𝐩_i}}=`$ (19)
$`=`$ $`\phi _0(\rho ,d_t,t)\epsilon (𝐫(t),t)`$ (20)
The equation (19) is analogous of the Liouville equation of the classical systems with the time defined on multifractal sets. The analog of a collision integral in a right member (19) is stipulated by interaction with the carrier of a measure of multifractal set $`M_t`$ and is equal’s to zero if sets of time $`s_t`$ is substitutes by sets with topological dimension equal to unity.
## VI Production of an entropy
Let’s consider a classical system which is be found in an equilibrium state at the usual describing of the time.The production of an entropy in such system is equal to zero. Let’s consider the production of the entropy $`S=\rho (X,t)\mathrm{}n\rho (X,t)𝑑X`$ of same classical system defined on multifractal set of time points $`S_t`$ (for $`d_t=1\epsilon <1)`$:
$$D_{+,t}^{d_t}S=D_{+,t}^{d_t}[\rho \mathrm{}n\rho ]𝑑X$$
(21)
Permissible, as well as earlier, that $`d_t=1+\epsilon (𝐫(t),t)`$, $`|\epsilon 1|`$. As $`\rho (X,t)`$ has for the equilibrium system at $`d_t1`$ the complete GFD which is non-equal zero, the right member (21) is not equal to zero. It means, that equilibrium systems does not exist in space with multifractal time. Really, as
$`D_{+,t}^{d_t}S={\displaystyle D_{+,t}^{d_t}[\rho \mathrm{}n\rho ]𝑑X}`$ $``$ (22)
$`{\displaystyle \frac{\epsilon S}{t^{d_t}}}\pm {\displaystyle \frac{1}{2}}{\displaystyle \frac{}{t}}(\epsilon S)+{\displaystyle \frac{S}{t}}0`$ (23)
that (22) is an inequality and in a case $`\frac{S}{t}=0`$. For $`\frac{\epsilon }{t}=0`$ the production of the entropy is positive. For $`\frac{\epsilon }{t}0`$, $`\frac{S}{t}=0`$ the production of the entropy can have any sign (in particular, the entropy can be decreasing too). Let’s mark in that circumstance, that all new results for behavior of the entropy are stipulated by the multifractality of the time and disappear after transition to time with topological dimension equal to unity.
## VII About connection FD $`d_t`$ with Lagrangians of physical fields
In the monograph the following approximating connection of fractional dimension of time $`d_t(𝐫(t),t)`$ with a Lagrangian density $`L`$ of all physical fields in a point $`𝐫(t)`$ in an instant $`t`$ (see also , ) is obtained:
$$d_t(𝐫(t),t)=1+\underset{i}{}\beta _iL_i(𝐫(t),t)$$
(24)
where $`\beta _i`$ is dimensional numerical factors ensuring a zero dimension of products $`\beta _iL_i`$. In it is shown, that for coinciding the results founded on the (24) with results of the theory general relativity (GR) it is necessary (for gravitational forces) to choose $`\beta =\frac{2}{c^2}`$ ($`c`$\- speed of light). For correspondence with results of a quantum mechanics, for electric fields the $`\beta _e`$ has an order of magnitude $`\beta _e=(2mc^2)^1`$ ($`m`$ -is mass of particle or body creating the electric charge). The small differences FD from unity is satisfied by condition
$$\underset{i}{}\beta _iL_i=\epsilon 1$$
(25)
The connection $`\epsilon `$ with density of Lagrangians adds physical sense GFD (more in detail about it see ) and renders concrete relations obtained in the previous paragraphs.
## VIII Conclusions
The present note is devoted to the appendix of idea of the multifractal time offered in , for researching of the problem of an irreversibility in large classical systems consisting from an identical objects. The following results, leading from this paper, on our sight, are essential:
1. The equations describing behavior of conservative systems (in usual time), are irreversible in space with multifractal time;
2. The neglecting by the fractionality of dimension of time and transition in space with topological dimension of time equal to unity allows to receive the known reversible equations of classical dynamics;
3. In space with the multifractal time there are no invariable objects, since the GFD with respect to stationary values are not equal to zero. From the physical point of view it is the reflection of non-stationary of the Universe with multifractal time. Last statement corresponds to mathematical exposition of behavior of physical objects and not contradict the exposition of the Einstein type of Universe , in which, in connection with its expansion, there are no invariable objects;
4. The quantity of the fractional additional to topological dimension of time member $`\epsilon `$ is determined by physical fields and depends on the density of energy that presents in the given moment in the given point of space . At the small densities of energy the corrections are very small. So, for the gravitational fields at distances more larger that gravitational radius (for example, for FD created by mass of the Earth on a surface of Earth) and for electric fields on atomic distances the value of $`\epsilon `$ is equal $`\epsilon 10^8`$. Therefore the multifractal nature of time ($`d_t1+\epsilon `$) does not contradicts an existing experimental data. |
warning/0002/astro-ph0002317.html | ar5iv | text | # Analytic solutions for coupled linear perturbations
## 1 Introduction
Several methods (e.g., Croft et al. 1998; Gnedin 1998; Nusser & Haehnelt 1999) have recently been proposed for extracting information on the mass density field from the Lyman-$`\alpha `$ forest. The underlying physical picture behind these method is that the absorbing Neutral hydrogen in the low density intergalactic medium (IGM) is tightly related to the mass density field on scales larger than the Jeans length. Below the Jeans length gas pressure segregates the baryons from the total mass fluctuations. On scales near the Jeans length, the evolution of the baryonic perturbation can affect estimates of the clustering amplitude from observations of the Lyman-$`\alpha `$ forest.
Hydrodynamical simulations (Petitjean, Mücket & Kates 1995; Zhang, Anninos & Norman 1995; Hernquist et al. 1996; Miralda-Escudé et al. 1996, Theuns et al. 1998) and semi-analytic models (e.g., Bi et. al. 1992, Gnedin & Hui 1996) of the IGM have been successful at explaining observations of the forest. Despite the success of the simulations, it is usually difficult to use them to study in detail the evolution of the gas below Jeans length (Theuns et al. 1998). The equations governing the evolution of baryons and dark matter in the nonlinear regime are extremely difficult to solve, even for special configurations like spherical collapse. Fortunately, since most of the IGM is of moderate density, linear analysis can be a suitable tool for understanding the evolution of the baryons (Gnedin & Hui 1998). Here we derive analytic solutions to the linear equations in a flat universe without a cosmological constant. Although the linear equations can readily be numerically integrated under a variety of conditions (Gnedin & Hui 1998), analytic treatment offers better understanding of the equations. Further, the paucity of analytic solutions makes their pursue worthwhile, even if tedious at times. The analytic solutions we derive here are subject to the condition that the baryonic and dark matter trace the same density and velocity fields before a sudden reionization epoch. After reionization the temperature of the IGM is assumed to be inversely proportional to the scale factor so that the comoving Jeans length is constant.
The paper is organized as follows. In section 2 we cast the equations in the form of linear differential equations with constant coefficients. In Section 3 we present the solutions to these equations for several cases. We conclude in Section 4.
## 2 The linear equations
Let $`\delta _x(t,k)`$ and $`\delta _b(t,k)`$ be , respectively, the Fourier modes of baryonic and dark matter density fluctuations. Let also $`f_X`$ and $`f_b=1f_x`$ be the the mean mass fractions of these two types of matter. We will restrict the analysis to perturbations in a flat universe without a cosmological constant. The linear equations governing the evolution of $`\delta _b`$ and $`\delta _x`$ are (e.g., Bi et. al. 1992, Padmanabhan 1993, Gnedin & Hui 1998),
$`{\displaystyle \frac{\mathrm{d}^2\delta _x}{\mathrm{d}t^2}}+2H{\displaystyle \frac{\mathrm{d}\delta _x}{\mathrm{d}t}}`$ $`=`$ $`{\displaystyle \frac{3}{2}}H^2\left(f_x\delta _x+f_b\delta _b\right)`$ (1)
$`{\displaystyle \frac{\mathrm{d}^2\delta _b}{\mathrm{d}t^2}}+2H{\displaystyle \frac{\mathrm{d}\delta _b}{\mathrm{d}t}}`$ $`=`$ $`{\displaystyle \frac{3}{2}}H^2\left(f_x\delta _x+f_b\delta _b\right){\displaystyle \frac{3}{2}}H^2\left({\displaystyle \frac{k}{k_J}}\right)^2\delta _b,`$ (2)
where $`a(t)t^{2/3}`$ is the scale factor, $`H(t)=2/(3t)`$ the Hubble function, and $`k_J`$ is the comoving Jeans wavenumber related to the speed of sound $`c_s`$ and the mean total (baryons plus dark) density, $`\overline{\rho }=3H^2/(8\pi G)`$, by
$$k_J=\frac{a}{c_s}\sqrt{4\pi G\overline{\rho }}=\sqrt{\frac{3}{2}}\frac{aH}{c_s}.$$
(3)
For $`f_x`$ close to unity the gravity of the baryons is negligible and $`\delta _xa`$. If also $`k_J=const`$ (i.e., $`c_s^21/a`$) and we impose, at some initial time $`t_i`$, the condition $`\delta _b(t_i)=\delta _x(t_i)/(1+(k/k_J)^2)`$, then $`\delta _b(t)=\delta _x(t)/(1+(k/k_J)^2)`$ at any later time, $`t`$. In the next section we will show how to solve these equations with $`0<lef_x`$ assuming that $`\delta _b(t)=\delta _x(t)`$ for $`tt_i`$, which is appropriate for a sudden reionization of the IGM at $`t_i`$.
In order to solve these equations we work with a new time variable $`\tau \mathrm{ln}(a)`$ instead of $`t`$ (Nusser & Colberg 1998). In terms of $`\tau `$ the equations become,
$`{\displaystyle \frac{\mathrm{d}^2\delta _x}{\mathrm{d}\tau ^2}}+{\displaystyle \frac{1}{2}}{\displaystyle \frac{\mathrm{d}\delta _x}{\mathrm{d}\tau }}`$ $`=`$ $`{\displaystyle \frac{3}{2}}\left(f_x\delta _x+f_b\delta _b\right)`$ (4)
$`{\displaystyle \frac{\mathrm{d}^2\delta _b}{\mathrm{d}\tau ^2}}+{\displaystyle \frac{1}{2}}{\displaystyle \frac{\mathrm{d}\delta _b}{\mathrm{d}\tau }}`$ $`=`$ $`{\displaystyle \frac{3}{2}}\left(f_x\delta _x+f_b\delta _b\right){\displaystyle \frac{3}{2}}\kappa ^2\delta _b,`$ (5)
where have defined $`\kappa =k/k_J`$.
When $`\kappa `$ is constant, the differential equations (5) are linear with constant coefficients and they can be solved by means of Laplace transformation. Since Laplace transforms are seldom used in cosmological studies, it seems prudent to briefly review their basic properties which are relevant to us. We refer the reader to Arfkin (1985) and references therein for mathematical details. The Laplace transform, $`f(s)`$, of a function $`F(t)`$, where $`t0`$, is defined as
$$f(s)\{F(t)\}=_0^{\mathrm{}}\mathrm{exp}\left(st\right)F(t)dt.$$
(6)
We will need the Laplace transforms of first and second derivatives of a function. Using (6) these transforms can be related to the $`f(s)`$ by
$`\{F^{}(t)\}`$ $`=`$ $`sf(s)F(0)`$ (7)
$`\{F^{\prime \prime }(t)\}`$ $`=`$ $`s^2f(s)sF(0)F^{}(0),`$ (8)
where the prime and double prime denote first and second order derivatives, respectively. The Bromwich integral expresses $`F(t)`$ in terms of $`f(s)`$ as
$$F(t)=\frac{1}{2\pi i}_{\gamma i\mathrm{}}^{\gamma +i\mathrm{}}\mathrm{exp}\left(st\right)f(s)ds$$
(9)
where $`i=\sqrt{1}`$ and $`\gamma `$ is a real number chosen so that all poles of $`f(s)`$ lie, in the complex plane, to the left of the vertical line defining the integration path. Therefore, by the residue theorem we have
$$F(t)=\left[\mathrm{residues}\mathrm{of}\mathrm{exp}\left(st\right)f(s)\right]$$
(10)
As an example consider $`f(s)=1/(ss_1)(ss_2)`$ which has two simple poles at $`s=s_1`$ and $`s_2`$. The residues of $`\mathrm{exp}(st)f(s)`$ at these poles are $`\mathrm{exp}(s_1t)/(s_1s_2)`$ and $`\mathrm{exp}(s_2t)/(s_1s_2)`$ so that, by (10), $`F(t)=[\mathrm{exp}(s_1t)\mathrm{exp}(s_2t)]/(s_1s_2)`$. If $`s_1=s_2`$, the function has a pole of order two at $`s_1`$. The residue in this case is $`\mathrm{d}[(ss_1)^2\mathrm{exp}(st)f(s))]/\mathrm{d}s`$ evaluated at $`s=s_1`$. Therefore, $`F(t)=t\mathrm{exp}(s_1t)`$.
## 3 The solutions
Denote by $`\mathrm{\Delta }_x`$ and $`\mathrm{\Delta }_x`$ the Laplace transforms of $`\delta _x`$ and $`\delta _b`$, respectively. By taking the Laplace transform of (5) we can obtain relations between $`\mathrm{\Delta }_x`$ and $`\mathrm{\Delta }_b`$. The initial conditions are contained in the Laplace transforms of the first and second derivatives of the densities. So first we have to specify in mathematical terms our choice for the initial conditions. For simplicity of notation we fix the initial conditions at $`\tau =1`$ assuming that before that time the temperature of the baryonic fluid is zero, i.e., $`\kappa =0`$. The initial conditions are fixed by the values of $`\delta _x`$ and $`\delta _b`$ and their first derivatives at $`\tau =1`$. Before $`\tau =1`$, we have $`\delta _x=\delta _b=\mathrm{exp}(\tau )`$, ignoring the decaying mode and setting arbitrarily $`\delta _x(\tau =1)=1`$. The first derivatives of $`\delta _x`$ and $`\delta _b`$ are therefore equal to unity at $`\tau =1`$. This fixes the initial conditions necessary for solving (5). Although we will present solutions satisfying only these initial conditions, we will, for completeness, write the Laplace transformation of (5) for $`\delta _b=\alpha \delta _x`$ and $`\mathrm{d}\delta _b/\mathrm{d}\tau =\alpha \mathrm{d}\delta _x/\mathrm{d}\tau `$, at $`\tau =1`$. Then the Laplace transformation of (5) yields
$`\left(s^2+{\displaystyle \frac{1}{2}}s\right)\mathrm{\Delta }_x`$ $`=`$ $`{\displaystyle \frac{3}{2}}\left(f_x\mathrm{\Delta }_x+f_b\mathrm{\Delta }_b\right)+s+{\displaystyle \frac{3}{2}}`$ (11)
$`\left(s^2+{\displaystyle \frac{1}{2}}s\right)\mathrm{\Delta }_b`$ $`=`$ $`{\displaystyle \frac{3}{2}}\left(f_x\mathrm{\Delta }_x+f_b\mathrm{\Delta }_b\right){\displaystyle \frac{3}{2}}\kappa ^2\mathrm{\Delta }_b+\alpha s+{\displaystyle \frac{3}{2}}\alpha ,`$ (12)
where have have used (8) to computed the transforms of the first and second derivatives of $`\delta _x`$ and $`\delta _b`$. For $`f_x=1`$ the first of these equations yields
$$\mathrm{\Delta }_x=\frac{1}{s1},$$
(13)
which is the Laplace transform of $`\mathrm{exp}(\tau )`$. If we take $`\alpha =(1+\kappa ^2)^1`$ and substitute (13) in the second equation of (12) we get
$$\mathrm{\Delta }_b=\frac{1}{s1}\frac{1}{1+\kappa ^2}=\frac{\mathrm{\Delta }_x}{1+\kappa ^2},$$
(14)
which leads to the well known solution $`\delta _b=\delta _x/(1+\kappa ^2)`$.
Subsequently we will present solutions only for $`\alpha =1`$. In this case, equations (12) yield
$$\mathrm{\Delta }_x=\frac{\left(s+\frac{3}{2}\right)\left(\frac{3}{2}\kappa ^2+s^2+\frac{1}{2}s\right)}{\left(s^2+\frac{1}{2}s\right)\left(\frac{3}{2}\kappa ^2+s^2+\frac{1}{2}s\frac{3}{2}\right)\frac{9}{4}f_x\kappa ^2},$$
(15)
and
$$\mathrm{\Delta }_b=\frac{s^2+\frac{1}{2}s}{\frac{3}{2}\kappa ^2+s^2+\frac{1}{2}s}\mathrm{\Delta }_x.$$
(16)
Before solving these equations for any value of $`f_x`$ in the range 0–1, it is instructive to examine the solutions for the special values $`f_x=1`$ and $`0`$.
### 3.1 Case I: $`f_x=1`$
In this case $`\mathrm{\Delta }_x=(s1)^1`$ and equation (16) can be written in the form
$$\mathrm{\Delta }_b=\frac{s^2+\frac{1}{2}s}{(s1)(ss_{_{}})(ss__+)}$$
(17)
where $`s__\pm `$ are the roots of $`3\kappa ^2/2+s^2+s/2`$. They are given by
$$s__\pm =\frac{1}{4}(1\pm \chi );\chi ^2=124\kappa ^2$$
(18)
We will deal with the case $`\chi ^2=0`$ at the end of this subsection. For $`\chi ^20`$, all three poles of $`\mathrm{\Delta }_b`$ are simple and so, by (10), its inverse transform is
$$\delta _b=\frac{\mathrm{exp}\left(\tau \right)}{1+\kappa ^2}+\frac{1}{s_{_{}}s__+}\frac{\kappa ^2}{1+\kappa ^2}\left[\left(s_{_{}}1\right)\mathrm{exp}\left(s__+\tau \right)\left(s__+1\right)\mathrm{exp}\left(s_{_{}}\tau \right)\right]$$
(19)
For $`\chi ^2>0`$ the roots $`s__\pm `$ are real and the result is
$$\delta _b=\frac{\mathrm{exp}\left(\tau \right)}{1+\kappa ^2}+\frac{1}{2\chi }\frac{\kappa ^2}{1+\kappa ^2}\left[\left(\chi 5\right)\mathrm{exp}\left(\frac{\chi }{4}\tau \right)+\left(5+\chi \right)\mathrm{exp}\left(+\frac{\chi }{4}\tau \right)\right]\mathrm{exp}\left(\frac{\tau }{4}\right)$$
(20)
The first term is the solution given in the previous section for $`\alpha =1/(1+\kappa ^2)`$. The maximum value $`\chi ^2`$ attains is unity, so the second and third terms are always decaying. If $`\chi ^2<0`$, then (19) gives the solution
$$\delta _b=\frac{\mathrm{exp}\left(\tau \right)}{1+\kappa ^2}+\frac{1}{\stackrel{~}{\chi }}\frac{\kappa ^2}{1+\kappa ^2}\left[5\mathrm{sin}\left(\frac{\stackrel{~}{\chi }}{4}\tau \right)+\chi \mathrm{cos}\left(\frac{\stackrel{~}{\chi }}{4}\tau \right)\right]\mathrm{exp}\left(\frac{\tau }{4}\right)$$
(21)
where $`\stackrel{~}{\chi }`$ is the imaginary part of $`\chi `$. The solution shows an oscillatory behavior with a period of $`16\pi /\stackrel{~}{\chi }`$. The envelope of these oscillations decays like $`\mathrm{exp}(\tau /4)t^{1/6}`$.
We deal now with the case $`\chi ^2=0`$, which occurs for $`\kappa ^2=1/24`$. Here special care is needed because $`s_{_{}}=s__+=1/4`$. However the contribution of the pole at $`s=1`$ to the Bromwich integral remains unchanged and the contribution of the second order pole at $`s=1/4`$ is simply the first derivative of $`\mathrm{exp}(s\tau )(s+1/4)^2\mathrm{\Delta }_b`$ at $`s=1/4`$. The result is
$$\delta _b=\frac{24}{25}\mathrm{exp}\left(\tau \right)+\frac{1}{20}\left(\tau +\frac{4}{5}\right)\mathrm{exp}\left(\frac{\tau }{4}\right).$$
(22)
The first term on the left is the familiar $`\delta _x/(1+\kappa ^2)`$ evaluated at $`\kappa ^2=1/24`$. The expression can also be derived by taking the limit $`\chi ^20`$ in either (20) or (21).
In the limit $`\tau \mathrm{}`$ the ratio $`\delta _b/\delta _x`$ is $`1/(1+\kappa ^2)`$.
### 3.2 Case II: $`f_x=0`$
This is equivalent to ignoring the gravity of the dark matter. Of course here only the behavior of the perturbation in the baryons is relevant since the dark matter plays no role. However, for the sake of completeness and comparison with other situations we will solve for the dark matter fluctuations as well.
We first find the solution for $`\delta _x`$. If $`f_x=0`$, we can express (15) in terms of $`s__\pm `$, the roots of $`3\kappa ^2/2+s^2+s/23/2`$, as
$$\mathrm{\Delta }_x=\frac{\left(s+\frac{3}{2}\right)\left(\frac{3}{2}\kappa ^2+s^2+\frac{1}{2}s\right)}{s\left(s+\frac{1}{2}\right)\left(ss_{_{}}\right)\left(ss__+\right)},$$
(23)
where
$$s__\pm =\frac{1}{4}(1\pm \chi );\chi ^2=2524\kappa ^2$$
(24)
If $`\kappa 1`$, then the function $`\mathrm{\Delta }_b`$ has three simple poles at $`s=0`$, $`s_{_{}}`$, and $`s__+`$. So for $`\kappa 1`$ and $`\chi ^2>0`$, Bromwich integral yields
$$\delta _x=\frac{\kappa ^2}{1\kappa ^2}\left[2\mathrm{exp}\left(\frac{\tau }{2}\right)3\right]+\frac{1}{2\chi }\frac{1}{1\kappa ^2}\left[\left(\chi 5\right)\mathrm{exp}\left(\frac{\chi }{4}\tau \right)+\left(\chi +5\right)\mathrm{exp}\left(+\frac{\chi }{4}\tau \right)\right]\mathrm{exp}\left(\frac{\tau }{4}\right).$$
(25)
The expression for $`\delta _b`$ is
$$\delta _b=\frac{1}{2\chi }\left[\left(\chi 5\right)\mathrm{exp}\left(\frac{\chi }{4}\tau \right)+\left(\chi +5\right)\mathrm{exp}\left(\frac{\chi }{4}\tau \right)\right]\mathrm{exp}\left(\frac{\tau }{4}\right).$$
(26)
When $`\kappa =1`$ the solution can be found either by taking the limit $`\kappa 1`$ in the previous two expressions or by direct evaluation of the Bromwich integral with a two poles of order two at $`s=0`$ and $`1/2`$. The result is
$$\delta _x=27+28\mathrm{exp}\left(\frac{\tau }{2}\right)+3\tau \left[3+2\mathrm{exp}\left(\frac{\tau }{2}\right)\right]\mathrm{and}\delta _b=32\mathrm{exp}\left(\frac{\tau }{2}\right).$$
(27)
This implies that $`\delta _x`$ grows linearly with $`\tau `$ at late times while $`\delta _b`$ reaches an asymptotic value of $`3`$.
The oscillatory behavior of $`\delta _b`$ and $`\delta _x`$ appears when $`\chi ^2<0`$, i.e., for $`\kappa ^2>25/24`$. The expressions in this case can be obtained by replacing $`\chi `$ with $`i\stackrel{~}{\chi }`$ in (25) and (26). For $`\kappa ^2>>25/24`$, the solution coincides with that given in Padmanabhan (1993)
In the limit $`\tau \mathrm{}`$ the ratio $`\delta _b/\delta _x`$ is $`1\kappa ^2`$ for $`\kappa 1`$ and, zero otherwise.
### 3.3 Case III: $`0<f_x<1`$
Again we first derive $`\delta _x`$. The denominator and numerator in (15) do not have any common roots. Then the poles of $`\mathrm{\Delta }_x`$ are the roots of the denominator. We find these roots as follows. Denote $`y(s)=s^2+s/23/2`$ and equate denominator to zero to obtain
$$y^2+\frac{3}{2}y\left(1+\kappa ^2\right)+\frac{9}{4}f_b\kappa ^2=0$$
(28)
where we have used $`1f_x=f_b`$. This equation is satisfied for the following values of $`y`$,
$$y_{p,m}=\frac{3}{4}(1+\kappa ^2)(1\pm \mathrm{\Xi });\mathrm{\Xi }^2=14f_b\frac{\kappa ^2}{\left(1+\kappa ^2\right)^2},$$
(29)
where the subscripts $`p`$ and $`m`$ correspond to the plus and minus sign, respectively. So the roots of the denominator in (15) are the values of $`s`$ which make $`y(s)=y_{p,m}`$. Let $`s_{p,\pm }`$ and $`s_{_{m,\pm }}`$ be the roots of $`y(s)y_p=0`$ and $`y(s)y_m=0`$, respectively. They are given by
$`s_{_{m,\pm }}={\displaystyle \frac{1}{4}}\left(1\pm \chi _m\right)`$ ; $`\chi _{m}^{}{}_{}{}^{2}=2512\left(1+\kappa ^2\right)\left(1+\mathrm{\Xi }\right),`$ (30)
$`s_{_{p,\pm }}={\displaystyle \frac{1}{4}}\left(1\pm \chi _p\right)`$ ; $`\chi _{p}^{}{}_{}{}^{2}=2512\left(1+\kappa ^2\right)\left(1\mathrm{\Xi }\right).`$ (31)
Excluding the values $`f_b=0`$ and $`1`$, which we have considered in the previous subsections, we have $`0<\mathrm{\Xi }^2<1f_b`$. This ensures that all four roots, $`s_{_{m,\pm }}`$ and $`s_{_{p,\pm }}`$, are distinct. Also, $`1<\chi _p^2<25`$ for any $`\kappa `$, so the roots $`s_{_{p,\pm }}`$ are real. On the other hand, $`0<\chi _m^2<1`$ when $`\kappa ^2<25/(600576f_b)`$, and negative otherwise. So $`s_{_{m,\pm }}`$ can be complex.
For $`\chi _m^2>0`$, the Bromwich integral yields
$$\delta _x=\frac{24\kappa ^2+\chi _p^21}{2\chi _p\left(\chi _p^2\chi _m^2\right)}[(\chi _p5)\mathrm{exp}(\frac{\chi _p}{4}\tau )+(\chi _p+5)\mathrm{exp}(+\frac{\chi _p}{4}\tau )]\mathrm{exp}(\frac{\tau }{4})+(\chi _p\chi _m),$$
(32)
with the second term on the r.h.s is obtained interchanging $`\chi _p`$ and $`\chi _m`$ in the first term. This result can be extended to $`\chi _m^2<0`$ by writing $`\chi _m=i\stackrel{~}{\chi }_m`$ where $`\stackrel{~}{\chi }_m`$ is real.
Using (16) we similarly obtain $`\delta _b`$
$$\delta _b=\frac{\chi _p^21}{2\chi _p\left(\chi _p^2\chi _m^2\right)}[(\chi _p5)\mathrm{exp}(\frac{\chi _p}{4}\tau )+(\chi _p+5)\mathrm{exp}(+\frac{\chi _p}{4}\tau )]\mathrm{exp}(\frac{\tau }{4})+(\chi _p\chi _m).$$
(33)
To visualize these solutions, we show in Fig.1 the density evolution for various values of $`f_x`$ and $`\kappa `$. The dark matter curve for $`f_x=0.9`$ is very close to $`\mathrm{exp}(\tau )`$. For $`\kappa =5`$ the baryonic perturbations for both values of $`f_x`$ show oscillations with similar period and amplitude. This is simply because for high $`\kappa `$, the evolution is mainly dictated by pressure forces.
In Fig.2 we the solid lines represent the ratio $`\delta _b(\kappa )/\delta _x(\kappa )`$ as a function of $`\kappa `$ at different $`\tau `$ designated in the plot by the redshift, $`z`$. The initial conditions are satisfied at $`z=6`$ and $`f_x=0.9`$ was taken. Also plotted, as the dotted line in each panel, the function $`1/(1+\kappa ^2)`$ which represents the limiting solution as $`\tau \mathrm{}`$. The analytic curves show more oscillations as they get closer to the limiting solution, when the redshift is decreased.
In many applications (e.g., Bi et al 1992, Bi & Davidsen 1997, Nusser & Haehnelt 2000) the limiting ratio $`1/(1+\kappa ^2)`$ is often used to filter the mass power spectrum in order to generate density fluctuations in the gas. As pointed out by Gnedin & Hui (1998) this may lead to a significant bias in statistics of the gas density. As an illustration of this bias we compute the rms values of the gas density by filtering a scale free mass power spectrum of slope $`n=2.5`$ with the ratio $`(\delta _b/\delta _x)^2`$ given from the analytic, and the limiting solutions, respectively. Fig.3 shows the ratio of the former to the latter rms value. As in the previous figure, $`f_x=0.9`$ and initial conditions satisfied at $`z=6`$. Using the filter $`1/(1+\kappa ^2)`$ can seriously underestimate the amplitude of gas density fluctuations. Only when we approach $`z=0`$ the ratio gets close to to unity.
## 4 Summary
We have found analytic solutions to the linear equations governing the evolution of baryonic and dark matter under four assumptions. First, the Universe is flat without cosmological constant. Second, sudden reionization of the IGM. Third, the temperature of the low density IGM drops like $`1/a`$ so that the comoving Jeans length is time independent. Fourth, before reionization the IGM is cold and the baryonic and dark matter trace the same density and velocity fields.
Of these assumptions, only the fourth has a physical basis, at least before any heating has occurred and when the IGM temperature is low. This is also the only assumption which if changed, the equations can still be readily solved by Laplace transformation. Unfortunately, relaxing any of the other assumptions complicates the analytic treatment of equations (5) by means of Laplace transformation. For example, suppose that the Jeans wavenumber changes with time according to $`a^\beta `$. Then the Laplace transformation of the term involving $`k_J`$ will yield $`\mathrm{\Delta }_b`$ at $`(s+\beta )`$ while other terms involve $`\mathrm{\Delta }_b(s)`$.
Yet the solutions can be useful for semi-analytic modeling of the IGM. They offer a convenient improvement over the commonly used filter $`1/[1+(k/k_J)^2]`$ for generating gas fluctuations associated with a given mass density field.
The analytic solutions presented here were verified by a comparison with the solutions obtained by numerical integration of equations (2). All numerical and analytic solutions agreed up to the numerical accuracy.
## 5 acknowledgement
This research was supported by a grant from the Israeli Science Foundation. |
warning/0002/hep-ph0002185.html | ar5iv | text | # References
1.Introduction
As known, the dynamics of the neutral fermions with anomalous magnetic moments are described by Dirac equation with non-minimal coupilngs of neutral fermions with electromagnetic field ,:
$$(\widehat{k}m+\mu (\stackrel{}{\mathrm{\Sigma }}\stackrel{}{B}i\stackrel{}{\alpha }\stackrel{}{E})+iq(\stackrel{}{E}\stackrel{}{B})\gamma _5)\psi (k)=0,$$
(1)
where $`\mu `$-is anomalous magnetic moment, $`\frac{1}{2}\stackrel{}{\mathrm{\Sigma }}`$ is spin operator,and defined by formula (21,21) of operator $`\stackrel{}{\alpha }`$ is defined by formula (21,20) of .The last term described by lagrangian:
$$L=iq(\stackrel{}{E}\stackrel{}{H})\overline{\psi }\gamma _5\psi $$
(2)
in Dirac equations is obtained in and depends only on $`r`$ in case of monopole (for monopoles see e.g. and references therein) which have both electric and magnetic fields.
In this article we consider cylindrically symmetric bound states and resonances of particles with anomalous magnetic moments.
It is of interest to consider several special cases in particular:
1)Particle energy levels in pure electric axial field case ($`\stackrel{}{E}=(E_r(r),0,0)`$) which created e.g. by homogeneously charged cylinder.
In accordance with we have:
$$E_r=\frac{2\sigma }{r}atr>R$$
(3)
and
$$E_r=\frac{2\sigma r}{R^2}atr<R$$
(4)
where $`\sigma `$ is density of charge of the unit of length of the cylinder.
Also we consider second order equations for axially symmetric $`Z_0(r)`$-boson field, radial axially symmetric magnetic field obtain second order equations for case of bound states of fermions in the radial electric and magnetic fields , and also obtain second order equations for case of pseudoscalars described by last term in Dirac equation (1).
Radial axially symmetric electric field
In ref. has been obtained the system of equations for radial functions (see below formulas (20)-(23)in which we add also radial magnetic field besides radial electric field).Below will be presented the second order equations which obtained after excluding two of four radial functions:
$$(\frac{1}{r}\frac{d}{dr}r\frac{d}{dr}+ϵ^2m^2p_z^24\pi \rho (\frac{l}{r}\mu E)^2)f_1(r)+2i\mu Ep_zf_2(r)=0$$
(5)
$$(\frac{1}{r}\frac{d}{dr}r\frac{d}{dr}+ϵ^2m^2p_z^24\pi \rho (\frac{(l+1)}{r}+\mu E)^2)f_2(r)2i\mu Ep_zf_1(r)=0$$
(6)
The presence of the $`i`$ mean that one of the radial functions must be purely imagine.
Analogous system of equations obtained if we esclude $`\varphi `$ and consider $`\chi `$ as $`\chi =(e^{il\varphi }f_3(r),e^{i(l+1)\varphi }f_4(r))`$:
$$(\frac{1}{r}\frac{d}{dr}r\frac{d}{dr}+ϵ^2m^2p_z^2+4\pi \rho (\frac{l}{r}\mu E)^2)f_3(r)+2i\mu Ep_zf_4(r)=0$$
(7)
$$(\frac{1}{r}\frac{d}{dr}r\frac{d}{dr}+ϵ^2m^2p_z^2+4\pi \rho (\frac{(l+1)}{r}+\mu E)^2)f_4(r)2i\mu Ep_zf_3(r)=0$$
(8)
In case of electric field created by charged line or charged cylinder at $`r>a`$ we obtain:
$$(\frac{1}{r}\frac{d}{dr}r\frac{d}{dr}+ϵ^2m^2p_z^2\frac{(l2\mu \sigma )^2}{r^2})f_1(r)+\frac{4i\mu \sigma p_z}{r}f_2(r)=0$$
(9)
$$(\frac{1}{r}\frac{d}{dr}r\frac{d}{dr}+ϵ^2m^2p_z^2\frac{(l+1+2\mu \sigma )^2}{r^2})f_2(r)\frac{4i\mu \sigma p_z}{r}f_1(r)=0$$
(10)
Inside homogeneously charged cylinder we obtain:
$$(\mathrm{\Omega }(\frac{(l)^2}{r^2}+\frac{4\mu ^2\sigma ^2r^2}{a^4}+\frac{4\mu \sigma l}{a^2}))f_1(r)+\frac{4i\mu \sigma p_zr}{a^2}f_2(r)=0$$
(11)
$$(\mathrm{\Omega }(\frac{(l+1)^2}{r^2}+\frac{4\mu ^2\sigma ^2r^2}{a^4}+\frac{4\mu \sigma (l+1)}{a^2}))f_2(r)\frac{4i\mu \sigma p_z}{r}f_1(r)=0$$
(12)
where $`\mathrm{\Omega }=\frac{1}{r}\frac{d}{dr}r\frac{d}{dr}+ϵ^2m^2p_z^2`$ we see that at $`p_z=0`$ both equations decouples and every of them is the same as for 2-dimensional harmonic oscillator and in accordance with we have the following energy levels:
$$ϵ^2=m^2+\frac{4\mu \sigma }{a^2}+\frac{8|\mu \sigma |}{a^2}(n_r+\frac{l+|l|+1}{2})$$
(13)
Analogously for second equation in result for energy levels we must replace:
$`l(l+1)`$
The particle is localized at distances $`r_H(\frac{4\mu \sigma }{a^2})^{\frac{1}{2}}`$ or smaller. Of course our consideration is available only if size of the particle localization is essentially smaller than the radius of the cylinder($`r_H<<a`$).The wave functions at large distances are suppressed by exponent $`e^{\frac{r^2}{r_H^2}}`$.
Above was considered energy levels of neutrons in case infinite radius of the cylinder. It is of interest to consider also case of the finite radius of the cylinder.
For this purpose we must find wave function inside and outside of the cylinder and in this case energy levels of neutrons will be defined from condition:
$$\frac{R_{r>a}^{}(a)}{R_{r>a}(a)}=\frac{R_{r<a}^{}(a)}{R_{r<a}(a)}.$$
(14)
Inside cylinder as we seen above we obtain equation which is equivalent to the 2-dimensional harmonic oscillator and radial wave functions are expresses through degenerate hypergeometric function:
$$R_{r<a}(r)=C_1x^{\frac{|l|}{2}}e^{\frac{x}{2}}F(\frac{l+|l|+1}{2}\frac{ϵ^2m^24\pi \mu \rho }{2m\omega },|l|+1,x)$$
(15)
where $`\omega =\frac{4|\mu \sigma |}{ma^2}`$,$`x=m\omega r^2`$. Outside of the cylinder radial wave functions are expresses through Macdonald’s function:
$$R_{r<a}(r)=C_2K_{|l+2\mu \sigma |}(\sqrt{|ϵ^2m^2|}r)$$
(16)
Below we consider case of the charged line.
In the limit:
$$\mu \sigma <<1$$
(17)
we obtain Coulomb-like spectrum for energy levels of neutrons in the field of charged line.
Indeed, if $`l`$ is not equal to the $`0,1`$ we can put in equations (5), (6) $`l2\mu \sigma l`$ in the limit $`\mu \sigma <<1`$.At $`l=0,1`$ must be
$$\frac{\mu ^2\sigma ^2}{r_B^2}<<\frac{\mu \sigma p_z}{r_B}$$
(18)
where $`r_B=\frac{1}{\mu \sigma p_z}`$ is Bohr radius (see below).Substituting $`r_B`$ in (18) we obtain again the condition $`\mu \sigma <<1`$.Thus, in the limit $`\mu \sigma <<1`$ at $`l=0,1`$ term $`=(0+\mu \sigma )^2r^2`$ must be neglected, and we obtain for all $`l`$ the following equations:
$$(\frac{1}{r}\frac{d}{dr}r\frac{d}{dr}+ϵ^2m^2p_z^2+\frac{l^2}{r^2})f_1(r)+\frac{4\mu \sigma p_z}{r}if_2(r)=0$$
(19)
$$(\frac{1}{r}\frac{d}{dr}r\frac{d}{dr}+ϵ^2m^2p_z^2+\frac{(l+1)^2}{r^2})if_2(r)+\frac{4\mu \sigma p_z}{r}f_1(r)=0$$
(20)
This equations are similar to equations obtained in where has been consider non-relativistic neutrons energy levels in the magnetic field $`\stackrel{}{H}=(0,H_\varphi (r)=2I/r,0)`$ if we making the following replacement:
$$I\frac{\sigma p_z}{m}$$
(21)
Thus, energy levels is also same and in accordance with result of and substitution (21) are defines by quantum number $`n`$:
$$ϵ_n^2=m^2+p_z^2\frac{\mu ^2\sigma ^2p_z^2}{n^2}$$
(22)
It must be noted, our consideration is relativistic,only assumption $`\mu \sigma <<1`$ has been used. In the near future we will present the solution where $`\mu \sigma `$ is not small.
Radial Axially Symmetric Magnetic +Radial Axially Symmetric electric field and bound states of neutral fermions with an anomalous magnetic moments
In this paper we consider also bound states of neutral fermions with an anomalous magnetic moments in radial axially symmetric magnetic field $`\stackrel{}{H}=(\frac{x}{\sqrt{x^2+y^2}}H(r),\frac{y}{\sqrt{x^2+y^2}}H(r),0)`$ which analogously to the above considered electric field in case of cylinder has the following form:
$$H_r(r)=\frac{2\sigma _m}{r}atr>R$$
(23)
and
$$H_r(r)=\frac{2\sigma _mr}{R^2}atr<R$$
(24)
where $`\sigma _m`$ is density of magnetic charge of the unit of length of the cylinder.
$$(ϵm)f_1+\mu Hf_2p_3f_3i(\frac{d}{dr}+\frac{l}{r}\mu E)f_4=0$$
(25)
$$\mu Hf_1+(ϵm)f_2i(\frac{d}{dr}\frac{l1}{r}\mu E)f_3+p_3f_4=0$$
(26)
$$p_3f_1+(\frac{d}{dr}+\frac{l}{r}+\mu E)f_2(ϵ+m)f_3+\mu Hf_4=0$$
(27)
$$i(\frac{d}{dr}\frac{l1}{r}+\mu E)f_1p_3f_2+\mu Hf_3(ϵ+m)f_4=0$$
(28)
It is interesting to notice that equations for radial magnetic field is similar to equations for magnetic field $`\stackrel{}{H}=(0,H_\varphi (r),0)`$ considered in .
In non-relativistic approximation we have the following system of equations (instead Dirac equation has been used Pauli equation with non-relativistic spinor $`\varphi =(f_1(r)e^{i(l1)}\varphi ,f_2(r)e^{il})`$):
$$\frac{1}{2m}\mathrm{\Omega }_1(l1)f_1(r)+\mu H(r)f_2(r)=0$$
(29)
$$\frac{1}{2m}\mathrm{\Omega }_1(l)f_2(r)+\mu H(r)f_1(r)=0$$
(30)
This equations are similar to equations obtained in where has been consider non-relativistic neutrons energy levels in the magnetic field $`\stackrel{}{H}=(0,H_\varphi (r)=2I/r,0)`$. It is seen from the following replacement in above derived equations (),():
$$f_2if_2,\sigma _mI$$
(31)
Thus, energy levels is also same and in accordance with result of are defines by quantum number $`n`$:
$$E_n=\frac{(\mu \sigma _m)^2m}{2n^2}$$
(32)
Fermions with anomalous magnetic moments in magnetic field $`\stackrel{}{H}=(0,0,H_z(r)=H(r))`$
In component form ($`\psi ^T=(f_1(r)e^{i(l1)\varphi },f_2(r)e^{i(l)\varphi },f_3(r)e^{i(l1)\varphi },f_4(r)e^{i(l)\varphi })`$) the equations has been obtained in :
$$(ϵm+\mu H)f_1p_3f_3p_{}f_4=0$$
(33)
$$(ϵm\mu H)f_2p_+f_3+p_3f_4=0$$
(34)
$$p_3f_1+p_{}f_2+(ϵm+\mu H)f_3=0$$
(35)
$$p_+f_1p_3f_2+(ϵm\mu H)f_4=0$$
(36)
where
$$p_{}=i(\frac{d}{dr}+\frac{l}{r})$$
(37)
$$p_+=i(\frac{d}{dr}\frac{l1}{r})$$
(38)
Excluding two of four components we obtain the system of two second order equations:
$$i(2\mu H(r)(\frac{d}{dr}\frac{l1}{r})+4\pi \mu j_\varphi (r))f_1(r)+((m+\mu H(r))^2+\mathrm{\Omega }_1(l))f_4(r)=0$$
(39)
$$i(2\mu H(r)(\frac{d}{dr}+\frac{l}{r})+4\pi \mu j_\varphi (r))f_4(r)+((m\mu H(r))^2+\mathrm{\Omega }_1(l1))f_1(r)=0$$
(40)
where $`\mathrm{\Omega }_1(l)=\frac{1}{r}\frac{d}{dr}r\frac{d}{dr}+\frac{l^2}{r}+p_z^2+m^2ϵ^2`$, $`\stackrel{}{j}=(0,j_\varphi (r),0)`$ is density of current.
Majorana neutrinos
In this paper we consider second order equations for Majorana neutrinos($`g_V=0`$).For this purpose it is necessary writing Dirac equation:
$$(\widehat{k}g_A\widehat{Z}\gamma _5m)\psi (k)=0,$$
(41)
in component form and excluding one of the component.
The final result for equations of the second order is following:
$$T(\kappa )f_1+g_A(Z_0^{}(r)f_2+2Z_0(r)(f_2^{}(r)+\frac{1\kappa }{r}f_2(r))=0$$
(42)
$$T(\kappa 1)f_22g_AZ_0(r)(f_1^{}(r)+\frac{1+\kappa }{r}f_1(r))=0$$
(43)
where
$$T(\kappa )=ϵ^2+\frac{1}{r^2}\frac{d}{dr}r^2\frac{d}{dr}\frac{\kappa (\kappa +1)}{r^2}g_A^2Z_0^2(r)m^2,$$
(44)
$$\kappa =l(l+1)j(j+1)\frac{1}{4},$$
(45)
We see that term $`g_A^2Z_0^2(r)`$ is always repulsive. In cylindrical symmetry case choosing two-component spinor $`\varphi `$ as $`\varphi =(f_1(r)e^{il}\varphi ,f_2(r)e^{i(l+1)})`$) we obtain the following system of the two second order equations:
$$(P(l)2g_AZ_0(r)p_z)f_1+g_A(2Z_0(r)(\frac{d}{dr}+\frac{1+l}{r})+Z_0^{}(r))if_2(r)=0$$
(46)
$$(P(l+1)+2g_AZ_0(r)p_z)if_2g_A(2Z_0(r)(\frac{d}{dr}\frac{1}{r})+Z_0^{}(r))f_1(r)=0$$
(47)
where $`P(l)=ϵ^2+\frac{1}{r}\frac{d}{dr}r\frac{d}{dr}\frac{l^2}{r^2}p_z^2g_A^2Z_0^2m^2`$.
As known, anomalous magnetic moments of the Majorana neutrino is equal to the zero. Thus , in case of Majorana neutrino the attraction is possible only via interaction (2) in radial electric and magnetic fields of monopole (see for detailes). The second order equations in this case (i.e. if only interaction (2) presented) are following:
$$(ϵ^2+\frac{1}{r^2}\frac{d}{dr}r^2\frac{d}{dr}\frac{\kappa (\kappa +1)}{r^2}A^2m^2)f_1+Aif_2(r)=0$$
(48)
$$(ϵ^2+\frac{1}{r^2}\frac{d}{dr}r^2\frac{d}{dr}\frac{\kappa (\kappa 1)}{r^2}A^2m^2)if_2+Af_1(r)=0,$$
(49)
where $`A=iqE(r)H(r)`$.We see that term $`A^2(r)r^8`$ is always repulsive, dominate at small distances and prevent fall down on the center.
In cylindrical symmetry case we have:
$$(ϵ^2+\frac{1}{r}\frac{d}{dr}r\frac{d}{dr}\frac{l^2}{r^2}p_z^2A^2m^2)f_1+Aif_2(r)=0,$$
(50)
$$(ϵ^2+\frac{1}{r}\frac{d}{dr}r\frac{d}{dr}\frac{(l+1)^2}{r^2}p_z^2A^2m^2)f_1+Aif_2(r)=0,$$
(51)
Bound states of fermions with anomalous magnetic moments in radial electric and magnetic field: second order equations
Although it is possible to exclude to of four radial functions from system of equations of the first order (see (12)-(16) in ) which defines energy levels of the neutral fermion with anomalous magnetic moment , much more convenient in order to obtain the second order equations to start from Dirac equation in component form:
$$(ϵm+\mu \stackrel{}{\sigma }\stackrel{}{H})\varphi \stackrel{}{\sigma }(\stackrel{}{p}+i\mu \stackrel{}{E})\chi =0$$
(52)
$$(ϵm+\mu \stackrel{}{\sigma }\stackrel{}{H})\chi +\stackrel{}{\sigma }(\stackrel{}{p}i\mu \stackrel{}{E})\varphi =0$$
(53)
Excluding e.g. component $`\chi `$ and presenting $`\varphi `$ as linear
combinations of spherical spinors with different $`P`$-parity (analogously , because in studied case P-parity violation take place ) <sup>1</sup><sup>1</sup>1it must be stressed that besides $`P`$-parity violation also presented $`T`$-parity violation (term $`\stackrel{}{\mathrm{\Sigma }}\stackrel{}{r}H`$ in Dirac equation is $`P`$\- and $`T`$ -parity violating) and purely imagine character of the two of four radial function in equations (12)-(16) of the is connected with this circumstances. Also, in (55),(56) e.g. $`f_2`$ must be purely imagine, $`f_1`$ must be real. :
$$\varphi ^T=f_1(r)\mathrm{\Omega }_{jlM}(\stackrel{}{n})+(1)^{\frac{1+ll^{}}{2}}f_2(r)\mathrm{\Omega }_{jl^{}M}(\stackrel{}{n})$$
(54)
we obtain the following system of the second order equations:
$$2m\mu Hif_2T_+(\kappa )f_1+Q(\kappa )(a_+S_+(\kappa )if_2b_+S_+(\kappa )f_1)$$
(55)
$$2m\mu Hf_1T_+(\kappa )if_2+Q(\kappa )(a_+S_+(\kappa )f_1b_+S_+(\kappa )if_2)$$
(56)
Analogously, excluding $`\varphi `$ and presenting $`\chi `$ as:
$$\chi =g_1(r)\mathrm{\Omega }_{jlM}(\stackrel{}{n})+(1)^{\frac{1+ll^{}}{2}}g_2(r)\mathrm{\Omega }_{jl^{}M}(\stackrel{}{n}))$$
(57)
we obtain the system for $`g_{1,2}`$:
$$2m\mu Hig_2T_{}(\kappa )g_1Q(\kappa )(a_{}S_{}(\kappa )ig_2b_{}S_{}(\kappa )g_1)$$
(58)
$$2m\mu Hg_1T_{}(\kappa )ig_2Q(\kappa )(a_{}S_{}(\kappa )g_1b_{}S_{}(\kappa )ig_2)$$
(59)
where:
$$Q(\kappa )=4\pi \mu \rho _m\frac{2\mu H(r)}{r}(1\pm \kappa ),$$
(60)
$$S_\pm (\kappa )=\frac{d}{dr}+\frac{1\pm \kappa }{r}\pm \mu E,$$
(61)
$$T_\pm (\kappa )=ϵ^2m^2+\frac{1}{r^2}\frac{d}{dr}r^2\frac{d}{dr}\frac{\kappa (\kappa +1)}{r^2}\mu ^2(E^2+H^2)\pm 4\mu \pi \rho \frac{2\mu E}{r}(1\pm \kappa )$$
(62)
$$a_\pm =\frac{ϵ\pm m}{(ϵ\pm m)^2\mu ^2H^2}$$
(63)
$$b_\pm =\frac{\pm \mu H}{(ϵ\pm m)^2\mu ^2H^2}$$
(64)
During derivation of this formulas we take into account that $`div\stackrel{}{E}=4\pi \rho `$,$`div\stackrel{}{H}=4\pi \rho _m`$.
In nonrelativistic limit terms which are proportional to the $`a_\pm ,b_\pm `$ are small and may be neglected.
We see, that terms $`(\mu E(r))^2r^4`$ $`(\mu H(r))^2r^4`$ are repulsive and prevent fall down on the center.
The author express his sincere gratitude to Zh.K.Manucharyan and E.B.Prokhorenko for helpful discussions. |
warning/0002/cond-mat0002423.html | ar5iv | text | # Nonlinear Effect of Transport Current on Response of Metals to Electromagnetic Radiation
## I Introduction
As well-known metals possess quite peculiar nonlinear electrodynamic properties (see., for example, Refs. rev1 ; rev2 ). Indeed, nonlinearity in a response of plasma or semiconductors to an electromagnetic excitation is usually associated to significant deviation of the electron subsystem from equilibrium. To the contrary, nonequilibrium in metals as a rule is weak due to large concentration of charge carriers. Nevertheless, a nonlinear regime in these media is rather easy to observe. Such a situation is possible because of the fact that nonlinearity in metals does not necessarily related to the nonequilibrium of the electron subsystem. Nonequilibrium is caused by the existence of a weak electric field, while nonlinearity is originates from strong magnetic fields. The Lorentz force, determined by the magnetic field of an electromagnetic wave or magnetic field of the transport current, affects the dynamics of charge carriers. As a result, the conductivity of a sample placed in an AC electromagnetic field depends on the spatial distribution of magnetic field of the wave. Such a magnetodynamic mechanism of nonlinearity is inherent to pure metals at low temperatures, if the mean free path of conducting electrons is rather large.
Magnetodynamic nonlinearity causes a number of nontrivial phenomena in the electrodynamics of metals. As an example, one can mention generation of the current states Dolg ; CS in a sample placed in a DC external magnetic field. The sample acquires a DC magnetic moment if irradiated by an additional strong AC electromagnetic field. The magnitude of the magnetic moment depends in a hysteretic manner on an external DC magnetic field. Under current states conditions, a hysteresis-like interaction of radiowaves is observed int as well as the appearance of electromagnetic dissipative structures diss . This specific mechanism of nonlinearity in metals causes a decrease of collisionless damping of helicons helic1 . Therefore, the spiral waves with large amplitudes can propagate even in conditions when there is no their linear electromagnetic excitations helic2 . Magnetoplasmic shock waves shock and soliton-like excitations waves are also predicted for the regime of strong magnetodynamic nonlinearity.
In the present paper, we study a novel manifestation of magnetodynamic nonlinearity, namely, interaction of an external electromagnetic wave and a strong DC transport current in a thin metal film, which is also displayed in a quite unusual way. The sample of thickness $`d`$ is assumed to be much smaller than electron mean free path $`l`$,
$$dl,$$
(1)
and electron scattering on a surface of the film is supposed to be diffuse. It is known Snapiro that in the static case (when external AC signal is absent), the magnetic field of a current can essentially affect the conductivity of a thin metal specimen and, thus, its current-voltage characteristics (CVC). In this situation the value $`I`$ of the current is rather small so that the typical radius of curvature $`R(I)`$ of electron trajectories in magnetic field is much greater than the film thickness,
$$dR(I),R(I)=cp_F/eH(I)I^1.$$
(2)
Here $`e`$ and $`p_F`$ represent electron charge and Fermi momentum, respectively. In Ref. Snapiro , it was shown that nonlinear peculiarities of CVC are connected to the antisymmetrical spatial distribution of the magnetic field of the DC current inside the sample. The magnetic field equals to zero at the middle of the film and takes the values $`H`$ to $`H`$ at the opposite boundaries of the sample, where
$$H=2\pi I/cD.$$
(3)
In this formula, $`c`$ denotes the speed of light in vacuum and $`D`$ is the sample width. Spatially alternating field of the DC current entraps a part of electrons in a potential well. Trajectories of such particles are flat curves winding around the plane of alternation of the magnetic field. The relative part of the trapped electrons in the order of magnitude is equal to the typical angle $`(d/R)^{1/2}1`$ of their crossing of that plane. Taking into account that trapped carriers do not collide with the film boundaries and interact with the electric field along their whole free path $`l`$, one can write the following estimating formula for their conductivity $`\sigma _{tr}`$:
$$\sigma _{tr}\sigma _0(d/R(I))^{1/2}I^{1/2}.$$
(4)
Here $`\sigma _0`$ represents the conductivity of the bulk sample. At the same time, there exist flying electrons which do collide with the boundaries of the specimen and, according to Ref. Fuks , have the conductivity of the order of $`\sigma _0(d/l)`$. Apparently in the range of rather strong currents, when the inequality
$$(dR(I))^{1/2}l,$$
(5)
holds, the conductivity of the film is determined by the group of the trapped carriers. As a result, we observe the deviation from the Ohm’s law: the voltage $`U`$ is proportional to the square root of current,
$$UI^{1/2}.$$
(6)
For the film with thickness $`d=10^3`$ cm, the electron mean free path $`l=10^1`$ cm, and Fermi momentum $`p_F=10^{19}`$ g$``$ cm/s the nonlinearity becomes noticeable ($`(dR)^{1/2}l`$) at values of the magnetic field $`H(I)`$ about $`1`$ Oe. The theory developed in Ref. Snapiro is in a good qualitative agreement with experimental data (see, for example, Ref. Fish1 ).
In an external magnetic field $`𝐡`$, collinear to the magnetic field of the current, the plane of the sign alternation for the magnetic field shifts to one of the two boundaries of the film (see Fig. 1). That in turn leads to appreciable diminution of the number of the trapped particles and, therefore, their conductivity. In particular, such a situation would take place under symmetrical irradiation of the film by the low-frequency electromagnetic wave of large amplitude. The frequency is supposed to be so small that the AC magnetic field $`𝐡(t)`$ of the wave is virtually uniform across the metal (i.e. the wave penetration depth $`\delta `$ is much greater than the sample thickness $`d`$). Then the conductivity of metal essentially depends on time and, therefore, strong nonlinear effects in the sample response to the AC electromagnetic excitation should appear. Being of interest from the both theoretical and experimental points of view this problem has not been investigated yet.
In the present paper we study theoretically the temporal dependence of electric field at the surface of the film, which carries a strong DC current of the fixed value $`I`$ satisfying inequalities (2) and (5). It is shown that with an increase of the amplitude $`h_m`$ of the AC magnetic field this dependence becomes anharmonic, turning into a series of sharp nonanalytic peaks. The case of sufficiently high amplitudes $`h_m>H`$, when the total magnetic field in the sample is spatially alternating during some part of the wave period and has constant sign during the other part, is of particular interest. In such a situation, the electric field is also characterized by kinks in its temporal dependence due to the periodical appearance and disappearance of the group of trapped carriers. The effect of amplification of the electric signal on the film surface is predicted as well. It turns out that, because of the presence of the strong DC transport current in the sample, the absolute value of the AC electric field of the wave is $`l/(dR)^{1/2}1`$ times as many as the corresponding magnitude in the absence of the DC current.
We also calculate the nonlinear surface impedance of the film, which turns out to be pure imaginary value in the main approximation in the parameter $`d/\delta 1`$, and show that its modulus monotonically diminishes with the growth of the AC amplitude decreasing in $`l/(dR)^{1/2}1`$ times. Simultaneously the conductivity of the trapped particles falls down and, consequently, $`\delta `$ increases.
## II Problem statement and geometry
Consider a metal film of the thickness $`d`$ with a DC current $`I`$ flowing along. The sample is irradiated from both sides by a monochromatic electromagnetic wave with a magnetic component collinear to the magnetic field of the current. The $`x`$-axis is oriented perpendicularly to the film boundaries. The plane $`x=0`$ corresponds to the middle of the sample (see Fig. 1).
The $`y`$-axis is directed along the current, and the $`z`$-axis is parallel to the vector $`(x,t)`$ of the total magnetic field which is a sum of the DC magnetic field of the current $`𝐇(x,t)`$ and the AC magnetic field of the wave $`𝐡(x,t)`$,
$$(x,t)=\{0,0,H(x,t)+h(x,t)\}.$$
(7)
The film length $`L`$ (the dimension along the $`y`$-axis) and its width $`D`$ ($`z`$-axis dimension) are much greater than the sample thickness $`d`$. We assume diffuse scattering of the electrons on the film boundaries. Maxwell’s equations in the assumed geometry can be written as
$$\frac{(x,t)}{x}=\frac{4\pi }{c}j(x,t),\frac{E(x,t)}{x}=\frac{1}{c}\frac{(x,t)}{t},$$
(8)
where $`j(x,t)`$ and $`E(x,t)`$ represent the $`y`$ components of the current density and the electric field. Boundary conditions for Eqs. (8) are
$$(\pm d/2,t)=h_m\mathrm{cos}\omega tH.$$
(9)
Let $`H`$ be the absolute value of the magnetic field on the surface of the metal film and $`h_m`$ denote a wave amplitude. According to Eq. (3), the field $`H`$ is determined by the DC current $`I`$. No special relation between magnitudes $`H`$ and $`h_m`$ is assumed.
We consider a quasistatic situation when the wave frequency $`\omega `$ is much less than the relaxation frequency $`\nu `$ of the charge carriers,
$$\omega \nu .$$
(10)
Here we suppose that the AC magnetic field inside the sample is quasiuniform and virtually does not differ from its value on the sample surface, $`h(x,t)h_m\mathrm{cos}\omega t`$. In other words, the typical spatial scale $`\delta (\omega )`$ of variation of the AC magnetic field is much greater than the film thickness $`d`$. Furthermore we assume that the curvature radius $`R(x,t)`$ of electron trajectories in the total magnetic field $`(x,t)`$ is also much greater than $`d`$,
$$d\delta (\omega ),dR(x,t),R(x,t)=cp_F/e|(x,t)|.$$
(11)
## III Electron dynamics, current density, and CVC of film
Let us consider electron dynamics in the nonuniform AC magnetic field $`(x,t)`$. We shall assume the following gauge of the vector potential:
$$𝐀(x,t)=\{0,A(x,t),0\},A(x,t)=^x𝑑x^{}(x^{},t).$$
(12)
It is suitable to choose the lower limit of integration in Eq. (12) depending on whether or not there exists the plane $`x=x_0(t)`$ of the sign alternation of the magnetic field $`(x,t)`$ at the present moment. This plane exists during the time intervals when $`h_m|\mathrm{cos}\omega t|<H`$ because the values $`h_m\mathrm{cos}\omega tH`$ and $`h_m\mathrm{cos}\omega t+H`$ of the total magnetic field at the film boundaries have opposite signs (see Eq. (9)). In this case, one should take $`x_0(t)`$ as the lower limit in integral (12). Then the vector potential $`A(x,t)`$ is negative. It reaches its maximum value (which equals to zero) at the point $`x=x_0(t)`$. Within other time intervals, when the inequality $`h_m|\mathrm{cos}\omega t|>H`$ holds, the magnetic field $`(x,t)`$ inside the sample is of a constant sign. In such a situation, one should choose $`\mathrm{sign}(\mathrm{cos}\omega t)d/2`$ ($`\mathrm{sign}(x)`$ is the sign function) as a lower limit of the integration. In this case, vector potential, also being negative, vanishes at one of the boundaries of the film.
The integrals of motion of electron in the field $`(x,t)`$ are the total energy (it equals to the Fermi energy) and the canonical momenta $`p_z=mv_z`$ and $`p_y=mv_yeA(x,t)/c`$ (m is the electron mass). The electron trajectory in a perpendicular to the direction of the magnetic field plane is determined by the velocities $`v_x(x,t)`$ and $`v_y(x,t)`$. In the case of Fermi-sphere with radius $`p_F=mv`$, we obtain
$$|v_x(x,t)|=(v_{}^2v_y^2)^{1/2},v_{}=(v^2v_z^2)^{1/2},v_y(x,t)=(p_y+eA(x,t)/c)/m.$$
(13)
Classically allowable regions of the electron motion along the $`x`$ axis are determined by the inequalities,
$$p_ymv_{}eA(x,t)/cp_y+mv_{}.$$
(14)
These inequalities provide the positivity of the radicand in Eq. (13) for $`|v_x(x,t)|`$.
The regions of the electron motion in the phase plane $`(x,p_y)`$ are described schematically in Fig. 2 for two cases: when there exists the plane $`x=x_0(t)`$ of the sign alternation of the magnetic field $`(x,t)`$ (Fig. 2, a) and when such plane is absent (Fig. 2, b).
For definiteness we have chosen the moment of time when the magnetic field of the wave is positive ($`\mathrm{cos}\omega t>0`$). The upper border on the phase plane is described by the curve $`p_y=mv_{}eA(x,t)/c`$ and the lower one is given by $`p_y=mv_{}eA(x,t)/c`$. The electrons are naturally divided in three groups depending on the sign and value of the integral of motion $`p_y`$. Below, we give inequalities determining the regions of their existence at an arbitrary moment of time.
1. Flying electrons:
$$p_y^{}mv_{}eA[\mathrm{sign}(\mathrm{cos}\omega t)d/2,t]/cp_ymv_{},|x|d/2.$$
(15)
These particles collide with the both boundaries of the film. Their trajectories do not twist significantly because of $`dR(x,t)`$. Flying electrons exist at every moment of time irrespective of the presence of the plane $`x=x_0(t)`$ (i.e irrespective of the relation between $`h_m\mathrm{cos}\omega t`$ and $`H`$).
2. Trapped electrons: They appear during the periods of time when $`h_m|\mathrm{cos}\omega t|<H`$ and the total magnetic field $`(x,t)`$ within the sample passes trough zero. Their states are bounded by the region (see. Fig. 2,a),
$`mv_{}p_yp_y^+mv_{}eA[\mathrm{sign}(\mathrm{cos}\omega t)d/2,t]/c,`$
$`x_{}(t)\mathrm{sign}(\mathrm{cos}\omega t)<x\mathrm{sign}(\mathrm{cos}\omega t)<d/2.`$ (16)
Here $`x_{}(t)`$ represents the breakpoint of the trapped electron most distant from the film boundary. One can find it from the equation,
$$A(x_{},t)=A[\mathrm{sign}(\mathrm{cos}\omega t)d/2,t].$$
(17)
According to Eq. (16), this electron group occupies the region $`x_{}(t)<x<d/2`$ when $`\mathrm{cos}\omega t>0`$ and the region $`d/2<x<x_{}(t)`$ if $`\mathrm{cos}\omega t<0`$. The trajectories of trapped particles are almost flat oscillating curves due to periodical motion of the particles along $`x`$-direction and the uniform motion along the $`y`$-and $`z`$-axes. The temporal period of oscillations with respect to the plane $`x=x_0`$ equals to $`2T`$, where
$$T=_{x_1(t)}^{x_2(t)}\frac{dx}{|v_x(x,t)|}.$$
(18)
The breakpoints $`x_1(t)`$ and $`x_2(t)`$ ($`x_1(t)<x_0(t)<x_2(t)`$) are the roots of the equation,
$$eA(x_{1,2},t)/c=mv_{}p_y.$$
(19)
3. Surface electrons:
These particles collide only with one of the boundaries of the film. In our case of diffuse scattering of the electrons on the surface, their influence on the nonlinear conductivity of metal is negligible Snapiro . Thus, we do not take them into account thereafter.
The current density of the flying and trapped particles can be deduced by means of solving the Boltzmann kinetic equation. One should linearize the kinetic equation with respect to the electric field $`E(x,t)`$, which can be represented as a sum,
$`E(x,t)`$ $`=`$ $`E_0+(x,t),`$
$`(x,t)`$ $`=`$ $`{\displaystyle \frac{1}{c}}\left({\displaystyle \frac{A(x,t)}{t}}{\displaystyle \frac{\overline{A}(t)}{t}}\right).`$ (20)
Here the first term, $`E_0`$, is a potential (uniform) component and $`(x,t)`$ is a rotational (nonuniform) field of the wave. Spatial averaging of the latter over the $`x`$-axis direction gives zero. The value $`\overline{A}(t)`$ represents a spatially averaged magnitude of the vector potential,
$$\overline{A}(t)=\frac{1}{d}_{d/2}^{d/2}A(x^{},t)𝑑x^{}.$$
(21)
The magnetodynamic nonlinearity is accounted for in the kinetic equation by means of terms which contain the total magnetic field $`(x,t)=H(x,t)+h(x,t)`$ entering the Lorentz force. We calculate the current density in the main approximation with respect to the small parameter $`d/\delta (\omega )`$ (see Eq. (11)). In this approximation, as it was mentioned above, the AC magnetic field $`h(x,t)`$ becomes spatially uniform and is equal to its boundary value, $`h(x,t)=h_m\mathrm{cos}\omega t`$. The electric field is also independent of the $`x`$-coordinate and coincides with the value $`E_0(t)`$. For the case of uniform electric and external magnetic fields, the current density was obtained in Ref. Snapiro . If the conditions (2) and (5) hold the following asymptotics for the current density of the flying and trapped electrons are valid:
$$j_{fl}(t)=\sigma _{fl}(t)E_0(t),$$
$$\sigma _{fl}(t)=\frac{3}{8}\sigma _0\frac{d}{l}\mathrm{ln}\frac{R_+(t)}{d},R_\pm (t)=cp_F/e|h_m|\mathrm{cos}\omega t|\pm H|,$$
(22)
$$j_{tr}(x,t)=\sigma _{tr}(x,t)E_0(t),$$
$$\sigma _{tr}(x,t)=\frac{36\pi ^{1/2}}{5\mathrm{\Gamma }^2(1/4)}\sigma _0\left\{\frac{e}{cp_F}\left[A(x,t)A(\mathrm{sign}(\mathrm{cos}\omega t)d/2,t)\right]\right\}^{1/2},$$
(23)
$$x_{}(t)\mathrm{sign}(\mathrm{cos}\omega t)<x\mathrm{sign}(\mathrm{cos}\omega t)<d/2.$$
In the limit $`\omega 0`$, Eqs. (22) and (23) transform into the corresponding formulae of Ref. Snapiro .
Let us substitute the current density in the first of Maxwell’s equations (8) for its asymptotic expressions (22) and (23) and introduce a dimensionless coordinate and vector potential,
$$\xi =2x\mathrm{sign}(\mathrm{cos}\omega t)/d,a(\xi ,t)=A(x,t)/A(\mathrm{sign}(\mathrm{cos}\omega t)d/2,t).$$
(24)
The equation for the quantity $`a(\xi ,t)`$ has the form,
$$\frac{^2a(\xi ,t)}{\xi ^2}=u\{\begin{array}{cc}r[1a(\xi ,t)]^{1/2}+1,\hfill & \xi _{}(t)\xi 1\text{,}\hfill \\ & \\ 1,\hfill & 1\xi \xi _{}(t),\hfill \end{array}$$
(25)
$$\xi _{}(t)=2x_{}(t)\mathrm{sign}(\mathrm{cos}\omega t)/d.$$
(26)
The dimensionless coordinate $`\xi _{}(t)`$ confines the region of existence of the trapped particles and, according to Eqs. (17) and (24), satisfies the equation, $`a(\xi _{},t)=1`$. The parameter $`r`$ represents the ratio of the maximum magnitude of the conductivity of the trapped electrons to the conductivity of the flying particles,
$$r=\frac{\sigma _{tr}(x_0)}{\sigma _{fl}}=\frac{96\pi ^{1/2}}{5\mathrm{\Gamma }^2(1/4)}\frac{l}{d}\left[\frac{e}{cp_F}|A(\mathrm{sign}(\mathrm{cos}\omega t)d/2,t)|\right]^{1/2}\mathrm{ln}^1(R_+/d).$$
(27)
The dimensionless quantity, $`u`$, is related to the voltage $`U=E_0L`$ on the sample,
$$u=\frac{U}{cL|A(\mathrm{sign}(\mathrm{cos}\omega t)d/2,t)|/\pi \sigma _{fl}d^2}.$$
(28)
Equation (25) should be solved together with the boundary conditions,
$$\frac{a(1,t)}{\xi }=\frac{d}{2}\frac{h_m|\mathrm{cos}\omega t|H}{A(\mathrm{sign}(\mathrm{cos}\omega t)d/2,t)},$$
$$\frac{a(1,t)}{\xi }=\frac{d}{2}\frac{h_m|\mathrm{cos}\omega t|+H}{A(\mathrm{sign}(\mathrm{cos}\omega t)d/2,t)},a(1,t)=1.$$
(29)
The first two of these expressions are dimensionless boundary conditions (9), and the third one is a consequence of normalization (24) of the vector potential.
Within the interval $`\xi _{}(t)\xi 1`$, the solution of Eq. (25) is symmetrical with respect to the point $`\xi _0(t)=(1+\xi _{}(t))/2`$, where the dimensionless vector potential reaches its minimum value (which equals to zero, $`a(\xi _0,t)=a(\xi _0,t)/\xi =0`$). This solution is described by the formula,
$$|\xi \xi _0(t)|=(3/4ru)^{1/2}_0^{a(\xi ,t)}𝑑\zeta [1(1\zeta )^{3/2}+3\zeta /2r]^{1/2}.$$
(30)
One can not obtain the field distribution and the current density within the region of existence of the trapped electrons in an explicit form. However, by means of Eq. (30), it is possible to calculate the average magnitude of the conductivity of the trapped carriers (23) within the interval (16),
$`{\displaystyle \frac{\overline{\sigma }_{tr}}{\sigma _{fl}}}`$ $`=`$ $`r{\displaystyle _0^1}𝑑\zeta (1\zeta )^{1/2}[1(1\zeta )^{3/2}+3\zeta /2r]^{1/2}`$ (31)
$`\times `$ $`\left({\displaystyle _0^1}𝑑\zeta [1(1\zeta )^{3/2}+3\zeta /2r]^{1/2}\right)^1.`$
The bar above $`\sigma _{tr}`$ denotes spatial averaging. In the remaining region of the sample ($`1\xi \xi _{}(t)`$), there exist only flying electrons, and the solution of Eq. (25) is given by the following formula:
$$a(\xi ,t)=1(2u)^{1/2}(1+2r/3)^{1/2}(\xi \xi _{}(t))+u(\xi \xi _{}(t))^2/2.$$
(32)
Expressions (30) and (32) and their derivatives are sewn together at the point $`\xi =\xi _{}(t)`$. The solution given by Eqs. (30) and (32) contains three parameters, $`\xi _0`$, $`u`$ and $`r`$, which should be found from boundary conditions (29). It is essential that the value $`A(\mathrm{sign}(\mathrm{cos}\omega t)d/2,t)`$ of the vector potential appearing in Eq. (29) is not an independent parameter due to its relation to $`r`$ via formula (27).
Adding term by term the first two boundary conditions in Eq. (29) and using Eqs. (30) and (32), and (28), we find the following expression for the drift of the plane $`x=x_0`$:
$$\xi _0=2x_0\mathrm{sign}(\mathrm{cos}\omega t)/d=\frac{Lh_m|\mathrm{cos}\omega t|}{2\pi U\sigma _{fl}d},h_m|\mathrm{cos}\omega t|H.$$
(33)
In order to determine the value of $`u`$ (i.e. the voltage $`U`$), let us integrate the left and right-hand sides of Eq. (25) from -1 to 1 taking into account the boundary conditions for the derivative $`a(\xi ,t)/\xi `$ in (29). The integral of the function $`[1a(\xi ,t)]^{1/2}`$ appearing in the right-hand side can be reduced to the product $`2(1\xi _0)\overline{\sigma }_{tr}/r\sigma _{fl}`$ with the use of the condition $`a(1,t)=1`$. Taking this into consideration as well as formulae (28) and (33) for the quantities $`u`$ and $`\xi _0`$, we have after some simple transformations,
$$U=\frac{cL}{2\pi d\sigma _{fl}(t)}\frac{H(I)+(\overline{\sigma }_{tr}/\sigma _{fl})h_m|\mathrm{cos}\omega t|}{1+\overline{\sigma }_{tr}/\sigma _{fl}},h_m|\mathrm{cos}\omega t|H.$$
(34)
According to Eq. (31), the ratio of conductivities, $`\overline{\sigma }_{tr}/\sigma _{fl}`$, depends on the parameter $`r`$. Using expression (28) for $`u`$, relation (27) between $`A(\mathrm{sign}(\mathrm{cos}\omega t)d/2,t)`$ and $`r`$, and solution (30), we obtain from the first boundary condition in Eq. (29) the algebraic equation for $`r`$,
$$r^2(1+2r/3)=\left(\frac{Hh_m|\mathrm{cos}\omega t|}{\stackrel{~}{H}}\right)^2\frac{\stackrel{~}{U}}{U\mathrm{ln}^3(R_+/d)},h_m|\mathrm{cos}\omega t|H.$$
(35)
Here we have introduced the following notations:
$$\stackrel{~}{H}=\frac{25\mathrm{\Gamma }^4(5/4)}{9\pi }\frac{cp_Fd}{el^2},\stackrel{~}{U}=\frac{4clL\stackrel{~}{H}}{3\pi \sigma _0d^2}.$$
(36)
The parameters $`\stackrel{~}{H}`$ and $`\stackrel{~}{U}`$ represent those magnitudes of the magnetic field and voltage for which the characteristic length $`(Rd)^{1/2}`$ of the arch of electron’s trajectory is of the order of the mean free path $`l`$.
Expressions (31), (34), and (35) define, in an implicit form, dependence of the voltage $`U`$ on the current $`I`$ for the case $`h_m|\mathrm{cos}\omega t|H`$. At these conditions there exists the plane of the alternation of sign of the total magnetic field within the sample. If the opposite inequality, $`h_m|\mathrm{cos}\omega t|H`$, is valid, the trapped electrons are absent ($`r=0`$, $`\xi _{}=1`$, $`\sigma _{tr}=0`$) and CVC is described by the formula,
$$U=\frac{cLH(I)}{2\pi d\sigma _{fl}(t)},h_m|\mathrm{cos}\omega t|H.$$
(37)
As seeing from formula (34), the voltage on the sample displays nonanalytical behavior vs. time: the dependence $`U(t)`$ has kinks at the moments when the AC magnetic field $`h_m\mathrm{cos}\omega t`$ vanishes. This is an essentially nonlinear effect caused by the contribution of a large group of trapped electrons into the electric current.
The temporal dependence of voltage (34) for the case when the wave amplitude is not too large ($`h_m<H`$) and there exist the trapped carriers during the whole period $`2\pi /\omega `$ is shown in Fig. 3, a. Fig. 3. b represents the dependence $`U(t)`$ for the opposite case $`h_m>H`$, in which during some part of the wave period (at $`h_m|\mathrm{cos}\omega t|H`$) the conductivity is caused by the flying particles only.
## IV Nonanalytical temporal dependence of electric field
Knowing the vector potential $`A(x,t)`$, one can calculate the rotational electric field $`(x,t)`$ as a correction to $`E_0(t)`$, (see Ref. (III)). We are interested in the difference $`\mathrm{\Delta }(t)=(d/2,t)(d/2,t)`$. This value is proportional to the rate of alteration of the magnetic flux trough the cross-sectional plane, which is perpendicular to the direction of the vector of the total field $`(x,t)`$, and thus can be measured in experiment.
From Eqs. (30) and (32), it follows that the difference $`a(1,t)a(1,t)`$ is connected to the derivatives $`a(1,t)/\xi `$ and $`a(1,t)/\xi `$ by the relations,
$`a(1,t)a(1,t)=\xi _0(t)\left[{\displaystyle \frac{a(1,t)}{\xi }}{\displaystyle \frac{a(1,t)}{\xi }}\right],h_m|\mathrm{cos}\omega t|H,`$ (38)
$`a(1,t)a(1,t)={\displaystyle \frac{a(1,t)}{\xi }}+{\displaystyle \frac{a(1,t)}{\xi }},h_m|\mathrm{cos}\omega t|H.`$ (39)
Let us now turn to the dimensional variables in Eqs. (38) and (39) using boundary conditions (29) and relation (27) between the values of $`A(\mathrm{sign}(\mathrm{cos}\omega t)d/2,t)`$ and $`r`$. After that one can obtain the following expression for the magnitudes of the vector potential at the film boundaries:
$$A(\mathrm{sign}(\mathrm{cos}\omega t)d/2,t)=\stackrel{~}{H}d\mathrm{ln}^2(R_+/d)r^2/4,$$
$$A(\mathrm{sign}(\mathrm{cos}\omega t)d/2,t)=\stackrel{~}{H}d\mathrm{ln}^2(R_+/d)r^2/42H|x_0(t)|$$
(40)
at
$$h_m|\mathrm{cos}\omega t|H$$
and
$$A(\mathrm{sign}(\mathrm{cos}\omega t)d/2,t)=0,A(\mathrm{sign}(\mathrm{cos}\omega t)d/2,t)=dh_m|\mathrm{cos}\omega t|$$
(41)
at
$$h_m|\mathrm{cos}\omega t|H.$$
Formulae (40) and (41) are sewn at the time moment when $`h_m|\mathrm{cos}\omega t|=H`$. The parameter $`r`$ in Eq. (27) vanishes, and the plane $`x=x_0(t)`$ coincides with one of the boundaries of the sample, $`|x_0(t)|=d/2`$. From relations (40) and (III), by means of formula (33) for $`\xi _0(t)`$, we derive the expression for the difference $`\mathrm{\Delta }(t)`$ of magnitudes of the electric field at the film boundaries,
$$\mathrm{\Delta }(t)=\frac{2H}{c}\frac{x_0(t)}{t}=\frac{H(I)Lh_m}{2\pi }\frac{}{t}\left[\frac{\mathrm{cos}\omega t}{\sigma _{fl}(t)U(t)}\right],h_m|\mathrm{cos}\omega t|H.$$
(42)
If the inequality $`h_mH`$ holds, the previous relation is valid during the whole period of the wave. However, in the case $`h_m>H`$, there exists a time interval when the plane $`x=x_0(t)`$ of alternation of the sign of the total magnetic field is absent. If such a situation takes place one should use formula (41) in order to obtain the dependence $`\mathrm{\Delta }(t)`$. Finally we come to the result below,
$$\mathrm{\Delta }(t)=\mathrm{\Delta }_L\mathrm{sin}\omega t,\mathrm{\Delta }_L=dh_m\omega /c,h_m|\mathrm{cos}\omega t|H.$$
(43)
From this, it follows that the difference $`\mathrm{\Delta }(t)`$ is a harmonic function of time, i.e. the response of the film on the external electromagnetic excitation turns out to be linear if there are no trapped electrons. It is obvious that formula (43) also describes the dependence $`\mathrm{\Delta }(t)`$ at small magnitudes of the current $`I`$ ($`H\stackrel{~}{H}`$), when the contribution of trapped particles to the conductivity is negligible during the whole period of the wave. Then the value $`\mathrm{\Delta }_L`$ represents the amplitude of a linear response.
The dependence $`\mathrm{\Delta }(t)`$ is shown in Fig. 4 for a wide range of the AC amplitudes $`h_m`$ and for the large magnitudes of the DC magnetic field $`H`$ of the current $`I`$, when the inequality $`H\stackrel{~}{H}`$ (or inequality (5)) is valid. It is obvious that the ratio of the amplitude $`\mathrm{\Delta }_m`$ to its linear value $`\mathrm{\Delta }_L`$ does not depend on $`h_m`$. From relations (42), (34), and (35) at $`\mathrm{cos}\omega t=0`$, we find the expression for $`\mathrm{\Delta }_m`$,
$$\frac{\mathrm{\Delta }_m}{\mathrm{\Delta }_L}=0.83\left(\frac{H}{\stackrel{~}{H}}\right)^{1/2}\frac{1}{\mathrm{ln}(R/d)},\left(\frac{H}{\stackrel{~}{H}}\right)^{1/2}\frac{\sigma _{tr}}{\sigma _{fl}}|_{\mathrm{cos}\omega t=0}\frac{l}{(Rd)^{1/2}}1.$$
(44)
The ratio $`\mathrm{\Delta }_m/\mathrm{\Delta }_L`$ is determined by the magnitude of the DC magnetic field $`H`$ and can be much greater than unity. In other words, there exists an effect of amplification of the electric signal at the film surface. For small AC amplitudes (curve $`1`$, $`h_m=H/300`$) the signal turns out to be quasi-harmonic. However, with the increase of $`h_m`$ the dependence $`\mathrm{\Delta }(t)`$ shows kinks. Curve $`2`$ has kinks at the points of extremum, i.e at the time moments when the AC magnetic field $`h_m\mathrm{cos}\omega t`$ vanishes. These singularities are related to the nonanalytical behavior of CVC of the film (see. Eq. (34) and Fig. 3). Curve $`3`$ corresponds to the case $`h_m=5H/3`$, in which the trapped electrons are absent during a part of the wave period. In such a situation, the dependence $`\mathrm{\Delta }(t)`$ contains additional kinks arising at the moments of appearance and disappearance of the plane $`x=x_0(t)`$ of the sign alternation of the total magnetic field. They are located symmetrically with respect to the points of extremum as shown in curve $`3`$. By means of formulae (42), (43),(34), and (35), we find the right and left derivatives of the function $`\mathrm{\Delta }(t)`$ at the point $`t_0=(1/\omega )\mathrm{arccos}(H/h_m)`$ of the first kink,
$$\frac{}{t}\frac{\mathrm{\Delta }(t)}{\mathrm{\Delta }_L}|_{t=t_00}=\frac{\omega H}{h_m},$$
(45)
$$\frac{}{t}\frac{\mathrm{\Delta }(t)}{\mathrm{\Delta }_L}|_{t=t_0+0}=\frac{\omega H}{h_m}\left[1\frac{\pi }{2\mathrm{ln}(R_+/d)}\left(\frac{H}{\stackrel{~}{H}}\right)^{1/2}\left(\frac{h_m^2}{H^2}1\right)\right].$$
(46)
According to Eq. (46), the right derivative is negative and has large absolute value even at $`[(h_m/H)^21]1`$.
## V Surface impedance of film
Let us analyze the dependence of the surface impedance at the film boundary $`x=d/2`$ on the AC amplitude $`h_m`$ under conditions of interaction of the transport current and the electromagnetic wave. The impedance is proportional to the ratio of the first Fourier harmonics of the electric $`_\omega `$ and magnetic $`h_\omega `$ fields at the surface of the sample,
$`Z={\displaystyle \frac{4\pi }{c}}{\displaystyle \frac{_\omega }{h_\omega }}`$ $`={\displaystyle \frac{8\pi }{c}}{\displaystyle \frac{_\omega }{h_m}},_\omega ={\displaystyle \frac{\omega }{2\pi c}}{\displaystyle _0^{2\pi /\omega }}\left({\displaystyle \frac{A(d/2,t)}{t}}{\displaystyle \frac{\overline{A}}{t}}\right)e^{\mathrm{i}\omega t}𝑑t`$ (47)
$`={\displaystyle \frac{\mathrm{i}\omega ^2}{2\pi c}}{\displaystyle _0^{2\pi /\omega }}\left(A(d/2,t)\overline{A}(t)\right)e^{\mathrm{i}\omega t}𝑑t.`$
Taking into account Eqs. (33) and (40), we deduce the boundary value of the vector potential for the periods of time given by the inequality $`h_m|\mathrm{cos}\omega t|H`$,
$$A(d/2,t)=\{\begin{array}{cc}\stackrel{~}{H}d\mathrm{ln}^2(R_+/d)r^2/4,at\mathrm{cos}\omega t>0,\hfill & \\ \stackrel{~}{H}\mathrm{ln}^2(R_+/d)r^2/4+cHLh_m\mathrm{cos}\omega t/2\pi U(t)\sigma _{fl}(t),at\mathrm{cos}\omega t<0.\hfill & \end{array}$$
(48)
In the case $`h_m|\mathrm{cos}\omega t|H`$, the following expression is valid (see. Eq. (41)):
$$A(d/2,t)=\{\begin{array}{cc}0,\hfill & at\mathrm{cos}\omega t>0\text{,}\hfill \\ dh_m\mathrm{cos}\omega t,\hfill & at\mathrm{cos}\omega t<0\text{.}\hfill \end{array}$$
(49)
Let us calculate the mean value of the vector potential $`\overline{A}(t)`$ for $`h_m|\mathrm{cos}\omega t|H`$, when there exists the plane of alternation of sign of the field. According to Eqs. (30) and (32), we have
$`{\displaystyle \frac{\overline{A}(t)}{A(\mathrm{sign}(\mathrm{cos}\omega t)d/2,t)}}={\displaystyle \frac{1}{2}}{\displaystyle _1^1}a(\xi ,t)𝑑\xi =\xi _0(t)+(2u(t))^{1/2}(1+r(t)/3)^{1/2}\xi _0^2(t)`$
$`+(2/3)u(t)\xi _0^3+\left({\displaystyle \frac{3}{4r(t)u(t)}}\right)^{1/2}{\displaystyle _0^1}{\displaystyle \frac{\zeta d\zeta }{\sqrt{1(1\zeta )^{3/2}+3\zeta /2r(t)}}}.`$ (50)
In the case $`h_m|\mathrm{cos}\omega t|H`$, one should use solution (32) with $`r=0`$, $`\xi _{}=1`$ in order to find $`\overline{A}(t)`$. Proceeding to dimensional variables and using Eqs. (28) and (37), one can easily obtain
$$\overline{A}(t)=\frac{dh_m|\mathrm{cos}\omega t|}{2}\frac{1}{6}Hd.$$
(51)
We draw reader’s attention to the fact that the mean value of the vector potential depends on time only via the term $`|\mathrm{cos}\omega t|`$, $`\overline{A}(t)=\overline{A}(|\mathrm{cos}\omega t|)`$. This follows from formulae (35), (28), and (33) for the values $`r`$ , $`u`$ , and $`\xi _0`$ as well as from the relation (27) between $`A(\mathrm{sign}(\mathrm{cos}\omega t)d/2,t)`$ and $`r`$. It also implies that the surface impedance in the main approximation with respect to $`d/\delta `$ has imaginary part (reactance) only. The latter is a consequence of the full transparency of the film.
We start calculation of the reactance $`X`$ with the case of relatively small amplitudes $`h_m<H`$, when the group of trapped electrons exists during the whole period of the wave. Let us substitute expressions (48) and (V) into Eq. (47). Then, the integrals containing $`\overline{A}(t)`$ and $`\stackrel{~}{H}d\mathrm{ln}^2(R_+/d)r^2/4`$ vanish since these functions depend on $`|\mathrm{cos}\omega t|`$ only. By means of formula (34) for the voltage $`U`$, the remaining integral can be transformed into the form,
$$X=\frac{8d\omega }{c^2}_0^{\pi /2}\frac{1+\overline{\sigma }_{tr}(\tau )/\sigma _{fl}(\tau )}{1+(\overline{\sigma }_{tr}(\tau )/\sigma _{fl}(\tau ))(h_m/H)\mathrm{cos}\tau }\mathrm{cos}^2\tau d\tau ,h_mH.$$
(52)
For the case of large amplitudes $`h_m>H`$, one should calculate the reactance using formulae (48), (49), (V), and (51). It represents a sum of two terms,
$`X`$ $`=`$ $`{\displaystyle \frac{8d\omega }{c^2}}[{\displaystyle _{\pi /2\mathrm{arcsin}H/h_m}^{\pi /2}}{\displaystyle \frac{1+\overline{\sigma }_{tr}(\tau )/\sigma _{fl}}{1+(\overline{\sigma }_{tr}(\tau )/\sigma _{fl}(\tau ))(h_m/H)\mathrm{cos}\tau }}\mathrm{cos}^2\tau d\tau `$ (53)
$`+`$ $`{\displaystyle _0^{\pi /2\mathrm{arcsin}H/h_m}}\mathrm{cos}^2\tau d\tau ],ath_m>H.`$
The first term corresponds to the temporal interval when the trapped electrons exist in the sample, and the second one is related to the interval when these particles are absent.
Let us calculate the asymptotics of the surface reactance for the case of rather large amplitudes $`h_mH`$. For this purpose, we rewrite integral (53) in another form,
$`X`$ $`=`$ $`{\displaystyle \frac{8d\omega }{c^2}}[{\displaystyle _0^{\pi /2}}\mathrm{cos}^2\tau d\tau `$ (54)
$`+`$ $`{\displaystyle _{\pi /2\mathrm{arcsin}H/h_m}^{\pi /2}}({\displaystyle \frac{1+\overline{\sigma }_{tr}(\tau )/\sigma _{fl}}{1+(\overline{\sigma }_{tr}(\tau )/\sigma _{fl}(\tau ))(h_m/H)\mathrm{cos}\tau }}1)\mathrm{cos}^2\tau d\tau ].`$
In the second integral, we substitute the variable of integration $`(h_m\mathrm{cos}\tau )/H=\eta `$ and expand the integrand in a power series in the ratio $`H/h_m`$. Then one finds,
$$\frac{X}{X_L}=1+\frac{4}{\pi }(H/h_m)^3_0^1\left[\frac{1+\overline{\sigma }_{tr}(\pi /2)/\sigma _{fl}(\pi /2)}{1+\overline{\sigma }_{tr}(\pi /2)/\sigma _{fl}(\pi /2)\eta }1\right]\eta ^2𝑑\eta ,$$
(55)
where
$$X_L=\frac{2\pi }{c^2}\omega d$$
(56)
is the same as the value of reactance in the absence of the DC transport current. The conductivities $`\overline{\sigma }_{tr}(\pi /2)`$, and $`\sigma _{fl}(\pi /2)`$ are taken at the moment of time when the AC magnetic field $`h_m\mathrm{cos}\omega t`$ turns into zero. Therefore, their ratio is much greater than unity due to inequality (44). Taking into account condition (44), we calculate integral (55) and obtain the following asymptotics for the reactance,
$$\frac{X}{X_L}=1+\frac{2}{3\pi }\left(\frac{H}{h_m}\right)^3,Hh_m.$$
(57)
Now we consider the case of the extremely small amplitudes described by the inequality $`h_mH\sigma _{fl}(\pi /2)/\sigma _{tr}(\pi /2)(H\stackrel{~}{H})^{1/2}`$. The integrand in Eq. (52) can be presented as a power series in $`h_m/(H\stackrel{~}{H})^{1/2}`$. As a result the asymptotic takes the form,
$`{\displaystyle \frac{X}{X_L}}`$ $`=`$ $`{\displaystyle \frac{4}{\pi }}{\displaystyle \frac{\overline{\sigma }_{tr}(\pi /2)}{\sigma _{fl}(\pi /2)}}{\displaystyle _0^{\pi /2}}\left[1{\displaystyle \frac{\sigma _{tr}(\pi /2)}{\sigma _{fl}(\pi /2)}}{\displaystyle \frac{h_m}{H}}\mathrm{cos}\tau \right]\mathrm{cos}^2\tau d\tau `$ (58)
$`=`$ $`{\displaystyle \frac{\sigma _{tr}(\pi /2)}{\sigma _{fl}(\pi /2)}}(1{\displaystyle \frac{8}{3\pi }}{\displaystyle \frac{\sigma _{tr}(\pi /2)}{\sigma _{fl}(\pi /2)}}{\displaystyle \frac{h_m}{H}}),ath_m(H\stackrel{~}{H})^{1/2}.`$
We notice that reactance (58) is $`\sigma _{tr}(\pi /2)\sigma _{fl}(\pi /2)1`$ times greater than that in the absence of the DC current. This is a direct consequence of the effect of amplification of electric signal at the film boundary which was treated in the previous section (see Eq. (44)). The presence of strong DC current in the sample also causes linear behaviour of the reactance in the region of small amplitudes. As shown in Fig 5, the reactance decreases monotonically within the region between asymptotics (58) and (57).
## VI Conclusion
Nonlinear interaction of electromagnetic waves with a strong DC transport current in a thin metal film leads to unusual physical effects due to specific, typical only for metals, magnetodynamic mechanism of nonlinearity. These effects have been studied by analyzing the nonlinear response of the film, which carries the DC current, irradiated bilaterally by electromagnetic wave. The interaction of the wave with the current results in nonanalytical behaviour of the AC electric field on the sample surface which is characterized by appearance of sharp kinks. The increase of the current is accompanied by a rise of the amplitude of oscillations of the electric field at the surface of the sample. This, in turn, causes to the growth of the imaginary part of the surface impedance of the conductor.
The results obtained in this treatment are valid under certain applicability conditions. Firstly, the AC electric field $`\mathrm{\Delta }(x,t)`$ must be small comparing to the potential electric field $`E_0(t)`$. It follows from formulae (III), (40), and (41) that the quantities $``$ and $`\mathrm{\Delta }_m`$, Eq.(44), are of the same order. Therefore, to ascertain the restrictions imposed by the condition $`E_0(t)`$, we can use quantity $`\mathrm{\Delta }_m`$ in the latter condition. The quantity $`\mathrm{\Delta }_m`$should be much less than the minimum value the function $`E_0(t)`$, i.e. the magnitude of potential field (34) for $`\mathrm{cos}\omega t=0`$. The desired inequality reads
$$d^2\frac{h_ml}{HR}\delta _n^2(\omega ),\delta _n^2(\omega )=\frac{c^2}{4\pi \sigma _0\omega },$$
(59)
where $`\delta _n(\omega )`$ represents the characteristic penetration depth of the AC field into a metal under the condition of normal skin effect. Secondly, the non-uniform component of magnetic field inside the film must necessarily be much less than $`h_m`$. This stems from the assumption that the AC magnetic field $`h(x,t)`$ should be quasi-uniform ($`h(x,t)h_m\mathrm{cos}\omega t`$) across the bulk of the film. The maximum value of the non-uniform correction can be estimated from the first of Maxwell’s equations (8) as $`(4\pi \sigma _{tr}\mathrm{\Delta }_md/c)h_m(d/\delta )^2`$, where an effective penetration depth $`\delta (\omega )`$ equals to $`\delta _n(\omega )(R/l)^{1/2}`$. As a result, we come to a requirement of the quasi-uniform property of the AC magnetic field which can be written in the following form:
$$d^2\frac{l}{R}\delta _n^2(\omega ).$$
(60)
Comparing the restrictions imposed by inequalities (59) and (60), it can be seen that condition (59) is more strict at large AC amplitudes, $`h_m>H`$, while for small values of $`h_m`$ one should use inequality (60).
For a sample with thickness $`d=10^3`$ cm, the electron free path $`l=10^1`$ cm, the concentration of electron $`N=10^{23}`$ cm<sup>-3</sup>, the Fermi momentum $`p_F=10^{19}`$ g$``$cm/sec and for magnetic fields $`h_m=H=100`$ Oe, we have $`\omega <10^5`$ sec<sup>-1</sup> using conditions (59) and (60). At such values of the parameters, conditions (59) and (60) are fulfilled as well as condition (5), which states that the value of the mean free path of electron should be large.
Unusual manifestation of specific magnetodynamic mechanism of nonlinearity discussed in the present paper calls for future investigation. In particular it would be very interesting to explore experimentally theoretical predictions made in this work. |
warning/0002/hep-ph0002041.html | ar5iv | text | # 1 Introduction
## 1 Introduction
In this theoretical summary lecture at the High Energy Physics 99 conference of the European Physical Society, I am charged to review some of the new conceptual developments presented at this conference. At the same time, I would like to review more generally the progress of high-energy physics over the past year, and to highlight areas in which our basic understanding has been affected by the new developments. There is no space here for a status report on the whole field. But I would like to give extended discussion to five areas that I think have special importance this year. These are (1) precision electroweak physics, which celebrates its tenth anniversary this summer; (2) CP violation, which entered a new era this summer with the inauguration of the SLAC and KEK B-factories; (3) QCD, which now branches into new lines of investigation; and two rapidly developing topics from physics beyond the Standard Model, (4) supersymmetry spectroscopy and (5) the experimental study of extra dimensions. A sixth important topic, that of neutrino masses and mixing, is covered in the experimental summary talk of Lorenzo Foa . The location of the conference in Finland makes it appropriate to end the lecture with a Lutheran sermon.
## 2 Precision Electroweak Physics
This summer marks the tenth anniversary of a watershed in high-energy physics that took place in the summer and fall of 1989. In that year, the UA2 and CDF experiments announced the first truly precision measurements of the $`W`$ boson mass. These experiments also pushed the mass of the top quark above 60 GeV, thus insuring that radiative corrections due to the top quark would play an important role in the interpretation of weak interaction measurements. SLC and LEP began their high-statistics study of $`Z^0`$ resonance in $`e^+e^{}`$ annihilation. The data from these machines rapidly produced a $`Z^0`$ mass accurate to four significant figures, limited the number of light neutrinos to 3, and began the program of precisely testing the weak-interaction couplings. In the spring of 1989, it was permissible to believe that the $`W`$ and $`Z`$ bosons were composite and that the gauge symmetry of the weak interactions was merely a low-energy approximation. Today, because of the experimental program set in motion that year, this is no longer an option.
It is interesting to contrast our knowledge of the weak interaction parameters in 1989 with our knowledge today. In Figure 1, I show the summary of the constraints on the weak interaction mixing angle $`\mathrm{sin}^2\theta _w`$ that were presented by Altarelli at the 1989 Lepton-Photon conference . The vertical axis shows $`\mathrm{sin}^2\theta _w`$ in Sirlin’s definition
$$\mathrm{sin}^2\theta _w|_{\text{Sirlin}}=1\frac{m_W^2}{m_Z^2}.$$
(1)
The horizontal axis shows the top quark mass, which enters the comparison of weak interaction processes through $`𝒪(\alpha )`$ vacuum polarization diagrams. The constraints shown were all new: the Ellis-Fogli fit to deep-inelastic neutrino and electron scattering , the UA2 and CDF measurements of $`m_W/m_Z`$ , and the SLC measurement of the $`Z^0`$ mass . Contrast the precision of this figure, remarkable at the time, with the small trace labeled ‘1999’. This dot represents our current knowledge of $`m_t`$ and $`\mathrm{sin}^2\theta _w`$.
Within a month after Altarelli showed this figure, the LEP collider began its high-statistics study of the $`Z^0`$ resonance. The precision of these experiments, and their remarkable agreement with the Standard Model predictions, has led to a major change in the way that we think about the weak interactions. Today, we regard $`m_Z`$ as a fundamental constant of Nature, determined to precision of five significant figures and thus standing on a par with $`\alpha `$ and $`G_F`$. The precise values of these three parameters fix the tree-level predictions of the Standard Model. Experiments can then focus on possible small deviations from these predictions, which might be due to the radiative corrections of the Standard Model, or to new physics.
There are several schemes for presenting relatively model-independent constraints on weak interaction radiative corrections. My favorite is to parametrize the $`W`$ and $`Z`$ boson vacuum polarization diagrams in terms of two parameters $`S`$ and $`T`$, scaled so that an effect of order 1 in these parameters corresponds to a correction of order $`\alpha `$ in the electroweak observables . (Other similar parameters sets are defined in .) The parameter $`S`$ is weak-isospin conserving and measures the overall size of a new physics sector; the parameter $`T`$ measures the extent of its weak-isospin violation. The zero point of $`S`$ and $`T`$ is fixed by convention; for my discussion here, I will fix it to correspond to the minimal Standard Model with a top quark mass of 175 GeV and a Higgs boson mass of 100 GeV.
In Figure 2, I present the $`S`$, $`T`$ fit to the corpus of precision electroweak data from the summer of 1989 and from the summer of 1999, both prepared by Morris Swartz . The bottom figure fits into the small dashed rectangle in the top figure. The current best values are
$$S=0.04\pm 0.10,T=0.06\pm 0.11.$$
(2)
What have we learned from this dramatic improvement in our experimental knowledge? I extract three morals:
First, we have learned that the predictions of the minimal Standard Model are amazingly successful! I remind you that the agreement of predictions to better than order-1 on the $`S,T`$ plot requires radiative corrections, and that the experimental success thus tests the Standard Model at the loop level. This aspect is discussed in more detail in .
This level of agreement causes deep difficulty for many schemes of physics beyond the Standard Model. Certain models of new physics have the property that they ‘decouple’ when the scale of new physics becomes large. In brief, this means that new particles of mass $`M`$ produce corrections to $`S`$ and $`T`$ that are of order
$$S,T\frac{1}{4\pi }\frac{m_Z^2}{M^2}$$
(3)
for $`Mm_Z`$. The precision electroweak results imply that, generically, any model of new physics that does not naturally decouple in this way is excluded. This is a severe setback for technicolor models, models with a fourth generation of quarks and leptons, and models in which quarks and leptons are composite. Models with decoupling are typically also models in which the Higgs boson is a fundamental weakly-coupled scalar particle, so the precision electroweak results support this hypothesis.
There are still some notable discrepancies in the picture. In his review at HEP99, Mnich presented a combined value of the $`Z`$ polarization asymmetry of $`b`$ quarks ,
$$A_b=0.893\pm 0.016(SM:0.935).$$
(4)
On the other hand, the value of the $`b`$ fraction of hadronic $`Z`$ decays has now settled down to
$$R_b=0.21642\pm 0.00073,$$
(5)
which agrees with the standard model to 0.3% accuracy. For comparison, technicolor models typically predict a 3% discrepancy . A theorist who wanted to pursue this matter could construct a model with a large deviation in $`A_b`$ and no deviation in $`R_b`$, but essentially all models constructed in advance of the data predicted the opposite pattern.
Further improvements in the precision of the comparison of electroweak data to the Standard Model will require a more precise determination of the renormalization of $`\alpha `$ from $`Q^2=0`$ to $`Q^2=m_Z^2`$. This requires a precise knowledge of the total cross section for $`e^+e^{}`$ annihilation to hadrons through this energy range. We are fortunate that the Beijing Electron Synchrotron has made this measurement a focus of its experimental program and expects to dramatically improve our knowledge of the cross section in the charm threshold region .
Second, we have determined the parameters of the Standard Model to remarkable precision. In particular, we now have accurate values for the fundamental Standard Model coupling constants. In terms of $`\overline{MS}`$ couplings at the scale $`m_Z`$,
$`\alpha ^{}`$ $`=`$ $`1/98.42\pm 0.27`$
$`\alpha _w`$ $`=`$ $`1/29.60\pm 0.08`$
$`\alpha _s`$ $`=`$ $`1/8.40\pm 0.14.`$ (6)
At the same time, the precision electroweak data constraints give us information on the mass of the Higgs boson. In a fit to the minimal Standard Model, one now finds
$$m_H<245\text{GeV}\text{(95\% conf.)}$$
(7)
These results generalize to models with multiple Higgs bosons. Assume, for example, that there are several Higgs bosons $`\varphi _i`$ with vacuum expectation values $`v_i`$, and set
$$v=246\text{GeV}.$$
(8)
Then there is a sum rule $`_iv_i^2=v^2`$ , and the precision electroweak data adds the information
$$\underset{i}{}\left(\frac{v_i^2}{v^2}\right)\mathrm{log}\frac{m_i}{\text{(}245GeV)}<0.$$
(9)
In principle, new interactions at high energy can contribute to the right-hand side of (9); however, in the simplest models, these contributions are negative (corresponding to positive contributions to $`S`$) .
Both of the new pieces of information are encouraging for supersymmetry. The precisely known values of the couplings are consistent with a grand unification of couplings if the renormalization group equations of supersymmetry are used in the comparison . Supersymmetric grand unified theories require a Higgs boson below 180 GeV. A useful reference value is the prediction of the minimal supersymmetric generalization of the Standard Model (MSSM), for large superparticle masses and reasonably large $`\mathrm{tan}\beta `$, $`m_h120`$ GeV. This should be compared to the new direct search limit on the Higgs boson mass announced at HEP99
$$m_h>95.2\text{GeV}\text{(95\% conf.)}.$$
(10)
I am still hoping that the Higgs boson will be found at LEP before its time runs out. If not, there is a new entrant into the race to discover the Higgs boson, the Run II Tevatron experiments. A new analysis of Tevatron capabilities takes account of many possible improvements from earlier studies . The Higgs is searched for in the decay mode $`h^0b\overline{b}`$ in the reaction $`p\overline{p}Wh`$, with a leptonic decay of the $`W`$, and $`p\overline{p}Zh`$, with a $`\nu \overline{\nu }`$ decay of the $`Z`$. The expected improvements in vertex identification and $`b\overline{b}`$ mass resolution are included, and the new ability to trigger on displayed vertices plays an important role. The expected sensitivity of the Tevatron experiments is shown in Figure 3. If the Tevatron can accumulate 20 fb<sup>-1</sup> of integrated luminosity, it should be able to find the Higgs boson in the whole region expected in the MSSM, and in most of the region expected for any model with a weakly-coupled Higgs boson.
Third, we have acquired a tactile appreciation for the Standard Model couplings to quarks and leptons. The beautiful experiments at the $`Z^0`$ resonance do not simply give the $`Z^0`$ couplings as outputs of a fit; they show directly how the Standard Model works. From the many remarkable plots that have come out of the LEP and SLC program, I show three of my favorites in Figure 4. The upper left shows the ALEPH determination of the $`\tau `$ polarization at the $`Z^0`$, in which three decay modes, each with its own characteristic physics, show a 14% excess of $`\tau _L`$ over $`\tau _R`$. The upper right shows the profound effect of electron beam polarization on the $`b`$ angular distribution, characteristic of the almost complete dominance of $`b_L`$ in $`Z^0`$ decays, as observed by the SLD experiment. The lower left shows the OPAL determination of the resonance line-shape for $`Z^0`$ decays to hadrons. The remarkable agreement shown reflects our understanding of all three of the fundamental interactions, weak, through the gross form and precision radiative corrections, strong, through the order $`\alpha _s`$ correction to the decay width to quarks, and electromagnetic, through the distortion of the line-shape by initial-state radiation. In the lower right, I add a new figure shown for the first time at HEP99, the measurement of the $`W`$ boson production and decay angular distributions by L3 . This shows the forward peak in the production angle expected from neutrino exchange, and the correct proportion of events with central values of the decay angle, characteristic of longitudinal $`W`$ polarization. All four plots speak directly to the basic underlying physics. It is no longer a tenable position that the Standard Model is a ‘social construct’; we see its reality before our eyes.
## 3 CP violation
Just as this year marks the completion of an era in electroweak physics, it marks the beginning of an era in the study of CP violation. We have seen one of the major questions about CP violation finally answered, and we have seen the first results from new facilities that will dramatically reshape our experimental knowledge.
To put both developments in perspective, I will begin my discussion with a capsule history of CP violation. The phenomenon was discovered in 1964, in the classic experiment of Christensen, Cronin, Fitch, and Turlay . Almost immediately thereafter, Wolfenstein asked a crucial question : Is CP violation a part of the weak interactions, or is it due to a new interaction at very small distances? Over the years, many models have been proposed in which CP violation arises from weak-interaction couplings of particles with masses of the order of $`m_W`$; the Kobayashi-Maskawa model , in which CP violation is due to quark mixing, the Weinberg model , in which CP violation is due to Higgs boson mixing, and other models in which CP violation comes from phases in the mixing of more exotic species. Behind all of these, though, lurked the possibility of a ‘superweak’ origin for CP violation, in which CP violation arose from a new hard coupling which affected only the $`K^0`$$`\overline{K}^0`$ mixing.
In 1979, Gilman and Wise proposed a crucial test of the weak-interaction origin of CP violation . They showed that such theories typically predict a small but nonzero influence of CP violation on the $`K^0`$ decay amplitudes through the parameter $`ϵ^{}`$. In 1988, the CERN NA31 experiment found a nonzero value for $`ϵ^{}`$ , but this result was not confirmed by the competing experiment E731 at Fermilab . Finally, this year, the new high precision experiments NA48 and KTeV agree that $`ϵ^{}`$ is nonzero and find quite similar values . The new world average presented at HEP99 is
$$ϵ^{}/ϵ=(21.2\pm 4.6)\times 10^4.$$
(11)
The nonzero result in (11) rules out a superweak origin of CP violation. The specific value is too small to be compatible with the original Weinberg model. It is an interesting question whether the value can be compatible with the Kobayashi-Maskawa model or whether it requires new particles with CP violating couplings. This topic was discussed at length at HEP99 , and I would like to give my impression of the current situation.
Though the complete formula for $`ϵ^{}/ϵ`$ in the Standard Model is very complicated, one can argue about the uncertainties in $`ϵ^{}/ϵ`$ by using the simplified approximate relation
$`ϵ^{}/ϵ`$ $`=`$ $`10^4`$ (12)
$`\left[16B_6^{(1/2)}8B_8^{(3/2)}\left({\displaystyle \frac{m_t}{165}}\right)^{2.5}\right]\left({\displaystyle \frac{110}{m_s}}\right)^2,`$
where $`m_t`$ is the $`\overline{MS}`$ top quark mass evaluated at $`m_t`$, $`m_s`$ is the $`\overline{MS}`$ strange quark mass evaluated at 2 GeV, and $`B_6^{(1/2)}`$ and $`B_8^{(3/2)}`$ are conventionally defined factors giving the matrix elements of penguin operators arising from the strangeness-changing weak interaction. The convention for the $`B`$ coefficients factors out the dependence on the strange quark mass, and one should keep in mind that it is the combination $`B/m_s^2`$ which corresponds to a physical matrix element. These matrix elements must be determined by a nonperturbative technique, for example, lattice QCD simulation. The first term in (12) is due to the strong-interaction penguin diagram, the second to the electroweak penguin (as in Figure 5). The strong cancellation between these two effects for large top quark mass is the reason that the observed value of $`ϵ^{}/ϵ`$ is much smaller than the original prediction of Gilman and Wise . The cancellation also amplifies the considerable uncertainties in the operator matrix elements.
New estimates of the parameter $`B_8^{(3/2)}`$ were reported at HEP99:
$$B_8^{(3/2)}\left(\frac{110}{m_s}\right)^2=\{\begin{array}{cc}0.91\pm 0.16\hfill & \text{[40]}\hfill \\ 0.86\pm 0.07\hfill & \text{[41]}\hfill \end{array}.$$
(13)
Both estimates are given in the chiral limit $`m_s0`$. (For the true value of $`m_s`$, one should multiply these estimates by 1.3.) The first of these estimates is based on perturbative QCD analysis of spectral-function sum rules, the second is derived from a lattice QCD calculation using the new technique of domain-wall chiral fermions . From the agreement, it seems that this part of the problem is now fairly well understood. Unfortunately, the situation for $`B_6^{(1/2)}`$ is much worse. This matrix element vanishes in the chiral limit and in the $`SU(3)`$ limit, making the usual techniques for both lattice gauge theory and QCD estimates awkward to apply. For perturbative QCD estimates, $`B_6^{(1/2)}`$ depends on the scalar and pseudoscalar spectral functions, which are poorly known. The operator which is responsible for the $`\mathrm{\Delta }I=\frac{1}{2}`$ rule enhancement of the $`K^0\pi \pi `$ matrix element has similar problems, and, indeed, to this day no lattice gauge theory calculation has been able to compute the $`\mathrm{\Delta }I=\frac{1}{2}`$ enhancement accurately. Thus, the value of $`B_6^{(1/2)}`$ is not known, and this uncertainty can easily allow one to reconcile the prediction (12) with the observed value (11).
We have now reached the situation in which we know that CP violation arises from weak-interaction couplings, but we do not have a sufficiently good theoretical understanding of the measured observables to know whether CP violation is accounted for by the Kobayashi-Maskawa model or whether new particles with CP violating couplings are required. Fortunately, we are entering a new era in which the SLAC, KEK, and Cornell $`e^+e^{}`$ B-factories, the HERA-B experiment, and measurements of $`B`$ decay at high-luminosity hadron colliders will provide measurements of new CP violation observables which can be interpreted with very small theoretical uncertainty. This new era offers us a remarkable opportunity either to put the conventional picture of CP violation on a firm footing or to overturn it and discover signal of new physics. In order to do this, however, we must change our view of what the important CP violation observables are and how we should compare them.
An example of the CP violation observables of the new era is the time-dependent asymmetry $`𝒜`$ in an exclusive $`B`$ decay, an observable first discussed by Carter and Sanda . For example,
$$\mathrm{\Gamma }(\overline{B}^0\text{or}B^0J/\psi K_S^0)e^{\mathrm{\Gamma }t}\left[1\pm 𝒜\mathrm{sin}\mathrm{\Delta }m_dt\right].$$
(14)
In this equation, $`\mathrm{\Delta }m_d`$ is the $`B_L^0`$$`B_S^0`$ mass difference, and the sign $`\pm `$ refers to the two possible initial states. The parameter $`𝒜`$ is manifestly CP violating and can be extracted with essentially no uncertainty from our knowledge of hadronic matrix elements. In the Kobayashi-Maskawa model, $`𝒜=\mathrm{sin}2\beta `$, where $`\beta `$ is simply related to the phase in the Cabibbo-Kobayashi-Maskawa (CKM) mixing matrix. In models with new CP violating couplings, $`𝒜`$ can obtain additional large contributions from these sources.
At HEP99, we had a first taste of the new era of CP violation with the report of the first significant measurement the CP asymmetry in $`B^0J/\psi K^0`$ by the CDF collaboration . The experiment observed the $`J/\psi `$ in its decay to $`\mathrm{}^+\mathrm{}^{}`$ and the $`K_S^0`$ in its decay to $`\pi ^+\pi ^{}`$. The initial flavor of the $`B^0`$ was determined either by the lepton charge or jet charge on the opposite side of the event, or by the charge of a pion accompanying the $`B^0`$ in the same jet. The figure of merit for such flavor tags, giving the fraction of the event sample that corresponds to the effective number of perfectly tagged $`B^0`$’s, is
$$ϵD^2,$$
(15)
where $`ϵ`$ is the efficiency of the tag and $`D`$ is the dilution, the difference between the probability of a correct tag and the probability of a wrong tag. (In the next decade, high-energy physicists will mutter ‘$`ϵD^2`$’ as often as, in last one, they were heard to mutter ‘$`\mathrm{sin}^2\theta _w`$’.) For the CDF measurement, the $`ϵD^2`$ for each of the three tagging methods is about 2%, so that the sample of 400 events corresponds effectively to 25 tagged $`B^0`$ decays. From this, one finds
$$𝒜=0.79_{0.44}^{+0.41},$$
(16)
roughly a $`2\sigma `$ determination that $`𝒜>0`$. I show the data binned as a function of $`t`$ in Figure 6 and leave it to you to judge the quality of the evidence.
Within a year or so, we should have the first accurate measurements of these new CP observables, and we will need a framework to use in comparing them. A useful pictorial device is the ‘unitarity triangle’ , the triangle in the complex plane which reflects the unitarity relation of CKM matrix elements
$$V_{ub}V_{ud}^{}+V_{cb}V_{cd}^{}+V_{tb}V_{td}^{}=0,$$
(17)
Using the approximations $`V_{ud}V_{tb}1`$, $`V_{cd}\mathrm{sin}\theta _C`$, we find the relation shown in Figure 7. The internal angles of this triangle are referred to as $`\alpha `$, $`\beta `$, $`\gamma `$, except in the Far East, where the notation $`\varphi _2`$, $`\varphi _1`$, $`\varphi _3`$ is used.
It is often said that the goal of the new CP violation measurements is to ‘check whether the unitarity triangle closes’. I would like to substitute for this a more precise idea.
Since CP violating phases can be redefined by convention, CP violation observables typically involve phase differences between two different amplitudes. Usually, these are $`B`$ or $`K`$ mixing amplitudes or other loop diagrams on one hand, and weak decay amplitudes on the other hand. I will assume that the phases of the decay amplitudes come only from the CKM matrix elements. This is correct unless the decay amplitudes also receive corrections from the tree-level exchange of light exotic particles such as charged Higgs bosons. On the other hand, a loop diagram which contribute to mixing can receive corrections from any new particles with masses in the range up to 1 TeV, and the couplings of these particles can bring new contributions to its phase. If we try to determine the unitarity triangle from a set of processes which all involve the same loop diagram, it is possible to get a consistently determined triangle which does not coincide with the true unitarity triangle of the CKM matrix. The way to test models of CP violation, then, is to compare the unitarity triangles determined from different classes of CP observables. This point of view, set out in the original work of Nir and Silverman , has been emphasized more recently by Cohen, Kaplan, Lepeintre, and Nelson and by Grossman, Nir, and Worah .
I would now like to distinguish four classes of CP violation measurements, corresponding to four different physical systems, such that each class would determine the unitarity triangle completely if the Kobayashi-Maskawa model were a complete description of CP violation. The test of the Kobayashi-Maskawa model will come from the comparison of these triangles. The four triangles that I will discuss are shown in Figure 8, with error boxes for the sides or angles that might be realized within the next decade.
Figure 8(a) shows the ‘non-CP triangle’. This triangle takes advantage of the fact that one can determine the unitarity triangle by measuring the absolute values of CKM matrix elements and thus show the existence of the phase through non-CP-violating observables. The left-hand side of the triangle is determined by the rate of $`bu`$ weak decays; the right-hand side is determined by ratio of $`B`$$`\overline{B}`$ mixing amplitudes for $`B_s`$ and $`B_d`$. The rate of $`bu`$ transitions depends only on the CKM matrix element $`V_{ub}`$ and is not affected by new physics. The $`B`$ mixing amplitudes involve box diagrams that might have large nonstandard contributions. However, in many models, including models with light supersymmetric particles in which squarks with the same electroweak quantum number are naturally degenerate, these contributions have the same ratio as the standard contributions . Thus, this ‘non-CP triangle’ is the most likely of the four to agree with the true unitarity triangle determined from the CKM matrix.
The expected accuracy that I have displayed in this figure—10% for the $`V_{ub}`$ side and 5% for the $`V_{td}`$ side—is surprisingly small, and I would like to defend these estimates now. I will begin with $`V_{td}`$. This parameter is determined by the relation
$$\frac{\mathrm{\Delta }m_d}{\mathrm{\Delta }m_s}=\frac{m_{Bd}f_{Bd}^2B_{Bd}}{m_{Bs}f_{Bs}^2B_{Bs}}\left|\frac{V_{td}}{V_{ts}}\right|^2=\xi ^1\left|\frac{V_{td}}{V_{ts}}\right|^2,$$
(18)
where $`f_{Bd}`$ is the $`B_d`$ decay constant and $`B_{Bd}`$ is the matrix element of a 4-fermion operator in the $`B_d`$ wavefunction. The $`B_d`$ mixing parameter $`\mathrm{\Delta }m_d`$ is now known to 3.5% accuracy . The $`B_s`$ mixing parameter can be determined by looking for a fast oscillation in tagged $`B^0`$ decays superimposed on the slow oscillation from $`B_d`$ mixing. There is suggestive evidence that such an oscillation appears in the $`B`$ vertex distribution at the $`Z^0`$, corresponding to an oscillation frequency $`\mathrm{\Delta }m_s16`$ ps<sup>-1</sup> ; I have used this value in constructing the figure. Once the oscillation is seen, the frequency can be determined to a few percent. The CDF experiment should be able to make this measurement early in Run II, even for $`\mathrm{\Delta }m_s`$ so large that the triangle collapses onto the real axis. Looking back at (18), the magnitude of $`V_{ts}`$ is constrained by unitarity to be very close to $`|V_{cb}|`$. Thus, the main source of uncertainty is in the estimation of the $`\xi `$. The ratio $`\xi ^1`$ is roughly equal to 0.8 and tends to 1 in the chiral limit or in the $`SU(3)`$ limit $`m_d=m_s`$. To achieve 5% accuracy in $`\xi `$, it is only necessary to compute the deviation of $`\xi `$ from 1 to 25% accuracy.
Lattice gauge theory should be up to the task. At HEP99, the CP-PACS collaboration reported a calculation
$$\left(\frac{f_{Bd}^2}{f_{Bs}^2}\right)=0.69(1\pm 7\%\pm 3\%{}_{7\%}{}^{+5\%}),$$
(19)
where the three errors come from Monte Carlo statistics, the determination of $`m_s`$, and the continuum extrapolation. It seems to me that, with further effort, a 5% determination of $`|V_{td}/V_{cb}|`$ is quite feasible. A useful review of the status of lattice gauge theory determinations of these and other heavy-quark matrix elements can be found in .
The situation is less clear for $`V_{ub}`$, but still there is reason for optimism . The best current measurement of $`V_{ub}`$ is based on the CLEO measurement of the rate of $`B\rho \mathrm{}\nu `$ ,
$$V_{ub}=(3.25\pm 0.14{}_{0.29}{}^{+0.21}\pm 0.55)\times 10^3,$$
(20)
where the third contribution to the error represents a 20% spread in the relations given by models between the underlying parameters and the observed rate. The experimental uncertainties are thus quite adequate, and they will decrease in the era of the B-factories. What is needed is a method for computing $`V_{ub}`$ that has less model uncertainty. Two methods have been proposed. The first is an inclusive technique based on the idea that in a decay $`BX\mathrm{}\nu `$, if $`m(X)<m_D`$, then the decay must be a $`bu`$ transition . The problem with this method is that energy from neutral particles cannot be unambiguously associated with a displaced vertex, so one must work with vertex masses based on charged particles and use models to estimate the background from $`bc`$ decays. The DELPHI collaboration has made a promising first application of this technique , obtaining
$$V_{ub}=(4.1\pm 0.5\pm 0.6\pm 0.3)\times 10^3,$$
(21)
where the last error indicates a 8% model uncertainty. It is a very interesting question how one defines the optimized vertex mass for this measurement applicable to the B-factory environment. The second method is to measure the spectrum of $`B\rho \mathrm{}\nu `$ decays as a function of $`m(\mathrm{}\nu )`$ and evaluate it at the ‘zero-recoil’ point where the heavy $`B`$ quark decays to a $`u`$ quark at rest. The value of the form factor at this point can be computed by lattice gauge theory simulations .
Figure 8(b) shows the ‘B triangle’. This triangle is constructed from the CP asymmetries in $`B^0/\overline{B}^0`$ decays. To draw the figure, I have used the asymmetry in $`BJ/\psi K_S^0`$ and the asymmetry in $`B\rho \pi `$. (I ignore the discrete ambiguities in determining the CKM angles from the measured asymmetries.) Both of these asymmetries involve the phase in the $`B^0`$$`\overline{B}^0`$ mixing amplitude and are sensitive to new physics through this source. For $`BJ/\psi K_S^0`$, at least four independent experiments (BaBar, BELLE, CDF, HERA-B) should determine $`\mathrm{sin}2\beta `$ an accuracy better than $`\pm 0.1`$ in the near future. LHC-B or BTeV should determine this parameter to the level of $`\pm 0.01`$. The constraint from $`B\rho \pi `$ is actually a measurement of $`\alpha `$ in the CKM picture, but I have moved the constraint to the lower vertex of the triangle for clarity. The process originally thought to best determine $`\alpha `$, $`B\pi \pi `$, is now disfavored due to potential large background contributions from strong and electromagnetic penguin diagrams. With sufficient statistics, one can fit to the Dalitz plot in $`B\rho \pi `$ to measure and remove these contributions. The BaBar collaboration has estimated an accuracy of $`7^{}`$ in $`\alpha `$ for a sample of 600 events . Such a large sample would require luminosities well above the design level. It is also possible to measure $`\gamma `$ by the comparison of rates for $`B^\pm K^\pm D`$ decays . This determination involves only tree-level decay amplitudes and so measures the true CKM unitarity triangle rather than the ‘B triangle’.
Figure 8(c) shows the ‘B<sub>s</sub> triangle’. The time-dependent CP asymmetry in $`B_sD_s^\pm K^{}`$ is connected to $`\mathrm{sin}\gamma `$. LHC-B or BTeV should measure $`\gamma `$ using this reaction to about 10. The $`B_s`$ system also allows an interesting null experiment. The time-dependent CP violation in $`B_sc\overline{c}s\overline{s}`$ decays is expected to be very small in the Standard Model. On the other hand, the phase in $`B_sJ/\psi \eta `$ and $`B_sJ/\psi \varphi `$ should be measurable to a few degrees by LHC-B or BTeV. These reactions will be a very sensitive indicator for new CP violating physics in the $`B_s`$$`\overline{B}_s`$ mixing amplitude. This constraint is shown, just for the purpose of illustration, as a constraint on the base of the unitarity triangle.
Figure 8(d) shows the ‘K triangle’. This is the triangle determined by the rare $`K`$ decays $`K^+\pi ^+\nu \overline{\nu }`$, which has an amplitude approximately proportional to $`V_{td}`$ in the Standard Model, and $`K_L^0\pi ^0\nu \overline{\nu }`$, which is a CP-violating process whose amplitude is proportional to Im\[$`V_{td}`$\] in the Standard Model. These decays proceed through box diagrams which could well have exotic contributions from new particles with masses of a few hundred GeV. The rare $`K`$ decays are frighteningly difficult to detect. Experiment E787 at Brookhaven has recently observed 1 event of the $`K^+`$ decay . There are preliminary plans for experiments that would run at Fermilab in the next decay and collect 100 events in each of these rare modes; I have drawn the triangle assuming that these experiments succeed and that their statistical errors are dominant. I have also included in this plot the lower bound of the constraint from $`ϵ`$, with conservative errors reflecting the uncertainty in a lattice determination of the overall normalization of hadronic matrix element. It should be noted that, while large deviations from the CKM prediction are possible in rare $`K`$ decays, broad classes of models give only a relatively small effect .
Figure 9 shows the four unitarity triangles superposed on one another. This could well indicate the state of our knowledge of CP violation ten years from now. If the agreement of the various triangles is as good as what is shown here, it will provide striking evidence that the Kobayashi-Maskawa model explains the observed CP violation in weak interactions. But keep in mind the possibility that these four triangles might disagree completely due to loop diagrams involving new heavy particles. We will soon find out which alternative is realized.
The HEP99 meeting also saw new theoretical developments in the theory of $`B`$ meson CP asymmetries. In my discussion of $`B\rho \pi `$ above, I mentioned the difficulties associated with penguin diagrams, which modify the current-current weak interaction and potentially add a different set of phases. It is an important issue in the theory of $`B`$ decays to provide methods to calculate these penguin contributions or to extract them from data. I would like to highlight three recent pieces of work along these lines. In a presentation at HEP99, Fleischer proposed using $`SU(3)`$ (more specifically, U-spin) to relate the decay amplitudes for $`B\pi ^+\pi ^{}`$ and $`B_sK^+K^{}`$. Using these relations, it is possible to solve for $`\beta `$ and $`\gamma `$ without assumptions about the size of the penguin effects. In another recent paper Neubert and Rosner , following up on ideas of Fleischer and Mannel , have shown how to extract $`\gamma `$ without assumptions on the size of the penguin contributions by fitting all partial rate differences among $`B^\pm \pi K`$ and $`B^\pm \pi ^\pm \pi ^0`$ decays. In the most ambitious of these projects, Beneke, Buchalla, Neubert, and Sachrajda reported a new factorization formula applicable to the decays of a $`B`$ meson to two pseudoscalars meson valid in the formal limit $`m_b\mathrm{}`$. In perturbative QCD, the leading term in this formula is the naive factorization in which one current from the weak interaction operator creates one final meson. However, the corrections to this term are finite and calculable. A typical result of their method is the formula
$`BR(\overline{B}^0\pi ^+\pi ^{})`$ $`=`$ $`6.5\times 10^6`$ (22)
$`\left|e^{i\gamma }+0.09e^{i(13^{})}\right|^2.`$
On the right-hand side, the prefactor should not be taken seriously, but it can eventually be well determined because it is computed from the same form factor that will be measured in the decay $`B\pi \mathrm{}\nu `$. The phase in the second term inside the bracket arises from the imaginary part of a QCD loop diagram. It would be very interesting to understand the accuracy of the formulae obtained by this method, since potentially they make many more processes available for the determination of CKM matrix elements.
Before leaving the subject of CP violation, I would like to remind you that there are many other possible probes which should be explored. Even within the realm of meson asymmetries, there is $`D`$$`\overline{D}`$ mixing, which could have a large CP violating component from sources beyond the Standard Model . Our knowledge of $`D`$$`\overline{D}`$ mixing will be greatly improved by the B factory experiments. Already, CLEO has a new and very impressive limit, which was presented at HEP99 . The neutron and electron electric dipole moments remain important constraints, especially on new angles in supersymmetry that couple to light flavors. CP violation might also be specifically associated with the top quark. The LHC experiments should be able to observe a $`10^3`$ energy asymmetry between leptons $`\mathrm{}^\pm `$ produced in top decays, and this is would be a significant test of CP violation models .
Finally, it is important to keep in mind that the Kobayashi-Maskawa model of CP violation cannot provide large enough CP asymmetries to create the baryon/antibaryon asymmetry in the universe . So, we must eventually find a new source of CP violation. It is possible that this source is the mass matrix of heavy leptons, or some other effect at extremely high energy. But it is also possible that the new mechanism of CP violation will appear in the experimental program that we are now beginning to carry out.
## 4 QCD
We turn next to Quantum Chromodynamics (QCD). In contrast to the previous two topics, the fundamental questions about QCD have been answered already some time ago. The experimental confirmation of QCD as the theory of the strong interactions is now very strong. QCD is now known to account for a wide variety of processes with large momentum transfer, in $`e^+e^{}`$ annihilation, $`ep`$ collisions, and $`p\overline{p}`$ collisions, to about the 10% level of accuracy, and with a common value of the coupling constant $`\alpha _s`$. At the same time, numerical lattice studies confirm that QCD explains the spectrum of light hadrons to about the same level of numerical precision. Wittig has reviewed this latter, less familiar, evidence for QCD at HEP99 . So, what are the important scientific issues for QCD today? I would like to highlight four of these, and give some examples of new work presented at HEP99.
The first issue is the precision determination of $`\alpha _s`$. At the moment, the $`\overline{MS}`$ coupling $`\alpha _s(m_Z)`$ is known to about 3% accuracy . It is important to reduce this error below 1%. This level of accuracy is needed as an input to the precision experiments of the next decade, for example, the study of the top quark at threshold. It is also already needed to assess the validity of grand unification. I have already noted that grand unification with the renormalization group equations of supersymmetry successfully relates the values of the Standard Model couplings given in (6). In particular, the prediction for $`\alpha _s`$ agrees with experiment at the 10% level, but it is subject at this level to uncertainties from threshold corrections at the scale of grand unification. With a more accurate $`\alpha _s`$, we could evaluate the needed threshold contribution and begin to test explicit models of grand unification. The primary barrier to a more accurate determination of $`\alpha _s`$ come not from experiment (though it would be good to have more precise data on multi-jet rates at energies well above the $`Z^0`$) but rather from theory. However difficult it may be, we need the order $`\alpha _s^2`$ corrections to the most important processes which determine $`\alpha _s`$, in particular, the rate for $`e^+e^{}3`$ jets.
The second issue is the determination of essential strong interaction parameters needed for high energy experiments. Here I mean especially the parton distributions in the proton. Though the quark distributions are well determined, the gluon distribution is not well constrained at moderate values of $`x`$. This is the freedom that was used to correct the discrepancy between the CDF measurement of the jet rate at large $`E_T`$ and QCD predictions . Gluons at moderate $`x`$ and low $`Q`$ evolve to the gluons at low $`x`$ and high $`Q`$ which are the dominant source of new particle production at the LHC. The study of high-energy $`e^+e^{}`$ collisions requires another set of input data, the total cross section for the process $`\gamma \gamma `$ hadrons. It is interesting in its own right to understand what part of this cross section comes from pointlike processes and what part from soft processes involving the hadronic constituents of the photon. The eventual theory should explain, as be constrained by, the data both for $`\sigma (\gamma \gamma )`$ and $`\sigma (\gamma p)`$.
The third issue is the study of the detailed structure of jets as predicted by QCD. QCD predicts that the hardest components of jets are built up by successive processes in which gluons or quarks split off from the hardest parton. This gives jets a fractal structure. On top of this backbone, hadrons are produced in a way that reflects the color pairings of the hard partons. If groups of partons are separated in phase space and are separately color neutral, we should find a phase space or rapidity gap in particle production. All of these features are just beginning to be understood from the data.
The fourth issue is one for theorists, the development of new techniques to compute higher-order and multiparton QCD amplitudes. This issue provides essential theoretical support to the first and third topics just listed. At HEP99, Uwer , Draggiotis , and Harlander set out new ideas for calculational programs whose results should be very interesting.
As an illustration of recent progress in these areas, I present in Figure 10 four of the new QCD results presented at HEP99. The upper left shows that event shape measurements from HERA are now contributing to the precision $`\alpha _s`$ determination. The upper right shows the new determination by L3 of the $`\gamma \gamma `$ total cross section at LEP. (A similar determination was presented for the $`\gamma p`$ total cross section at HERA .) The large uncertainty, reflected in the difference between the two fits, comes from the fact that about 40% of the total cross section is unobserved, and that theoretical models differ on the size of the contribution from these very soft events. This is a problem that must be addressed. The last two figures show new studies of QCD event shapes. The lower left shows the narrowing of jets with increasing $`Q^2`$ in $`ep`$ collisions in the H1 event sample. The lower right, from the D0 experiment , shows that the particle production in $`W+`$ jet events reflects the color flow expected for a color singlet $`W`$ recoiling against a colored parton.
## 5 Supersymmetry
From the areas of current experimental interest, we now turn to the future. What should be the main topic of discussion at HEP09? What new era in experimental high energy physics will be opening up at that time?
In this section, I would like to take very seriously the second moral I drew in Section 2 from the precision electroweak data. The precision study of the $`Z^0`$ points us toward a world in which the interactions responsible for electroweak symmetry breaking are weakly coupled and the Higgs boson is an elementary scalar particle. I have already explained that some specific aspects of the data are compatible with the idea of supersymmetry at the weak interaction scale. But there is a much stronger argument for the presence of supersymmetry in the fundamental description of Nature. While it is possible in principle that there is no explanation for the negative value of the Higgs (mass)<sup>2</sup> and the instability to symmetry breaking in the Higgs potential, I insist that we must find a way to explain this instability on the basis of physics. For this, we must have a theoretical framework for the Higgs field in which its potential energy function is calculable. A part of this requirement is that some symmetry must forbid the addition of a Higgs mass term by hand. If the Higgs boson is treated as an elementary field, the only known symmetry with this power is supersymmetry. Thus, to the extent that the precision electroweak data excludes models in which the Higgs is composite or strongly coupled, we should expect to see not a light Higgs boson but also the new particles predicted by supersymmetry.
The idea that the data drives us to a weak-coupling picture of the Higgs boson was controversial at HEP99. Among the people arguing vocally on the other side were Holger Nielsen and Gerard ’t Hooft. So if you do not wish to accept this argument, you are in good company (but still wrong). In any event, for the rest of this lecture I will take this conclusion very seriously and use it to map out future questions for high-energy experimentation.
Among the many theoretical problems connected with supersymmetry, I would like to focus on the spectrum of supersymmetric particles. Supersymmetry predicts that every particle of the Standard Model has a partner with the opposite statistics. That is, the chiral fermions of the Standard Model have scalar partners, and the gauge bosons have spin-$`\frac{1}{2}`$ (‘gaugino’) partners. What are the masses of these particles?
The interest of this question goes beyond the issue of where or when these particles will be found. To produce a reasonable phenomenology, supersymmetry must itself be spontaneously broken. The supersymmetrized Standard Model cannot directly break supersymmetry, because this hypothesis leads to unwanted very light superpartners . In most models, the description of the superparticle masses involve two ingredients, a sector in which supersymmetry is broken and a ‘mediator’ which connects the symmetry-breaking to the Standard Model fields. The identity of the mediator is typically connected to the very short distance physics of the model, the connection of the Standard Models fields to grand unification or to gravity. And, the nature of the mediator is reflected in the detailed pattern of the masses of superparticles.
If supersymmetry explains the phenomenon of electroweak symmetry breaking, the superparticle masses must have masses of a few hundred GeV; that is, they should be accessible to the experiments of the next decade and possibly even to LEP and the Tevatron. By measuring this mass spectrum, we should, ten years from now, have a wealth of new data which speaks directly to the physics at this fundamental level. It is an exciting prospect. To help you think about it more clearly, I would like to offer a first lesson in superspectroscopy. I will contrast three paradigmatic theories of the mediation of supersymmetry breaking. Though supersymmetry phenomenology has been studied for almost twenty years and went through a period of complacency in the early 1990’s, two of these paradigms were discovered only recently. The second was invented in 1995, and the third only in the past year. Presumably, there are more theoretical insights into this subject that are waiting to be uncovered.
The three paradigms for the superspectrum that I would like to discuss are those of ‘gravity mediation’, ‘gauge mediation’, and ‘anomaly mediation’. In gravity mediation, as introduced in , the mediator is supergravity itself. One imagines a sector that spontaneously breaks supersymmetry. Let $`m_{\text{Pl}}`$ be the Planck scale. Then the gaugino masses arise from direct order $`1/m_{\text{Pl}}`$ couplings of this sector to the Standard Model Yang-Mills Lagrangian, and scalar masses arise from direct order $`1/m_{\text{Pl}}^2`$ couplings of this sector into the Standard Model potential. If the superparticle masses are of the order of the weak scale, the mass of the gravitino is of the same order. The spectrum acquires additional structure from the renormalization group evolution of the mass parameters from the Planck scale to the weak scale.
In gauge mediation, introduced in , the mediator is the supersymmetric version of the Standard Model gauge interactions. One imagines that supersymmetry is broken by a new sector which includes heavy particles with Standard Model gauge couplings. Then gaugino masses arise from one-loop diagrams involving these heavy particles, and scalar masses arise from two-loop diagrams in which the loop of heavy particles appears in a scalar self-energy diagram. Because what is computed for scalars is the (mass)<sup>2</sup>, both gauginos and scalars acquire masses of order $`\alpha `$ with respect to the heavy scale. The gravitino has a mass of order eV and is typically the lightest superparticle. Other superparticles may be observed to decay to the gravitino if the rate of this decay is not suppressed by a very large value of the heavy mass.
Anomaly mediation, introduced in represents the opposite extreme pole. To realize this possibility, one considers models in which the supergravity couplings needed to work gravity mediation vanish. Then the partners of Standard Model particles acquire no mass at tree level. In this case, it turns out that the first mass contribution is universal in character and is connected to the breaking of scale invariance by the running of the Standard Model coupling constants. The gravitino mass does arise at the tree level, and so in this scheme the gravitino mass is about 100 TeV if the Standard Model superparticles acquire weak scale masses.
In Table 1, I have collected the basic formulae for the gaugino and scalar masses in these three paradigms. For clarity, I have eliminated the underlying parameters in terms of the mass $`m_2`$ of the $`W`$ boson superpartner. In the case of gravity mediation, there is another independent parameter $`m_0`$. For gauge mediation and anomaly mediation, the superparticle masses naturally depend only the Standard Model quantum numbers, a feature that suppresses new flavor-changing neutral current effects from supersymmetry loop diagrams. This property also holds in gravity mediation in the ‘no-scale’ limit $`m_00`$. Away from this limit, there is no clear reason why $`m_0`$ should not depend on flavor except that this could lead to unwanted neutral current effects.
Anomaly mediation makes two specific predictions that deserve comment. First, the lightest gaugino are expected to be an almost degenerate $`SU(2)`$ triplet $`\stackrel{~}{w}^\pm `$, $`\stackrel{~}{w}^0`$, with the charged states only a few hundred MeV above the neutral one . This gives rise to a distinctive phenomenology, discussed in . Second, the scalars partners of leptons are computed to have negative (mass)<sup>2</sup>, a disaster. Some cures for this problem are given in .
In Figure 11, I show a comparison of the spectra for the three paradigms; for gravity mediation, I give both the ‘no-scale’ case $`m_0=0`$ and the case $`m_0/m_2=5`$ and universal.
I emphasize that the formulae in Table 1 represent only the first lesson in superspectroscopy. They omit possible mass mixings and effects of large top, bottom, and $`\tau `$ Yukawa couplings, and they omit higher-order corrections . Nevertheless, these formulae and Figure 11 already give a feeling for the complexity of the spectrum that might be found when superparticles appear in experiments.
If supersymmetry is the explanation of electroweak symmetry breaking, it is likely that the LHC will be able to sample the whole superparticle mass spectrum, including the heaviest states. In addition, as members of the ATLAS collaboration have recently demonstrated, the LHC experiments have the ability to make precision measurements of superparticle masses in a number of different scenarios . Nevertheless, the richness of the phenomena calls for the exploration of these particles also in $`e^+e^{}`$ annihilation. It is worth remembering that an $`e^+e^{}`$ linear collider of the next generation will provide not only a relatively clean environment with kinematic constraints that aid in particle mass measurements, but also the availability of beam polarization, which is very useful in resolving questions of particle mixing . The production cross sections for superparticles are electroweak and can be computed precisely, allowing an unambiguous determination of the quantum numbers of each new particle.
The superspectrum is complex, but the LHC and linear collider are powerful instruments. With these two facilities, with their complementary strengths, we could fully explore the supersymmetry spectrum of particles and mine the information it contains for information about a truly fundamental level of physics. This is already cause for optimism about the future of experimental particle physics. But, there is more.
## 6 New space dimensions
Many people say that the key problem of quantum physics has nothing to do with what we do at accelerators. Rather, they say, it is the problem of the compatibility of quantum mechanics with general relativity. To solve this problem, one must do two things, first, remove the divergences from the quantum theory of gravity, and, second, unify gravity with the microscopic particle interactions.
Most people who recite this litany do not realize that we have at least one possible solution already in hand. It is string theory. String theory has not been proved to be the correct theory of Nature, but it does demonstrably solve these two problems. It is the only known approach to these problems which has no glaring weaknesses. Therefore, we must take it very seriously.
There has been tremendous progress in string theory since the 1995 discovery of string dualities by Hull and Townsend and Witten . Just in the past year, two new and very profound ideas have been validated: The first is the idea of ‘t Hooft and Susskind that quantum gravity is ‘holographic’, in the sense that its physical degrees of freedom are those of a manifold with one lower dimension that the dimension of space time. The second is an explicit realization of this relation, due to Maldacena , a duality linking supersymmetric Yang-Mills theory in 4 dimensions with supergravity in 5-dimensional anti-de Sitter space. I do not have space here to do justice to these ideas, but they are described clearly in Bachas’ lecture at HEP99 .
I would like to concentrate instead on another consequence of the new understanding of string dualities. These developments have led to new classes of models in which quantum gravity and string physics is much more accessible to experiment and may even appear directly in the realm of the LHC and the linear collider.
String theory requires that we live in a world which has 11 dimensions. Until recently, it was thought that this could only be compatible with our observations if seven of these dimensions were compact and very small, of the order of the Planck scale. It is interesting, though, to think about new space dimensions that are not so small. In 4 dimensions, the gravitational force falls off as
$$F\frac{m_1m_2}{r^2};$$
(23)
in $`(4+n)`$ dimensions, it falls off as
$$F\frac{m_1m_2}{r^{2+n}}.$$
(24)
Consider the $`n`$ extra dimensions to be periodic with period $`2\pi R`$. Then two masses separated by a distance much larger than $`R`$ would feel a gravitational force of the form (23), while two masses separated by a distance much less than $`R`$ would feel a gravitational force of the form (24). The short-distance force is the more fundamental. We can define the fundamental quantum gravity scale $`M`$ by writing the dimensionful constant of proportionality in (24) as a numerical constant times $`M^{(n+2)}`$. The constant of proportionality in (23) is just Newton’s constant. If we insist that the forces match at the distance scale $`R`$, we obtain the relation
$$(4\pi G_N)^1=R^nM^{n+2}.$$
(25)
This equation has a surprising implication. If we fix $`G_N`$ to its observed value and imagine larger values of $`R`$, then the true fundamental quantum gravity scale becomes smaller.
How large could $`R`$ be? In principle, $`R`$ could vary continuously. But there are four natural choices represented in the literature as explicit classes of models. Using a standard American nomenclature, these sizes are:
1. micro-: In this case, all three quantities in (25) are of the order of $`M_{\text{Pl}}`$. This is the original proposal of Scherk and Schwarz for the size of extra dimensions in string theory .
2. mini-: In this case, $`M`$ is taken to be of the order of the grand unification scale, $`2\times 10^{16}`$ GeV. In fact, Hořava and Witten have argued that there is a solution in which all fundamental scales in Nature are of order the grand unification scale, with the scale of the 11th dimension only a small amount larger. This theory includes the unification of Standard Model couplings provided by supersymmetry, and a unification with gravity as well.
3. midi-: In this case, $`R`$ is taken to be at the TeV scale. From (25), for $`n=6`$ for example, the fundamental gravity scale $`M`$ would be 8000 TeV. This choice was first advocated by Antoniadis , who showed how it could lead to superparticle masses of the order of the weak interaction scale. It is possible to arrange a unification of Standard Model couplings at the scale $`M`$ . Recently, Randall and Sundrum have presented a 5-dimensional model with curvature in the extra dimension which gives a novel way to relate the Planck scale and the TeV scale. The phenomenology of this model is quite similar to that of models with flat extra dimensions of TeV-scale size .
4. maxi-: In this case, $`M`$ is taken to be at the TeV scale. Then $`R`$ would be at some scale from millimeters ($`n=2`$) to fermi ($`n=7`$). These large distances, which are huge on the scale of high-energy physics, would seems to violate common sense. But Arkani-Hamed, Dimopoulos, and Dvali have argued that this aggressive choice is not excluded. In this case, the quantum gravity scale $`M`$, the shortest possible distance in Nature, might already be accessible to accelerator experiments.
The maxi- case requires one extra condition. From tests of Bhabha scattering and fermion pair production at LEP and quark-antiquark scattering at the Tevatron, we know that the strong, weak, and electromagnetic interactions follow the force laws predicted for four dimensions up to momentum transfers of about 1 TeV. This means that the quarks, leptons, and gauge bosons must be confined to a 4-dimensional submanifold of thickness less than 1/TeV inside the new large dimensions. This is known to be possible in string theory. Indeed, string theory contains a classical solution called a ‘D-brane’, which can have fermions, bosons, and gauge fields bound to its surface . In the maxi- case, the Standard Model particles would live on a D-brane, whose thickness would be of order $`1/M`$, while gravity and perhaps other light fields could propagate in the full space, out to distances of order $`R`$. A picture of this construction is given in Figure 12.
The micro- case above is famously difficult to test and has given rise to unfortunate statements that string theory is not physics. But I would like to argue now that the other three cases are amenable to experimental test, and that in fact they can be tested at the LHC and the linear collider. Indeed, ten years from now, we could be arguing from experimental data about the true number of space dimensions in Nature.
I will discuss the three cases in turn, from large to small. Consider first the maxi-scale case. One might think that such large extra dimensions are excluded by Cavendish experiments, but actually the best current limit is only $`R<0.8`$ mm ($`M>940`$ GeV) . More significant constraints come from searches for quantum gravity effects at accelerators. Two methods have been proposed. The first is to search for processes in which a collision causes a graviton to be radiated off the brane, carrying with it unobserved momentum . The simplest processes of this kind are
$$e^+e^{}\gamma G,q\overline{q}gG,$$
(26)
where $`G`$ is a graviton, and these can be observed as missing energy processes at $`e^+e^{}`$ and hadron colliders. I will discuss the experimental status of this search in a moment. The second is to search for a contact interaction in fermion-fermion reactions due to graviton exchange . The coefficient of the induced contact interaction is model-dependent, to one cannot use this effect to set strict limits on $`M`$. But if the effect were there, it would be striking, causing the cross sections in $`e^+e^{}f\overline{f}`$ to bend upward as a function of energy, and also modifying the production angular distributions. There is new data from LEP on the pair-production total cross sections, but, unfortunately for this purpose, it is in remarkable agreement with the Standard Model prediction. The new data from DELPHI is shown in Figure 13 .
On the other hand, the cross sections for missing-energy processes can be computed absolutely in terms of the gravity scale defined by (25) and the number of extra dimensions $`n`$, so that bounds on these processes allow us to place lower bounds on $`M`$. In Table 2, taken from , I give the best current limits on $`M`$ (at 95% confidence) from LEP and the Tevatron and the sensitivity expected at LHC and at a 1 TeV linear collider. The LEP results correspond to new limits announced at HEP99 . The first line of the table gives a set of bounds from an astrophysical source, the constraint that supernova 1987A did not radiate away most of its energy in gravitons . This bound is very strong for $`n=2`$ but is unimportant for larger $`n`$. I exclude cosmological bounds that are really constraints on the cosmological scenario.
In Table 2, the sensitivity to missing-energy processes expected at the LHC is quite remarkable. These values cannot be completely trusted, for the usual reason that the LHC cross sections integrate over very large momentum transfer processes. However, it is argued in that the values in the Table are most likely to underestimate the LHC sensitivity. It was the original idea of Arkani-Hamed, Dimopoulos, and Dvali that the size of the weak interaction scale should be set by the scale of $`M`$. The LHC search for missing energy processes should provide a sensitive test of this hypothesis.
I should note that the maxi- case throws away the grand unification scale and all of the physics associated with it. This includes the unification of coupling constants through the renormalization group. It also includes the suppression of neutrino masses and proton decay matrix elements by the factor $`m_W/M_{\text{GUT}}`$, a factor that arises naturally in the standard picture from the fact that these effects are mediated by dimension 5 operators. New suppression mechanisms are needed if we have large extra dimensions. Actually, this may be less a problem than an opportunity to discover new physical mechanisms; see for an example.
We turn next to the midi- case. In this class of models, the Standard Model fields can consistently explore the extra dimensions. Direct quantum gravity effects are inaccessible, but we should expect to see the excitation of states of the photon, $`Z^0`$, and gluon with nonzero momentum in the extra dimensions. These states appear in experiment as massive vector resonances, called ‘Kaluza-Klein recurrences’. If the extra dimensions are flat and have periodicity $`2\pi R`$, the masses of the these states are $`|\stackrel{}{N}|/R`$, where $`\stackrel{}{N}`$ is a vector with integer components. The spectrum of Kaluza-Klein recurrences is a Fourier transform of the shape of the extra dimensions . Figure 14 shows the effect of these recurrences in producing resonances in the dilepton invariant mass distribution that would be observed at the LHC.
Finally, we come to the mini- case. In this class of models, the direct effects of the extra dimensions occur only at the grand unification scale and cannot be observed experimentally. A test of the hypothesis would have to be based on a characteristic set of Lagrangian parameters following from the geometry. These parameters would provide a boundary condition for the renormalization group equations, and we would compare the results of integrating those equation to the weak interaction scale with the results of our experiments.
Here is a concrete picture of how such a comparison could be made. In the original model of Hořava and Witten , the geometry of Nature is effectively five-dimensional and has the form shown in Figure 15. The fifth dimension is bounded, and gauge bosons, fermions, and scalars are bound to four-dimensional walls at the boundary of the space. On one wall, we would have the supersymmetric Standard Model. What is on the other wall? Hořava and Witten proposed that this would be a natural place to put the hidden sector responsible for supersymmetry breaking.
We must now ask, what supersymmetry spectrum follows from this hypothesis? Two answers have been given in the literature. Hořava has argued that one should find the spectrum of gravity mediation in the no-scale limit. However, this result has been criticized by Nilles, Olechowski, and Yamaguchi , who have found large contributions to $`m_0`$ in his picture. Randall and Sundrum have argued that one should find the spectrum of anomaly mediation. However, we have already seen that that spectrum is not self-consistent and requires correction to produce positive slepton masses. Despite these problems with the answers that have been proposed up to now, I believe that the question I have asked has a definite answer, and many theorists are now working to find it. If someone succeeds, the result will give a remarkable and concrete goal for the experimental studies of supersymmetry spectroscopy that we are soon to undertake.
## 7 Lutheran sermon
Tampere has a number of beautiful stone churches, and, in touring them, we learned that sermons play an important role in the Finnish Lutheran traditions. So I will conclude with a sermon.
For me, the most memorable part of the HEP99 meeting was a formal ceremony conducted by four young physicists representing the four LEP collaborations—Fabio Cerutti, Magali Grüwe, Simonetta Gentile, and Mario Pimenta. The title of the ceremony was: ‘Any sign of New Physics in the 1999 LEP data?’. The speakers were thorough, precise, and extremely well-informed. The answer to the question in the title was, no.
It is wrong to be cynical about such an exercise, but it is correct to be disappointed. These speakers stood at the apex of a huge superstructure, representing more than a billion dollars of investment in equipment and training, all focused on the goal of breaking through to the next layer of physics beneath the strong, weak, and electromagnetic interactions. This time, we did not succeed. What moral should we draw from this?
Most of the people in my audience for this lecture were still in grammar school in the 1970’s. This was a very different era in high energy physics, with surprising discoveries and puzzles coming from experiment, forming a cloudy picture in which one struggled to the see the final outcome. I was a graduate student in that period, and the excitement drew me in, away from a perhaps more sensible career in the physics of materials. We do not feel this sort of excitement in high-energy physics today, and many people now ask if it will ever return.
On the other hand, it is important to recognize that the experimental progress we have made in the 1990’s is remarkable in another way. It was often said in the early 1970’s that the experimental picture was necessarily unclear, because we were exploring a realm very remote from human experience. Here I do not refer to the requirement that high energy phenomena need to be observed by complex detectors, but to the conceptual problem of visualizing the basic objects that were used to construct theories—quarks, gluons, heavy bosons, and the like. It was thought that we could view these objects only indirectly, by matching experimental results to abstruse theoretical predictions.
In the 1990’s we learned that this attitude is hogwash. Quarks and leptons may be unimaginably small, but with the right experiments, we can reveal all of the fine details of their behavior. Especially through the program of precision experiments at the $`Z^0`$ resonance, we have been able to examine the quarks or leptons of each individual species, shake their hands, and watch their dances. However remote this microworld is, we understand it pictorially, and with certainty. It is important to add that the means by which we have achieved this understanding is that of using accelerators to go to the basic scale on which these objects act, and then looking and seeing what is manifest there.
In the process, we have learned that physics at this microscopic scale has a basis as rational as chemistry. Quarks and leptons move, not by magic, but because there is a mechanism at work. Experimentally, we have found the moving parts and exhibited their properties. This understanding is very encouraging for the major unsolved aspects of the behavior of the Standard Model, the breaking of electroweak gauge symmetry. This phenomenon must also have a cause, and our experience in the 1990’s tells us that, if we can patiently continue our investigations to higher energy, we can find it out.
The idea that there are mechanisms and reasons for physical phenomena, and that we can find the next one by searching to smaller distances, is an article of faith. As our experimental devices become more complex and expensive, and as the time required to realize them stretches out, it becomes harder and harder to keep the faith. The public wants results on the nightly news. Our allies in government work on the time scale of an election cycle; our colleagues in industry measure progress in ‘Internet time’. It requires continuing effort to persuade them that, though our enterprise moves much more slowly, it is in motion, and toward important goals. But our hardest struggle is with ourselves and our community, to press on to the next great period of discovery which is still over the horizon.
In talking to many people in the experimental community, I sense a pessimism, not about whether there is a next scale of physics, but about what we will find there. Just the Standard Model, they say, just the Higgs boson, just the familiar pattern that some theorist has set out. The excitement of the 1970’s has receded very far into the past, so far that it is difficult to imagine that it will come again.
This is the reason that I have given so much attention in this lecture to the possibility of new space dimensions. Though I have tried to motivate this idea, I think that its major importance comes not so much because it must be true as because it gives an example of how much we could have to learn, and how profoundly different the deep structure of the universe could be from what we now conceive. There could really be unguessed secrets in the laws of Nature. And, these secrets are not hiding in the cosmos or on the large scales of the universe, and not in rare materials or the organization of matter, but only at very small distances.
This is where our accelerators will take us, if we can marshall our resources and our intellectual strength. Above all, we have to keep our belief in our joint enterprise, the belief that Nature has more wonders, beyond our imagination, which wait patiently for our tools to reach them.
I am grateful to Profs. Wolfgang Kummer and Matts Roos and to the members of the HEP99 organizing committee for giving me the opportunity to present this lecture, and to colleagues too numerous to mention, both at SLAC and at HEP99, who have helped me to understand the topics discussed here. I thank particularly Adam Falk, Yuval Grossmann, and Zoltan Ligeti for discussions of CP violation. My work is supported by the US Department of Energy under contract DE–AC03–76SF00515. |
warning/0002/astro-ph0002529.html | ar5iv | text | # Spherical, Oscillatory 𝜶^𝟐-Dynamo Induced by Magnetic Coupling Between a Fluid Shell and an Inner Electrically Conducting Core: Relevance to the Solar Dynamo
## 1 Introduction
$`\alpha ^2`$\- and $`\alpha \omega `$-dynamo models have been studied for decades (Braginsky, 1964; Steenbeck and Krause, 1966; Roberts, 1972; Moffatt, 1978; Gubbins and Roberts, 1987; Baryshnikova and Shukurov, 1987). The $`\alpha ^2`$-dynamo is usually stationary (e.g., Roberts, 1972; Gubbins and Roberts, 1987; Hollerbach, 1996), although oscillatory $`\alpha ^2`$-dynamos have been found to occur in the special circumstance wherein $`\alpha `$ changes rapidly in boundary layers (Radler and Brauer, 1987; Baryshnikova and Shukurov, 1987). In such cases, the period of the $`\alpha ^2`$-dynamo depends strongly on the location of the $`\alpha `$-boundary layer and is typically an order of magnitude or more smaller than the magnetic diffusion time across the dynamo generation region. In general, oscillatory dynamo behavior has been produced by combination of the $`\alpha `$\- and $`\omega `$-effects. Kinematic models of the solar dynamo, which is inherently oscillatory, have been of the $`\alpha \omega `$-type (e.g., Roald and Thomas, 1997).
In this paper, we report that spherical oscillatory $`\alpha ^2`$-dynamos can be simply induced by the magnetic coupling between an electrically conducting outer fluid shell and a conducting inner spherical core even when $`\alpha `$ in the outer shell is a constant. The period of oscillation is of the same order of magnitude as the magnetic diffusion time across the outer shell and depends largely on the electrical conductivity of the inner core. Oscillatory behavior occurs when the outer region of dynamo action surrounds a large, less magnetically diffusive core. The radiative interior of the Sun is a large region with a smaller magnetic diffusivity than the overlying convection zone wherein dynamo action occurs. The oscillatory character of the Sun’s magnetic field, as expressed in the 22 year periodicity of the sunspot cycle, could then be related to the electromagnetic coupling of the region of magnetic field generation in the convection zone with the radiative core of the Sun. The importance of an inner electrically conducting core to the problem of magnetic field generation in an overlying spherical shell has been emphasized in the $`\alpha ^2`$-type models of the geodynamo by Hollerbach and Jones (1993, 1995), Hollerbach (1996), and Gubbins (1999). Nevertheless, the effects of an inner core, with magnetic diffusivity different from that of the overlying fluid convecting shell in which dynamo action takes place, have not been fully elucidated.
The problem of inner core-fluid shell coupling is made difficult by complicated electromagnetic matching conditions at the interface between the regions. For this reason, and to facilitate understanding of the physical effects, we consider the simplest type of $`\alpha ^2`$-dynamo model consisting of a spherical shell with $`\alpha =\text{constant}`$ surrounding a core with $`\alpha =0`$. We derive the appropriate matching conditions for a core and shell of arbitrary magnetic diffusivity. These matching conditions do not appear to have been considered in previous studies of spherical $`\alpha ^2`$-dynamos and they result in oscillatory dynamo solutions.
## 2 Model, Equations and Boundary Conditions
The model consists of a turbulent fluid spherical shell of inner radius $`r_i`$ and outer radius $`r_o`$ with constant (turbulent) magnetic diffusivity $`\lambda _o`$. A magnetic field is generated in the shell by the $`\alpha `$-effect (Steenbeck and Krause, 1966; Roberts, 1972). For $`r>r_o`$, we assume there is a non-conductor; for $`r<r_i`$ we assume that there is a conductor with magnetic diffusivity $`\lambda _i`$. The kinematics of the $`\alpha ^2`$-dynamo in the spherical shell is governed by the non-dimensional linear equations for the magnetic field $`𝐁_o`$
$$\frac{𝐁_o}{t}=R\left(1\eta \right)\times \alpha 𝐁_o+^2𝐁_o$$
(1)
$$𝐁_o=0$$
(2)
In the inner sphere the magnetic field $`𝐁_i`$ is governed by
$$\frac{𝐁_i}{t}=\beta ^2𝐁_i$$
(3)
$$𝐁_i=0$$
(4)
Equations (1)–(4) are scaled by the thickness of the shell $`\left(r_or_i\right)`$ and by the magnetic diffusion timescale $`\left(r_or_i\right)^2/\lambda _o`$. The scaling of the linear system of equations for the magnetic field is arbitrary. The non-dimensional parameters in the above equations, $`\beta `$, $`\eta `$, and the magnetic Reynolds number $`R`$ are defined as
$$\beta =\frac{\lambda _i}{\lambda _o},\eta =\frac{r_i}{r_o},R=\frac{r_o\alpha }{\lambda _o}$$
(5)
Since the main purpose of this paper is to understand the effect of an electrically conducting inner core, we adopt the simplest possible model and take $`\alpha `$ constant in the spherical shell $`r_i<r<r_o`$; $`\alpha `$ is zero outside the shell. With this assumption, spherical harmonics are decoupled and the problem is reduced to a one-dimensional problem with complicated boundary conditions.
At the interface between the shell and the perfectly insulating exterior, i.e., at $`r=r_o`$, the magnetic field must be continuous
$$𝐁_o=𝐁^{(e)}\text{at}r=r_o$$
(6)
where $`𝐁^{(e)}=\varphi `$ is the magnetic field in the insulating exterior $`r>r_o`$, and $`^2\varphi =0`$. At the interface between the shell and the conducting inner sphere, i.e., at $`r=r_i`$, both the magnetic field $`𝐁`$ and tangential components of the electric field $`𝐄`$ must be continuous
$$𝐁_o=𝐁_i,\widehat{𝒓}\times 𝐄_o=\widehat{𝒓}\times 𝐄_i\text{at}r=r_i$$
(7)
where $`\widehat{𝒓}`$ is the unit radial vector, $`𝐄_o`$ is the electric field in the outer shell, and $`𝐄_i`$ is the electric field in the inner core.
Conditions (2) and (4) allow us to express the magnetic fields as a sum of poloidal and toroidal vectors
$`𝑩_o`$ $`=`$ $`\times \times 𝒓h_o+\times 𝒓g_o`$ (8)
$`𝑩_i`$ $`=`$ $`\times \times 𝒓h_i+\times 𝒓g_i`$ (9)
where $`𝒓`$ is the position vector. Use of Equation (8) in boundary condition (6) and expansion of $`h_o`$ and $`g_o`$ in terms of spherical harmonics give
$$g_o=0,\frac{h_o}{r}+\frac{(l+1)h_o}{r}=0\text{at}r=r_o$$
(10)
where $`l`$ is the degree of the spherical harmonic $`Y_l^m`$.
Extra care must be taken for the magnetic boundary conditions at the interface $`r=r_i`$. There are four different cases that we have studied:
(I). The limit $`\beta \mathrm{}`$ for both stationary and oscillatory dynamos. In this case, the boundary condition for the magnetic field is simply
$$g_o=0,\frac{h_o}{r}\frac{lh_o}{r}=0\text{at}r=r_i$$
(11)
(II). The limit $`\beta 0`$ for a stationary dynamo. In this case, boundary conditions (7) require
$$h_o=0,R\left(1\eta \right)r\frac{h_o}{r}\frac{(rg_o)}{r}=0\text{at}r=r_i$$
(12)
(III). The limit $`\beta 0`$ for an oscillatory dynamo. In this case, boundary conditions (7) require
$$\frac{h_o}{r}\frac{lh_o}{r}=0,R\left(1\eta \right)\left(l+1\right)h_o\frac{(rg_o)}{r}=0\text{at}r=r_i$$
(13)
(IV). The general case for $`\beta `$ not tending toward 0 or $`\mathrm{}`$, for both stationary and oscillatory dynamos. In this case, boundary conditions (7) require
$$g_o=g_i,h_o=h_i,\frac{h_o}{r}=\frac{h_i}{r},R\left(1\eta \right)\frac{(rh_o)}{r}\frac{(rg_o)}{r}+\beta \frac{(rg_i)}{r}=0\text{at}r=r_i$$
(14)
The last case is evidently the most complicated one. The solutions presented below show that for $`\beta 10`$, case I provides a good approximation to case IV, while for $`\beta 0.1`$, cases II and III provide a good approximation to case IV. In cases II through IV, $`g`$ and $`h`$ are coupled by boundary conditions (7). The solutions are invariant to a change in the sign of $`R`$.
## 3 Solution Method
In all cases, solutions are expanded in terms of spherical harmonics, implicit in the forms of the boundary and interface conditions given above. The spherical harmonics are decoupled and only the lowest one $`\left(\mathrm{}=1\right)`$ is used in the analysis. The $`\mathrm{}=2`$ mode, not discussed here, behaves similarly to the $`\mathrm{}=1`$ mode.
The time dependence of the solutions is written as $`\mathrm{exp}\left(\sigma _r+i\omega t\right)`$ and onset of dynamo action $`\left(\sigma _r=0\right)`$ is sought. As discussed below, dynamos are either stationary $`\left(\omega =0\right)`$ or oscillatory $`\left(\omega 0\right)`$. The frequency $`\omega `$ is dimensionless with respect to the timescale $`\left(r_or_i\right)^2/\lambda _o`$.
In case I, a perfectly insulating core, and in cases II and III, a perfectly conducting core, it is only necessary to solve for $`g_o`$ and $`h_o`$ subject to the above boundary conditions at $`r_i`$ and $`r_o`$. Exact analytic solutions for $`g_o`$ and $`h_o`$ in these cases can be found in terms of the spherical Bessel functions of the first and second kind. The solutions reduce to finding the eigenvalues of a $`4\times 4`$ matrix. The eigenvalues give critical values of the magnetic Reynolds number $`R`$ as a function of $`\eta =r_i/r_o`$, for which steady or oscillatory dynamos are possible (i.e., $`\sigma _r=0`$).
In case IV, a core of arbitrary $`\beta `$, solutions must be obtained in both the shell and the core subject to the above matching conditions, i.e., $`g_o`$, $`h_o`$, $`g_i`$, and $`h_i`$ must be determined. In principle, analytic solutions are possible, but it is computationally more efficient to seek numerical solutions. We do this by employing a spectral-Tau method which expands solutions in terms of Chebyshev polynomials. The numerical solutions for arbitrary $`\beta `$, and the analytic solutions determined independently in cases I, II, and III, provide a mutual validation of the separate methods. For appropriate values of $`\beta `$, the solutions of the separate methods agree essentially exactly.
## 4 Results
The principal results of this study are summarized in Figure 1
which gives the critical magnetic Reynolds number $`R_{cr}`$ for the onset of dynamo action in the dipole $`(l=1)`$ mode as a function of $`\eta `$, the ratio of the inner radius of the shell to its outer radius. The critical value of $`R`$ increases with increasing $`\eta `$. For an insulating core ($`\beta >>1`$, dashed curve) dynamo solutions are always steady $`\left(\omega 0\right)`$. For a perfectly conducting core ($`\beta <<1`$, solid curve) there are two branches of dynamo solutions depending on $`\eta `$; for $`\eta 0.55`$ the dynamo is steady, but for $`\eta 0.55`$ the dynamo is oscillatory. There is a jump in $`R_{cr}`$ at the transition from steady dynamos to oscillatory dynamos near $`\eta =0.55`$. The dimensionless frequency $`\omega `$ of the oscillatory dynamos, also shown in Figure 1 as a function of $`\eta `$, varies between about 2.5 and 3 for all values of $`\eta `$ considered.
The values of $`R_{cr}`$ for arbitrary $`\beta `$ lie in the narrow space between the solid and dashed curves of Figure 1. Importantly, it is found that $`\beta `$ need not in fact be very small compared with unity for oscillatory dynamos to exist. For example, when $`\eta =0.8`$, the value appropriate to the solar dynamo, oscillatory dynamo solutions are found for $`\beta `$ less than 2 to 3.
## 5 Discussion
The problem solved above is a classically simple one of the type considered by Steenbeck and Krause (1966) and Roberts (1972) decades ago. Yet the effects of an inner electrically conducting core on the $`\alpha ^2`$-dynamo are subtle and not heretofore appreciated. They enter through the complicated electromagnetic matching conditions at the interface between the core and the surrounding shell in which dynamo action occurs. The main effect of the core is to introduce time dependence into the dynamo solutions for cores whose radii are greater than about 0.55 of the outer radius of the shell. An additional requirement for time dependence is that the core be a reasonably good electrical conductor; in terms of the magnetic diffusivity ratio $`\beta =\lambda _i/\lambda _o`$, $`\beta O(1)`$ suffices for oscillatory dynamo behavior.
The importance of all this to the solar dynamo is that the parameters of the solar dynamo satisfy the requirements of oscillatory $`\alpha ^2`$-dynamo solutions. For the solar dynamo $`\eta `$ is about 0.8 and $`\beta `$ is about $`10^3`$ (Moffatt, 1978). In addition, if $`\omega `$ (from Figure 1) is made dimensional using the time scale $`\left(r_or_i\right)^2/\lambda _o`$ with $`r_or_i=1.4\times 10^5\text{km}`$ and $`\lambda _o=O\left(10^2\text{ km}^2\text{ s}^1\right)`$ (eddy magnetic diffusivity), then the period of the oscillatory dynamo solution is comparable to the 22 year period of the sunspot cycle. Thus, $`\alpha ^2`$-dynamo action alone could be responsible for the observed time dependence of the large scale solar magnetic field. It is not our intent to suggest that the $`\omega `$-effect is not significant in dynamo action in general, or in the solar dynamo in particular, because it represents a physically important process. Our purpose is only to clarify some physics and demonstrate the potential importance of a hitherto overlooked effect, that of the oscillatory $`\alpha ^2`$-dynamo. A detailed analysis of the cases $`\alpha =\alpha (r)`$ and $`\alpha \mathrm{cos}\theta (\theta =\text{polar angle})`$ is in progress.
G. S. acknowledges support from NASA’s Planetary Atmospheres Program. K. Z. is supported by PPARC and NATO grants. |
warning/0002/math0002081.html | ar5iv | text | # Entropy in type I algebras
## 1 Introduction
In the theory of non-commutative entropy the attention has almost exclusively been concentrated on non type I algebras. We shall in the present paper remedy this situation by proving the basic facts on entropy of automorphisms of type I C\- and von Neumann-algebras. The results are as nice as one can hope. The CNT-entropy of an automorphism of a von Neumann algebra of type I with respect to an invariant normal state is the classical entropy of the restriction of the automorphism to the center of the algebra. If one factor of a tensor product of two von Neumann algebras is of type I and the other injective, then the entropy of a tensor product automorphism with respect to an invariant product state is the sum of the entropies. The results have obvious corollaries to type I C-algebras. The main idea behind our proofs is the use of conditional expectations of finite index, as employed in \[GN\].
We shall use the notation $`h_\varphi (\alpha )`$ for the CNT-entropy of a C-dynamical system as defined by Connes, Narnhofer and Thirring in \[CNT\], and $`h_\varphi ^{}(\alpha )`$ for the ST-entropy defined by Sauvageot and Thouvenot in \[ST\].
## 2 Main results
We first prove a general result for the Sauvageot-Thouvenot entropy for the restriction of an automorphism to a globally invariant C-subalgebra of finite index. Recall the definition of ST-entropy and its connection with CNT-entropy.
A stationary coupling of a C-dynamical system $`(A,\varphi ,\alpha )`$ with a commutative system $`(C,\mu ,\beta )`$ is an $`\alpha \beta `$-invariant state $`\lambda `$ on $`AC`$ such that $`\lambda |_A=\varphi `$ and $`\lambda |_C=\mu `$. Given such a coupling and a finite-dimensional subalgebra $`P`$ of $`C`$ with atoms $`p_1,\mathrm{},p_n`$, consider the quantity
$$H_\mu (P|P^{})H_\mu (P)+\underset{i=1}{\overset{n}{}}\mu (p_i)S(\varphi ,\varphi _i),$$
where $`\varphi _i(a)=\frac{1}{\mu (p_i)}\lambda (ap_i)`$. By definition, the ST-entropy $`h_\varphi ^{}(\alpha )`$ of the system $`(A,\varphi ,\alpha )`$ is the supremum of these quantities.
By \[ST, Proposition 4.1\], ST-entropy coincides with CNT-entropy for nuclear C-algebras. In fact, the proof of the inequality $`h_\varphi (\alpha )h_\varphi ^{}(\alpha )`$ does not use any assumptions on the algebra. On the other hand, given a coupling $`\lambda `$ and an algebra $`P`$ as above, for each $`m`$ we can form the decomposition
$$\varphi =\underset{i_1,\mathrm{},i_m=1}{\overset{n}{}}\varphi _{i_1\mathrm{}i_m},\varphi _{i_1\mathrm{}i_m}(a)=\lambda (ap_{i_1}\beta (p_{i_2})\mathrm{}\beta ^{m1}(p_{i_m})).$$
If $`\gamma `$ is a unital completely positive mapping of a finite-dimensional C-algebra into $`A`$, we can use these decompositions in computing the mutual entropy $`H_\varphi (\gamma ,\alpha \gamma ,\mathrm{},\alpha ^{m1}\gamma )`$ \[CNT\]. Indeed, since the atoms in $`\beta ^j(P)`$ are $`\beta ^j(p_1),\mathrm{},\beta ^j(p_n)`$ we have by \[CNT, III.3\]
$`H_\varphi (\gamma ,\alpha \gamma ,\mathrm{},\alpha ^{m1}\gamma )S\left(\mu |{\displaystyle \underset{0}{\overset{m1}{}}}\beta ^j(P)\right){\displaystyle \underset{j=0}{\overset{m1}{}}}S\left(\mu |\beta ^j(P)\right)`$
$$+\underset{j}{}\underset{i}{}\mu (\beta ^j(p_i))S(\varphi \alpha ^j\gamma ,\frac{\lambda ((\alpha ^j\gamma )()\beta ^j(p_i))}{\mu (\beta ^j(p_i))}).$$
Hence by invariance of $`\varphi `$, $`\mu `$ and $`\lambda `$ with respect to $`\alpha `$, $`\beta `$ and $`\alpha \beta `$ respectively
$$\frac{1}{m}H_\varphi (\gamma ,\alpha \gamma ,\mathrm{},\alpha ^{m1}\gamma )\frac{1}{m}H_\mu \left(\underset{0}{\overset{m1}{}}\beta ^j(P)\right)H_\mu (P)+\underset{i}{}\mu (p_i)S(\varphi \gamma ,\varphi _i\gamma ).$$
It follows that
$$h_\varphi (\alpha )H_\mu (P|P^{})H_\mu (P)+\underset{i=1}{\overset{n}{}}\mu (p_i)S(\varphi \gamma ,\varphi _i\gamma ).$$
Thus what is really necessary for the coincidence of the entropies, is the existence of a net of unital completely positive mappings $`\gamma _i`$ of finite-dimensional C-algebras into $`A`$ such that $`S(\varphi ,\psi )=lim_iS(\varphi \gamma _i,\psi \gamma _i)`$ for any positive linear functional $`\psi `$ on $`A`$, $`\psi \varphi `$. In particular, $`h_\varphi (\alpha )=h_\varphi ^{}(\alpha )`$ if $`A`$ is an injective von Neumann algebra and $`\varphi `$ is a normal state on it.
###### Proposition 1
Let $`(A,\varphi ,\alpha )`$ be a unital C-dynamical system. Let $`BA`$ be an $`\alpha `$-invariant C-subalgebra (with $`1B`$). Suppose there exists a conditional expectation $`E:AB`$ such that $`E\alpha =\alpha E`$, $`\varphi E=\varphi `$ and $`E(x)cx`$ for all $`xA^+`$ for some $`c>0`$. Then $`h_\varphi ^{}(\alpha )=h_\varphi ^{}(\alpha |_B)`$.
Proof. Let $`(C,\mu ,\beta )`$ be a C-dynamical system with $`C`$ abelian. Using $`E`$ we can lift any stationary coupling on $`BC`$ to a stationary coupling on $`AC`$. This, together with the property of monotonicity of relative entropy, shows that $`h_\varphi ^{}(\alpha )h_\varphi ^{}(\alpha |_B)`$.
Conversely, suppose $`\lambda `$ is a stationary coupling of $`(A,\varphi ,\alpha )`$ with $`(C,\mu ,\beta )`$, $`P`$ a finite-dimensional subalgebra of $`C`$ with atoms $`p_1,\mathrm{},p_n`$, and $`\varphi _i(a)=\frac{1}{\mu (p_i)}\lambda (ap_i)`$ for $`aA`$. Since $`\varphi _i\frac{1}{\mu (p_i)}\varphi `$, $`\varphi _i`$ is normal in the GNS-representation of $`\varphi `$. Since $`E`$ is $`\varphi `$-invariant, it extends to a normal conditional expectation of the closure of $`A`$ in the GNS-representation onto the closure of $`B`$. Thus we can apply \[OP, Theorem 5.15\] to $`\varphi `$ and $`\varphi _i`$, and (as in the proof of Lemma 1.5 in \[GN\]) get
$$\underset{i=1}{\overset{n}{}}\mu (p_i)S(\varphi ,\varphi _i)=\underset{i=1}{\overset{n}{}}\mu (p_i)(S(\varphi |_B,\varphi _i|_B)+S(\varphi _iE,\varphi _i))\underset{i=1}{\overset{n}{}}\mu (p_i)S(\varphi |_B,\varphi _i|_B)\mathrm{log}c.$$
It follows that $`h_\varphi ^{}(\alpha )h_\varphi ^{}(\alpha |_B)\mathrm{log}c`$. Then for each $`m`$
$$h_\varphi ^{}(\alpha )=\frac{1}{m}h_\varphi ^{}(\alpha ^m)\frac{1}{m}h_\varphi ^{}(\alpha ^m|_B)\frac{1}{m}\mathrm{log}c=h_\varphi ^{}(\alpha |_B)\frac{1}{m}\mathrm{log}c.$$
Thus $`h_\varphi ^{}(\alpha )h_\varphi ^{}(\alpha |_B)`$.
###### Corollary 2
If in the above proposition $`A`$ and $`B`$ are injective von Neumann algebras and $`\varphi `$ is normal then $`h_\varphi (\alpha )=h_\varphi (\alpha |_B)`$.
To prove our main result we need also two simple lemmas. The first lemma is more or less well-known.
###### Lemma 3
Let $`(M,\varphi ,\alpha )`$ be a W-dynamical system. Then
(i) if $`p`$ is an $`\alpha `$-invariant projection in $`M`$ such that $`\mathrm{supp}\varphi p`$, then $`h_\varphi (\alpha )=h_\varphi (\alpha |_{M_p})`$;
(ii) if $`\{p_i\}_{iI}`$ is a set of mutually orthogonal $`\alpha `$-invariant central projections in $`M`$, $`_ip_i=1`$, then
$$h_\varphi (\alpha )=\underset{i}{}\varphi (p_i)h_{\varphi _i}(\alpha _i),$$
where $`\varphi _i=\frac{1}{\varphi (p_i)}\varphi `$ is the normalized restriction of $`\varphi `$ to $`Mp_i`$, and $`\alpha _i=\alpha |_{Mp_i}`$.
Proof. (i) easily follows from the definitions; (ii) follows from \[CNT, VII.5(iii)\], (i) and \[SV, Lemma 3.3\] applied to the subalgebras $`M(p_{i_1}+\mathrm{}+p_{i_n})+(1p_{i_1}\mathrm{}p_{i_n})`$.
The proof of the following lemma is left to the reader.
###### Lemma 4
Let $`T`$ be an automorphism of a probability space $`(X,\mu )`$, $`fL^{\mathrm{}}(X,\mu )`$ a $`T`$-invariant function such that $`f0`$ and $`_Xf𝑑\mu =1`$. Let $`\mu _f`$ be the measure on $`X`$ such that $`d\mu _f/d\mu =f`$. Then $`h_{\mu _f}(T)f_{\mathrm{}}h_\mu (T)`$.
###### Theorem 5
Let $`(M,\varphi ,\alpha )`$ be a W-dynamical system with $`M`$ a von Neumann algebra of type I. Let $`Z`$ denote the center of $`M`$. Then $`h_\varphi (\alpha )=h_\varphi (\alpha |_Z)`$.
Proof. By Lemma 3(i) we may suppose that $`\varphi `$ is faithful. Then $`M`$ is a direct sum of homogeneous algebras of type I<sub>n</sub>, $`n\{\mathrm{}\}`$. By Lemma 3(ii) we may assume that $`M`$ is homogeneous of type I<sub>n</sub>. We first assume that $`n`$. Then $`Z=L^{\mathrm{}}(X,\mu )`$, where $`(X,\mu )`$ is a probability space and $`\varphi |_Z=\mu `$. Thus
$$MZ\mathrm{Mat}_n()=L^{\mathrm{}}(X,\mathrm{Mat}_n()),\varphi =_X^{}\varphi _x𝑑\mu (x),$$
where $`\varphi _x=\mathrm{Tr}(Q_x)`$ is a state on $`\mathrm{Mat}_n()`$, $`\mathrm{Tr}`$ the canonical trace on $`\mathrm{Mat}_n()`$. We first assume $`Q_xc>0`$ for all $`x`$.
If $`sM^+`$, $`s`$ is a function in $`L^{\mathrm{}}(X,\mathrm{Mat}_n())`$. Define the $`\varphi `$-preserving conditional expectation $`E:MZ`$ by $`E(s)(x)=\varphi _x(s(x))`$. Then
$$E(s)(x)=\mathrm{Tr}(s(x)Q_x)c\mathrm{Tr}(s(x))cs(x),$$
so $`E(s)cs`$, and it follows from Corollary 2 that $`h_\varphi (\alpha )=h_\varphi (\alpha |_Z)`$.
If there is no $`c>0`$ such that $`Q_xc`$ for all $`x`$, let $`X_c=\{xX|Q_xc\}`$, ($`c>0`$),
$$N_c=L^{\mathrm{}}(X_c,\mathrm{Mat}_n())\text{and}M_c=N_c+\chi _{X\backslash X_c},$$
where $`\chi _{X\backslash X_c}`$ is the characteristic function of $`X\backslash X_c`$. Since $`\varphi `$ is $`\alpha `$-invariant so is $`M_c`$, so by the above argument and Lemma 3, letting $`\varphi _c=\frac{1}{\mu (X_c)}\varphi |_{N_c}`$ and $`\mu _c=\frac{1}{\mu (X_c)}\mu |_{X_c}`$, we obtain
$$h_\varphi (\alpha |_{M_c})=\mu (X_c)h_{\varphi _c}(\alpha |_{N_c})=\mu (X_c)h_{\mu _c}(T|_{X_c})h_\mu (T),$$
where $`T`$ is the automorphism of $`(X,\mu )`$ induced by $`\alpha `$. Letting $`c0`$ and using \[SV, Lemma 3.3\] we obtain the Theorem when $`M`$ is finite.
If $`M`$ is homogeneous of type I, we have $`ML^{\mathrm{}}(X,\mu )B(H)`$, where $`H`$ is a separable Hilbert space. Let $`\mathrm{Tr}`$ denotes the canonical trace on $`B(H)`$. Write again
$$\varphi =_X^{}\varphi _xd\mu (x),\varphi _x=\mathrm{Tr}(Q_x),$$
and let $`E_x(U)`$ denote the spectral projection of $`Q_x`$ corresponding to a Borel set $`U`$. Let $`P_cM=L^{\mathrm{}}(X,B(H))`$ be the projection defined by $`P_c(x)=E_x([c,+\mathrm{}))`$, where $`c>0`$. Then $`P_c`$ is an $`\alpha `$-invariant finite projection. Let
$$M_c=P_cMP_c+(1P_c).$$
Then $`M_c`$ is a finite type I von Neumann algebra. Its center is isomorphic to $`L^{\mathrm{}}(X_c,\mu _c)`$, and the restriction of $`\varphi `$ to it is $`\varphi (P_c)\mu _c\varphi (1P_c)`$, where $`X_c=\{xX|P_c(x)0\}`$ and
$$_{X_c}f(x)𝑑\mu _c(x)=\frac{1}{\varphi (P_c)}_{X_c}f(x)\varphi _x(P_c(x))𝑑\mu (x).$$
So we can apply the first part of the proof to $`M_c`$. Since $`d\mu _c/d\mu \frac{1}{\varphi (P_c)}`$, applying Lemma 4 we get
$$h_\varphi (\alpha |_{M_c})=\varphi (P_c)h_{\mu _c}(T|_{X_c})h_\mu (T).$$
Now letting $`c0`$ we conclude that $`h_\varphi (\alpha )=h_\mu (T)`$.
It should be remarked that in a special case the above theorem was proved in \[GS, Proposition 2.4\].
If $`A`$ is a C-algebra and $`\varphi `$ a state on $`A`$, the central measure $`\mu _\varphi `$ of $`\varphi `$ is the measure on the spectrum $`\widehat{A}`$ of $`A`$ defined by $`\mu _\varphi (F)=\varphi (\chi _F)`$, where $`\varphi `$ is regarded as a normal state on $`A^{\prime \prime }`$, see \[P, 4.7.5\]. Thus by Theorem 5 and \[P, 4.7.6\] we have the following
###### Corollary 6
Let $`(A,\varphi ,\alpha )`$ be a C-dynamical system with $`A`$ a separable unital type I C-algebra. Then $`h_\varphi (\alpha )=h_{\mu _\varphi }(\widehat{\alpha })`$, where $`\widehat{\alpha }`$ is the automorphism of the measure space $`(\widehat{A},\mu _\varphi )`$ induced by $`\alpha `$.
Since inner automorphisms act trivially on the center we have
###### Corollary 7
If $`(M,\varphi ,\alpha )`$ is a W-dynamical system with $`M`$ of type I and $`\alpha `$ an inner automorphism then $`h_\varphi (\alpha )=0`$.
Note that in the finite case the above corollary also follows from a result of N. Brown \[Br, Lemma 2.2\].
The next result was shown in \[S\] when $`\varphi `$ is a trace.
###### Corollary 8
Let $`R`$ denote the hyperfinite II<sub>1</sub>-factor. Let $`A`$ be a Cartan subalgebra of $`R`$ and $`u`$ a unitary operator in $`A`$. If $`\varphi `$ is a normal state such that $`u`$ belongs to the centralizer of $`\varphi `$ then $`h_\varphi (\mathrm{Ad}u)=0`$.
Proof. As in \[S\], it follows from \[CFW\] that there exists an increasing sequence of full matrix algebras $`N_1N_2\mathrm{}`$ with union weakly dense in $`R`$ such that $`AA_nB_n`$, where $`A_n=N_nA`$ and $`B_n=(N_n^{}R)A`$ for all $`n`$. Let $`M_n=N_nB_n`$. Then $`M_n`$ is of type I and contains $`u`$. Hence $`h_\varphi (\mathrm{Ad}u|_{M_n})=0`$. Since $`(_nM_n)^{}=R`$, $`h_\varphi (\mathrm{Ad}u)=0`$ by \[SV, Lemma 3.3\].
If $`(A,\varphi ,\alpha )`$ and $`(B,\psi ,\beta )`$ are C-dynamical systems we always have
$$h_{\varphi \psi }(\alpha \beta )h_\varphi (\alpha )+h_\psi (\beta ),$$
see \[SV, Lemma 3.4\]. The equality does not always hold, see \[NST\] or \[Sa\]. However, we have
###### Theorem 9
Let $`(A,\varphi ,\alpha )`$ and $`(B,\psi ,\beta )`$ be W-dynamical systems. Suppose that $`A`$ is of type I, and $`B`$ is injective. Then
$$h_{\varphi \psi }(\alpha \beta )=h_\varphi (\alpha )+h_\psi (\beta ).$$
Proof. We shall rather prove that $`h_{\varphi \psi }(\alpha \beta )=h_\varphi (\alpha |_{Z(A)})+h_\psi (\beta )`$. For this it suffices to consider the case when $`A`$ is abelian; the general case will follow by the same arguments as in the proof of Theorem 5. (Note that the mapping $`x\mathrm{Tr}(x)x`$ on $`\mathrm{Mat}_n()`$ is not completely positive, but the mapping $`x\mathrm{Tr}(x)\frac{1}{n}x`$ is by the Pimsner-Popa inequality. Thus replacing $`M`$ with $`MB`$ and $`Z`$ with $`ZB`$ in the proof of Theorem 5 we have to replace the inequality $`E(s)cs`$ in the proof with $`E(s)\frac{c}{n}s`$.)
So suppose that $`A`$ is abelian. It is clear that it suffices to prove that if $`A_1,\mathrm{},A_n`$ are finite-dimensional subalgebras of $`A`$, and $`B_1,\mathrm{},B_n`$ are finite-dimensional subalgebras of $`B`$, then
$$H_{\varphi \psi }(A_1B_1,\mathrm{},A_nB_n)=H_\varphi (A_1,\mathrm{},A_n)+H_\psi (B_1,\mathrm{},B_n).$$
We always have the inequality ”$``$”, \[SV, Lemma 3.4\]. To prove the opposite inequality consider a decomposition
$$\varphi \psi =\underset{i_1,\mathrm{},i_n}{}\omega _{i_1\mathrm{}i_n}.$$
Let $`H_{\{\varphi \psi ={\scriptscriptstyle \omega _{i_1\mathrm{}i_n}}\}}(A_1B_1,\mathrm{},A_nB_n)`$ be the entropy of the corresponding abelian model, so
$`H_{\{\varphi \psi ={\scriptscriptstyle \omega _{i_1\mathrm{}i_n}}\}}(A_1B_1,\mathrm{},A_nB_n)=`$
$$=\underset{i_1,\mathrm{},i_n}{}\eta \omega _{i_1\mathrm{}i_n}(1)+\underset{k=1}{\overset{n}{}}\underset{i}{}S(\varphi \psi |_{A_kB_k},\underset{i_k=i}{}\omega _{i_1\mathrm{}i_n}|_{A_kB_k}).$$
Set $`C=_{k=1}^nA_k`$. Let $`p_1,\mathrm{},p_r`$ be those atoms $`p`$ of $`C`$ for which $`\varphi (p)>0`$. Define positive linear functionals $`\psi _{m,i_1\mathrm{}i_n}`$ on $`B`$,
$$\psi _{m,i_1\mathrm{}i_n}(b)=\frac{\omega _{i_1\mathrm{}i_n}(p_mb)}{\varphi (p_m)}.$$
Let also $`\varphi _m`$ be the linear functional on $`C`$ defined by the equality $`\varphi _m(a)=\varphi (ap_m)`$. Then
$$\omega _{i_1\mathrm{}i_n}=\underset{m=1}{\overset{r}{}}\varphi _m\psi _{m,i_1\mathrm{}i_n}\text{on}CB,$$
and
$$\psi =\underset{i_1,\mathrm{},i_n}{}\psi _{m,i_1\mathrm{}i_n}\text{for}m=1,\mathrm{},r.$$
Since the supports of the states $`\varphi _m`$ are mutually orthogonal minimal projections in $`C`$, we have $`{\displaystyle \underset{k=1}{\overset{n}{}}}{\displaystyle \underset{i}{}}S(\varphi \psi |_{A_kB_k},{\displaystyle \underset{i_k=i}{}}\omega _{i_1\mathrm{}i_n}|_{A_kB_k})`$
$``$ $`{\displaystyle \underset{k=1}{\overset{n}{}}}{\displaystyle \underset{i}{}}S(\varphi \psi |_{CB_k},{\displaystyle \underset{i_k=i}{}}\omega _{i_1\mathrm{}i_n}|_{CB_k})`$
$`=`$ $`{\displaystyle \underset{k=1}{\overset{n}{}}}{\displaystyle \underset{i}{}}S(\varphi \psi |_{CB_k},{\displaystyle \underset{m=1}{\overset{r}{}}}\varphi _m\left({\displaystyle \underset{i_k=i}{}}\psi _{m,i_1\mathrm{}i_n}\right)|_{CB_k})`$
$`=`$ $`{\displaystyle \underset{k=1}{\overset{n}{}}}{\displaystyle \underset{i}{}}{\displaystyle \underset{m=1}{\overset{r}{}}}\varphi (p_m)S(\psi |_{B_k},{\displaystyle \underset{i_k=i}{}}\psi _{m,i_1\mathrm{}i_n}|_{B_k}).`$
If $`a_i0`$ and $`\underset{i}{}a_i1`$ then $`\eta (\underset{i}{}a_i)\underset{i}{}\eta (a_i)`$. Hence we have
$`{\displaystyle \underset{i_1,\mathrm{},i_n}{}}\eta \omega _{i_1\mathrm{}i_n}(1)`$ $``$ $`{\displaystyle \underset{m=1}{\overset{r}{}}}{\displaystyle \underset{i_1,\mathrm{},i_n}{}}\eta (\varphi _m\psi _{m,i_1\mathrm{}i_n})(1)`$
$`=`$ $`{\displaystyle \underset{m=1}{\overset{r}{}}}\eta \varphi (p_m){\displaystyle \underset{i_1,\mathrm{},i_n}{}}\psi _{m,i_1\mathrm{}i_n}(1)+{\displaystyle \underset{m=1}{\overset{r}{}}}\varphi (p_m){\displaystyle \underset{i_1,\mathrm{},i_n}{}}\eta \psi _{m,i_1\mathrm{}i_n}(1)`$
$`=`$ $`{\displaystyle \underset{m=1}{\overset{r}{}}}\eta \varphi (p_m)+{\displaystyle \underset{m=1}{\overset{r}{}}}\varphi (p_m){\displaystyle \underset{i_1,\mathrm{},i_n}{}}\eta \psi _{m,i_1\mathrm{}i_n}(1).`$
Thus
$`H_{\{\varphi \psi ={\scriptscriptstyle \omega _{i_1\mathrm{}i_n}}\}}(A_1B_1,\mathrm{},A_nB_n)`$
$`{\displaystyle \underset{m=1}{\overset{r}{}}}\eta \varphi (p_m)+{\displaystyle \underset{m=1}{\overset{r}{}}}\varphi (p_m)H_{\{\psi ={\scriptscriptstyle \psi _{m,i_1\mathrm{}i_n}}\}}(B_1,\mathrm{},B_n).`$
Since $`_m\eta \varphi (p_m)=H_\varphi (C)=H_\varphi (A_1,\mathrm{},A_n)`$, we conclude that
$$H_{\varphi \psi }(A_1B_1,\mathrm{},A_nB_n)H_\varphi (A_1,\mathrm{},A_n)+H_\psi (B_1,\mathrm{},B_n),$$
completing the proof of the Theorem. |
warning/0002/nlin0002023.html | ar5iv | text | # Coupled higher-order nonlinear Schrödinger equations: a new integrable case via the singularity analysis
## 1 Introduction
In this paper, we study the integrability of the following system of two symmetrically coupled higher-order nonlinear Schrödinger equations:
$$\begin{array}{c}q_t=hq_{xxx}+aq\overline{q}q_x+bq^2\overline{q}_x+cr\overline{r}q_x+dq\overline{r}r_x+eqr\overline{r}_x+\hfill \\ \mathrm{i}(sq_{xx}+fq^2\overline{q}+gqr\overline{r}),\\ r_t=hr_{xxx}+ar\overline{r}r_x+br^2\overline{r}_x+cq\overline{q}r_x+dr\overline{q}q_x+erq\overline{q}_x+\hfill \\ \mathrm{i}(sr_{xx}+fr^2\overline{r}+grq\overline{q}),\end{array}$$
(1)
where $`h,a,b,c,d,e,s,f,g`$ are real parameters, $`h0`$, and the bar denotes the complex conjugation.
Section 2 is devoted to the singularity analysis of (1), and there, under certain simplifying assumptions, we find the following new case when the system (1) passes the Painlevé test for integrability well:
$$a0,b=0,c=d=e=a,f=g=\frac{as}{3h}.$$
(2)
In Section 3, we prove the integrability of (1) with (2), constructing a corresponding $`4\times 4`$ Lax pair, and then obtain a multi-field generalization of the new integrable system.
Section 4 contains some concluding remarks.
## 2 Singularity analysis
Let us apply the Weiss-Kruskal algorithm of the singularity analysis , to the system (1) (we set $`h=1`$ w.l.g.). With respect to $`q,\overline{q},r,\overline{r}`$, which should be considered as mutually independent during the Painlevé test, the system (1) is a normal system of four third-order equations, of total order twelve. A hypersurface $`\varphi (x,t)=0`$ is non-characteristic for (1) if $`\varphi _x0`$, and we set $`\varphi _x=1`$. Then we substitute the expansions
$$\begin{array}{c}q=q_0(t)\varphi ^\alpha +\mathrm{}+q_n(t)\varphi ^{n+\alpha }+\mathrm{},\hfill \\ \overline{q}=\overline{q}_0(t)\varphi ^\beta +\mathrm{}+\overline{q}_n(t)\varphi ^{n+\beta }+\mathrm{},\hfill \\ r=r_0(t)\varphi ^\gamma +\mathrm{}+r_n(t)\varphi ^{n+\gamma }+\mathrm{},\hfill \\ \overline{r}=\overline{r}_0(t)\varphi ^\delta +\mathrm{}+\overline{r}_n(t)\varphi ^{n+\delta }+\mathrm{}\hfill \end{array}$$
(3)
(the bar does not mean the complex conjugation now) into (1) and obtain four algebraic equations for $`\alpha ,\beta ,\gamma ,\delta ,q_0\overline{q}_0,r_0\overline{r}_0`$, which determine the dominant behavior of solutions near $`\varphi =0`$, as well as one twelfth-degree algebraic equation with respect to $`n`$, which determines the positions of resonances in the expansions. The perfect analysis of those five equations is very complicated, and we will publish it later on. But now we impose certain restrictions on the expansions (3) in order to reach a new integrable case of (1) by a short way.
For this purpose, we set $`\alpha =\beta =\gamma =\delta =1`$ and require that exactly two of twelve resonances lie in the position $`n=0`$. Under these simplifying assumptions, we find from (1) and (3) that $`q_0\overline{q}_0=r_0\overline{r}_0=\mathrm{constant}0`$ (we set $`\mathrm{constant}=1`$ w.l.g.), and that
$$a=6bcde,$$
(4)
$$\begin{array}{c}(n+1)n^2(n3)(n4)\times \\ (n^26n2b2d+5)(n^26n2b2e+5)\times \\ (n^36n^2+(52d2e)n+4(c+d+e+3))=0.\end{array}$$
(5)
Due to (5), five resonances lie in the positions $`n=1,0,0,3,4`$. Denoting the positions of other seven resonances as $`n_1,n_2,\mathrm{},n_7`$, we find from (5) that
$$\begin{array}{c}n_2=6n_1,n_4=6n_3,n_7=6n_5n_6,\\ d=\frac{1}{2}(52b6n_1+n_1^2),e=\frac{1}{2}(52b6n_3+n_3^2),\\ b=\frac{1}{4}(56n_1+n_1^26n_3+n_3^2+6n_5n_5^2+6n_6n_6^2n_5n_6),\\ c=\frac{1}{4}(22+12n_52n_5^2+12n_62n_6^28n_5n_6+n_5^2n_6+n_5n_6^2).\end{array}$$
(6)
We require that the considered branch is generic, i.e. eleven resonances lie in nonnegative positions. Taking into account the admissible multiplicity of resonances, we have to study 23 distinct cases listed in Table 1.
At the next step of the analysis, we find from (1) and (3) the recursion relations for $`q_n,\overline{q}_n,r_n,\overline{r}_n`$, $`n=0,1,2,\mathrm{}`$, and then check the consistency of those relations at the resonances, using the *Mathematica* system for computations. The results are listed in Table 1, in its column $`n_{log}`$, where $`n_{log}`$ denotes the position in which some logarithmic terms should be introduced into the expansions (3) due to the following reasons: either the actual number of arbitrary functions is less than the multiplicity of the resonance, or the compatibility conditions at the resonance cannot be satisfied identically. As we see, all the cases but one, #14, have already failed to pass the Painlevé test. In the case #14, where
$$h=1,a=c=d=e=\frac{3}{2},b=0$$
(7)
due to our assumption and formulae (4) and (6), we have to set
$$f=g=\frac{1}{2}s$$
(8)
for the compatibility conditions at $`n=2,3`$ to become identities.
The system (1) with (7) and (8) admits many branches, i.e. kinds of expansions (3). We have already studied the generic branch. Other branches either are Taylor expansions governed by the Cauchy-Kovalevskaya theorem, or are related to the following two:
1. $`\alpha =1`$, $`\beta =1`$, $`\gamma =2`$, $`\delta =2`$, $`q_0\overline{q}_0=4`$, $`r_0,\overline{r}_0`$, positions of resonances are $`n=4,1,0,0,0,1,1,3,4,4,5,5`$;
2. $`\alpha =2`$, $`\beta =0`$, $`\gamma =3`$, $`\delta =2`$, $`q_0\overline{q}_0=8`$, $`r_0,\overline{r}_0`$, positions of resonances are $`n=5,2,1,0,0,0,2,4,5,5,6,7`$.
Compatibility conditions at all resonances of these branches turn out to be identities. The Painlevé test is completed.
Since we consider $`q`$ and $`\overline{q}`$, $`r`$ and $`\overline{r}`$ as mutually independent, evident scale transformations of $`q,\overline{q},r,\overline{r}`$ and $`t`$ relate the conditions (7) and (8) with the more general condition (2). On the other hand, the transformation
$$\begin{array}{c}q^{}=\frac{1}{2}q\mathrm{exp}\omega ,\overline{q}^{}=\frac{1}{2}\overline{q}\mathrm{exp}(\omega ),\\ r^{}=\frac{1}{2}r\mathrm{exp}\omega ,\overline{r}^{}=\frac{1}{2}\overline{r}\mathrm{exp}(\omega ),\\ \omega =\frac{\mathrm{i}s}{3}x+\frac{2\mathrm{i}s^3}{27}t,x^{}=x+\frac{s^2}{3}t,t^{}=t\end{array}$$
(9)
changes the system (1) with (7) and (8) into the system of four coupled mKdV equations
$$\begin{array}{c}q_t+q_{xxx}+6q\overline{q}q_x+6(qr\overline{r})_x=0,\hfill \\ \overline{q}_t+\overline{q}_{xxx}+6q\overline{q}\overline{q}_x+6(\overline{q}r\overline{r})_x=0,\hfill \\ r_t+r_{xxx}+6r\overline{r}r_x+6(q\overline{q}r)_x=0,\hfill \\ \overline{r}_t+\overline{r}_{xxx}+6r\overline{r}\overline{r}_x+6(q\overline{q}\overline{r})_x=0,\hfill \end{array}$$
(10)
where the prime of $`x,t,q,\overline{q},r,\overline{r}`$ is omitted, and the bar does not mean (but may mean) the complex conjugation. This form (10) is useful for obtaining a Lax pair for the new integrable case (2) of (1).
## 3 Lax pair and generalization
Let us consider the linear problem
$$\mathrm{\Psi }_x=U\mathrm{\Psi },\mathrm{\Psi }_t=V\mathrm{\Psi }$$
(11)
with the matrices $`U`$ and $`V`$ given in the following block form :
$$U=\mathrm{i}\zeta \left(\begin{array}{cc}I_1& 0\\ 0& I_2\end{array}\right)+\left(\begin{array}{cc}0& Q\\ R& 0\end{array}\right),$$
(12)
$$\begin{array}{c}V=\mathrm{i}\zeta ^3\left(\begin{array}{cc}4I_1& 0\\ 0& 4I_2\end{array}\right)+\zeta ^2\left(\begin{array}{cc}0& 4Q\\ 4R& 0\end{array}\right)+\\ \mathrm{i}\zeta \left(\begin{array}{cc}2QR& 2Q_x\\ 2R_x& 2RQ\end{array}\right)+\left(\begin{array}{cc}Q_xRQR_x& Q_{xx}+2QRQ\\ R_{xx}+2RQR& R_xQRQ_x\end{array}\right),\end{array}$$
(13)
where $`I_1`$ and $`I_2`$ are unit matrices, $`\zeta `$ is a parameter. The compatibility condition of the linear problem (11),
$$U_tV_x+UVVU=0,$$
(14)
becomes the system of two matrix mKdV equations :
$$\begin{array}{c}Q_t+Q_{xxx}3Q_xRQ3QRQ_x=0,\hfill \\ R_t+R_{xxx}3R_xQR3RQR_x=0.\hfill \end{array}$$
(15)
If we substitute
$$Q=\left(\begin{array}{cc}q& r\\ \overline{r}& \overline{q}\end{array}\right),R=\left(\begin{array}{cc}\overline{q}& r\\ \overline{r}& q\end{array}\right)$$
(16)
into (15), we obtain exactly the new system (10). This proves that the new case (2) of the coupled higher-order nonlinear Schrödinger equations (1) possesses a parametric Lax pair.
A multi-field generalization of the system (10) is obtained by choosing
$$\begin{array}{c}Q=\left(\begin{array}{cc}u_0II+_{k=1}^{2m1}u_ke_kI& v_0II+_{k=1}^{2m1}v_kIe_k\\ v_0II_{k=1}^{2m1}v_kIe_k& u_0II_{k=1}^{2m1}u_ke_kI\end{array}\right),\hfill \\ R=\left(\begin{array}{cc}u_0II_{k=1}^{2m1}u_ke_kI& v_0II+_{k=1}^{2m1}v_kIe_k\\ v_0II_{k=1}^{2m1}v_kIe_k& u_0II+_{k=1}^{2m1}u_ke_kI\end{array}\right),\hfill \end{array}$$
(17)
where $`I`$ is the $`2^{m1}\times 2^{m1}`$ unit matrix, and $`\{e_1,\mathrm{},e_{2m1}\}`$ are $`2^{m1}\times 2^{m1}`$ anti-commutative and anti-Hermitian matrices:
$$\{e_i,e_j\}_+=2\delta _{ij}I,e_k^{}=e_k.$$
(18)
Then the compatibility condition (14) becomes the system
$$\begin{array}{c}u_{j,t}+u_{j,xxx}+6_{k=0}^{2m1}u_k^2u_{j,x}+6\left(_{k=0}^{2m1}v_k^2u_j\right)_x=0,\hfill \\ v_{j,t}+v_{j,xxx}+6_{k=0}^{2m1}v_k^2v_{j,x}+6\left(_{k=0}^{2m1}u_k^2v_j\right)_x=0,\hfill \\ \hfill j=0,1,\mathrm{},2m1.\end{array}$$
(19)
If we assume that $`u_k`$ and $`v_k`$ are real and set
$$\begin{array}{cc}\hfill \begin{array}{c}u_{2j2}+\mathrm{i}u_{2j1}=q_j,\\ v_{2j2}+\mathrm{i}v_{2j1}=r_j,\end{array}& j=1,2,\mathrm{},m,\end{array}$$
(20)
the system (19) is expressed as
$$\begin{array}{c}q_{j,t}+q_{j,xxx}+6_{k=1}^m\left|q_k\right|^2q_{j,x}+6\left(_{k=1}^m\left|r_k\right|^2q_j\right)_x=0,\hfill \\ r_{j,t}+r_{j,xxx}+6_{k=1}^m\left|r_k\right|^2r_{j,x}+6\left(_{k=1}^m\left|q_k\right|^2r_j\right)_x=0,\hfill \\ \hfill j=1,2,\mathrm{},m.\end{array}$$
(21)
## 4 Conclusion
In the literature, the following three integrable cases of coupled higher-order nonlinear Schrödinger equations (1) are known and studied (sometimes in a form of coupled mKdV equations):
$$a0,b=e=0,c=d=\frac{1}{2}a,f=g=\frac{as}{3h},$$
(22)
$$a0,b=d=e=0,c=a,f=g=\frac{as}{3h},$$
(23)
$$a0,b=d=e=\frac{1}{3}a,c=\frac{2}{3}a,f=g=\frac{2as}{9h};$$
(24)
they were introduced in , , , respectively.
The new integrable case (2) of the system (1), obtained in this paper by means of the singularity analysis, admits the reduction $`r0`$ to the Hirota equation and the reduction $`rq`$ to the Sasa-Satsuma equation . Its soliton solutions and conservation laws deserve further investigation.
Acknowledgments. The work of S. Yu. S. was supported in part by the Fundamental Research Fund of Belarus, grant $`\mathrm{\Phi }`$98-044. The work of T. T. was supported by a JSPS Research Fellowship for Young Scientists. |
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