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warning/0002/nlin0002043.html | ar5iv | text | # Dynamics of Kinks in One- and Two- Dimensional Hyperbolic Models with Quasi-Discrete Nonlinearities
## 1 Introduction
In the last few years, partial differential equations with discrete nonlinearities have been used to model phenomena in fields ranging from physics to biology, including the study of pinning of dislocation motions in crystals, breathers in nonlinear crystal lattices, Josephson junction arrays and the biophysical description of calcium release waves \[kn:flakla1, kn:flawil1, kn:flomaz1, kn:fuklee1, kn:leeric1, kn:cop1, kn:bishor1, kn:mitkla1, kn:mitkla2, kn:mit1, kn:klamit1, kn:ponkei1, kn:keismi1, kn:peapon1, kn:gru1, kn:gru2, kn:kee1, kn:kee2\]. Recently, the discrete one-dimensional stationary version of the Klein-Gordon equation
$$\varphi _{tt}+\gamma \varphi _t=D\varphi _{xx}+\alpha \underset{k}{}\delta (xx_k)[f(\varphi )+h],$$
(1)
with $`\gamma =0`$ has been analyzed by Flach and Kladko \[kn:flakla1\] (see also references therein). In (1) $`\varphi `$ is an order parameter, the non-negative constant $`\gamma `$ is the dissipation coefficient and the positive constants $`D`$ and $`\alpha `$ are the diffusion coefficient and the amplitude of the discrete nonlinearity, respectively. The function $`f`$ is a bistable function (the derivative of a double well potential having the two equal minima); i.e., a real odd function with positive maximum equal to $`\varphi ^{}`$, negative minimum equal to $`\varphi ^{}`$ and precisely three zeros in the closed interval $`[a_{},a_+]`$ located at $`a_{}`$, $`a_0`$ and $`a_+`$. For simplicity and without lost of generality we will consider in our analysis $`a_{}=1`$, $`a_0=0`$ and $`a_+=1`$. The prototype example is $`f(\varphi )=(\varphi \varphi ^3)/2`$. The constant $`h`$, assumed to be small in absolute value, specifies the difference of the potential minima of the system; i.e., $`f(\varphi )+h`$ is the derivative of a double well potential with one local minimum and one global minimum. Note that (1) reduces to the Klein-Gordon equation when $`_k\delta (xx_k)`$ is replaced by a constant with appropriate rescaling. In \[kn:flakla1\] a first order perturbation calculation for the heteroclinic orbits of the corresponding stationary kink solution for (1) was presented. Kink solutions, connecting the two local minima of the double well potential, were also obtained for the sine-Gordon case, $`f(\varphi )=\mathrm{sin}(\varphi )`$, and the Klein-Gordon case, $`f(\varphi )=(\varphi \varphi ^3)/2`$. Both are particular cases of the function $`f(\varphi )`$ as defined above. Note that the sine-Gordon case is equivalent to the derivative of a double well potential in a restricted domain of definition.
In this manuscript, we study the dynamics of kinks for a quasi-discrete version of the Klein-Gordon equation
$$\varphi _{tt}+\gamma \varphi _t=D\mathrm{\Delta }\varphi +\alpha \beta (x,y)[f(\varphi )+h],$$
(2)
in a bounded region $`\mathrm{\Omega }R^n`$, $`n=1,2`$ with smooth boundary $`\mathrm{\Omega }`$ for Neumann boundary conditions on $`\mathrm{\Omega }`$. Equation (2) reduces to (1) when $`\beta `$ is one-dimensional and $`\beta (x)=_k\delta (xx_k)`$. Altough the analysis presented below will be valid for a general class of positive differentiable functions $`\beta `$, we have in mind some particular cases which approximate a distribution of discrete nonlinearities for large $`\eta `$, a positive constant defined below.
Case 1) There is a sequence of points on the real line, $`x_k`$, $`k=1,\mathrm{},N`$, with $`N`$ finite or infinite, where the function $`\beta `$ reaches a maximum,
$$\beta (x)=\underset{k=1}{\overset{N}{}}e^{\eta (xx_k)^2}.$$
(3)
Case 2) There is a sequence of lines in the plane, $`y_k`$, $`k=1,\mathrm{},N`$, with $`N`$ finite or infinite, where the function $`\beta `$, independent of $`x`$, reaches a maximum,
$$\beta (x,y)=\underset{k=1}{\overset{N}{}}e^{\eta (yy_k)^2}.$$
(4)
Case 3) There is a sequence of points in the plane, $`(x_k,y_j)`$, $`k=1,\mathrm{},N`$, $`j=1,\mathrm{},M`$ with $`N`$ and $`M`$ finite or infinite, where the function $`\beta `$ reaches a maximum,
$$\beta (x,y)=\underset{k=1}{\overset{N}{}}\underset{j=1}{\overset{M}{}}\sigma (xx_k,yy_j;\eta ),\text{w}here\sigma (x,y;\eta )=e^{\eta (x^2+y^2)}.$$
(5)
Case 4) There is a sequence of circles in the plane, $`\rho =\rho _k`$, $`k=1,\mathrm{},N`$, with $`N`$ finite or infinite, and where $`\rho `$ represents the radial polar coordinate, where the function $`\beta `$ reaches a maximum,
$$\beta (\rho )=\underset{k=1}{\overset{N}{}}e^{\eta (\rho \rho _k)^2}.$$
(6)
We refer to the points $`x_k`$ and $`(x_k,y_j)`$,$`k=1,\mathrm{},N`$, $`j=1,\mathrm{},M`$ as quasi-discrete (QD) sites and to the stripes $`y=y_k`$ and circles $`\rho =\rho _k`$ $`k=1,\mathrm{},N`$ as quasi-semi-discrete (QS) sites. We define $`d`$ to be the minimum distance between two adjacent QD or QS sites. Note that the function $`\beta `$ can be chosen to depend not only on the spatial variable but also on $`t`$. The specific form of $`\beta (x,y,t)`$ will depend on the particular model. One might, for example, have the product of a spatially dependent function $`\beta (x,y)`$ with a probabilistic time-dependent function.
For (2) we define the following dimensionless variables and parameters
$$\widehat{x}=\frac{x}{d},\widehat{y}=\frac{y}{d},\widehat{t}=\frac{\sqrt{D}t}{d}$$
(7)
and
$$ฯต=\sqrt{\frac{D}{\alpha }}\frac{1}{d}\widehat{\gamma }=\frac{\gamma d}{\sqrt{D}}\widehat{\eta }=\eta d^2,\widehat{h}=\frac{h}{ฯต}.$$
(8)
Substituting (7) and (8) into (2) and dropping the $`\widehat{}`$ from the variables and parameters
$$ฯต^2\varphi _{tt}+ฯต^2\gamma \varphi _t=ฯต^2\mathrm{\Delta }\varphi +\beta (x,y)[f(\varphi )+ฯตh].$$
(9)
We will consider the case $`0<ฯต1`$; i.e., when diffusion is slow, $`d`$ is large or $`\alpha `$ is large, and there is a small dissipation.
The homogeneous version of (9),
$$ฯต^2\varphi _{tt}+ฯต^2\gamma \varphi _t=ฯต^2\mathrm{\Delta }\varphi +f(\varphi )+h,$$
(10)
possesses a travelling kink solution. The point on the line (for $`n=1`$) or the set of points in the plane (for $`n=2`$) for which the order parameter $`\varphi `$ vanishes are called the interface or front of the system. For (10) the front moves according to an extended version of the Born-Infeld equation \[kn:neu1, kn:rotnep1\]
$$(1s_t^2)s_{xx}+2s_xs_ts_{xt}(1+s_x^2)s_{tt}\gamma s_t(1+s_x^2s_t^2)\widehat{h}(1+s_x^2s_t^2)^{\frac{3}{2}}=0,$$
(11)
where $`\widehat{h}`$, proportional to $`h`$, will be defined later. Planar fronts moving according to (11) with $`\gamma =\widehat{h}=0`$ (no dissipation and both phases with equal potential) move with a constant velocity equal to their initial velocity. For other values of $`\gamma `$ or $`\widehat{h}`$, fronts move with a velocity that asymptotically approaches $`\widehat{h}/(\gamma ^2+\widehat{h}^2)^{\frac{1}{2}}`$ as long as the initial velocity is bounded from above by $`1`$ in absolute value. Linear perturbations to these planar fronts decay, either in a monotonic or an oscillatory way, to zero as $`t\mathrm{}`$ \[kn:rotnep1\]. Circular interfaces moving according to (11) with $`h>0`$ shrink to a point in finite time \[kn:neu1, kn:rotnep1\]. If $`h<0`$, then circles shrink to points for some initial conditions and for others they grow unboundedly. Neu \[kn:neu1\] showed that for $`\gamma =h=0`$, closed kinks can be stabilized against collapse by the appearance of short wavelength, small amplitude waves. For the more general case, two situations are possible. Either linear perturbations to a circle decay and curves shrink to a point in finite time or they are still present at the shrinkage point of the circle. Note that Equation (11) expressed in terms of its kinematic and geometric properties reads \[kn:rotnep1\]
$$\frac{dv}{dt}+\gamma v(1v^2)\kappa (1v^2)+\widehat{h}(1v^2)^{\frac{3}{2}}=0,$$
(12)
where $`\kappa `$ is the curvature of the front and $`dv/dt`$ is the โLagrangianโ time derivative of $`v`$ which is calculated along the trajectory fo the interfacial point moving with the normal velocity $`v`$ \[kn:rotnep1\].
One of the goals of this papers is to determine whether the dynamic behavior of kinks in (9) differs from its homogeneous (discrete) nonlinearity counterpart (10). For the overdamped version of equation (2), which is a parabolic bistable equation, it has been shown that there are essential differences between the homogeneous and non-homogeneous (discrete) cases, in that latter exhibits propagation failure \[kn:kee1, kn:kee2, kn:mit1, kn:mitkla1, kn:roteps1\].
In Section 2 we present an equation of motion for the front in equation (9), and we describe briefly the method by which it was derived. This equation generalizes equation (11) with the strong nonlinearity accounting for the influence of the function $`\beta `$ on the front motion. In Section 3 we study the evolution of one-dimensional fronts. We show that for $`h=0`$ the function $`\beta `$ acts as a โpotential functionโ for the motion of the front; i.e., a front initially placed between two maxima of $`\beta `$ asymptotically approaches the intervening minimum. When $`h0`$, fronts that start between two maxima of $`\beta `$ asymptotically approach an equilibrium point determined by $`h`$ and $`\beta `$, producing a kink propagation failure. In Section 4 we study the evolution of two-dimensional fronts with radial symmetry. We find that when there is no dissipation circles can shrink to a point in finite time, grow unboundedly or their radius can oscillate, depending on the initial conditions. When dissipation effects are present, the oscillations decay spirally or not depending of the value of $`\gamma `$. The final result is the stabilization of the circular domain of one phase inside the other phase. Our conclusions appear in Section 5.
## 2 Front Dynamics: The Equation of Motion
For (9) and the law of motion of the interface in two dimensions is given by
$$(1s_t^2)s_{xx}+2s_xs_ts_{xt}(1+s_x^2)s_{tt}\gamma s_t(1+s_x^2s_t^2)$$
$$\frac{\beta _y(x,s)\beta _x(x,s)s_x+\beta _t(x,s)s_t}{2\beta (x,s)}(1+s_x^2s_t^2)\widehat{h}\beta ^{\frac{1}{2}}(x,s)(1+s_x^2s_t^2)^{\frac{3}{2}}=0,$$
(13)
where $`y=s(x,t)`$ is the Cartesian description of the interface and $`\widehat{h}`$ is proportional to $`h`$ as will be explained later. Equation (13) was obtained by carrying out a non-rigorous but self-consistent singular perturbation analysis for $`ฯต1`$, treating the interface as a moving internal layer of width $`O(ฯต)`$. We focused on the dynamics of the fully developed layer, and not on the process by which it was generated. The method that we applied is similar to that used in \[kn:rotnep1\] and \[kn:roteps1\] for the study of the evolution of kinks in both the nonlinear wave equation (10) and the Allen-Cahn equation with quasi-discrete sources of reaction (the overdamped version of (2)). The basic assumptions made were:
\- For small $`ฯต0`$ and all $`t[0,T]`$, the domain $`\mathrm{\Omega }`$ can be divided into two open regions $`\mathrm{\Omega }_+(t;ฯต)`$ and $`\mathrm{\Omega }_{}(t,ฯต)`$ by a curve $`\mathrm{\Gamma }(t;ฯต)`$, which does not intersect $`\mathrm{\Omega }`$. This interface, defined by $`\mathrm{\Gamma }(t;ฯต):=\{x\mathrm{\Omega }:\varphi (x,t;ฯต)=0\}`$, is assumed to be smooth, which implies that its curvature and its velocity are bounded independently of $`ฯต`$.
\- There exists a solution $`\varphi (x,t;ฯต)`$ of (2), defined for small $`ฯต`$, for all $`x\mathrm{\Omega }`$ and for all $`t[0,T]`$ with an internal layer. As $`ฯต0`$ this solution is assumed to vary continuously through the interface, taking the value $`1`$ when $`x\mathrm{\Omega }_+(t;ฯต)`$, $`1`$ when $`x\mathrm{\Omega }_{}(t,ฯต)`$, and varying rapidly but smoothly through the interface.
\- The curvature of of the front is small compared to its width.
As a first stage in the derivation of equation (13) we define near the interface a new variable $`z=(ys)/ฯต`$ which is $`๐ช(1)`$ as $`ฯต\mathrm{}`$ and then express equation (9) in terms of this new variable. Next we expand $`\varphi `$ and $`\beta `$ asymptotically in a power series in $`ฯต`$ and substitute these expansions into the differential equation. After equating the coefficients of corresponding powers of $`ฯต`$, we obtain two equations. The first can be reduced to an equation of the type $`\mathrm{\Phi }_{zz}^0+f(\mathrm{\Phi }^0)=0`$ which has to satisfy $`\mathrm{\Phi }^0(0)=0`$ and $`\mathrm{\Phi }^0(\pm 1)=\pm 1`$, giving a kink solution. Here $`\mathrm{\Phi }^0`$ represents the leading order term of the order parameter $`\varphi `$ in terms of $`z`$. The second problem is a linear non-homogeneous second order ODE. Equation (13) is obtain by applying the solvability condition (Fredholm alternative) after defining $`\widehat{h}:=h[\mathrm{\Phi }^0(+\mathrm{})\mathrm{\Phi }^0(\mathrm{})]/_{\mathrm{}}^{\mathrm{}}(\mathrm{\Phi }_z^0)^2๐z`$. Note that for $`f(\varphi )=(\varphi \varphi ^3)/2`$ (Ginsburg-Landau theory), $`\mathrm{\Phi }^0(z)=\mathrm{tanh}\frac{z}{2}`$ and $`\widehat{h}=3h`$ whereas for $`f(\varphi )=\mathrm{sin}\varphi `$ (sine-Gordon), $`\mathrm{\Phi }^0(z)=4\text{tan}^1e^z\pi `$ and $`\widehat{h}=\frac{\pi }{4}h`$.
## 3 Front Motion in 1D
For a one-dimensional system, equation (13) reads
$$s_{tt}+\gamma s_t(1s_t^2)+\frac{\beta ^{}(s)}{2\beta (s)}(1s_t^2)+\widehat{h}\beta ^{\frac{1}{2}}(s)(1s_t^2)^{\frac{3}{2}}=0.$$
(14)
We concentrate on functions $`\beta `$ of the form (3), altough the same analysis can be done for a general differentiable function. We define $`u=s`$ and $`v=s_t`$ obtaining
$$\{\begin{array}{c}u_t=v,\hfill \\ v_t=\gamma v(1v^2)\frac{\beta ^{}(u)}{2\beta (u)}(1v^2)\widehat{h}\beta ^{\frac{1}{2}}(u)(1v^2)^{\frac{3}{2}}.\hfill \end{array}$$
(15)
The fixed points of (15) are $`(u_0,0)`$, where $`u_0`$ satisfies $`g(u)=\beta ^{}(u)+2\widehat{h}\beta ^{\frac{3}{2}}(u)=0`$. The trace, $`\tau `$, and determinant, $`\mathrm{\Delta }`$, of matrix of the linearized system are $`\tau =\gamma `$ and $`\mathrm{\Delta }=[\beta ^{\prime \prime }(u_0)\beta (u_0)\beta ^2(u_0)]/2\beta ^2(u_0)+\widehat{h}\beta ^{}(u_0)/[2\beta ^{\frac{1}{2}}(u_0)]`$, respectively. If $`\widehat{h}=0`$ then the fixed points are the maxima (unstable) and minima (stable) of $`\beta (u)`$. Thus, a front initially between two maxima of $`\beta `$ will move and asymptotically approach the intervening minimum. When there is no dissipation, this behavior is in contrast with the homogeneous case (11) where, as was pointed out in the introduction, fronts move with a constant velocity equal to their initial velocity. In the non-homogeneous case, we can predict the final position of the front from the structure of $`\beta `$. In order to understand the behavior of $`g(u)`$ as $`\widehat{h}`$ increases above zero we consider a function $`\beta (u)`$ with a single peak at $`0`$; i.e., $`\beta (u)=e^{\eta u^2}`$. This function will approximate the more general (3) if $`\eta 1`$, so that the influence of peaks on one another is very small. In this case $`g(u)=2e^{\eta u^2}[\eta u\widehat{h}e^{\frac{\eta u^2}{2}}]`$. For $`\widehat{h}=0`$, $`g(u)`$ vanishes at $`u=0`$ and it is positive for $`u>0`$ and negative for $`u<0`$. As $`\widehat{h}`$ moves away from zero, $`\widehat{u}`$, the root of $`g(u)`$, will be given by the solution of $`\eta u\widehat{h}e^{\frac{\eta u^2}{2}}=0`$, an equation that always has a solution. If $`\widehat{h}>0`$, then $`\widehat{x}>0`$, and $`g(u)`$ is positive for $`x>\widehat{x}`$ and negative for $`x<\widehat{x}`$. If $`\widehat{h}<0`$, then $`\widehat{x}<0`$. As an illustration, We can see the shape of $`g(u)`$ as $`\widehat{h}`$ increases in Figure 1. In summary, as $`\widehat{h}`$ increases or decreases the behavior of the front is similar to the case $`\widehat{h}=0`$, in contrast to the classical homogeneous case (11) where, as noted in Section 1, fronts with an initial velocity whose absolute value is bounded from above by $`1`$, move with a velocity that asymptotically approaches $`\widehat{h}/(\gamma ^2+\widehat{h}^2)^{\frac{1}{2}}`$.
## 4 Front Motion in 2 D
The analysis of front motion in two dimensions governed by (13) with a function $`\beta `$ of type (4) reduces to the analysis of one-dimensional front motion, and we shall not consider this case further. For radially symmetric functions, $`\beta =\beta (\rho )`$, and radially symmetric fronts, equation (13) for the radial coordinate $`\rho `$ of the front reads
$$\rho _{tt}+(\gamma \rho _t+\frac{1}{\rho })(1\rho _t^2)+\frac{\beta ^{}(\rho )}{2\beta (\rho )}(1\rho _t^2)+\beta ^{\frac{1}{2}}(\rho )\widehat{h}(1\rho _t^2)^{\frac{3}{2}},$$
(16)
We define $`u=\rho `$ and $`v=\rho _t`$ obtaining
$$\{\begin{array}{c}u_t=v,\hfill \\ v_t=[\gamma v+\frac{1}{u}+\frac{\beta ^{}(u)}{2\beta (u)}+\widehat{h}\beta ^{\frac{1}{2}}(u)(1v^2)^{\frac{1}{2}}](1v^2).\hfill \end{array}$$
(17)
The lines $`v=\pm 1`$ are trajectories of (17) in the corresponding phase plane. They define a region $`D`$ with the property that every curve starting in this region remains inside it for all future time. We confine our analysis to $`u>0`$. We analyze here the case $`\widehat{h}=0`$. The fixed points of (17) are $`(u_0,0)`$, where $`u_0`$ are solutions (if they exist) of $`2\beta (u)+u\beta ^{}(u)=0`$. The trace, $`\tau `$, and the Determinant, $`\mathrm{\Delta }`$, of the matrix of the linearized system are given by $`\tau =\gamma `$ and $`\mathrm{\Delta }=1/u_o^2+[\beta ^{\prime \prime }(u_0)\beta (u_0)\beta ^2(u_0)]/2\beta ^2(u_0)`$ respectively. The simplest case is $`\beta (\rho )=e^{\eta (\rho \rho _1)^2}`$ for a given $`\rho _1>0`$. For this case $`u_0=\left(\eta \rho _1+\sqrt{\eta ^2\rho _1^2+4\eta }\right)/2\eta `$ and $`(u_0,0)`$ is a saddle point. For (4) with $`\eta =50`$ and $`N=2`$, $`\rho _1=0.5`$ and $`\rho _2=1.5`$, we calculated the fixed points of (17) using the Newton-Raphson method with a tolerance of $`0.0001`$. They are $`z_1=0.537228`$, $`z_2=0.999165`$ and $`z_3=1.513217`$. The corresponding values of $`\mathrm{\Delta }`$ are $`\mathrm{\Delta }(z_1)=53.464832`$, $`\mathrm{\Delta }(z_2)=1196.824463`$ and $`\mathrm{\Delta }(z_3)=50.436714`$. Then $`z_1`$ and $`z_3`$ are saddle points, and $`z_2`$ is stable. Since the discriminant, $`\mathrm{\Lambda }=\tau ^24\mathrm{\Delta }`$ of $`z_2`$, is $`4787.297852<0`$, $`z_2`$ is a neutrally stable center for $`\gamma =0`$, a stable spiral for $`0<\gamma \gamma _0`$ and a stable node for $`\gamma >\gamma _0`$, where $`\gamma _0`$ is the value of $`\gamma `$ for which $`\mathrm{\Lambda }=0`$.
For the case $`\gamma =0`$, dividing the second equation in (17) by the first and solving one obtains
$$c^2u^2\beta (u)+v^2=1,$$
(18)
where $`c^2=(1v_i^2)/(u_i^2\beta (u_i))`$ and $`(u_i,v_i)`$ are the initial conditions. In Figure 2 we present a graph of (18) for $`\beta `$ given by (4) with $`N=2`$, $`\rho _1=0.5`$, $`\rho _2=1.5`$ $`\eta =50`$ (a) and $`\eta =10`$ (b). We observe that there are shrinking trajectories, oscillatory trajectories and growing trajectories in contrast with the homogeneous case where all trajectories are shrinking trajectories given by $`c^2u^2+v^2=1`$ with $`c^2=(1v_i^2)/u_i^2`$ \[kn:rotnep1\]. As $`\eta `$ decreases, the oscillatory trajectories dissapear, leaving growing trajectories, which will ultimately vanish as $`\eta 0`$.
In order to study more generally the behavior of the system away from the fixed points we can see in Figure 3 the phase plane for $`\gamma =0`$ (a) and $`\gamma =1`$ (b). The dashed lines are the nullclines of the system and the โoโ are its fixed points. The trajectories were calculated solving (17) using a Runge-Kutta method of order four. We observe there that there are different situations according to the initial conditions. In Figure 3-a (no dissipation), trajectories $`A`$, $`B`$ and $`C`$ correspond to circles that shrink to a point in finite time. If their initial velocity is positive, then their radius grows initially to a value bounded by $`z_1`$ before shrinkage takes place. Trajectories $`D`$ and $`J`$ also correspond to a circles that finally shrink to a point in finite time. In the case of $`D`$, although the initial conditions are close to those of trajectory $`C`$, the dynamics is very different. In addition to shrinkage, trajectories can display unbounded growth, represented by trajectory $`G`$, and periodic behaviour, represented by trajectories $`E`$ and $`F`$. Trajectories $`H`$ and $`I`$ also correspond to circles growing unboundedly, but if the initial velocity is negative they shrink to a valued bounded from below by $`z_3`$ and then they start growing. In Figure 3-b ($`\gamma =1`$) we see that trajectories $`A`$, $`B`$ and $`C`$ correspond to circles that shrink to points in finite time after growing to a radius bounded by $`z_1`$. Trajectory $`D`$ also shrinks to a point in finite time, but it grows initially to a radius bounded from above by $`z_3`$ and from below by $`z_2`$. As we pointed out before, as a consequence of dissipation ($`\gamma 0`$) $`z_2`$ is a stable spiral. We see that trajectory $`E`$ spirals into $`z_2`$, and there are no longer periodic trajectories. For $`N>2`$ we expect the phase plane analysis to be similar to that presented here. In contrast with the homogeneous nonlinear wave equation, where any circular front shrinks to a point in finite time, the nonhomogeneous version (13) presents a very rich dynamics with periodic motion and stabilization of circular domains of one phase inside the other.
In the absense of dissipation there are two โforcesโ responsible for the motion of the front: the curvature of the circular front, $`1/\rho `$, and the โpotential functionโ $`\beta `$. For initial conditions near enough to the minimum of $`\beta `$ the two โforcesโ balance and oscillations are possible. When dissipation is present that โbalanceโ is lost, and the oscillations decay.
## 5 Conclusions
In this manuscript we have presented equation (13) as governing the evolution of a fully developed front in a nonhomogeneous version of the nonlinear wave equation, (9), when $`ฯต1`$. This equation generalizes the damped version of the Born-Infeld equation (11) to include the effects of stronger nonlinearities and accounts for the influence of the nonhomogeneous nonlinear term on the motion of the front. The motion of interfaces according to (13) is qualitatively different from that of the homogeneous counterpart given by (11). This difference arises primarily from the fact that the function $`\beta `$ acts as a โpotential functionโ for the motion of the front. For the one dimensional case, an initial front initially placed between two maxima of $`\beta `$ (which for a homogeneous nonlinear term will move with a velocity that asymptotically approaches $`\widehat{h}/(\gamma ^2+\widehat{h}^2)^{\frac{1}{2}}`$ as long as the initial velocity bounded from above by $`1`$ in absolute value asymptotically approaches a point depending on $`\widehat{h}`$ and on the structure of $`\beta `$. For the radially symmetric two-dimensional case, the dynamics is richer than in the homogeneous counterpart, where for $`\widehat{h}=0`$ circles shrink to point in finite time. In the absence of dissipation, circles can shrink to a point in finite time, grow unboundedly or their radius oscillates, depending on the initial conditions. When dissipation effects are present, the oscillations decay, spirally or not, depending on the value of $`\gamma `$. The final result is the stabilization of a circular domain of one phase inside the other phase.
The evolution of circular interfaces in more complicated arrangement of QD sites and the evolution of more complicated fronts like convex closed curves calls for further research. We hope to address this questions in a forthcoming paper.
Acknowledgement: We thank the Chemistry Division of the National Science Foundation for support of this work.
Captions
Figure 1:
a) Graph of $`\beta (u)`$ for $`\eta =1000`$, $`x_1=2`$, $`x_2=1`$, $`x_3=0`$ and $`x_4=1`$, $`x_5=2`$.
b) Graph of $`g(u)`$ for $`\eta =1000`$, $`h=0`$, $`x_1=2`$, $`x_2=1`$, $`x_3=0`$ and $`x_4=1`$, $`x_5=2`$.
c) Graph of $`g(u)`$ for $`\eta =1000`$, $`h=10`$, $`x_1=2`$, $`x_2=1`$, $`x_3=0`$ and $`x_4=1`$, $`x_5=2`$.
d) Graph of $`g(u)`$ for $`\eta =1000`$, $`h=20`$, $`x_1=2`$, $`x_2=1`$, $`x_3=0`$ and $`x_4=1`$, $`x_5=2`$.
The points $`x_k`$ are the maxima of $`\beta (u)`$.
Figure 2:
Graph of (18) with $`\beta `$ given by (4) with $`N=2`$, $`\rho _1=0.5`$, $`\rho _2=1.5`$ and
a) $`\eta =50`$.
b) $`\eta =10`$.
Figure 3:
Phase plane for (17) with $`\beta `$ given by (4) with $`N=2`$, $`\rho _1=0.5`$, $`\rho _2=1.5`$, $`\eta =50`$ and $`\widehat{h}=0`$. The dashed lines are the nullclines of the system and the โoโ are its fixed points.
a) $`\gamma =0`$.
b) $`\gamma =1`$. |
warning/0002/hep-th0002135.html | ar5iv | text | # Implications of ๐ฉ = 2 Superconformal Symmetry
## 1 Introduction
This note is a brief review of the results obtained in our recent paper where we analysed the general structure of two- and three- point functions of the supercurrent and the flavor current of $`๐ฉ=2`$ superconformal field theories. Our research was inspired by (i) similar results obtained by Osborn for $`๐ฉ=1`$ superconformal field theories; (ii) $`๐ฉ`$โextended superconformal kinematics due to Park , in particular the existence of nilpotent superconformal invariants of three points; (iii) the conjecture of Maldacena (see for a review), which relates superconformal gauge theories in four dimensional Minkowski space to extended gauge supergravities in five dimensional anti-de-Sitter space.
In $`๐ฉ=1`$ superconformal field theory, the conserved currents are contained in two different supermultiplets: (i) the supercurrent $`J_{\alpha \dot{\alpha }}`$ containing the energy-momentum tensor $`\mathrm{\Theta }_{mn}`$, the supersymmetry currents $`j_{m\widehat{\alpha }}`$ ($`\widehat{\alpha }=\alpha ,\dot{\alpha }`$) and the axial current $`j_m^{(R)}`$; (ii) the flavor current multiplet $`L^{\overline{a}}`$ containing the conserved flavor current $`v_m^{\overline{a}}`$ among its components. Both $`J_{\alpha \dot{\alpha }}`$ and $`L^{\overline{a}}`$ are real $`๐ฉ=1`$ superfields, and they satisfy the conservation equations
$`\overline{D}^{\dot{\alpha }}J_{\alpha \dot{\alpha }}`$ $`=`$ $`D^\alpha J_{\alpha \dot{\alpha }}=0,`$ (1)
$`\overline{D}^2L^{\overline{a}}`$ $`=`$ $`D^2L^{\overline{a}}=0.`$ (2)
In $`๐ฉ=2`$ superconformal field theory, the conserved currents are contained in two different supermultiplets: (i) the supercurrent $`๐ฅ`$ whose components include the energy-momentum tensor $`\mathrm{\Theta }_{mn}`$, the SU(2) $`R`$-current $`j_m^{(ij)}`$ ($`i,j=\underset{ยฏ}{1},\underset{ยฏ}{2}`$), the axial current $`j_m^{(R)}`$ and the $`๐ฉ=2`$ supersymmetry currents $`j_{m\widehat{\alpha }}^i`$; (ii) the flavor current multiplet $`_{ij}^{\overline{a}}`$ containing the conserved flavor current $`v_m^{\overline{a}}`$ among its components. Both $`๐ฅ`$ and $`_{ij}^{\overline{a}}`$ are real $`๐ฉ=2`$ superfields ($`\overline{_{ij}}=^{ij}`$), and they satisfy the conservation equations
$`D^{ij}๐ฅ`$ $`=`$ $`\overline{D}^{ij}๐ฅ=0,`$ (3)
$`D_\alpha ^{(i}^{jk)}`$ $`=`$ $`\overline{D}_{\dot{\alpha }}^{(i}^{jk)}=0,`$ (4)
where $`D^{ij}=D^{\alpha (i}D_\alpha ^{j)}`$, $`\overline{D}^{ij}=\overline{D}_{\dot{\alpha }}^{(i}\overline{D}^{j)\dot{\alpha }}`$.
Any $`๐ฉ=2`$ superconformal field theory is a special $`๐ฉ=1`$ superconformal model. Therefore, it is useful to know the decomposition of $`๐ฅ`$ and $`_{ij}`$ into $`๐ฉ=1`$ multiplets. For that purpose we introduce the $`๐ฉ=1`$ spinor covariant derivatives $`D_\alpha D_\alpha ^{\underset{ยฏ}{1}}`$, $`\overline{D}^{\dot{\alpha }}\overline{D}_{\underset{ยฏ}{1}}^{\dot{\alpha }}`$ and define the $`๐ฉ=1`$ projection $`U|U(x,\theta _i^\alpha ,\overline{\theta }_{\dot{\alpha }}^j)|_{\theta _{\underset{ยฏ}{2}}=\overline{\theta }^{\underset{ยฏ}{2}}=0}`$ of an arbitrary $`๐ฉ=2`$ superfield $`U`$. It follows from (3) that $`๐ฅ`$ is composed of three independent $`๐ฉ=1`$ multiplets
$`J`$ $``$ $`๐ฅ|,J_\alpha D_\alpha ^{\underset{ยฏ}{2}}๐ฅ|,`$ (5)
$`J_{\alpha \dot{\alpha }}`$ $``$ $`{\displaystyle \frac{1}{2}}[D_\alpha ^{\underset{ยฏ}{2}},\overline{D}_{\dot{\alpha }\underset{ยฏ}{2}}]๐ฅ|{\displaystyle \frac{1}{6}}[D_\alpha ^{\underset{ยฏ}{1}},\overline{D}_{\dot{\alpha }\underset{ยฏ}{1}}]๐ฅ|,`$
while the $`๐ฉ=1`$ flavor current multiplet is identified as follows
$$L\mathrm{i}^{\underset{ยฏ}{12}}|.$$
(6)
Here $`J`$ and $`J_\alpha `$ satisfy the conservation equations
$`\overline{D}^2J`$ $`=`$ $`D^2J=0,`$ (7)
$`D^\alpha J_\alpha `$ $`=`$ $`\overline{D}^2J_\alpha =0.`$ (8)
The spinor object $`J_\alpha `$ contains the second supersymmetry current and two of the three SU(2) currents, namely those which correspond to the symmetries belonging to SU(2)$`/`$U(1). Finally, the scalar $`J`$ contains the current corresponding to the special combination of the $`๐ฉ=2`$ U(1) $`R`$-transformation and SU(2) $`\sigma _3`$-rotation which leaves $`\theta _{\underset{ยฏ}{1}}`$ and $`\overline{\theta }^{\underset{ยฏ}{1}}`$ invariant.
## 2 Superconformal building blocks
In $`๐ฉ`$โextended global superspace $`๐^{4|4๐ฉ}`$ parametrised by $`z^A=(x^a,\theta _i^\alpha ,\overline{\theta }_{\dot{\alpha }}^i)`$, an infinitesimal superconformal transformation
$`z^Az^A+\xi z^A,`$ (9)
$`\xi =\overline{\xi }=\xi ^a(z)_a+\xi _i^\alpha (z)D_\alpha ^i+\overline{\xi }_{\dot{\alpha }}^i(z)\overline{D}_i^{\dot{\alpha }}`$
is generated by a superconformal Killing vector $`\xi `$ defined to satisfy
$$[\xi ,D_\alpha ^i]D_\beta ^j.$$
(10)
From here it follows
$$\xi _i^\alpha =\frac{\mathrm{i}}{8}\overline{D}_{\dot{\beta }i}\xi ^{\dot{\beta }\alpha },\overline{D}_{\dot{\beta }j}\xi _i^\alpha =0$$
(11)
while the vector component of $`\xi `$ is constrained by
$`D_{(\alpha }^i\xi _{\beta )\dot{\beta }}=\overline{D}_i^{(\dot{\alpha }}\xi ^{\dot{\beta })\beta }=0,`$ (12)
$`_a\xi _b+_b\xi _a={\displaystyle \frac{1}{2}}\eta _{ab}_c\xi ^c.`$
For $`๐ฉ<4`$, the algebra of superconformal Killing vectors is isomorphic to the $`๐ฉ`$โextended superconformal algebra, su$`(2,2|๐ฉ)`$.
Let us introduce the parameters of generalized Lorentz $`\omega _{(\alpha \beta )}`$, scaleโchiral $`\sigma `$ and SU$`(๐ฉ)`$ transformations $`\mathrm{\Lambda }_i^j`$ ($`\mathrm{\Lambda }^{}\mathrm{\Lambda }=\mathrm{tr}\mathrm{\Lambda }=0`$) generated by $`\xi `$
$`[\xi ,D_\alpha ^i]`$ $`=`$ $`(D_\alpha ^i\xi _j^\beta )D_\beta ^j`$ (13)
$`=`$ $`\omega _\alpha {}_{}{}^{\beta }D_{\beta }^{i}\mathrm{i}\mathrm{\Lambda }_j{}_{}{}^{i}D_{\alpha }^{j}`$
$``$ $`{\displaystyle \frac{1}{๐ฉ}}\left((๐ฉ2)\sigma +2\overline{\sigma }\right)D_\alpha ^i.`$
A primary superfield $`๐ช(z)`$, carrying some number of undotted and dotted spinor indices and transforming in some representation of the $`R`$โsymmetry SU$`(๐ฉ)`$, satisfies the following infinitesimal transformation law under the superconformal group
$`\delta ๐ช`$ $`=`$ $`\xi ๐ช+{\displaystyle \frac{1}{2}}\omega ^{ab}M_{ab}๐ช+\mathrm{i}\mathrm{\Lambda }_j{}_{}{}^{i}R_{i}^{}{}_{}{}^{j}๐ช`$ (14)
$`2\left(q\sigma +\overline{q}\overline{\sigma }\right)๐ช.`$
Here $`M_{ab}`$ are the Lorentz generators, and $`R_i^j`$ are the generators of SU$`(๐ฉ)`$. The constant parameters $`q`$ and $`\overline{q}`$ determine the dimension $`(q+\overline{q})`$ and U(1) $`R`$โsymmetry charge $`(q\overline{q})`$ of the superfield, respectively.
In $`๐ฉ=1`$ superconformal theory, the supercurrent $`J_{\alpha \dot{\alpha }}`$ and the flavor current $`L`$ are primary superfields with the superconformal transformations
$`\delta J_{\alpha \dot{\alpha }}`$ $`=`$ $`\xi J_{\alpha \dot{\alpha }}3\left(\sigma +\overline{\sigma }\right)J_{\alpha \dot{\alpha }}`$ (15)
$`+(\omega _\alpha {}_{}{}^{\beta }\delta _{\dot{\alpha }}^{}{}_{}{}^{\dot{\beta }}+\overline{\omega }_{\dot{\alpha }}{}_{}{}^{\dot{\beta }}\delta _{\alpha }^{}{}_{}{}^{\beta })J_{\beta \dot{\beta }},`$
$`\delta L`$ $`=`$ $`\xi L2\left(\sigma +\overline{\sigma }\right)L.`$ (16)
In $`๐ฉ=2`$ superconformal theory, the supercurrent $`๐ฅ`$ and the flavor current $`_{ij}`$ are primary superfields with the superconformal transformations
$`\delta ๐ฅ`$ $`=`$ $`\xi ๐ฅ2\left(\sigma +\overline{\sigma }\right)๐ฅ,`$ (17)
$`\delta _{ij}`$ $`=`$ $`\xi _{ij}2\left(\sigma +\overline{\sigma }\right)_{ij}`$ (18)
$`+\mathrm{\hspace{0.17em}2}\mathrm{i}\mathrm{\Lambda }_{(i}{}_{}{}^{k}_{j)k}^{}.`$
Correlation functions of primary superfields, $`๐ช_1(z_1)๐ช_2(z_2)\mathrm{}๐ช_n(z_n)`$, involve some universal building blocks which we are going to describe briefly. Associated with any two points $`z_1`$ and $`z_2`$ in superspace are (anti-)chiral combinations $`x_{\overline{1}2}`$, $`\theta _{12}`$ and $`\overline{\theta }_{12}`$:
$`x_{\overline{1}2}^a`$ $`=`$ $`x_{2\overline{1}}^a=x_1^ax_{2+}^a+2\mathrm{i}\theta _{2i}\sigma ^a\overline{\theta }_1^i,`$
$`x_\pm ^a`$ $``$ $`x^a\pm \mathrm{i}\theta _i\sigma ^a\overline{\theta }^i;`$
$`\theta _{12}`$ $`=`$ $`\theta _1\theta _2,\overline{\theta }_{12}=\overline{\theta }_1\overline{\theta }_2,`$ (19)
which are invariant under Poincarรฉ supersymmetry transformations (the notation โ$`x_{\overline{1}2}`$โ indicates that $`x_{\overline{1}2}`$ is antichiral with respect to $`z_1`$ and chiral with respect to $`z_2`$) but transforms semi-covariantly with respect to the superconformal group (see, e.g. ). In extended supersymmetry, there exist primary superfields with isoindices, and their correlation functions generically involve a conformally covariant $`๐ฉ\times ๐ฉ`$ unimodular matrix<sup>1</sup><sup>1</sup>1We use the notation adopted in . When the spinor indices are not indicated explicitly, the following matrix-like conventions are used : $`\psi =(\psi ^\alpha )`$, $`\stackrel{~}{\psi }=(\psi _\alpha )`$, $`\overline{\psi }=(\overline{\psi }^{\dot{\alpha }})`$, $`\stackrel{~}{\overline{\psi }}=(\overline{\psi }_{\dot{\alpha }})`$, $`x=(x_{\alpha \dot{\alpha }})`$, $`\stackrel{~}{x}=(x^{\dot{\alpha }\alpha })`$; but $`x^2x^ax_a=\frac{1}{2}\mathrm{tr}(\stackrel{~}{x}x)`$, and hence $`\stackrel{~}{x}^1=x/x^2`$.
$`\widehat{u}_i{}_{}{}^{j}(z_{12})`$ $`=`$ $`\left({\displaystyle \frac{x_{\overline{2}1}^2}{x_{\overline{1}2}^2}}\right)^{1/๐ฉ}u_i{}_{}{}^{j}(z_{12}),`$
$`u_i{}_{}{}^{j}(z_{12})`$ $`=`$ $`\delta _i{}_{}{}^{j}4\mathrm{i}{\displaystyle \frac{\theta _{12i}x_{\overline{1}2}\overline{\theta }_{12}^j}{x_{\overline{1}2}^2}},`$ (20)
with the basic properties
$`\widehat{u}^{}(z_{12})\widehat{u}(z_{12})=\mathrm{๐},`$
$`\widehat{u}^1(z_{12})=\widehat{u}(z_{21}),`$
$`det\widehat{u}(z_{12})=1,`$ (21)
and the transformation rule
$`\delta \widehat{u}_i{}_{}{}^{j}(z_{12})`$ $`=`$ $`\mathrm{i}\mathrm{\Lambda }_i{}_{}{}^{k}(z_1)\widehat{u}_k{}_{}{}^{j}(z_{12})`$ (22)
$``$ $`\mathrm{i}\widehat{u}_i{}_{}{}^{k}(z_{12})\mathrm{\Lambda }_k{}_{}{}^{j}(z_2).`$
Given three superspace points $`z_1,z_2`$ and $`z_3`$, one can define superconformally covariant bosonic and fermionic variables $`๐_1,๐_2`$ and $`๐_3`$, where $`๐_1=(๐_1,\mathrm{\Theta }_1^i,\overline{\mathrm{\Theta }}_{1i})`$ are
$`๐_1`$ $`=`$ $`\stackrel{~}{x}_{1\overline{2}}{}_{}{}^{1}\stackrel{~}{x}_{\overline{2}3}^{}\stackrel{~}{x}_{3\overline{1}}{}_{}{}^{1},`$
$`\stackrel{~}{\mathrm{\Theta }}_1^i`$ $`=`$ $`\mathrm{i}\left(\stackrel{~}{x}_{\overline{2}1}{}_{}{}^{1}\overline{\theta }_{12}^{i}\stackrel{~}{x}_{\overline{3}1}{}_{}{}^{1}\overline{\theta }_{13}^{i}\right),`$
$`\overline{๐}_1`$ $`=`$ $`๐_1^{}=๐_14\mathrm{i}\stackrel{~}{\mathrm{\Theta }}_1^i\stackrel{~}{\overline{\mathrm{\Theta }}}_{1i},`$
$`\stackrel{~}{\overline{\mathrm{\Theta }}}_{1i}`$ $`=`$ $`(\stackrel{~}{\mathrm{\Theta }}_1^i)^{}`$ (23)
and $`๐_2,๐_3`$ are obtained from here by cyclically permuting indices. These structures possess remarkably simple superconformal transformation rules:
$`\delta ๐_{1\alpha \dot{\alpha }}`$ $`=`$ $`(\omega _\alpha {}_{}{}^{\beta }(z_1)\delta _\alpha {}_{}{}^{\beta }\sigma (z_1))๐_{1\beta \dot{\alpha }}`$
$`+`$ $`๐_{1\alpha \dot{\beta }}(\overline{\omega }^{\dot{\beta }}{}_{\dot{\alpha }}{}^{}(z_1)\delta ^{\dot{\beta }}{}_{\dot{\alpha }}{}^{}\overline{\sigma }(z_1)),`$
$`\delta \mathrm{\Theta }_{1\alpha }^i`$ $`=`$ $`\omega _\alpha {}_{}{}^{\beta }(z_1)\mathrm{\Theta }_{1\beta }^i\mathrm{i}\mathrm{\Theta }_{1\alpha }^j\mathrm{\Lambda }_j{}_{}{}^{i}(z_1)`$
$``$ $`{\displaystyle \frac{1}{๐ฉ}}\left((๐ฉ2)\sigma (z_1)+2\overline{\sigma }(z_1)\right)\mathrm{\Theta }_{1\alpha }^i`$
and turn out to be essential building blocks for correlations functions of primary superfields. The variables $`๐`$ with different labels are related to each other, in particular:
$`\stackrel{~}{x}_{\overline{1}3}๐_3\stackrel{~}{x}_{\overline{3}1}`$ $`=`$ $`\overline{๐}_1{}_{}{}^{1},`$
$`\stackrel{~}{x}_{\overline{1}3}\stackrel{~}{\mathrm{\Theta }}_3^iu_i{}_{}{}^{j}(z_{31})`$ $`=`$ $`๐_1{}_{}{}^{1}\stackrel{~}{\mathrm{\Theta }}_{1}^{j}.`$ (24)
With the aid of the matrices $`u(z_{rs})`$, $`r,s=1,2,3`$, defined in (20), one can construct unitary matrices $`\widehat{๐ฎ}(๐_s)`$ , in particular
$`\widehat{๐ฎ}(๐_3)`$ $`=`$ $`\widehat{u}(z_{31})\widehat{u}(z_{12})\widehat{u}(z_{23})`$
$`=`$ $`\left({\displaystyle \frac{\overline{๐}_3^2}{๐_3^2}}\right)^{1/๐ฉ}\left(\delta _i{}_{}{}^{j}4\mathrm{i}\stackrel{~}{\overline{\mathrm{\Theta }}}_{3i}๐_3{}_{}{}^{1}\stackrel{~}{\mathrm{\Theta }}_{3}^{j}\right)`$
transforming at $`z_3`$ only. Their properties are
$$\widehat{๐ฎ}^{}(๐_3)=\widehat{๐ฎ}^1(๐_3),det\widehat{๐ฎ}(๐_3)=1.$$
(26)
The above general formalism has specific features in the case $`๐ฉ=2`$ that is of primary interest for us. Here we have at our disposal the SU(2)โinvariant tensors $`\epsilon _{ij}=\epsilon _{ji}`$ and $`\epsilon ^{ij}=\epsilon ^{ji}`$, normalized to $`\epsilon ^{\underset{ยฏ}{12}}=\epsilon _{\underset{ยฏ}{21}}=1`$. They can be used to raise and lower isoindices: $`C^i=\epsilon ^{ij}C_j`$, $`C_i=\epsilon _{ij}C^j`$. For $`๐ฉ=2`$, the condition of unimodularity of the matrix defined in (20) can be written as
$$\widehat{u}_{ji}(z_{21})=\widehat{u}_{ij}(z_{12}).$$
(27)
The importance of this relation is that it implies that the two-point function
$`A_{i_1i_2}(z_1,z_2)`$ $``$ $`{\displaystyle \frac{\widehat{u}_{i_1i_2}(z_{12})}{(x_{\overline{1}2}{}_{}{}^{2}x_{\overline{2}1}^{}{}_{}{}^{2})^{\frac{1}{2}}}}`$ (28)
$`=`$ $`{\displaystyle \frac{\widehat{u}_{i_2i_1}(z_{21})}{(x_{\overline{1}2}{}_{}{}^{2}x_{\overline{2}1}^{}{}_{}{}^{2})^{\frac{1}{2}}}}`$
is analytic in $`z_1`$ and $`z_2`$ for $`z_1z_2`$,
$`D_{1\alpha (j_1}A_{i_1)i_2}(z_1,z_2)`$ $`=`$ $`0,`$
$`\overline{D}_{1\dot{\alpha }(j_1}A_{i_1)i_2}(z_1,z_2)`$ $`=`$ $`0.`$ (29)
As we will see later, $`A_{i_1i_2}(z_1,z_2)`$ is a building block of correlation functions of analytic primary superfields like the $`๐ฉ=2`$ flavor currents. For $`๐ฉ=2`$, the fact that $`\widehat{๐ฎ}(๐_3)`$ is unimodular and unitary, implies
$`\mathrm{tr}\widehat{๐ฎ}^{}(๐_3)`$ $`=`$ $`\mathrm{tr}\widehat{๐ฎ}(๐_3),`$
$`\widehat{๐ฎ}_{ji}^{}(๐_3)`$ $`=`$ $`\widehat{๐ฎ}_{ij}(๐_3).`$ (30)
## 3 Correlation functions of $`๐ฉ`$ = 2 currents
According to the general prescription of , the two-point function of a primary superfield $`๐ช_{}`$, which is a Lorentz scalar and transforms in a representation $`T`$ of the $`R`$โsymmetry group SU$`(๐ฉ)`$, with its conjugate $`\overline{๐ช}^๐ฅ`$ reads
$`๐ช_{}(z_1)\overline{๐ช}^๐ฅ(z_2)`$ $`=`$ $`C_๐ช{\displaystyle \frac{T_{}{}_{}{}^{๐ฅ}\left(\widehat{u}(z_{12})\right)}{(x_{\overline{1}2}{}_{}{}^{2})^{\overline{q}}(x_{\overline{2}1}{}_{}{}^{2})^q}},`$
where $`C_๐ช`$ is a normalization constant.
For the $`๐ฉ=2`$ supercurrent $`๐ฅ`$ and the flavor current $`_{ij}^{\overline{a}}`$, the above prescription gives
$`๐ฅ(z_1)๐ฅ(z_2)=c_๐ฅ{\displaystyle \frac{1}{x_{\overline{1}2}{}_{}{}^{2}x_{\overline{2}1}^{}^2}},`$ (31)
$`_{i_1j_1}^{\overline{a}_1}(z_1)^{\overline{a}_2i_2j_2}(z_2)=2c_{}\delta ^{\overline{a}_1\overline{a}_2}`$
$`\times {\displaystyle \frac{\widehat{u}_{i_1}{}_{}{}^{(i_2}(z_{12})\widehat{u}_{j_1}{}_{}{}^{j_2)}(z_{12})}{x_{\overline{1}2}{}_{}{}^{2}x_{\overline{2}1}^{}^2}}.`$ (32)
The relevant conservation equations prove to be satisfied at $`z_1z_2`$,
$`D_1{}_{}{}^{ij}๐ฅ(z_1)๐ฅ(z_2)`$ $`=`$ $`0,`$
$`D_{1\alpha (k_1}_{i_1j_1)}(z_1)^{i_2j_2}(z_2)`$ $`=`$ $`0.`$ (33)
According to the general prescription of , the three-point function of primary superfields $`๐ช__1^{(1)}`$, $`๐ช__2^{(2)}`$ and $`๐ช__3^{(3)}`$ reads
$`๐ช__1^{(1)}(z_1)๐ช__2^{(2)}(z_2)๐ช__3^{(3)}(z_3)`$
$`={\displaystyle \frac{T^{(1)}{}_{_1}{}^{}{}_{}{}^{๐ฅ_1}\left(\widehat{u}(z_{13})\right)T^{(2)}{}_{_2}{}^{}{}_{}{}^{๐ฅ_2}\left(\widehat{u}(z_{23})\right)}{(x_{\overline{1}3}{}_{}{}^{2})^{\overline{q}_1}(x_{\overline{3}1}{}_{}{}^{2})^{q_1}(x_{\overline{2}3}{}_{}{}^{2})^{\overline{q}_2}(x_{\overline{3}2}{}_{}{}^{2})^{q_2}}}`$
$`\times H_{๐ฅ_1๐ฅ_2_3}(๐_3).`$
Here $`H_{๐ฅ_1๐ฅ_2_3}(๐_3)`$ transforms as an isotensor at $`z_3`$ in the representations $`T^{(1)},T^{(2)}`$ and $`T^{(3)}`$ with respect to the indices $`๐ฅ_1,๐ฅ_2`$ and $`_3`$, respectively, and possesses the homogeneity property
$`H_{๐ฅ_1๐ฅ_2_3}(\mathrm{\Delta }\overline{\mathrm{\Delta }}๐,\mathrm{\Delta }\mathrm{\Theta },\overline{\mathrm{\Delta }}\overline{\mathrm{\Theta }})`$
$`=\mathrm{\Delta }^{2p}\overline{\mathrm{\Delta }}^{2\overline{p}}H_{๐ฅ_1๐ฅ_2_3}(๐,\mathrm{\Theta },\overline{\mathrm{\Theta }}),`$
$`p2\overline{p}=\overline{q}_1+\overline{q}_2q_3,`$
$`\overline{p}2p=q_1+q_2\overline{q}_3.`$
In general, the latter equation admits a finite number of linearly independent solutions, and this can be considerably reduced by taking into account the symmetry properties, superfield conservation equations and, of course, the superfield constraints (such as chirality or analyticity ). Below we shall present the most general expressions for three-point functions of the $`๐ฉ=2`$ supercurrent $`๐ฅ`$ and the flavor current $`_{ij}^{\overline{a}}`$, which are compatible with all physical requirements. Details can be found in .
The three-point function of the $`๐ฉ=2`$ supercurrent is
$`๐ฅ(z_1)๐ฅ(z_2)๐ฅ(z_3)`$ (34)
$`={\displaystyle \frac{1}{x_{\overline{1}3}{}_{}{}^{2}x_{\overline{3}1}^{}{}_{}{}^{2}x_{\overline{2}3}^{}{}_{}{}^{2}x_{\overline{3}2}^{}^2}}`$
$`\times \{A({\displaystyle \frac{1}{๐_3^2}}+{\displaystyle \frac{1}{\overline{๐}_3^2}})`$
$`+B{\displaystyle \frac{\mathrm{\Theta }_3^{\alpha \beta }๐_{3\alpha \dot{\alpha }}๐_{3\beta \dot{\beta }}\overline{\mathrm{\Theta }}_3^{\dot{\alpha }\dot{\beta }}}{(๐_3{}_{}{}^{2})^2}}\},`$
where
$`\mathrm{\Theta }_3^{\alpha \beta }`$ $`=`$ $`\mathrm{\Theta }_3^{(\alpha \beta )}=\mathrm{\Theta }_3^{\alpha i}\mathrm{\Theta }_{3i}^\beta ,`$
$`\overline{\mathrm{\Theta }}_3^{\dot{\alpha }\dot{\beta }}`$ $`=`$ $`\overline{\mathrm{\Theta }}_3^{(\dot{\alpha }\dot{\beta })}=\overline{\mathrm{\Theta }}_{3i}^{\dot{\alpha }}\overline{\mathrm{\Theta }}_3^{\dot{\alpha }i},`$ (35)
and $`A,B`$ are real parameters. The second structure is nilpotent and real.
The three-point function of the $`๐ฉ=2`$ flavor current reads
$`_{i_1j_1}^{\overline{a}}(z_1)_{i_2j_2}^{\overline{b}}(z_2)_{i_3j_3}^{\overline{c}}(z_3)`$
$`=`$ $`{\displaystyle \frac{\widehat{u}_{i_1}{}_{}{}^{k_1}(z_{13})\widehat{u}_{j_1}{}_{}{}^{l_1}(z_{13})\widehat{u}_{i_2}{}_{}{}^{k_2}(z_{23})\widehat{u}_{j_2}{}_{}{}^{l_2}(z_{23})}{x_{\overline{3}1}{}_{}{}^{2}x_{\overline{1}3}^{}{}_{}{}^{2}x_{\overline{3}2}^{}{}_{}{}^{2}x_{\overline{2}3}^{}^2}}`$
$`\times `$ $`f^{\overline{a}\overline{b}\overline{c}}\{{\displaystyle \frac{\epsilon _{i_3(k_1}\widehat{๐ฎ}{}_{l_1)(l_2}{}^{}(๐_3)\epsilon _{k_2)j_3}}{(๐_3{}_{}{}^{2}\overline{๐}_{3}^{}{}_{}{}^{2})^{\frac{1}{2}}}}+(i_3j_3)\}`$
with $`f^{\overline{a}\overline{b}\overline{c}}=f^{[\overline{a}\overline{b}\overline{c}]}`$ a completely antisymmetric real tensor being proportional to the structure constants of the flavor group.
For mixed correlation functions of the $`๐ฉ=2`$ supercurrent and the flavor current, we get
$`๐ฅ(z_1)๐ฅ(z_2)_{ij}^{\overline{a}}(z_3)=0,`$
$`_{i_1j_1}^{\overline{a}}(z_1)_{i_2j_2}^{\overline{b}}(z_2)๐ฅ(z_3)=d\delta ^{\overline{a}\overline{b}}`$
$`\times `$ $`{\displaystyle \frac{\widehat{u}_{i_1}{}_{}{}^{k_1}(z_{13})\widehat{u}_{j_1}{}_{}{}^{l_1}(z_{13})\widehat{u}_{i_2}{}_{}{}^{k_2}(z_{23})\widehat{u}_{j_2}{}_{}{}^{l_2}(z_{23})}{x_{\overline{3}1}{}_{}{}^{2}x_{\overline{1}3}^{}{}_{}{}^{2}x_{\overline{3}2}^{}{}_{}{}^{2}x_{\overline{2}3}^{}^2}}`$
$`\times `$ $`{\displaystyle \frac{\epsilon _{k_2(k_1}\widehat{๐ฎ}{}_{l_1)l_2}{}^{}(๐_3)+\epsilon _{l_2(k_1}\widehat{๐ฎ}{}_{l_1)k_2}{}^{}(๐_3)}{(๐_3{}_{}{}^{2}\overline{๐}_{3}^{}{}_{}{}^{2})^{\frac{1}{2}}}},`$
with $`d`$ a real parameter which can be related, via supersymmetric Ward identities, to the parameter $`c_{}`$ in the two-point function (32),
$$d=\frac{1}{4\pi ^2}c_{}.$$
(39)
It is worth pointing out that eq. (3) is one of the important consequences of $`๐ฉ=2`$ superconformal symmetry and has no direct analog in the $`๐ฉ=1`$ case. In a generic $`๐ฉ=1`$ superconformal theory with a flavor current $`L`$, the correlation function $`J_{\alpha \dot{\alpha }}J_{\beta \dot{\beta }}L`$ is not restricted by $`๐ฉ=1`$ superconformal symmetry to vanish .
## 4 Reduction to $`๐ฉ`$ = 1 superfields
From the point of view of $`๐ฉ=1`$ superconformal symmetry, any $`๐ฉ=2`$ primary superfield consists of several $`๐ฉ=1`$ primary superfields. Having computed the correlation functions of $`๐ฉ=2`$ primary superfields, one can read off all correlators of their $`๐ฉ=1`$ superconformal components. Since any $`๐ฉ=2`$ superconformal theory is a particular $`๐ฉ=1`$ superconformal theory, one can then simply make use of $`๐ฉ=1`$ superconformal Ward identities to relate the coefficients of various correlators.
Using the explicit form (31) of the $`๐ฉ=2`$ supercurrent two-point function, one can read off the two-point functions of the $`๐ฉ=1`$ primary superfields contained in $`๐ฅ`$, in particular<sup>2</sup><sup>2</sup>2Here and below, all building blocks are expressed in $`๐ฉ=1`$ superspace.
$`J(z_1)J(z_2)`$ $`=`$ $`c_๐ฅ{\displaystyle \frac{1}{x_{\overline{1}2}{}_{}{}^{2}x_{\overline{2}1}^{}^2}},`$ (40)
$`J_{\alpha \dot{\alpha }}(z_1)J_{\beta \dot{\beta }}(z_2)`$ $`=`$ $`{\displaystyle \frac{64}{3}}c_๐ฅ{\displaystyle \frac{(x_{1\overline{2}})_{\alpha \dot{\beta }}(x_{2\overline{1}})_{\beta \dot{\alpha }}}{(x_{\overline{1}2}{}_{}{}^{2}x_{\overline{2}1}^{}{}_{}{}^{2})^2}}.`$
Similarly, the two-point function of the $`๐ฉ=1`$ flavor current follows from (32)
$$L^{\overline{a}_1}(z_1)L^{\overline{a}_2}(z_2)=c_{}\frac{\delta ^{\overline{a}_1\overline{a}_2}}{x_{\overline{1}2}{}_{}{}^{2}x_{\overline{2}1}^{}^2}.$$
(41)
We now present several $`๐ฉ=1`$ three-point functions which are encoded in that of the $`๐ฉ=2`$ supercurrent, given by eq. (34).
$`J(z_1)J(z_2)J(z_3)={\displaystyle \frac{A}{x_{\overline{1}3}{}_{}{}^{2}x_{\overline{3}1}^{}{}_{}{}^{2}x_{\overline{2}3}^{}{}_{}{}^{2}x_{\overline{3}2}^{}^2}}`$
$`\times \left({\displaystyle \frac{1}{๐_3^2}}+{\displaystyle \frac{1}{\overline{๐}_3^2}}\right),`$ (42)
$`J(z_1)J(z_2)J_{\alpha \dot{\alpha }}(z_3)={\displaystyle \frac{1}{12}}(8A3B)`$
$`\times {\displaystyle \frac{1}{x_{\overline{1}3}{}_{}{}^{2}x_{\overline{3}1}^{}{}_{}{}^{2}x_{\overline{2}3}^{}{}_{}{}^{2}x_{\overline{3}2}^{}^2}}`$
$`\times \{{\displaystyle \frac{2(๐_3๐_3)๐_{3\alpha \dot{\alpha }}+๐_3{}_{}{}^{2}๐_{3\alpha \dot{\alpha }}^{}}{(๐_3{}_{}{}^{2})^2}}`$
$`+(๐_3\overline{๐}_3)\},`$ (43)
$`J_{\alpha \dot{\alpha }}(z_1)J_{\beta \dot{\beta }}(z_2)J(z_3)={\displaystyle \frac{4}{9}}(8A+3B)`$
$`\times {\displaystyle \frac{(x_{1\overline{3}})_{\alpha \dot{\gamma }}(x_{3\overline{1}})_{\gamma \dot{\alpha }}(x_{2\overline{3}})_{\beta \dot{\delta }}(x_{3\overline{2}})_{\delta \dot{\beta }}}{(x_{\overline{1}3}{}_{}{}^{2}x_{\overline{3}1}^{}{}_{}{}^{2}x_{\overline{2}3}^{}{}_{}{}^{2}x_{\overline{3}2}^{}{}_{}{}^{2})^2}}`$
$`\times \{{\displaystyle \frac{๐_3{}_{}{}^{\gamma \dot{\gamma }}๐_{3}^{}^{\delta \dot{\delta }}}{(๐_3{}_{}{}^{2})^3}}+{\displaystyle \frac{1}{2}}{\displaystyle \frac{\epsilon ^{\gamma \delta }\epsilon ^{\dot{\gamma }\dot{\delta }}}{(๐_3{}_{}{}^{2})^2}}`$
$`+(๐_3\overline{๐}_3)\},`$ (44)
with $`๐_a`$ defined by
$$\overline{๐}_a๐_a=\mathrm{i}๐_a,๐_a=2\mathrm{\Theta }\sigma _a\overline{\mathrm{\Theta }}.$$
(45)
The most interesting correlator and by far the most laborious to compute is
$`J_{\alpha \dot{\alpha }}(z_1)J_{\beta \dot{\beta }}(z_2)J_{\gamma \dot{\gamma }}(z_3)`$ (46)
$`={\displaystyle \frac{(x_{1\overline{3}})_{\alpha \dot{\sigma }}(x_{3\overline{1}})_{\sigma \dot{\alpha }}(x_{2\overline{3}})_{\beta \dot{\delta }}(x_{3\overline{2}})_{\delta \dot{\beta }}}{(x_{\overline{1}3}{}_{}{}^{2}x_{\overline{3}1}^{}{}_{}{}^{2}x_{\overline{2}3}^{}{}_{}{}^{2}x_{\overline{3}2}^{}{}_{}{}^{2})^2}}`$
$`\times H^{\dot{\sigma }\sigma ,\dot{\delta }\delta }{}_{\gamma \dot{\gamma }}{}^{}(๐_3,\overline{๐}_3),`$
$`H^{\dot{\sigma }\sigma ,\dot{\delta }\delta }{}_{\gamma \dot{\gamma }}{}^{}(๐_3,\overline{๐}_3)=h^{\dot{\sigma }\sigma ,\dot{\delta }\delta }{}_{\gamma \dot{\gamma }}{}^{}(๐_3,\overline{๐}_3)`$
$`+h^{\dot{\delta }\delta ,\dot{\sigma }\sigma }{}_{\gamma \dot{\gamma }}{}^{}(\overline{๐}_3,๐_3),`$
where
$`h^{abc}(๐,\overline{๐}){\displaystyle \frac{1}{8}}(\sigma ^a)_{\alpha \dot{\alpha }}(\sigma ^a)_{\beta \dot{\beta }}(\stackrel{~}{\sigma }^c)^{\dot{\gamma }\gamma }`$
$`\times h^{\dot{\alpha }\alpha ,\dot{\beta }\beta }{}_{\gamma \dot{\gamma }}{}^{}(๐,\overline{๐})`$ (47)
$`={\displaystyle \frac{16}{27}}(26A{\displaystyle \frac{9}{4}}B){\displaystyle \frac{\mathrm{i}}{(๐^2)^2}}`$
$`\times \left(๐^a\eta ^{bc}+๐^b\eta ^{ac}๐^c\eta ^{ab}+\mathrm{i}\epsilon ^{abcd}๐_d\right)`$
$`{\displaystyle \frac{8}{27}}(8A9B){\displaystyle \frac{1}{(๐^2)^3}}`$
$`\times \{2(๐^a๐^b+๐^b๐^a)๐^c`$
$`3๐^a๐^b\left(๐^c+2{\displaystyle \frac{(๐๐)}{๐^2}}๐^c\right)`$
$`(๐๐)\left(3(๐^a\eta ^{bc}+๐^b\eta ^{ac})2๐^c\eta ^{ab}\right)`$
$`+{\displaystyle \frac{1}{2}}๐^2(๐^a\eta ^{bc}+๐^b\eta ^{ac}+๐^c\eta ^{ab})\}.`$
Our final relations (46) and (47) perfectly agree with the general structure of the three-point function of the supercurrent in $`๐ฉ=1`$ superconformal field theory .
Using the results of , one may express $`A`$ and $`B`$ in terms of the anomaly coefficients
$`a`$ $`=`$ $`{\displaystyle \frac{1}{24}}(5n_V+n_H),`$
$`c`$ $`=`$ $`{\displaystyle \frac{1}{12}}(2n_V+n_H),`$ (48)
where $`n_V`$ and $`n_H`$ denote the number of free $`๐ฉ=2`$ vector multiplets and hypermultiplets, respectivley. We get<sup>3</sup><sup>3</sup>3Our definition of the $`๐ฉ=1`$ supercurrent corresponds to that adopted in and differs in sign from Osbornโs convention .
$`A`$ $`=`$ $`{\displaystyle \frac{3}{64\pi ^6}}(4a3c),`$
$`B`$ $`=`$ $`{\displaystyle \frac{1}{8\pi ^6}}(4a5c).`$ (49)
In $`๐ฉ=1`$ supersymmetry, a superconformal Ward identity relates the coefficient in the two-point function of the supercurrent (40) to the anomaly coefficient $`c`$ as follows
$$c_๐ฅ=\frac{3}{8\pi ^4}c.$$
(50)
In terms of the coefficients $`A`$ and $`B`$ this relation reads
$$\frac{2}{\pi ^2}c_๐ฅ=8A3B.$$
(51)
Let us turn to the three-point function of the $`๐ฉ=2`$ flavor current given by eq. (3). From it one reads off the three-point function of the $`๐ฉ=1`$ flavor current
$`L^{\overline{a}}(z_1)L^{\overline{b}}(z_2)L^{\overline{c}}(z_3)`$ (52)
$`={\displaystyle \frac{1}{4}}f^{\overline{a}\overline{b}\overline{c}}{\displaystyle \frac{\mathrm{i}}{x_{\overline{1}3}{}_{}{}^{2}x_{\overline{3}1}^{}{}_{}{}^{2}x_{\overline{2}3}^{}{}_{}{}^{2}x_{\overline{3}2}^{}^2}}`$
$`\times \left({\displaystyle \frac{1}{\overline{๐}_3^2}}{\displaystyle \frac{1}{๐_3^2}}\right).`$
It is worth noting that the Ward identities allow one to represent $`f^{\overline{a}\overline{b}\overline{c}}`$ as a product of $`c_{}`$ and the structure constants of the flavor symmetry group, see for more details.
In $`๐ฉ=1`$ superconformal field theory, the three-point function of the flavor current superfield $`L`$ contains, in general, two linearly independent forms :
$`L^{\overline{a}}(z_1)L^{\overline{b}}(z_2)L^{\overline{c}}(z_3)`$
$`={\displaystyle \frac{1}{x_{\overline{1}3}{}_{}{}^{2}x_{\overline{3}1}^{}{}_{}{}^{2}x_{\overline{2}3}^{}{}_{}{}^{2}x_{\overline{3}2}^{}^2}}`$
$`\times \{\mathrm{i}f^{[\overline{a}\overline{b}\overline{c}]}({\displaystyle \frac{1}{๐_3^2}}{\displaystyle \frac{1}{\overline{๐}_3^2}})`$
$`+d^{(\overline{a}\overline{b}\overline{c})}({\displaystyle \frac{1}{๐_3^2}}+{\displaystyle \frac{1}{\overline{๐}_3^2}})\}.`$
The second term, involving a completely symmetric group tensor $`d^{\overline{a}\overline{b}\overline{c}}`$, reflects the presence of chiral anomalies in the theory. The field-theoretic origin of this term is due to the fact that the $`๐ฉ=1`$ conservation equation $`\overline{D}^2L=D^2L=0`$ admits a non-trivial deformation
$`\overline{D}^2L^{\overline{a}}d^{\overline{a}\overline{b}\overline{c}}W^{\overline{b}\alpha }W_\alpha ^{\overline{c}}`$
when the chiral flavor current is coupled to a background vector multiplet. Eq. (52) tells us that the flavor currents are anomaly-free in $`๐ฉ=2`$ superconformal theory. This agrees with the facts that (i) $`๐ฉ=2`$ super Yang-Mills models are non-chiral; (ii) the $`๐ฉ=2`$ conservation equation (4) does not possess non-trivial deformations.
Acknowledgements Support from DFG, the German National Science Foundation, from GIF, the German-Israeli foundation for Scientific Research and from the EEC under TMR contract ERBFMRX-CT96-0045 is gratefully acknowledged. This work was also supported in part by the NATO collaborative research grant PST.CLG 974965, by the RFBR grant No. 99-02-16617, by the INTAS grant No. 96-0308 and by the DFG-RFBR grant No. 99-02-04022. |
warning/0002/hep-th0002111.html | ar5iv | text | # 1 Introduction
## 1 Introduction
The proposed duality between string theory on anti-de Sitter space and lower-dimensional conformal field theory provides a non-perturbative definition of string theory, and could thus, subject to the restriction on the asymptotic boundary conditions, cast a bright light on many dark corners of quantum gravity. In particular, the field theory description encompasses arbitrary fluctuations of the metric and other fields in the interior, and should provide a fully quantum description of the formation and evaporation of a black hole. One of the major barriers to studying conceptual questions in quantum gravity using this theory is our poor understanding of how an approximately local classical (or semi-classical) spacetime description of the physics emerges from the fundamental gauge theory description, and the consequent absence of any intuition about how this approximate locality breaks down under extreme conditions. (A related problem is that in the regime where a classical spacetime description is a good approximation, we donโt have any other quantitative description; see for a recent attempt to construct calculationally useful approximations.)
The connection between asymptotic behavior of the spacetime fields and the field theory was one of the first subjects of study , and it was subsequently shown that the map between states in the field theory and states in spacetime identifies the asymptotic behavior of the fields with the expectation values of local operators in the field theory . This was used to show a โscale-radius dualityโ for a variety of bulk sources, and for wavepackets of supergravity fields โ the radial position of a bulk probe is encoded in the scale size of the dual expectation values. Dynamical sources for supergravity fields were studied in , where the radial position of a source particle following a bulk geodesic was reflected in the size and shape of an expectation-value bubble in the CFT. The expectation values of the operators produced by spacetime sources were further studied in .
However, the simple scale-radius relationship seen in these studies is a consequence of an isometry in pure AdS space which is dual to a scale transformation in the conformal field theory, under which the vacuum remains invariant. For situations describing black holes, which break the symmetries, the relationship between bulk position and boundary observables will be more complicated . The same phenomenon is apparent in the collision of two massless particles to form a black hole in ; after the particles collide, their radial position is fixed, but the scales in the boundary expectation values continue to evolve.
Furthermore, the asymptotic values of the fields are not sufficient to reproduce the whole spacetime. Since asymptotic values of fields in AdS space are dual to the expectation values of local operators in the CFT, it follows that such expectation values describe only a small piece of the physical information. A number of authors have studied spacetime sources which do not change the asymptotic values of the fields, such as particles in AdS<sub>3</sub> and spherical shells, and found that the location of the shell or particles is encoded in non-local operator expectation values in the field theory, such as the two-point function and Wilson loops . Similar work is described in . Thus, non-local operators must be included in any understanding of the bulk-boundary connection. A particularly striking case is asymptotically AdS<sub>3</sub> spaces, where we can describe a wide range of dynamics without relying on perturbations around some background , and the asymptotic metric only encodes the total mass and angular momentum of the system.
In , Balasubramanian and Ross used a stationary phase approximation to obtain predictions for the propagator in the gauge theory from the geodesics of supergravity solutions in which a black hole was formed. This propagator appeared to be sensitive to events in the interior of the black hole. Now, while the CFT may well encode information about the black hole interior, the particular CFT propagator studied in is in fact the restriction to the boundary of AdS space of a propagator associated with the bulk quantum field theory. This raises certain issues about causality<sup>1</sup><sup>1</sup>1These issues were brought to our attention by Lenny Susskind through his comments at the Val Morin workshop on Black Holes, June 1999. which we wish to clarify in the work below. Some general arguments are presented in section 2. In short, we argue that the propagator studied in is in fact a causal object, but that the stationary phase approximation is valid only in appropriately analytic spacetimes and not in the actual spacetime considered in . However, even without the stationary-phase approximation, the path-integral definition of the propagator used in should generally lead to a result which depends on the region inside the black hole; we argue that this should be interpreted as an object which is defined by a mixture of past and future boundary conditions.
We then proceed to explore the propagator in two analytic spacetimes in order to see more precisely what sort of object it represents. The spacetimes that we consider contain black holes, but are static outside the Killing horizon. In those cases, the stationary phase approximation is expected to be valid, and a computation of the propagator reduces to a study of various geodesics in the bulk spacetime. We show that, in such cases, the propagator of is in fact the boundary limit of a time-ordered expectation value of a product of local bulk fields. Our spacetimes are the spinless BTZ black hole and the associated $`^2`$ geon . We find that the propagator in each case is associated with a natural vacuum state for linearized quantum fields on the spacetime, and that geodesics passing behind the black hole horizon play an important role in determining the structure of this state. The states are analogues of the Hartle-Hawking state, and are defined by boundary conditions at past and future infinity. The propagators in these cases are known to be Greenโs functions of a (causal) wave equation, and sensitivity to โeventsโ behind the event horizon would once again seem to contradict this causality. In this case, the resolution is that the analyticity of these spacetimes implies that much of the information about the region inside the event horizon is in fact โstored outsideโ. Note however that knowledge of the region outside the Killing horizon is not enough to determine what happens inside the (future) event horizon; we also need access to the โwhite holeโ region, inside the past event horizon.
The next section is devoted to a short commentary on the AdS/CFT correspondence and to general arguments concerning the nature of the calculations in . Section 3 then reviews the BTZ and geon spacetimes and determines the propagators on these spacetimes given by the path integral of . In section 4, these calculations are compared to the propagator in the dual CFT. We discuss the extension of the propagator calculation to the rotating BTZ black hole spacetime in an appendix.
## 2 The setting and the approximations
We use this section to set the stage for our later calculations. The most relevant elements of the AdS/CFT correspondence are briefly reviewed in section 2.1. This allows us to discuss the particular regime in which we use the correspondence and to comment on certain subtleties. We then address the stationary phase approximation and the issue of causality in section 2.2. Section 2.3 includes a few further comments on the interpretation of the propagator.
### 2.1 The correspondence in the bulk classical limit
While the AdS/CFT correspondence is conjectured to relate the full quantum theories associated with bulk string theory and the CFT, it is fair to say that this correspondence is best understood in the neighborhood of the vacuum. In that region, a useful way to describe the correspondence is in terms of the partition functions $`Z_{CFT}`$ and $`Z_{bulk}`$, which in both cases are functions of external sources that may be coupled to the CFT and to the boundary of the AdS space. Recall that the CFT lives on a spacetime which may be identified with the boundary of AdS<sub>3</sub>. The partition functions are equal and, by differentiating them, we may arrive at relations between propagators and correlators in the two theories. For example, differentiating twice yields the relation
$$๐ช_{}(b),๐ช_{}(b^{})_{}=\underset{ฯต0}{lim}ฯต^{2\mathrm{\Delta }}๐ช_B(b_ฯต)๐ช_B(b^{}{}_{ฯต}{}^{})_B$$
(1)
between the propagators in the boundary and bulk, where the bulk operators $`๐ช_B`$ are at points $`b_ฯต`$, $`b^{}_ฯต`$ in the bulk that approach the points $`b,b^{}`$ in a certain way as $`ฯต0`$. (also see ). This is a relation between the Euclidean propagators or, via analytic continuation, between the Feynman propagators in the respective vacuum states. Since we are in the vacuum state, operators on the right-hand side may be viewed as fields on AdS space.
In the work below, we again wish to consider a propagator or correlator. However, we wish to work in a regime that is rather far from the vacuum state. We consider a state in which the bulk string theory is nearly classical and contains, or is in the process of forming, a large black hole. Since the bulk string theory is nearly classical, quantum fluctuations are infinitesimal and are well approximated by linear fields. In terms of the CFT, this is the limit of large โt Hooft coupling. While this is not the classical limit of the CFT, it is a limit in which we again expect certain kinds of classical behavior (such as factorization of correlation functions with infinitesimal corrections).
Now, by acting on the vacuum with a sufficient set of local operators, we should be able to reach any state in the Hilbert space. Thus, the relation between the partition functions implies that any state $`|\mathrm{\Phi }_{}`$ in the CFT will be associated with some state $`|\mathrm{\Phi }_B`$ in the bulk. Unfortunately, it is difficult to describe this relationship in detail. Nonetheless, given any bulk state and an associated state in the CFT, it follows that correlation functions in the CFT state are given much as above by the limit of correlation functions in the bulk as the points are moved to the boundary of the spacetime.
As stated above, the regime of interest here is the limit in which the bulk spacetime is nearly classical and in which the quantum fluctuations are effectively linear. This is just the usual setting of (free) quantum field theory in curved spacetime. As a result, it is clear that a given classical geometry does not determine a unique quantum state, but rather determines an entire space of states for the linearized fluctuations. For globally static spacetimes, one can identify a preferred vacuum state, though this is not generally possible. For example, in the familiar asymptotically flat black hole spacetimes, the โnaturalโ choices of state for the linear quantum fields include the Hartle-Hawking vacuum as well as the Unruh vacuum, and more complicated choices of state are possible as well. The particular choice of quantum state may be associated with initial and/or final conditions satisfied by the linearized fluctuations.
Now used the relation (1) to link a CFT object to a bulk propagator. As a result, some particular choice of state, or perhaps several states or a class of states, for the linearized bulk quantum fields must have been made implicitly. We note that in it was explicitly assumed that the โpropagatorโ for a scalar field $`\varphi `$ in the bulk was given by the path integral expression
$$\varphi (x)\varphi (x^{})_{\mathrm{FPI}}=๐๐ซe^{i\mathrm{\Delta }L(๐ซ)},$$
(2)
where $`L(๐ซ)`$ denotes the length of the path $`๐ซ`$. The measure $`d๐ซ`$ was not specified in detail as the intention of was to use the expression (2) only in the semiclassical approximation. The subscript FPI reminds us that this is the object defined by a Feynman path integral, to distinguish it from other two-point functions that we may wish to discuss. The conventions are set here so that spacelike paths have positive imaginary length, while timelike paths have real length. The question we wish to explore is whether this is in fact the 2-point function of any bulk quantum state and, if so, just which state it represents.
Now, the two-point function alone does not uniquely determine the quantum state. However, for linear fields there is the notion of a quasi-free state (see, e.g., ), also known as a Gaussian state, in which the higher connected n-point functions vanish, and all of the structure is in fact determined by the two-point function. It is therefore natural to attempt to associate the calculations of with a quasi-free state of the linearized bulk fields. We will show below that, on the BTZ black hole spacetime, the expression (2) does in fact yield the 2-point function of the Hartle-Hawking vacuum state. Similarly, on the $`^2`$ geon spacetime, it yields the 2-point function of the so called geon-vacuum, the analogue of the state discussed in for asymptotically flat geons. These states are in fact quasi-free. In subsection 2.3 below, we will discuss to what extent we can draw the same conclusion in more general spacetimes.
### 2.2 Causality and the stationary phase approximation
After calculating the propagator (2) using a stationary phase method, it was found in that this propagator was sensitive to events happening behind a black holeโs event horizon. This raises certain issues about causality. Stated most simply, we have noted (see eq. 1) that the correlation functions in the CFT are (up to a rescaling) the boundary limits of correlation functions in the bulk. However, in the current context of bulk correlators for linear quantum field theory in curved spacetime, it is well known that the evolution is causal. An operator at any point in the spacetime can be expressed purely in terms of operators in its past light cone. How, therefore, are we to interpret the results of which suggest that correlation functions of such operators near the boundary are sensitive to the interior of the black hole?
In order to address this question, we first provide a few words on the general interpretation of the propagator (2). Let us first note that there are at least two natural ways that we might try and interpret this object. The first is as a (time-ordered) correlation function in some quantum state. For definiteness, let us use the word โstateโ in the sense of algebraic quantum field theory. This means that a โstateโ $`\rho `$ may be either a pure state or a mixed state and that we would try to interpret (2) as $`Tr\left(\rho T(\varphi (x)\varphi (y))\right)`$ for some $`\rho `$. The second natural choice is to try to interpret the propagator as the time-ordered version of a transition amplitude: $`\alpha |T(\varphi (x)\varphi (y))|\beta `$. In either case, however, the propagator would be a Greenโs function for the wave operator and thus a causal object.
Thus, we need to know whether the propagator (2) does in fact yield a Greenโs function for the wave operator $`^2`$. That this is the case may be argued as follows. Let us consider the spacetime as the configuration space of a โnon-relativistic particleโ and take $`H=^2`$ to be its Hamiltonian<sup>2</sup><sup>2</sup>2Which, in this case, is unbounded from below due to the Lorentzian signature of the spacetime.. As usual, we may write
$$\frac{i}{H}=_0^{\mathrm{}}e^{iN(Hiฯต)}๐N,$$
(3)
so that the object on the right hand side defines a Greenโs function for the wave operator. By the usual path integral skeletonization arguments, one can write this as
$$x|\frac{i}{H}|y=_0^{\mathrm{}}๐NDxDp\mathrm{exp}i_0^1[\dot{x}pN(p^2+m^2)]๐\lambda ,$$
(4)
where $`\dot{}=d/d\lambda `$. We will see in a moment that (4) is just the path integral (2) in another form. Alternatively, (4) could be taken as the definition of the path integral (2.
The path integral above contains the action for a free relativistic particle. Note, however, that while such particles are typically associated with a time reparametrization invariance, there is no such explicit invariance above. We may thus consider (4) to be a gauge-fixed path integral, using in particular the gauge $`\dot{N}=0`$:
$$x|\frac{i}{H}|y=_0^{\mathrm{}}DNDxDp\delta (\dot{N})\mathrm{exp}i_0^1[\dot{x}pN(p^2+m^2)]๐\lambda .$$
(5)
The argument below will be more transparent if we change the gauge fixing scheme to use a gauge condition that depends only on the path $`x(\lambda )`$ through position space<sup>3</sup><sup>3</sup>3A complete such gauge fixing cannot be a smooth function of the path $`x(\lambda )`$, but this need not concern us here.. Thus, we write
$$x|\frac{i}{H}|y=_0^{\mathrm{}}DNDxDp\mathrm{\Delta }(x)\mathrm{exp}i_0^1[\dot{x}pN(p^2+m^2)]๐\lambda $$
(6)
where $`\mathrm{\Delta }(x)`$ contains both the gauge fixing condition and the associated Faddeev-Popov determinant. Note that $`\mathrm{\Delta }(x)`$ will depend only on $`x(\lambda )`$.
Now, to lowest order in the WKB approximation, performing an integral over some variable is equivalent to solving the associated classical equation of motion and inserting the result back into the action. Thus, we can do the integrals over $`N`$ and $`p`$ and write the result as follows:
$$x|\frac{i}{H}|y=Dx\mathrm{\Delta }^{}(x)\mathrm{exp}iL(x(\lambda )),$$
(7)
where $`L(x(\lambda ))`$ is the length of the path $`x(\lambda )`$ with exactly the same conventions as in (2).
The factor $`\mathrm{\Delta }^{}(x)`$ denotes $`\mathrm{\Delta }(x)`$ together with the various path-dependent measure factors arising from the corrections to the WKB approximation in integrating over $`N`$ and $`p`$. Identifying $`d๐ซ=\mathrm{\Delta }^{}(x)Dx`$, we find that our Greenโs function is just the propagator (2). Note that solving the equation of motion for $`N`$ involves taking a square root. For the timelike segments of path the restriction $`N>0`$ was used to choose the appropriate branch. For the spacelike segments, the appropriate branch is determined by the details of the measure as discussed in . Note that the action is an analytic function of both $`N`$ and $`p`$ so that we expect no problems with the use of stationary phase methods here. Thus, the propagator (2) is indeed a Greenโs function for a wave operator. That (2) satisfies Dirichlet boundary conditions on the smooth part of the boundary at infinity can be seen from the arguments of .
At this point, we can now reduce our physical question about causality in the setting of to a mathematical question about solutions of the wave equation. In further stationary phase methods were used to argue that, to leading order, the propagator was in fact determined by the shortest geodesic connecting the points $`x`$ and $`y`$. The authors considered a spacetime that was pure AdS before a certain spacelike hypersurface $`\mathrm{\Sigma }`$ on which two massless point particles entered through the boundary at infinity. From this it is clear that two points sufficiently far in the past of $`\mathrm{\Sigma }`$ can only be connected by geodesics that lie in the pure AdS part of the spacetime. Thus, the geodesic approximation leads to the conclusion that, to the past of some hypersurface $`\mathrm{\Sigma }^{}`$, the propagator is just as it would be in pure AdS space<sup>4</sup><sup>4</sup>4In that case, as we will discuss below, it is known to be the time ordered 2-point function in the AdS vacuum..
Nonetheless, at sufficiently late times, it was shown in that there are points outside the black hole such that the shortest geodesic connecting them runs through the interior of the black hole. It was therefore concluded that the propagator (2) outside the black hole was sensitive to events occurring inside the black hole.
In order to eliminate certain technical worries, let us consider a family of generalizations of the spacetime constructed in . Imagine replacing the singular null particles with a distribution of null fluid of compact support. Since there is no local gravitational dynamics in 2+1 dimensions, the resulting spacetime is easily made identical to that of outside of the region occupied by the null fluid. Until the formation of the black hole singularity, the resulting spacetime is then smooth<sup>5</sup><sup>5</sup>5It is not, however, asymptotically AdS where the null fluid enters the spacetime. We shall assume that this does not cause any further complications.. If the field $`\varphi `$ for which we compute the propagator does not couple to the null fluid, then the definition of the propagator on this spacetime remains just (2). Thus, we have a complete specification of the propagator, up to issues associated with the black hole singularity<sup>6</sup><sup>6</sup>6Such issues certainly exist. For example, if we take (4) as the definition of (2), the black hole singularity will imply that $`H`$ is not essentially self-adjoint and that some particular self-adjoint extension should be chosen. Here, we simply assume that some such choice has been made..
Suppose now that we arrange things such that the two bits of null fluid actually collide inside the black hole. That is, suppose that at some event the supports of the two distributions of fluid overlap. Note that, depending on the sort of null fluid used, various outcomes are possible. Some sorts of fluid would interpenetrate readily while other sorts would bounce solidly off of each other. The outcome should affect some of the geodesics mentioned above that connect two points near infinity by passing through the interior of the black hole.
Now we see that we have a real contradiction at hand. On the one hand, we have the statement that the propagator at early times is the AdS vacuum correlator โ independent of what goes on in the black hole interior. Also, we know that the propagator satisfies the wave equation and so evolves in a causal fashion. Thus, the propagator at points outside the black hole can be expressed in terms of initial data on an early hypersurface in a manner that is independent of what goes on in the black hole interior. Thus, the propagator outside the black hole cannot in fact depend on events inside the black hole. This is in direct contradiction to the conclusion of the previous paragraph.
The resolution seems to be that the geodesic โapproximationโ is not in fact a valid approximation<sup>7</sup><sup>7</sup>7It is also a logical possibility that the approximation is valid, but simply unstable in a manner that causes higher order effects at early times to evolve into lower order effects at late times.. In retrospect, it seems quite likely that this approximation fails for such a spacetime. Note that to arrive at the geodesic approximation, one would use a stationary phase argument to solve the classical equations of motion corresponding to the action $`m\sqrt{\dot{x}^2}`$. While the stationary point (the spacelike geodesic) does indeed lie on the original contour of integration (real values of $`x`$), this contour is not a steepest descent contour through the stationary point. In particular, in a Lorentzian signature spacetime, a spacelike geodesic is not a path of minimal length. As a result, if one wishes to argue that the stationary point dominates, one must first analytically continue the action to complex values of the coordinates and attempt to deform the original contour to the contour of steepest descent.
Now, the action involves the metric $`g_{ab}(x)`$. To avoid the issue of the singularities, let us consider the smoothed spacetimes with null fluid sources. Since the fluid density vanishes in an open region, but not in the entire spacetime, it is clear that such spacetimes are not analytic and that continuation is problematic. Thus, it is not at all clear that steepest descent methods should succeed in this case, and we are happy to associate their failure with nonanalyticities of the spacetime.
While this seems to settle the issue nicely, we should mention for completeness that, if one excises the region of non-zero fluid density from the spacetime, the resulting spacetime does have a real analytic atlas and can be continued. Presumably, excising the region occupied by the fluid prevents one from deforming the contours as one would like.
### 2.3 Interpreting the Propagator
Having ruled out the use of the geodesic approximation in general, what are we to conclude about the full propagator (2)? In principle, picking any two points $`x`$ and $`y`$ in the spacetime, the path integral includes contributions from paths connecting them that explore arbitrarily far into the future. As a result, even in the spacetime studied in , it is far from clear that the propagator at early times is independent of events in the interior of the black hole. It seems likely that the propagator does not correspond to a fixed initial condition, but instead to some mixture of initial and final conditions. In this case, the propagator between points near the boundary at late times may depend on events inside the black hole as well. That is, it may be possible to choose a state or states in such a way that the two-point function reproduces the important qualitative features found in the calculation in . Such a state (or states) will involve a mixture of initial and final conditions, reflecting the fact that the form of the two-point function in depends on the formation of the black hole in the future.
Let us return to the two natural interpretations of the propagator mentioned above: as a time ordered expectation value in some quasi-free state, and as a time ordered transition amplitude between two states. We note that either is compatible with the above observations. In the case of the expectation value, it may simply be the case that the quantum state itself is one that is naturally defined by a combination of retarded and advanced boundary conditions, and so is free to depend on events in the interior of the black hole. We note that the Hartle-Hawking state for an asymptotically flat black hole is an example of such a state that is naturally associated with boundary conditions in both past and future, while the Unruh state is associated only with boundary conditions in the past. In the case of the transition amplitude, both states may involve such โmixedโ boundary conditions, or perhaps one is defined by retarded boundary conditions and one by advanced boundary conditions.
In spacetimes that are asymptotically flat at both timelike and spacelike infinity, the propagator (2) can be shown to define a transition amplitude . On the other hand, the work of Wald effectively shows that (2) defines an expectation value for globally static spacetimes (without horizons). It is also known to give the expectation value of time-ordered fields in the Hartle-Hawking state on the Kruskal spacetime, though the status of this question on a general black hole spacetime is not yet understood . We will see that an expectation value is once again obtained on the spinless BTZ spacetime and the associated $`^2`$ geon.
## 3 The geodesics in AdS<sub>3</sub> and quotient spacetimes
We have argued in section 2 that stationary phase methods do not in general yield a valid approximation to the FPI propagator (2). Nevertheless, one may ask if there are cases in which it does provide a valid approximation and, if so, whether geodesics passing behind the horizon play any important role. We shall see in this section and the next that the answer to both of these questions is in the affirmative.
In the present section, we consider the lengths of spacelike geodesics in the AdS<sub>3</sub>, spinless BTZ, and $`^2`$ geon spacetimes. As these spacetimes are real Lorentzian sections of holomorphic complex manifolds, one may expect the geodesic approximation to succeed in these cases. Indeed, it is known to succeed in yielding the vacuum correlator on AdS<sub>3</sub>. In the following section, we consider the propagators obtained through this approximation, and compare to what we know about the field theory. This will allow us to explicitly check the agreement with certain CFT calculations and to trace the role of geodesics passing through the interior of the black hole. The final agreement provides additional confirmation of the accuracy of the geodesic approximation in these cases.
In fact, these calculations are not truly independent. Since the spinless BTZ and $`^2`$ geon spacetimes are quotients of AdS<sub>3</sub>, a method of images argument together with analytic continuations and the uniqueness of the Euclidean Greenโs functions shows that the success of the geodesic approximation to (2) in reproducing the vacuum correlator on AdS<sub>3</sub> implies that it must also approximate the Hartle-Hawking correlation function for the spinless BTZ hole and the related geon correlation function (see ) on the $`^2`$ geon. Thus, in these cases the FPI propagator gives the expectation value of a time-ordered product of fields in a quasi-free state.
### 3.1 Geodesics of AdS<sub>3</sub>
The AdS<sub>3</sub> spacetime can be constructed as the hyperboloid
$$(T^1)^2+(T^2)^2(X^1)^2(X^2)^2=1$$
(8)
in a flat embedding space with metric
$$ds^2=(dT^1)^2(dT^2)^2+(dX^1)^2+(dX^2)^2.$$
(9)
Here, we are choosing units so that the AdS length scale $`\mathrm{}`$ (related to the cosmological constant) is one. A set of intrinsic coordinates on AdS<sub>3</sub> is given in terms of these embedding coordinates by
$$T^1=\mathrm{cosh}\chi \mathrm{cos}\tau ,T^2=\mathrm{cosh}\chi \mathrm{sin}\tau ,X^1=\mathrm{sinh}\chi \mathrm{sin}\phi ,X^2=\mathrm{sinh}\chi \mathrm{cos}\phi ,$$
(10)
where $`\phi `$ has period $`2\pi `$, and $`0\chi \mathrm{}`$. For the hyperboloid, $`\tau `$ is also periodic with period $`2\pi `$, but we pass to the covering space, and take $`\tau `$ to run between $`\pm \mathrm{}`$. In terms of these coordinates, the metric is
$$ds^2=d\chi ^2+\mathrm{sinh}^2\chi d\phi ^2\mathrm{cosh}^2\chi d\tau ^2=\left(\frac{2}{1\rho ^2}\right)^2(d\rho ^2+\rho ^2d\phi ^2)\left(\frac{1+\rho ^2}{1\rho ^2}\right)^2d\tau ^2.$$
(11)
In the second equality, we have defined a new radial coordinate $`\rho =\mathrm{tanh}(\chi /2)`$, so $`0\rho 1`$. Fixed $`\tau `$ surfaces have the Poincarรฉ disc geometry, and the dual CFT is defined on a cylinder isomorphic to the $`\rho =1`$ boundary.
We will need the length of the unique geodesic traveling between $`(\tau ,\chi _\mathrm{m},\pm \phi _\mathrm{m})`$. Now, since the metric at fixed $`\tau `$ is that of the Poincarรฉ disc, equal-time geodesics of (11) are circle segments obeying the equation
$$\mathrm{tanh}\chi \mathrm{cos}(\phi \alpha )=\mathrm{cos}(\beta ),$$
(12)
where the geodesic reaches the $`\chi =\mathrm{}`$ boundary at $`\phi =\alpha \pm \beta `$. Setting $`\alpha =0`$, the unique geodesic between the boundary points $`(\tau ,\pm \beta )`$ intersects $`\chi =\chi _m`$ at $`\phi _m^\pm `$ which are fixed by
$$\mathrm{tanh}\chi _\mathrm{m}\mathrm{cos}\phi _\mathrm{m}^\pm =\mathrm{cos}(\pm \beta ),$$
(13)
which implies that
$$\phi _m^{}=\phi _m^+\phi _m.$$
(14)
Integrating (11) yields the length of the geodesic connecting $`(\tau ,\chi _\mathrm{m},\pm \phi _\mathrm{m})`$:
$$L(\phi _\mathrm{m},\phi _\mathrm{m})=2\mathrm{ln}\left[\mathrm{sinh}\chi _\mathrm{m}\mathrm{sin}\phi _\mathrm{m}+(\mathrm{sinh}^2\chi _\mathrm{m}\mathrm{sin}^2\phi _\mathrm{m}+1)^{1/2}\right].$$
(15)
### 3.2 Spacelike geodesics on the spinless BTZ hole
The spinless BTZ hole is obtained by taking the quotient of the region $`T^1>|X^1|`$ of AdS<sub>3</sub> by the isometry $`\mathrm{exp}(2\pi r_+\xi )`$, where $`\xi `$ is the Killing vector
$$\xi =X^1\frac{}{T^1}+T^1\frac{}{X^1}.$$
(16)
To express this geometry in the Schwarzschild-like coordinates of the original papers , we introduce on the region $`X^2>|T^2|`$, $`T^1>0`$ of AdS<sub>3</sub> the coordinates $`(t,r,\varphi )`$ by
$`T^1`$ $`=`$ $`{\displaystyle \frac{r}{r_+}}\mathrm{cosh}(r_+\varphi ),`$
$`X^1`$ $`=`$ $`{\displaystyle \frac{r}{r_+}}\mathrm{sinh}(r_+\varphi ),`$
$`T^2`$ $`=`$ $`\left({\displaystyle \frac{r^2}{r_+^2}}1\right)^{1/2}\mathrm{sinh}(r_+t),`$
$`X^2`$ $`=`$ $`\left({\displaystyle \frac{r^2}{r_+^2}}1\right)^{1/2}\mathrm{cosh}(r_+t).`$ (17)
$`t`$ and $`\varphi `$ take all real values, $`r>r_+`$, and the metric takes the form
$$ds^2=N^2dt^2+r^2d\varphi ^2+\frac{1}{N^2}dr^2;N^2=r^28GM,$$
(18)
where $`M=r_+^2/(8G)`$. The identification by $`\mathrm{exp}(2\pi r_+\xi )`$ amounts to $`(t,r,\varphi )(t,r,\varphi +2\pi )`$, and with this identification the coordinates $`(t,r,\varphi )`$ cover one exterior region of the BTZ hole.
We are interested in geodesics between two points, $`x_1`$ and $`x_2`$, in the exterior region of the hole. We take the value of $`r`$ at both points to be the same. To parametrize the locations of the points, let $`y_1`$ and $`y_2`$ be two points in AdS<sub>3</sub>, respectively at $`(t_1,r,\varphi _1)`$ and $`(t_2,r,\varphi _2)`$, and let $`x_1`$ (respectively $`x_2`$) be the equivalence class of $`y_1`$ ($`y_2`$). We write $`\mathrm{\Delta }\varphi =\varphi _2\varphi _1`$ and $`\mathrm{\Delta }t=t_2t_1`$, and we assume that $`|\mathrm{\Delta }\varphi +2\pi n|>|\mathrm{\Delta }t|`$ for all integers $`n`$. For fixed $`t_1`$, $`t_2`$, $`\varphi _1`$, and $`\varphi _2`$, it is then straightforward to show that for sufficiently large $`r`$ there are countably many spacelike geodesics connecting $`x_1`$ and $`x_2`$ in the BTZ hole.
To calculate the lengths of these geodesics, we exploit the symmetries to argue that the geodesic distance between $`y_1`$ and $`y_2`$ in AdS<sub>3</sub> is a function only of the chordal distance $`D`$ in the embedding space,
$`D`$ $`=`$ $`(\mathrm{\Delta }T^1)^2(\mathrm{\Delta }T^2)^2+(\mathrm{\Delta }X^1)^2+(\mathrm{\Delta }X^2)^2`$ (19)
$`=`$ $`{\displaystyle \frac{4r^2}{r_+^2}}\mathrm{sinh}^2\left({\displaystyle \frac{r_+\mathrm{\Delta }\varphi }{2}}\right)4\left({\displaystyle \frac{r^2}{r_+^2}}1\right)\mathrm{sinh}^2\left({\displaystyle \frac{r_+\mathrm{\Delta }t}{2}}\right).`$
By considering a simple example of a spacelike geodesic, we can show that the relation between chordal distance and proper length $`L`$ is
$$\mathrm{sinh}^2(L/2)=\frac{D}{4}.$$
(20)
It then follows from the quotient construction that the lengths, $`L_n(x_1,x_2)`$, of the geodesics connecting $`x_1`$ and $`x_2`$ in the BTZ hole have the large $`r`$ expansion
$$\mathrm{exp}\left[L_n(x_1,x_2)\right]=\frac{2r^2}{r_+^2}\left\{\mathrm{cosh}\left[r_+(\mathrm{\Delta }\varphi +2\pi n)\right]\mathrm{cosh}(r_+\mathrm{\Delta }t)\right\}+O\left(1\right),$$
(21)
where $`n`$.
### 3.3 Spacelike geodesics on the $`^2`$ geon
Recall that the $`^2`$ geon is obtained by taking the quotient of the region $`T^1>|X^1|`$ of AdS<sub>3</sub> by the isometry that is the composition of $`J_1:\mathrm{exp}(\pi r_+\xi )`$ and the involution $`J_2:(T^1,T^2,X^1,X^2)(T^1,T^2,X^1,X^2)`$. The resulting spacetime is not orientable, but one can construct a related orientable spacetime from the product of the BTZ spacetime with $`T^4`$. If the moduli of the $`T^4`$ are chosen so that there is an orientation-reversing involution $`J_4`$ of the torus, then one obtains an orientable spacetime by taking the quotient with respect to $`J_1J_2J_4`$.
Now, let $`y_1`$ and $`y_2`$ be points on AdS<sub>3</sub> as above, respectively at $`(t_1,r,\varphi _1)`$ and $`(t_2,r,\varphi _2)`$, and suppose that $`|\mathrm{\Delta }\varphi +2\pi n|>|\mathrm{\Delta }t|`$ for all integers $`n`$. Let $`x_1`$ and $`x_2`$ be two points in the exterior region of the geon, such that $`x_1`$ (respectively $`x_2`$) is the equivalence class of $`y_1`$ ($`y_2`$). For sufficiently large $`r`$, one class of spacelike geodesics connecting $`x_1`$ and $`x_2`$ is then obtained precisely as for the BTZ hole, with the result (21) for their lengths. The second class of geodesics arises from the AdS<sub>3</sub> geodesics connecting $`y_1`$ to the points $`\stackrel{~}{y}_{2;n}`$, located at
$`T^1=(r/r_+)\mathrm{cosh}[r_+(\varphi _2+\pi +2\pi n)],`$
$`X^1=(r/r_+)\mathrm{sinh}[r_+(\varphi _2+\pi +2\pi n)],`$
$`T^2=\sqrt{(r/r_+)^21}\mathrm{sinh}(r_+t_2),`$
$`X^2=\sqrt{(r/r_+)^21}\mathrm{cosh}(r_+t_2),`$ (22)
where $`n`$. As
$`D(y_1,\stackrel{~}{y}_{2;n})`$ $`=`$ $`2(r/r_+)^2\left\{\mathrm{cosh}\left[r_+(\varphi _2\varphi _1+\pi +2\pi n)\right]1\right\}`$ (23)
$`+2\left[(r/r_+)^21\right]\left\{\mathrm{cosh}\left[r_+(t_2+t_1)\right]+1\right\},`$
the lengths $`\stackrel{~}{L}_n(x_1,x_2)`$ of these geodesics have the large $`r`$ expansion
$$\mathrm{exp}\left[\stackrel{~}{L}_n(x_1,x_2)\right]=\frac{2r^2}{r_+^2}\left\{\mathrm{cosh}\left[r_+(\varphi _2\varphi _1+\pi +2\pi n)\right]+\mathrm{cosh}\left[r_+(t_2+t_1)\right]\right\}+O\left(1\right).$$
(24)
It is precisely this class of geodesics that pass through the black hole interior. We note that all such geodesics are longer than the shortest geodesic connecting $`x_1`$ and $`x_2`$ through the exterior region. Thus, at first sight one might think that geodesics passing through the interior cannot be relevant to leading order. Nonetheless, we shall see in section 4.3 that they do provide the leading contribution to the two-particle correlations in the geon vacuum, and that (24) reproduces expectations based on the dual CFT.
## 4 Matching to the CFT
It turns out that, due to difficulties in performing the various mode sums, there are few exact results for the bulk correlators in the spinless BTZ Hartle-Hawking state and in the geon vacuum. We will therefore proceed by comparing the limiting behaviors of (30) and (24) with expectations based on toy models of the dual CFT. We shall see that the agreement is surprisingly good. This supports both the accuracy of the bulk geodesic approximation in these cases and the ability of the toy models to capture much of the physics of the CFT. We first review the calculation showing that the geodesic approximation in AdS<sub>3</sub> reproduces the vacuum propagator, and then show that the asymptotic behavior of (24) reproduces the expected two-particle correlations in BTZ and the geon.
### 4.1 The propagator in AdS<sub>3</sub>
We will now review the calculation of the equal time correlation functions in the dual field theory for the AdS<sub>3</sub> geometry using the (bulk) WKB approximation. A scalar field of mass $`m`$ in a spacetime which is asymptotically AdS<sub>3</sub> is dual to an operator $`๐ช`$ of conformal dimension $`\mathrm{\Delta }=1+\sqrt{1+m^2}`$. The fiducial metric for the CFT on the cylinder is related to the induced metric obtained from (11) by a diverging Weyl factor. To relate operators to expectation values, we need to regulate this behavior by cutting off the spacetime at a boundary defined by
$$\rho _\mathrm{m}(\tau ,\phi )=1ฯต(\tau ,\phi ),ฯต(\tau ,\phi )=ฯต(\tau ,\phi ),$$
(25)
where $`ฯต`$ is some smooth function of the boundary coordinates. The symmetry of $`ฯต`$ under $`\phi \phi `$ is chosen for simplicity. For the calculations relating to the BTZ black hole and geon, we will take the cutoff surface to be at constant $`r`$ in the BTZ coordinates. According to , the Feynman propagator for $`๐ช`$ in the dual CFT is obtained from the spacetime propagator between the corresponding points on the cutoff boundary at $`\rho _\mathrm{m}`$ (also see ),
$$G_{}((\tau ,\phi ),(\tau ^{},\phi ^{}))=ฯต^{2\mathrm{\Delta }}G_B(๐,๐^{}),$$
(26)
where $`๐=(\tau ,\phi ,\rho _\mathrm{m}(\tau ,\phi ))`$. We will only need the propagator when $`\tau =\tau ^{}`$. For $`๐,๐^{}`$ causally unrelated, the Greenโs function $`G_B(๐,๐^{})`$ in the leading order semi-classical approximation is given by a sum over geodesics:
$$G(๐,๐^{})=\underset{g}{}e^{\mathrm{\Delta }L_g(๐,๐^{})}.$$
(27)
Here $`L_g`$ is the (real) geodesic length between the boundary points and only spacelike geodesics contribute since $`\tau =\tau ^{}`$.
By rotational invariance, it is sufficient to perform the calculation for $`\phi =\phi ^{}`$. For the particular case $`1ฯต(\tau ,\phi )=\mathrm{tanh}(\chi _m/2)=const`$, the length of the geodesic connecting $`๐`$ and $`๐^{}`$ is given by (15). In fact, the symmetry $`ฯต(\tau ,\phi )=ฯต(\tau ,\phi )`$ guarantees that a corresponding result holds for any such symmetric choice of $`ฯต`$. So, to leading order in $`ฯต`$, the geodesic length between the points $`๐,๐^{}`$ is
$$L(๐,๐^{})=2\mathrm{ln}\left(\frac{2\mathrm{sin}\phi }{ฯต}\right).$$
(28)
The bulk propagator is thus
$$G(๐,๐^{})=\left(\frac{2\mathrm{sin}\phi }{ฯต}\right)^{2\mathrm{\Delta }}$$
(29)
in the $`ฯต0`$ limit, where the boundary metric is $`ds^2=(1/ฯต(\tau ,\phi )^2)(d\tau ^2+d\phi ^2)`$. This correctly reproduces the CFT two-point correlator of for $`\mathrm{\Delta }\tau =0`$ and $`\mathrm{\Delta }\phi =2\phi `$, since the CFT is defined on the Weyl rescaled cylinder with metric $`ds^2=d\tau ^2+d\phi ^2`$.
### 4.2 The propagator in BTZ
We now apply the bulk geodesic approximation method of to the Greenโs function on the boundary of the spinless BTZ hole, using the geodesic length (21). The geodesic approximation to the path integral (2) reads
$`\varphi (x_1)\varphi (x_2)_{\mathrm{FPI}}`$ $`={\displaystyle ๐๐ซe^{i\mathrm{\Delta }L(๐ซ)}}{\displaystyle \underset{n}{}}\mathrm{exp}\left[\mathrm{\Delta }L_n(x_1,x_2)\right]`$ (30)
$`=\left({\displaystyle \frac{r_+^2}{2r^2}}\right)^\mathrm{\Delta }{\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{\left\{\mathrm{cosh}\left[r_+(\mathrm{\Delta }\varphi +2\pi n)\right]\mathrm{cosh}(r_+\mathrm{\Delta }t)\right\}^\mathrm{\Delta }}}`$
$`+O\left(\left({\displaystyle \frac{r_+^2}{r^2}}\right)^{\mathrm{\Delta }+1}\right).`$
This is the bulk propagator; to relate it to the boundary propagator, we observe that at large $`r`$, the boundary metric in the BTZ spacetime is $`ds^2=r^2(dt^2+d\varphi ^2)`$. Since we want the CFT to live on the same cylinder as above, the boundary propagator is given by $`G_{}=r^{2\mathrm{\Delta }}G_B`$ in this case. The rescaling thus precisely cancels the $`r^{2\mathrm{\Delta }}`$ in the prefactor.
We note that this Greenโs function is manifestly periodic in the global time $`t`$ in the imaginary direction, and the period $`2\pi /r_+`$ is the inverse of the spinless BTZ temperature. It is thus plausibly identified with the propagator in the analogue of the Hartle-Hawking state for this black hole. We also note that the geodesics used in this calculation lie entirely outside the black hole. This calculation successfully reproduces the result given in for the spinless BTZ black hole. A similar agreement is obtained for the rotating BTZ black hole in appendix A.
### 4.3 The propagator in the single-exterior black hole
We now proceed to address the dual CFT propagator associated with the bulk FPI propagator $`x_1x_2_{\mathrm{geon}}`$ on the $`^2`$ geon. Now, the corresponding path integral can be written as a sum of two contributions:
$$x_1x_2_{\mathrm{geon}}=x_1x_2_{\mathrm{BTZ}}+x_1J(x_2)_{\mathrm{BTZ}}$$
(31)
where $`x_1x_2_{\mathrm{BTZ}}`$ represents the bulk FPI propagator on the spinless BTZ hole. Here, we take $`x_1`$ and $`x_2`$ to lie outside the geon horizon so that we may naturally associate them with two points in an asymptotic region of the BTZ black hole. The first term ($`x_1x_2_{\mathrm{BTZ}}`$) was calculated in the geodesic approximation in section 3.2 while the second term ($`x_1J(x_2)_{\mathrm{BTZ}}`$) is given in the geodesic approximation by (24). The geodesics that contribute to this second term are longer than the shortest geodesic contributing to $`x_1x_2_{\mathrm{BTZ}}`$, so that one might at first think that $`x_1J(x_2)_{\mathrm{BTZ}}`$ can be neglected. However, let us now Fourier transform this result in order to compute the two-particle correlations in the geon state. Since the energies of the two particles cannot add to zero, the time translation invariance of the BTZ hole is enough to guarantee that the contributions from the first term (with both points in the same asymptotic region) vanish. However, the contribution of the second term need not vanish, corresponding to the fact that the geon does not itself have a time translation invariance. Thus, we see that the two-particle correlations in the geon state can be directly tied to the second term above, which results only from geodesics that pass through the interior of the BTZ black hole. In terms of the geon spacetime, the result is again that only geodesics passing behind the horizon can account for the two-particle correlations.
Thus, we might try to match these correlations to those computed in for a toy model of the CFT state $`|\mathrm{geon}`$ dual to the $`^2`$ geon, presumably with linearized quantum fluctuations in the geon vacuum state. The toy model replaced the CFT by a free scalar field and found, in the case of a nontwisted field, the correlations
$$\mathrm{geon}|d_{n,ฯต}d_{n^{},ฯต^{}}|\mathrm{geon}=\frac{(1)^n\delta _{n,n^{}}\delta _{ฯต,ฯต^{}}}{2\mathrm{sinh}(\pi n/r_+)},$$
(32)
where $`d_{n,ฯต}`$ is the annihilation operator for the mode with frequency quantum number $`n`$, $`n=1,2,\mathrm{}`$, and the index $`ฯต`$ takes the value $`1`$ for right-movers and $`1`$ for left-movers (see (33) below). Here and below we ignore issues involving the zero mode ($`n=0`$). As a consequence of rotational invariance, the correlations are between a right-mover and a left-mover with the same frequency. We note that this nontwisted free scalar field has conformal weight $`\mathrm{\Delta }=0`$.
We now show that this result can be obtained from the geodesic approximation (24) to the bulk Greenโs function. It is clear, however, that due to the simplified nature of the toy model, one should not expect to be able to Fourier transform the asymptotic values of the propagator and obtain (32) directly. In particular, the bulk propagator will not be built from only the discrete mode spectrum of (32). In the full interacting CFT, the correlator will similarly not be periodic in time, and so will not be a simple combination of these discrete modes. However, we may attempt to extract information analogous to (32) by modifying the Fourier transform of (24) to take into account the fact that the bulk propagator is not periodic in time. We shall see that the agreement is impressive.
In the toy (free) CFT, the oscillator modes that correspond to the annihilation operators $`d_{n,ฯต}`$ in (32) are
$$u_{n,ฯต}=\frac{1}{\sqrt{4\pi n}}e^{in(tฯต\varphi )}$$
(33)
where $`n=1,2,\mathrm{}`$ and $`ฯต=\pm 1`$. If the Greenโs function $`G(t_1,\varphi _1;t_2,\varphi _2)`$ had the periodicity of the oscillator modes, we would thus have
$`d_{n,ฯต}d_{n^{},ฯต^{}}`$ $`=`$ $`{\displaystyle \frac{\sqrt{nn^{}}}{4\pi ^3}}{\displaystyle _0^{2\pi }}๐t_1{\displaystyle _0^{2\pi }}๐t_2{\displaystyle _0^{2\pi }}๐\varphi _1{\displaystyle _0^{2\pi }}๐\varphi _2`$ (34)
$`\times \mathrm{exp}[i(nt_1+n^{}t_2nฯต\varphi _1n^{}ฯต^{}\varphi _2)]G(t_1,\varphi _1;t_2,\varphi _2).`$
We shall modify (34) to take into account the lack of periodicity shortly.
As discussed above, the part of $`G(t_1,\varphi _1;t_2,\varphi _2)`$ coming from the geodesics that do not pass through the geon does not contribute to $`d_{n,ฯต}d_{n^{},ฯต^{}}`$. The part of $`G(t_1,\varphi _1;t_2,\varphi _2)`$ coming from the geodesics that do pass through the geon is, from (24),<sup>8</sup><sup>8</sup>8The asymptotic metric in the geon spacetime is the same as in the BTZ spacetime, so the rescaling relating the bulk and boundary propagators is the same as in the previous subsection.
$$\left(\frac{r_+^2}{2}\right)^\mathrm{\Delta }\underset{k=\mathrm{}}{\overset{\mathrm{}}{}}\frac{1}{\left\{\mathrm{cosh}\left[r_+(\varphi _2\varphi _1+\pi +2\pi k)\right]+\mathrm{cosh}\left[r_+(t_2+t_1)\right]\right\}^\mathrm{\Delta }}.$$
(35)
As (35) depends on $`\varphi _1`$ and $`\varphi _2`$ only through the combination $`\varphi _2\varphi _1`$, integrating over $`\varphi _2+\varphi _1`$ in (34) is immediate. Next, we observe that each term in (35) depends on $`\varphi _2\varphi _1`$ and $`k`$ only through the combination $`\varphi _2\varphi _1+2\pi k`$. Integrating $`\varphi _2\varphi _1`$ from zero to $`2\pi `$ and summing over $`k`$ is thus equivalent to integrating any one term in (35) in $`\varphi _2\varphi _1`$ from negative infinity to positive infinity. Writing the integration in terms of the variable $`y:=ฯต(\varphi _2\varphi _1+\pi )`$, we obtain
$`d_{n,ฯต}d_{n^{},ฯต^{}}={\displaystyle \frac{\sqrt{nn^{}}}{2\pi ^2}}\delta _{nฯต,n^{}ฯต^{}}(1)^n\left({\displaystyle \frac{r_+^2}{2}}\right)^\mathrm{\Delta }{\displaystyle _0^{2\pi }}๐t_1{\displaystyle _0^{2\pi }}๐t_2`$
$`\times {\displaystyle _{\mathrm{}}^{\mathrm{}}}dy{\displaystyle \frac{\mathrm{exp}[i(nt_1+n^{}t_2+ny)]}{\left\{\mathrm{cosh}\left(r_+y\right)+\mathrm{cosh}\left[r_+(t_2+t_1)\right]\right\}^\mathrm{\Delta }}}.`$ (36)
We must now face the fact that the integrand in (4.3) is not periodic in $`t_1`$ and $`t_2`$. We reinterpret (4.3) by hand so that $`t_2+t_1:=\alpha `$ is integrated over $``$ but $`t_2t_1`$ over $`4\pi `$. The integral over $`t_2t_1`$, combined with the Jacobian that arises from the change of variables, yields then $`2\pi \delta _{n,n^{}}`$. The factor $`\delta _{nฯต,n^{}ฯต^{}}`$ can thus be replaced by $`\delta _{ฯต,ฯต^{}}`$, and we find
$`d_{n,ฯต}d_{n^{},ฯต^{}}`$ $`=`$ $`{\displaystyle \frac{n}{\pi }}\delta _{n,n^{}}\delta _{ฯต,ฯต^{}}(1)^n\left({\displaystyle \frac{r_+^2}{2}}\right)^\mathrm{\Delta }`$ (37)
$`\times {\displaystyle _{\mathrm{}}^{\mathrm{}}}d\alpha {\displaystyle _{\mathrm{}}^{\mathrm{}}}dy{\displaystyle \frac{\mathrm{exp}[in(\alpha +y)]}{\left[\mathrm{cosh}\left(r_+\alpha \right)+\mathrm{cosh}\left(r_+y\right)\right]^\mathrm{\Delta }}}.`$
Changing variables to $`u=\alpha y`$, $`v=\alpha +y`$, gives finally
$`d_{n,ฯต}d_{n^{},ฯต^{}}={\displaystyle \frac{n}{2\pi }}\delta _{n,n^{}}\delta _{ฯต,ฯต^{}}(1)^n\left({\displaystyle \frac{r_+}{2}}\right)^{2\mathrm{\Delta }}{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \frac{du}{\left[\mathrm{cosh}\left(r_+u/2\right)\right]^\mathrm{\Delta }}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \frac{dv\mathrm{exp}(inv)}{\left[\mathrm{cosh}\left(r_+v/2\right)\right]^\mathrm{\Delta }}}`$
$`={\displaystyle \frac{n}{2\pi }}\delta _{n,n^{}}\delta _{ฯต,ฯต^{}}(1)^n(r_+)^{2(\mathrm{\Delta }1)}\left({\displaystyle \frac{\mathrm{\Gamma }(\mathrm{\Delta }/2)}{\mathrm{\Gamma }(\mathrm{\Delta })}}\right)^2\mathrm{\Gamma }\left({\displaystyle \frac{\mathrm{\Delta }}{2}}+{\displaystyle \frac{in}{r_+}}\right)\mathrm{\Gamma }\left({\displaystyle \frac{\mathrm{\Delta }}{2}}{\displaystyle \frac{in}{r_+}}\right).`$ (38)
In the limit $`\mathrm{\Delta }0_+`$, (38) reduces to
$`d_{n,ฯต}d_{n^{},ฯต^{}}`$ $`=`$ $`{\displaystyle \frac{2n}{\pi }}\delta _{n,n^{}}\delta _{ฯต,ฯต^{}}(1)^n(r_+)^2\mathrm{\Gamma }\left({\displaystyle \frac{in}{r_+}}\right)\mathrm{\Gamma }\left({\displaystyle \frac{in}{r_+}}\right)`$ (39)
$`=`$ $`{\displaystyle \frac{4}{r_+}}\times {\displaystyle \frac{(1)^n\delta _{n,n^{}}\delta _{ฯต,ฯต^{}}}{2\mathrm{sinh}(\pi n/r_+)}},`$
which agrees with (32) up to the factor $`4/r_+`$. This factor may be a consequence of our having neglected any pre-exponential factors in the bulk Greenโs function, or from our by-hand reinterpretation of the $`dt_1dt_2`$ integrals in (34).
This result verifies the importance of geodesics passing behind the horizon in obtaining the proper 2-particle correlations, and shows that the toy free CFT does indeed match well with the bulk spacetime results. As discussed earlier, it is only in special spacetimes which are appropriately analytic than we can expect the geodesic approximation to hold. As a result, the fact that our calculation relies on geodesics passing behind the horizon of the black hole is consistent with the causal nature of the FPI propagator and with the ideas of that one must look beyond simple products of local operators in the CFT to encode useful information about the interior of a black hole.
Acknowledgements
We have enjoyed discussions with Vijay Balasubramanian, Sumit Das, Ted Jacobson, Bernard Kay, Per Kraus, Finn Larsen, Emil Martinec, Rafael Sorkin, Lenny Susskind, Sandip Trivedi, Bob Wald, and others at the Val Morin workshop on Black Holes (June, 1999) and the ICTP conference on black hole physics (July, 1999). The work of S.F.R. was partially supported by NSF grant PHY95-07065. D.M. is an Alfred P. Sloan Research Fellow and was supported in part by funds from Syracuse University, and from NSF grant PHY-9722362.
## Appendix A The propagator for the rotating BTZ hole
In this appendix we generalize the treatment of sections 3.2, 4.2 to show that the bulk geodesic approximation method of reproduces the Greenโs function in the Poincarรฉ vacuum (see and the references therein) on a single boundary component of the rotating nonextremal BTZ hole.
The generalization of equations (17) to the rotating case is the rotating exterior BTZ coordinate transformation
$`T^1=\sqrt{\alpha }\mathrm{cosh}(r_+\varphi r_{}t),`$
$`X^1=\sqrt{\alpha }\mathrm{sinh}(r_+\varphi r_{}t),`$
$`T^2=\sqrt{\alpha 1}\mathrm{sinh}(r_+tr_{}\varphi ),`$
$`X^2=\sqrt{\alpha 1}\mathrm{cosh}(r_+tr_{}\varphi ),`$ (40)
with
$$\alpha =\frac{r^2r_{}^2}{r_+^2r_{}^2},$$
(41)
where $`r>r_+`$, $`\mathrm{}<t<\mathrm{}`$, and $`\mathrm{}<\varphi <\mathrm{}`$, and the parameters $`r_\pm `$ satisfy $`0r_{}<r_+`$. For $`r_{}=0`$, this transformation reduces to the spinless transformation (17).
Introducing the points $`y_1`$ and $`y_2`$ in AdS<sub>3</sub> as in section 3.2, respectively at $`(t_1,r,\varphi _1)`$ and $`(t_2,r,\varphi _2)`$, we find
$`D(y_1,y_2)`$ $`=`$ $`2\alpha \left[\mathrm{cosh}(r_+\mathrm{\Delta }\varphi r_{}\mathrm{\Delta }t)1\right]`$ (42)
$`2(\alpha 1)\left[\mathrm{cosh}(r_+\mathrm{\Delta }tr_{}\mathrm{\Delta }\varphi )1\right].`$
When $`t_1`$, $`t_2`$, $`\varphi _1`$, and $`\varphi _2`$ are fixed, and such that $`|\mathrm{\Delta }\varphi |>|\mathrm{\Delta }t|`$, equation (42) shows that $`D(y_1,y_2)>0`$ for sufficiently large $`r`$. $`y_1`$ and $`y_2`$ can then be joined by a spacelike geodesic, and the length $`L(y_1,y_2)`$ of this geodesic has the large $`r`$ expansion
$$\mathrm{exp}\left[L(y_1,y_2)\right]=\frac{2r^2}{r_+^2}\left[\mathrm{cosh}(r_+\mathrm{\Delta }\varphi r_{}\mathrm{\Delta }t)\mathrm{cosh}(r_+\mathrm{\Delta }tr_{}\mathrm{\Delta }\varphi )\right]+O\left(1\right).$$
(43)
Now, in the region of AdS<sub>3</sub> covered by the exterior BTZ coordinates, the rotating BTZ quotient construction amounts to the identification $`(t,r,\varphi )(t,r,\varphi +2\pi )`$. Let again $`x_1`$ (respectively $`x_2`$) be the equivalence class of the point $`y_1`$ ($`y_2`$). Assuming $`|\mathrm{\Delta }\varphi +2\pi n|>|\mathrm{\Delta }t|`$ for all integers $`n`$, and proceeding as in section 4.2, we find that the geodesic approximation to the path integral (2) reads
$`\varphi (x_1)\varphi (x_2)_{\mathrm{FPI}}={\displaystyle ๐๐ซe^{i\mathrm{\Delta }L(๐ซ)}}{\displaystyle \underset{n}{}}\mathrm{exp}\left[\mathrm{\Delta }L_n(x_1,x_2)\right]`$
$`=\left({\displaystyle \frac{r_+^2}{2r^2}}\right)^\mathrm{\Delta }{\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{\left\{\mathrm{cosh}\left[r_+(\mathrm{\Delta }\varphi +2\pi n)r_{}\mathrm{\Delta }t\right]\mathrm{cosh}\left[r_+\mathrm{\Delta }tr_{}(\mathrm{\Delta }\varphi +2\pi n)\right]\right\}^\mathrm{\Delta }}}`$
$`+O\left(\left({\displaystyle \frac{r_+^2}{r^2}}\right)^{\mathrm{\Delta }+1}\right).`$ (44)
The boundary-dependent factor in the leading term in (44) at $`r\mathrm{}`$ can be rewritten as
$$\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}\frac{1}{\left\{\mathrm{sinh}\left[\frac{1}{2}(r_+r_{})(\mathrm{\Delta }\varphi +\mathrm{\Delta }t+2\pi n)\right]\mathrm{sinh}\left[\frac{1}{2}(r_++r_{})(\mathrm{\Delta }\varphi \mathrm{\Delta }t+2\pi n)\right]\right\}^\mathrm{\Delta }},$$
(45)
which is recognized as the dominant factor in the Greenโs function in the Poincarรฉ vacuum on the boundary of the rotating BTZ hole (see and the references therein). |
warning/0002/physics0002036.html | ar5iv | text | # Dielectronic recombination of lithium-like Ni25+ ions โ high resolution rate coefficients and influence of external crossed E and B fields
## I Introduction
Dielectronic recombination (DR) is an electron-ion collision process which is well known to be important in astrophysical and fusion plasmas. In DR the initially free electron is transferred to a bound state of the ion via a doubly excited intermediate state which is formed by an excitation of the core and a simultaneous attachment of the incident electron. This two step process
$$e^{}+A^{q+}[A^{(q1)+}]^{}A^{(q1)+}+h\nu $$
(1)
involves dielectronic capture (time-inverse Auger process) as the first step with a subsequent stabilization of the lowered charge state by radiative decay to a state below the ionization limit. This second step competes with autoionization which would transfer the ion back into its initial charge state $`q`$ with the net effect being resonant elastic or inelastic electron scattering. Another recombination process, which in contrast to DR is non-resonant, is radiative recombination (RR)
$$e^{}+A^{q+}A^{(q1)+}+h\nu $$
(2)
where the initially free electron is transferred to a bound state of the ion and a photon is emitted simultaneously. The cross section for RR diverges at zero electron energy and decreases rapidly towards higher energies. In the present investigation we regard RR as a continuous background on top of which DR resonances are emerging.
In the case of narrow non-overlapping DR resonances the DR cross section due to an intermediate state labelled $`d`$ can be well approximated by
$$\sigma _d(E_{\mathrm{cm}})=\overline{\sigma }_dL_d(E_{\mathrm{cm}})$$
(3)
with the electron-ion center-of-mass (c. m.) frame energy $`E_{\mathrm{cm}}`$, the Lorentzian line shape $`L_d(E)`$ normalized to $`L_d(E)๐E=1`$ and the resonance strength
$`\overline{\sigma }_d`$ $`=`$ $`4.95\times 10^{30}\mathrm{cm}^2\mathrm{eV}^2\mathrm{s}`$ (5)
$`\times {\displaystyle \frac{1}{E_d}}{\displaystyle \frac{g_d}{2g_i}}{\displaystyle \frac{A_\mathrm{a}(di)_fA_\mathrm{r}(df)}{_kA_\mathrm{a}(dk)+_f^{}A_\mathrm{r}(df^{})}}`$
where $`E_d`$ is the resonance energy, $`g_i`$ and $`g_d`$ are the statistical weights of the initial ionic core $`i`$ and the doubly excited intermediate state $`d`$, $`A_\mathrm{a}(di)`$ and $`A_\mathrm{r}(df)`$ denote the rate for an autoionizing transition from $`d`$ to $`i`$ and the rate for a radiative transition from $`d`$ to states $`f`$ below the first ionization limit, respectively. The summation indices $`k`$ and $`f^{}`$ run over all states which from $`d`$ can either be reached by autoionization or by radiative transitions, respectively.
Soon after the establishment of DR as an important process governing the charge state balance of ions in the solar corona, Burgess and Summers and Jacobs et al. realized that DR cross sections should be sensitive to external electric fields present in virtually any plasma environment. The electric field enhancement of DR rates was subsequently reproduced in a number of theoretical calculations. Briefly, the effect arises from the Stark mixing of $`\mathrm{}`$ states and the resulting influence on the autoionization rates which, by detailed balancing, determine the capture of the free electron. Autoionization rates strongly decrease with increasing $`\mathrm{}`$ and, therefore, only low $`\mathrm{}`$ states significantly contribute to DR. Electric fields mix low and high $`\mathrm{}`$ states and thereby increase the autoionization rates of the high $`\mathrm{}`$ states and consequently also the contribution of high Rydberg states to DR.
The control of external fields in experiments using intense ion and electron beams is a challenge. Results from early recombination experiments could only be brought in agreement with theory under the assumption that external electric fields had been present in the interaction region. The first experiment where external fields were applied under well controlled conditions was performed by Mรผller and coworkers who investigated DR in the presence of external fields (DRF) of singly charged Mg<sup>+</sup> ions. They observed an increase of the measured DR cross section by a factor of about 1.5 when increasing the motional $`\stackrel{}{v}\times \stackrel{}{B}`$ electric field from 7.2 to 23.5 V/cm. The agreement of these results with theoretical predictions was at the 20% level. Further DRF experiments with multiply charged C<sup>3+</sup> ions also revealed drastic DR rate enhancements by electric fields; however, the large uncertainties of these measurements left ambiguities.
The first DRF experiment using highly charged ions at a storage ring was carried out with Si<sup>11+</sup> ions by Bartsch et al.. It produced results with an unprecedented accuracy, enabling a detailed comparison with theory. Whereas the overall agreement between experiment and theory for the magnitude of the effect (up to a factor of 3 when increasing the field from 0 V/cm to 183 V/cm) was fair, discrepancies remained in the functional dependence of the rate enhancement on the electric field strength. This finding stimulated theoretical investigations of the role of the additional magnetic field which is always present in storage ring DR experiments, since it is needed to guide and confine the electron beam within the electron cooler. In a model calculation Robicheaux and Pindzola found that in a configuration of crossed $`\stackrel{}{E}`$ and $`\stackrel{}{B}`$ fields indeed the magnetic field influences through the mixing of $`m`$ levels the rate enhancement generated by the electric field. More detailed calculations confirmed these results. It should be noted that in theoretical calculations by Huber and Bottcher no influence of a pure magnetic field of at least up to 5 T on DR was found.
Inspired by these predictions we previously performed storage ring DRF experiments using Li-like Cl<sup>14+</sup> and Ti<sup>19+</sup> ions and crossed $`\stackrel{}{E}`$ and $`\stackrel{}{B}`$ fields where we clearly discovered a distinct effect of the magnetic field strength on the magnitude of the DR rate enhancement. The electric field effect decreased monotonically with the $`\stackrel{}{B}`$ field increasing from 30 mT to 80 mT. A decrease of the electric field enhancement by a crossed magnetic field is also predicted by the model calculation of Robicheaux and Pindzola for magnetic fields larger than approximately 20 mT where, due to a dominance of the magnetic over the electric interaction energy, the $`\mathrm{}`$-mixing weakens and consequently the number of states participating in DR decreases. At lower magnetic fields $`m`$-mixing yields an increase of the DR rate with increasing $`\stackrel{}{B}`$ field. A corresponding experimental observation has been made recently by Klimenko and coworkers who studied recombination of Ba<sup>+</sup> ions from a continuum of finite bandwidth which they had prepared by laser excitation of neutral Ba atoms. For a given electric field strength of 0.5 V/cm, they find that the recombination rate is increasingly enhanced by crossed magnetic fields up to about 20 mT. However, there is no effect of the magnetic field when it is directed parallel to the electric field vector. For the $`m`$-mixing to occur the crossed $`\stackrel{}{E}`$ and $`\stackrel{}{B}`$ arrangement is essential. In the case of parallel $`\stackrel{}{B}`$ and $`\stackrel{}{E}`$ fields $`m`$ remains a good quantum number and no influence of the magnetic field is expected.
The aim of the present investigation with Li-like Ni<sup>25+</sup> is to extend the previous studies to an ion with even higher nuclear charge $`Z`$. Because of the $`Z^4`$ scaling of radiative rates it is expected that with higher $`Z`$ less $`\mathrm{}`$ states of a given Rydberg $`n`$ level take part in DR and therefore the sensitivity to $`\mathrm{}`$-mixing decreases. Results of Griffin and Pindzola who calculated decreasing DR rate enhancements for increasing charge states of iron ions point into the same direction. Another aspect of going to higher $`Z`$ is that even high lying $`2p_jn\mathrm{}`$ Rydberg resonances are more separated in energy and therefore easier to resolve. This has been demonstrated by Brandau et al. who resolved $`2p_{1/2}n\mathrm{}`$ DR resonances up to $`n=41`$ in the recombination spectrum of Li-like Au<sup>76+</sup> ions. When choosing Ni<sup>25+</sup> we hoped to be able to study the field enhancement effect on a single $`2p_jn\mathrm{}`$ resonance. This would enable a quantitative comparison with theory which at present is limited to a single low value of $`n`$ when explicitly treating all $`n\mathrm{}m`$ levels required for a realistic description of DR in crossed $`\stackrel{}{E}`$ and $`\stackrel{}{B}`$ fields.
In our DRF studies we have chosen Li-like ions as test systems since on the one hand their electronic structure is simple enough to be treated theoretically on a high level of sophistication and on the other hand provides strong DR channels connected to the $`2s2p_j`$ core excitations.
## II Experiment
The measurements have been performed at the heavy ion storage ring TSR of the Max-Planck-Institut fรผr Kernphysik in Heidelberg. For a general account of experimental techniques at heavy ion storage rings the reader is referred to an article by Mรผller and Wolf. Recombination measurements at storage rings have been reviewed recently by Mรผller, Schippers and Wolf et al.. Detailed descriptions of the experimental procedure for field free DR measurements have been given by Kilgus et al. and more recently by Lampert et al.. Therefore, we here only describe more explicitly experimental aspects pertaining especially to the present investigation.
The <sup>58</sup>Ni<sup>25+</sup> ion beam was supplied by the MPI tandem booster facility and injected into the TSR with an energy of 343 MeV. Using multiturn injection and e-cool stacking ion currents of up to 3300 $`\mu `$A were stored in the TSR. At these high ion currents, however, intra beam scattering heated the ion beam during DR measurements resulting in a considerable loss of energy resolution. To avoid this and to limit the recombination count rate to below 1 MHz, i. e. to below a count rate where dead time effects are still negligible, we kept ion currents below 1 mA during all DR measurements. In the storage ring the circulating Ni<sup>25+</sup> ions were merged with the magnetically guided electron beam of the electron cooler. In the present experiment the electron density was $`5.4\times 10^6`$ cm<sup>-3</sup> at cooling energy. Generally the electron density varies with the cathode voltage $`U_\mathrm{c}`$ thereby following a $`U_\mathrm{c}^{3/2}`$ dependence. The distribution of collision velocities in the electron-ion center of mass frame can be described by the anisotropic Maxwellian
$`f(\stackrel{}{v},v_{\mathrm{rel}})`$ $`=`$ $`{\displaystyle \frac{m_\mathrm{e}}{2\pi k_\mathrm{B}T_{}}}\mathrm{exp}\left({\displaystyle \frac{m_\mathrm{e}v_{}^2}{2k_\mathrm{B}T_{}}}\right)`$ (7)
$`\times \left[{\displaystyle \frac{m_\mathrm{e}}{2\pi k_\mathrm{B}T_{||}}}\right]^{1/2}\mathrm{exp}\left({\displaystyle \frac{m_\mathrm{e}(v_{||}v_{\mathrm{rel}})^2}{2k_\mathrm{B}T_{||}}}\right)`$
characterized by the longitudinal and transverse temperatures $`T_{||}`$ and $`T_{}`$. In Eq. (7) $`m_\mathrm{e}`$ is the electron mass, $`k_\mathrm{B}`$ is the Boltzmann constant, and $`v_{\mathrm{rel}}`$ is the detuning of the average longitudinal electron velocity from that at cooling, which determines the relative energy $`E_{\mathrm{rel}}m_\mathrm{e}v_{\mathrm{rel}}^2/2`$ between the electron and the ion beam. The longitudinal temperature, inferred from the experimental resolution for relative energies $`E_{\mathrm{rel}}k_\mathrm{B}T_{}`$, was $`k_\mathrm{B}T_{||}0.25`$ meV. It implies an energy resolution given by $`\mathrm{\Delta }E`$(FWHM)$`=4\sqrt{\mathrm{ln}(2)k_\mathrm{B}T_{||}E_{\mathrm{rel}}}`$. The longitudinal velocity spread of the stored ion beam yields a considerable contribution to this temperature, while the velocity spread of the electron beam alone, after acceleration, is estimated to be $`<0.1`$ meV. In the transverse direction the electron beam was adiabatically expanded from a diameter $`d_c9.5`$ mm at the cathode to a diameter $`d_\mathrm{e}=29.5`$ mm in the interaction region; the reduction of its transverse velocity spread by this expansion determines the low value of the transverse temperature of $`k_\mathrm{B}T_{}10`$ meV.
Before starting a measurement, the ion beam was cooled for 5 seconds until the beam profiles reached their equilibrium widths. This can be monitored online by employing beam profile monitors based on residual gas ionization. The cooled ion beam had a diameter $`d_\mathrm{i}2`$ mm. During the measurement the electron cooler voltage was stepped through a preset range of values different from the cooling voltage, thus introducing non-zero mean relative velocities between ions and electrons. Recombined Ni<sup>24+</sup> ions were counted as a function of the cooler voltage with a CsI-scintillation detector located behind the first dipole magnet downstream of the electron cooler. The dipole magnet bends the circulating Ni<sup>25+</sup> ion beam onto a closed orbit and separates the recombined Ni<sup>24+</sup> ions from that orbit.
Two different measurement schemes were applied for the measurement of i) a high resolution โfield-freeโ DR spectrum (residual stray electric fields $`5`$ V/cm) and ii) DRF spectra with motional electric fields ranging up to 300 V/cm. In view of the result of Huber and Bottcher who calculated that purely magnetic fields below 5 T do not influence DR, the use of the term โfield-freeโ seems justified in case i) even with the magnetic guiding field (up to 80 mT) still present in the electron cooler.
### A Procedure for a field free high resolution measurement
In between two measurement steps for different values of $`E_{\mathrm{rel}}`$ the cooler voltage was first set back to the cooling value in order to maintain the ion beam quality and then set to a reference value which is chosen to lead to a relative velocity where the electron-ion recombination signal is very small, favourably being only due to a negligible RR contribution (reference relative energy $`E_{\mathrm{ref}}`$). Under this condition the recombination rate measured at the reference point monitors the background signal due to electron capture from residual gas molecules. Choosing short time intervals of the order of only 10 ms for dwelling on the measurement, cooling and reference voltages ensured that the experimental environment did not change significantly between the signal and the background measurements. An additional interval of 1.5 ms after each change of the cooler voltage allowed the power supplies to reach the preset values before data taking was started.
The electron-ion recombination coefficient
$$\alpha (E_{\mathrm{rel}})=d^3\stackrel{}{v}\sigma (v)vf(\stackrel{}{v},v_{\mathrm{rel}})$$
(8)
is obtained from the background corrected recombination count rate $`R(E_{\mathrm{rel}})R(E_{\mathrm{ref}})`$, the detection efficiency $`\eta `$, the electron density $`n_\mathrm{e}`$, the number of stored ions $`N_\mathrm{i}`$, the nominal length $`L=1.5`$ m of the interaction zone and the ring circumference $`C=55.4`$ m using the relation
$$\alpha (E_{\mathrm{rel}})=\frac{R(E_{\mathrm{rel}})R(E_{\mathrm{ref}})}{\gamma _i^2\eta n_\mathrm{e}(E_{\mathrm{rel}})N_\mathrm{i}L/C}+\alpha (E_{\mathrm{ref}})\frac{n_\mathrm{e}(E_{\mathrm{ref}})}{n_\mathrm{e}(E_{\mathrm{rel}})}$$
(9)
where $`\gamma _\mathrm{i}=1+E_\mathrm{i}/(m_\mathrm{i}c^2)`$ is the relativistic Lorentz factor for the transformation between the c. m. and the laboratory frames where the ions of mass $`m_\mathrm{i}`$ have the kinetic energy $`E_\mathrm{i}`$. The detection efficiency of the CsI-scintillation detector used to detect the recombined ions is very close to unity for count rates up to 2.5 MHz. The second term in Eq. (9) is to be added in case of a non-negligible electron-ion recombination rate at the reference energy. We insert the theoretical RR rate at $`E_{\mathrm{ref}}=131.5eV`$ which we have calculated to be $`\alpha (E_{\mathrm{ref}})=1.39\times 10^{11}`$ cm<sup>3</sup>/s using a semi-classical formula for the radiative recombination cross section
$`\sigma _{RR}(E_{\mathrm{rel}})`$ $`=`$ $`2.1\times 10^{22}\mathrm{cm}^2`$ (11)
$`\times {\displaystyle \underset{n_{\mathrm{min}}}{\overset{n_{\mathrm{cut}}}{}}}k_nt_n{\displaystyle \frac{q^4^2}{nE_{\mathrm{rel}}(q^2+n^2E_{\mathrm{rel}})}}`$
with $``$ denoting the Rydberg constant and $`k_n`$ being correction factors given by Andersen and Bolko. This expression for recombination on bare nuclei is used to approximately describe RR on a lithium-like core by introducing the lowest quantum number $`n_{\mathrm{min}}=2`$ and weight factors $`t_n`$ accounting for partial occupation of $`n`$-shells. In our calculation we use $`t_2=7/8`$ and $`t_n=1`$ for $`n>2`$. For the maximum (cut off) quantum number we use $`n_{\mathrm{cut}}=150`$ as explained below.
After the generation of a recombination spectrum from the experimental data via Eq. (9) a correction procedure accounting for non-perfect beam overlap in the merging sections of the cooler is applied which in our case only slightly redistributes the DR resonance strengths, resulting in DR peaks narrower and taller by small amounts. The systematic uncertainty in the absolute recombination rate coefficient is due to the ion and electron current determination, the corrections accounting for the merging and demerging sections of the electron and ion beams, and the detection efficiency. It is estimated to be $`\pm 15\%`$ of the measured recombination rate coefficient. The statistical uncertainty of the results presented below amounts to less than $`1\%`$ of the rate coefficient maximum.
### B Procedure for DRF measurements
The geometry of the magnetic and electric fields present in the merging section of the electron cooler is sketched in Fig. 1. We choose the $`z`$-axis to be defined by the ion beam direction. The magnetic guiding field $`\stackrel{}{B}`$ defines the electron beam direction. The field strength $`B`$ is limited both towards low and high values. Only fields $`B>25`$ mT guarantee a reliable operation of the electron cooler. The maximum tolerable current through the generating coils limits $`B`$ to at most 80 mT. Correction coils allow the steering of the electron beam in the $`x`$-$`y`$ plane. In the first place these are used to minimize the transverse field components $`B_x`$ and $`B_y`$ with respect to the ion beam, such that the two beams are collinear and centered to each other. The collinearity is inferred indirectly from beam profile measurements of the cooled ion beam with an accuracy of $`0.2`$ mrad; i. e. the transverse magnetic field components caused by imperfections in the beam alignment amounts at most to 2$`\times `$10$`{}_{}{}^{4}B`$. Residual fields which may vary in size and direction along the overlap length, are expected to be also of this magnitude. Since the settings of the various steering magnets result from a rather tedious beam optimization process, they are not exactly reproduced after each optimization procedure that is required e. g. after a change of the magnetic guiding field $`B_z`$. This means that the residual transverse magnetic fields for the collinear geometry may also slightly vary from one set of cooler settings to another. All uncertainties in the transverse magnetic field translate into an uncertainty in the motional electric field of less than $`\pm `$10 V/cm in our present experiment.
In DRF measurements we offset the current through the correction coils to generate additional magnetic field components $`B_x`$ and $`B_y`$. Their influence on the stored ion beam is negligible, i. e. the ion beam is still travelling with velocity $`v_\mathrm{i}`$ in $`z`$-direction. However, in the frame of the ion beam the magnetic field components $`B_x`$ and $`B_y`$ generate a motional electric field $`E_{}=\sqrt{E_y^2+E_x^2}=v_\mathrm{i}\sqrt{B_x^2+B_y^2}`$ in the $`x`$-$`y`$ plane rotated out of the $`y`$ direction by the azimuthal angle $`\varphi =\mathrm{arctan}(B_y/B_x)=\mathrm{arctan}(E_x/E_y)`$, i. e. $`E_{}=E_y`$ for $`\varphi =0`$. The electrons (due to their much lower mass) follow the resulting magnetic field vector $`\stackrel{}{B}`$ which now crosses the ion beam at the angle $`\theta =\mathrm{arctan}(\sqrt{B_x^2+B_y^2}/B_z)`$. Two consequences are to be dealt with: i) The ion beam now probes different portions of the space charge well of the electron beam. This reduces the energy resolution. In order to minimize this effect we used a rather small electron density of only $`5.4\times 10^6`$ cm<sup>-3</sup> at cooling, i. e. one order of magnitude smaller than in the Si<sup>11+</sup> experiment of Bartsch et al.. ii) The angle $`\theta `$ between electron beam and ion beam explicitly enters the formula for the transformation from the laboratory system to the c. m. system which is easily derived from the conservation of four-momentum. It reads
$$E_{\mathrm{rel}}=m_\mathrm{i}c^2(1+\mu )\left[\sqrt{1+\frac{2\mu }{(1+\mu )^2}(\mathrm{\Gamma }1)}1\right]$$
(12)
with the mass ratio $`\mu =m_\mathrm{e}/m_\mathrm{i}`$,
$$\mathrm{\Gamma }=\gamma _\mathrm{i}\gamma _\mathrm{e}\sqrt{(\gamma _\mathrm{i}^21)(\gamma _\mathrm{e}^21)}\mathrm{cos}\theta $$
(13)
and $`\gamma _\mathrm{e}=1+E_\mathrm{e}/m_\mathrm{e}c^2`$; $`E_\mathrm{e}`$ denotes the electron laboratory energy. It is obvious that the cooling condition $`E_{\mathrm{rel}}=0`$ can only be reached for $`\gamma _\mathrm{i}=\gamma _\mathrm{e}`$ and $`\theta =0`$, i. e. for $`\mathrm{\Gamma }=1`$. In regular field-free measurements a scheme of intermittent cooling is used during data taking, i. e. after each measured energy a cooling interval ($`E_{\mathrm{rel}}=0`$) is inserted. For DRF measurements this would require rapid switching from $`\theta 0`$ to cooling with $`\theta =0`$. It turned out that such a procedure heavily distorts the electron beam mainly because of the slow response of the power supplies controlling the steering coils. Under such conditions useful measurements could not be performed. Therefore, we omitted the intermittent cooling and reference measurement intervals, thereby losing resolution. After each injection into the ring and an appropriate cooling time the correction coils were set to produce a defined $`E_{}`$ and the cathode voltage was ramped very quickly through a preset range with a dwell-time of only 1 ms per measurement point. In such a manner a spectrum for one $`E_{}`$ setting was collected within only 4 s. After termination of the voltage ramp the correction coils were set back to $`\theta =0`$ and the whole cycle started again with the injection of ions into the ring. In subsequent cycles a range of typically 30 preset $`E_{}`$ values was scanned. Each spectrum was measured as many times as needed for reaching a satisfying level of statistical errors. This whole procedure was repeated for different settings of the guiding field strength $`B_z`$.
In order to compare only contributions from DR to the measured spectra we subtracted an empirical background function $`\alpha _{\mathrm{BG}}(E_{\mathrm{rel}})=a_0+a_1E_{\mathrm{rel}}+a_2/(1+a_3E_{\mathrm{rel}}+a_4E_{\mathrm{rel}}^2)`$ with the coefficients $`a_i`$ determined by fitting $`\alpha _{\mathrm{BG}}(E_{\mathrm{rel}})`$ to those parts of the spectrum which do not exhibit DR resonances. One should note that a proper calculation of the RR rate coefficient is hampered by the fact that for $`\theta 0`$ the electron velocity distribution probed by the ion beam cannot be described by Eq. (7).
## III Results and Discussion
### A Recombination at zero electric field
#### 1 DR cross section
For the $`\mathrm{\Delta }n=0`$ DR channels of Li-like Ni<sup>25+</sup>, i. e. for DR involving excitations which do not change the main quantum number of any electron in the $`1s^22s`$ core, Eq. (1) reads more explicitly
$`e^{}`$ $`+`$ $`\mathrm{Ni}^{25+}(1s^22s_{1/2})\mathrm{Ni}^{24+}(1s^22p_jn\mathrm{})`$ (14)
$``$ $`\{\begin{array}{cc}\mathrm{Ni}^{24+}(1s^22s_{1/2}n\mathrm{})+h\nu \hfill & \text{(type I)}\hfill \\ \mathrm{Ni}^{24+}(1s^22p_jn^{}\mathrm{}^{})+h\nu ^{}\hfill & \text{(type II)}\hfill \end{array}`$ (17)
The lowest Rydberg states which are energetically allowed are $`n=13`$ and $`n=11`$ for $`2s_{1/2}2p_{1/2}`$ and $`2s_{1/2}2p_{3/2}`$ core excitations, respectively.
The Ni$`{}_{}{}^{25+}(1s^22s_{1/2})`$ recombination spectrum has been measured for $`0|E_{\mathrm{rel}}|131.5`$ eV. The result is shown in Fig. 2. At $`E_{\mathrm{rel}}=0`$ a sharp rise of the recombination rate due to RR is observed. At higher energies DR resonances due to $`\mathrm{\Delta }n=0`$ $`2s_{1/2}2p_j`$ transitions occur, the lowest resonance appearing at $`E_{\mathrm{rel}}=0`$ eV. Individual $`2p_jn\mathrm{}`$ resonances are resolved for $`n32`$. Their resonance strengths have been extracted from the measured spectrum by first subtracting the theoretical recombination rate coefficient due to RR (cf. section II A) where the c. m. velocity spread can be neglected since $`E_{\mathrm{rel}}`$ is very large compared to $`kT_{}`$ and $`kT_{}`$. In principle the resulting rate coefficient should be zero at off-resonance energies. However, we find that probably due to our approximate treatment of RR (cf. Eq. (11)) small non-zero rate coefficients remain after subtraction of the calculated RR rate coefficient. These are removed by further subtracting a smooth background before the observed DR resonance structures are fitted by Gaussians. (Details of the observed smooth RR rate were not further investigated in the present work.) The resulting values for resonance positions and strengths are listed in Table I. The $`2p_{1/2}13\mathrm{}`$ and $`2p_{3/2}11\mathrm{}`$ resonances at about 2.5 eV and 4.5 eV, respectively, exhibit a splitting due to the interaction of the $`n\mathrm{}`$-Rydberg electron with the $`1s^22p_j`$ core. For higher $`n`$ resonances this splitting decreases and cannot be observed because of the finite experimental energy spread which increases as $`\sqrt{E_{\mathrm{rel}}}`$. The $`2p_j`$ Rydberg series limits, $`E_{\mathrm{}}`$, are obtained from a fit of the resonance positions $`E_n`$ with $`n16`$ to the Rydberg formula
$$E_n=E_{\mathrm{}}\left(\frac{q}{n\delta }\right)^2$$
(18)
with the quantum defect $`\delta `$ as a second fit parameter. The fit results are listed in Table II, where spectroscopic values for the series limits are also given. Our values agree with the spectroscopic values within $`0.6\%`$, i. e. within the experimental uncertainty of the energy scale. The result that the fitted quantum defects are almost zero reflects the fact that the interaction between the core electrons and the Rydberg electron is weak.
The measured DR rate decreases already below the $`2p_{1/2}`$ and $`2p_{3/2}`$ series limits as obtained from the fit to the peak positions. (cf. Fig. 2). This discrepancy results from field ionization of high Rydberg states with $`n>n_\mathrm{f}(3.2\times 10^8\mathrm{V}/\mathrm{cm}q^3/E_{\mathrm{dip}})^{1/4}`$ in the charge analyzing dipole magnet (see Sec. II) with magnetic field strength $`B_{\mathrm{dip}}=0.71`$ T where the moving ion experiences the motional electric field $`E_{\mathrm{dip}}=v_\mathrm{i}B_{\mathrm{dip}}`$ ($`n_\mathrm{f}`$ is the classical field ionization limit). A more realistic value for the cut-off has to account for Stark splitting and tunnelling effects. In Ref. an approximate value $`n_{\mathrm{cut}}(7.3\times 10^8\mathrm{V}/\mathrm{cm}q^3/E_{\mathrm{dip}})^{1/4}`$ was found. For the calculation of the actual cut-off quantum number $`n_{\mathrm{cut}}`$ relevant in this experiment one has to take into account that on the way from the cooler to the dipole magnet states with $`n>n_\mathrm{f}`$ may radiatively decay to states below $`n_f`$. An approximate calculation of this effect yields $`n_{\mathrm{cut}}=150`$ for the present case.
An estimate of the DR line strength escaping detection because of field ionization can be made by extrapolating the measured DR line strength to $`n=\mathrm{}`$ employing the $`n^3`$ scaling of the autoionization and the type II (cf. Eq. (17)) radiative rates. For autoionization rates we make the ansatz $`A_\mathrm{a}(n\mathrm{})=A_\mathrm{a}/n^3`$ for $`0\mathrm{}\mathrm{}_{\mathrm{max}}`$ and $`A_\mathrm{a}(n\mathrm{})=0`$ for $`\mathrm{}_{\mathrm{max}}<\mathrm{}<n`$. Rates for the sum of type I and type II radiative transitions we represent as $`A_\mathrm{r}(n\mathrm{})=A_\mathrm{r}^{(\mathrm{I})}+A_\mathrm{r}^{(\mathrm{II})}/n^3`$. The same representations of the relevant rates have already been used by Kilgus et al. in a recombination study of isoelectronic Cu<sup>26+</sup> ions. After summation over all $`\mathrm{}`$ substates Eq. (5) simplifies to
$$\overline{\sigma }_nE_n=S_0\frac{A_\mathrm{a}[A_\mathrm{r}^{(\mathrm{I})}+n^3A_\mathrm{r}^{(\mathrm{II})}]}{A_\mathrm{a}+n^3A_\mathrm{r}^{(\mathrm{I})}+A_\mathrm{r}^{(\mathrm{II})}}$$
(19)
with $`S_0=2.475(2j_\mathrm{c}+1)(\mathrm{}_{\mathrm{max}}+1)^2\times 10^{30}\mathrm{cm}^2\mathrm{eV}^2\mathrm{s}`$. For the statistical weights in Eq. (5) we have used $`g_i=2`$ and $`g_d=2(2\mathrm{}+1)(2j_\mathrm{c}+1)`$. Here, $`j_c`$ is the total angular momentum quantum number of the core excited state i. e. $`j_c=1/2`$ and 3/2 in the present case. The first step of the extrapolation procedure consists of adjusting the model parameters such that Eq. (19) fits the measured DR line strengths. Here the values for $`A_\mathrm{r}^{(\mathrm{I})}`$ have been taken from atomic structure calculations while $`S_0`$, $`A_a`$ and $`A_\mathrm{r}^{(\mathrm{II})}`$ were allowed to vary during the fit. It turned out that the fit is not very sensitive to large variations of the Auger rates $`A_a`$. In this situation we kept also the Auger rates fixed at values which have been inferred from atomic structure calculations and which are meant to be order of magnitude estimates only. In Fig. 3 we have plotted the measured and fitted resonance strengths (multiplied by the resonance energy) as a function of the main quantum number. The actual parameters used for drawing the fit curves are listed in Table II. The fit has been restricted to $`n20`$ for both the $`2p_{1/2}`$ and the $`2p_{3/2}`$ series because additional Coster-Kronig decay channels $`2p_{3/2}n\mathrm{}2p_{1/2}ฯต\mathrm{}^{}`$ open up when $`E_n`$ crosses the $`2p_{1/2}`$ series limit. A corresponding discontinuous decrease of the $`2p_{3/2}n\mathrm{}`$ DR resonance strength can be clearly discerned in Fig. 3. Since the additional autoionizing channels are included for the $`2p_{3/2}n\mathrm{}`$ series of resonances but not for the $`2p_{1/2}n\mathrm{}`$ series, $`A_\mathrm{a}`$ for the $`2p_{3/2}n\mathrm{}`$ series is more than a factor of 2 higher than $`A_\mathrm{a}`$ for the $`2p_{1/2}n\mathrm{}`$ series. A factor of 2 would just correspond to the ratio of statistical weights.
Inserting the fit parameters listed in Table II into Eqs. (18) and (19) now allows an extrapolation of the DR resonance positions and strengths, respectively, to be obtained for arbitrary high $`n`$. In order to check the quality of the extrapolation we have convoluted the extrapolated DR cross section with the experimental electron energy distribution using the electron beam temperatures $`k_\mathrm{B}T_{||}=0.25`$ meV and $`k_\mathrm{B}T_{}=10`$ meV. After adding the semiclassically calculated RR rate coefficient โ as described above โ the resulting extrapolated DR+RR recombination rate coefficient is plotted in Fig. 4 together with the experimental one. Despite the very simple model assumptions the calculated recombination rate agrees with the measured one also over the energy intervals covered by the $`2p_jn\mathrm{}`$ resonances with $`n>31`$ which are not resolved individually and therefore have not been used for the fits. Deviations for the Ni$`{}_{}{}^{24+}(1s^22p_{3/2}n\mathrm{})`$ resonances with $`n19`$ stem from the fact that for these resonances the Coster-Kronig decay channels to Ni$`{}_{}{}^{25+}(1s^22p_{1/2})`$ are closed whereas the fit has been made to resonances where they are open. At energies close to the series limits slight deviations from the model rate occurs even when the expected value of $`n_{\mathrm{cut}}=150`$ is inserted into the model (dashed line in Fig. 4). The origin of this discrepancy has probably to be searched for in the approximations made in both the field ionization model and in the model rate descriptions, in particluar regarding the dependence of the angular momentum $`\mathrm{}`$ which may be affected by even the small residual electric fields in the interaction region.
#### 2 Maxwellian plasma rate coefficient
The comparison between the measured data and the calculated extrapolation to $`n\mathrm{}`$ (full line in Fig. 4) suggests that only a minor part of the total Ni<sup>25+</sup> $`\mathrm{\Delta }n=0`$ DR resonance strength has not been measured. This enables us to derive from our measurement the Ni<sup>25+</sup> DR rate coefficient in a plasma. To this end the experimental DR rate coefficient is substituted by the extrapolated one at energies $`3`$ eV below the series limits. Compared to using the experimental result without extrapolation this results in a correction of the plasma rate coefficient of at most $`5\%`$. The experimental DR rate coefficient including the high $`n`$-extrapolation is convoluted by an isotropic Maxwellian electron velocity distribution characterized by the electron temperature $`T_\mathrm{e}`$. The resulting $`\mathrm{\Delta }n=0`$ DR plasma rate coefficient is displayed in Fig. 5 (thick full line). Summing experimental and extrapolation errors, the total uncertainty of the DR rate coefficient in plasmas determined in this work amounts to $`\pm 20\%`$.
A convenient representation of the plasma DR rate coefficient is provided by the following fit formula
$$\alpha (T_\mathrm{e})=T_\mathrm{e}^{3/2}\underset{i}{}c_i\mathrm{exp}(E_i/k_\mathrm{B}T_\mathrm{e})$$
(20)
It has the same functional dependence on the plasma electron temperature as the widely used Burgess formula, where the coefficients $`c_i`$ and $`E_i`$ are related to oscillator strengths and excitation energies, respectively. The results for the fit to the experimental Ni<sup>25+</sup> $`\mathrm{\Delta }n=0`$ DR rate coefficient in a plasma (thick full line in Fig. 5) are summarized in Table III. The fitted curve cannot be distinguished from the experimental plasma rate coefficient in a plot as presented in Fig. 5.
In Fig. 5 we also compare our results with theoretical results for $`\mathrm{\Delta }n=0`$ DR by Mewe et al. (dashed line), Romanik (dash-dotted line) and Teng et al. (dashed-dot-dotted line) who interpolated DR calculations performed by Chen for selected lithiumlike ions. At temperatures $`k_\mathrm{B}T_\mathrm{e}>1`$ eV the rate of Mewe et al., which is based on the Burgess formula, overestimates our experimental result by up to a factor of $`5`$. At lower temperatures the experimental result is underestimated by factors up to 10. Above an electron temperature of 30 eV Romanikโs theoretical result agrees with our $`\mathrm{\Delta }n=0`$ DR rate coefficient to within $`15\%`$, which is within the $`20\%`$ experimental accuracy. In this energy range the interpolation result of Teng et al. underestimates the experimental rate coefficient by 20โ30$`\%`$. It should be noted that neither the calculation of Romanik nor that of Teng et al. covers temperatures below 10 eV. When compared with our RR calculation (thin full line in Fig. 5) our experimental result shows that in the temperature range of 1 to 10 eV, where Ni<sup>25+</sup> ions may exist in photoionized plasmas, DR is still significant. The importance of DR in low temperature plasmas has been pointed out recently by Savin et al. who measured DR of fluorine-like Fe<sup>17+</sup> ions. At higher temperatures Romanikโs calculation suggests that above 100 eV $`\mathrm{\Delta }n=1`$ DR contributions become significant (upper dash-dotted line in Fig. 5).
### B DRF measurements
Fig. 6 shows a series of Ni<sup>25+</sup> recombination spectra measured in the presence of external electric fields $`E_{}`$ ranging from 0 to 270 V/cm. The magnetic field on the axis of the cooler has been $`B_z=80`$ mT. Due to the altered measurement scheme that leaves out the intermittent cooling of the ion beam, the energy resolution is reduced compared to Fig. 2. Now individual $`2p_jn\mathrm{}`$ DR resonances are resolved only up to $`n=21`$. Two features in the series of spectra are to be noted. Firstly, the strength of the DR resonances occurring below 47 eV does not depend on $`E_{}`$. Secondly, the strength of the unresolved high-$`n`$ DR resonances increases with increasing field strength. This can been seen more clearly from the close-up presented in Fig. 7. At energies of more than 10 eV below the $`2p_j`$ series limits the different DR spectra lie perfectly on top of each other whereas at higher energies (i. e. $`n>30`$) an increase of the DR intensities by up to a factor of 1.5 at $`E_{}=270`$ V/cm is observed. The degraded resolution of the DR spectrum does not allow us to resolve $`n`$ levels in the range of the electric field enhancement. In order to quantify this DR rate enhancement we consider integrated recombination rates with the integration intervals chosen as marked in Fig. 7.
The integration intervals 44.0โ53.5 eV and 66.2โ76.0 eV include all $`2p_jn\mathrm{}`$ resonances with $`n31`$ for $`j=1/2`$ and $`j=3/2`$, respectively. In the following we denote the resulting integrals by $`I_{1/2}`$ and $`I_{3/2}`$. The energy range 44.0โ53.5 eV also contains $`2p_{3/2}n\mathrm{}`$ resonances with $`16n19`$. These resonances, however, are not affected by the electric field strengths used in our experiment. Consequently, any change in the magnitude of $`I_{1/2}`$ as a function of $`E_{}`$ we attribute to field effects on $`2p_{1/2}n\mathrm{}`$ resonances. As a check of the proper normalization of the DR spectra we additionally monitor the integral $`I_0=_{2\mathrm{e}\mathrm{V}}^{18\mathrm{e}\mathrm{V}}\alpha _{\mathrm{DR}}(E_{\mathrm{rel}})๐E_{\mathrm{rel}}`$ which comprises the strengths of the $`2p_{3/2}11\mathrm{}`$, $`2p_{3/2}12\mathrm{}`$ and $`2p_{1/2}n\mathrm{}`$ DR resonances with $`13n15`$. Since these low $`n`$ resonances are not affected by $`E_{}`$ we expect $`I_0`$ to be constant. Any deviation of $`I_0`$ from a constant value would indicate a reduction of beam overlap due to too large a tilting angle $`\theta `$ of the electron beam. The maximum angle $`\theta _{\mathrm{max}}`$ to which the overlap of the electron beam with the ion beam is ensured over the full interaction length $`L`$ is given by $`\mathrm{tan}\theta _{\mathrm{max}}=(d_\mathrm{e}d_\mathrm{i})/L`$. With the geometrical values given above one obtains $`\theta _{\mathrm{max}}1^{}`$. Apart from the highest $`E_{}`$ at $`\varphi =180^{}`$ and at the lowest magnetic guiding field, i. e. at $`B_z=41.8`$ mT, the condition $`\theta <\theta _{\mathrm{max}}`$ was always met. This is exemplified in the upper panel of Fig. 8 where $`I_0`$, $`I_{1/2}`$ and $`I_{3/2}`$ are shown for -270 V/cm $`E_{}`$ 220 V/cm. There, positive (negative) field strength indicates $`\varphi =0^{}`$ ($`\varphi =180^{}`$). While $`I_{1/2}`$ and $`I_{3/2}`$ clearly exhibit a field effect which is nearly symmetric about $`E_{}=0`$ V/cm, $`I_0`$ is independent of $`E_{}`$.
The fact that the increase of the integrated recombination rate coefficient is independent of the sign of the electric field vector is expected from the cylindrical symmetry of the merged beams arrangement in the electron cooler. However, for the entire experimental setup this symmetry is broken by the charge analyzing dipole magnet with an electric field vector lying in the bending plane (the $`x`$-$`y`$ plane of the coordinate frame defined in the interaction region). This, in principle, could lead to redistribution of population between different $`m`$ substates in the dipole magnet and a resulting field ionization probability depending on the azimuthal angle $`\varphi `$ of the motional electric field vector in the cooler. In order to clarify this question we took a series of DRF spectra with $`\varphi `$ ranging from 5 to 175. At the same time the electric and magnetic fields were kept fixed at $`E_{}=100`$ V/cm and $`B_z=80`$ mT. A scan around a full circle was prohibited by the limited output of the power supplies used for steering the electron beam in the particular arrangement of this experiment. As shown in the lower panel of Fig. 8 no significant dependence of the integrated recombination rates on the azimuthal angle $`\varphi `$ was found.
As a measure for the magnitude of the field enhancement we introduce the field enhancement factor
$$r_j(E_{},B_z)=C_j(B_z)\frac{I_j(E_{},B_z)}{I_0(E_{},B_z)}$$
(21)
for $`j=1/2`$ or $`3/2`$ and the constant $`C_j(B_z)`$ chosen such that $`r_j(0,B_z)=1.0`$ (see below). Plots of enhancement factors as a function of $`E_{}`$ are shown in Fig. 9 for different values of $`B_z`$. Since the field effect is independent of the orientation of the electric field in the $`x`$-$`y`$ plane data points for $`\varphi =0^{}`$ and $`\varphi =180^{}`$ are plotted together. The enhancement factor exhibits a linear dependence on the electric field. Exceptions occur for $`B_z=41.8`$ mT and $`\varphi =180^{}`$ at $`E_{}`$ values where $`\theta `$ becomes maximal, and around $`E_{}=0`$ V/cm. In the former case the complete overlap of the ion beam and the electron beam over the full length of the interaction region is lost. This is indicated by a reduction of $`I_0`$, and apparently a consistent normalization cannot be carried out. In the latter case residual electric and magnetic field components resulting e. g. from a non-perfect alignment of the beams prevent us from reaching $`E_{}=0`$ eV. After excluding all data points with $`E_{}10`$ V/cm and those with $`\varphi =180^{}`$, $`E_{}200`$ V/cm for $`B_z=41.8`$ mT, we were able to fit straight lines to the measured field enhancement factors as a function of $`E_{}`$. The constants $`C_j(B_z)`$ in Eq. (21) have been chosen such that the fitted straight lines yield $`r_j^{(\mathrm{fit})}(0,B_z)=1.0`$.
As a measure for the electric-field enhancement of the DR rate enhancement we now consider the slopes
$$s_j(B_z)=\frac{dr_j^{(\mathrm{fit})}(E_{},B_z)}{dE_{}}$$
(22)
of the fitted straight lines, which are displayed as a function of the magnetic field strength in Fig. 10. The error bars correspond to statistical errors only. Systematic uncertainties e. g. due to residual fields are difficult to estimate. Nevertheless, their order of magnitude can be judged from the $`10\%`$ difference between the two data points at $`B_z=80`$ mT which have been measured with different cooler settings.
For all magnetic fields used the electric field dependence of the rate enhancement factor is steeper for the $`2p_{3/2}n\mathrm{}`$ series of Rydberg resonances than that for the $`2p_{1/2}n\mathrm{}`$ series. This can be understood from the fact that the multiplicity of states and consequently the number of states which can be mixed is two times higher for $`j=3/2`$ than for $`j=1/2`$. Following this argument one would expect a ratio of 2 for the respective incremental integrated recombination rates $`I_j(E_{},B_z)I_j(0,B_z)`$. From our measurements we find lower values scattering around 1.5 indicating a somewhat reduced number of states available for field mixing within the series of Ni$`{}_{}{}^{24+}(1s^22p_{3/2}n\mathrm{})`$ DR resonances. In calculations for Li-like Si<sup>11+</sup> and C<sup>3+</sup> ions ratios, even less than 1 have been found. This has been attributed to the electrostatic quadrupole-quadrupole interaction between the $`2p`$ and the $`n\mathrm{}`$ Rydberg electron in the intermediate doubly excited state, which more effectively lifts the degeneracy between the $`2p_{3/2}n\mathrm{}`$ than between the $`2_{1/2}n\mathrm{}`$ levels. Another reason for the reduced number of $`2p_{3/2}n\mathrm{}`$ states participating in DRF might be the existence of additional Coster-Kronig decay channels for these resonances with $`n20`$ to Ni$`{}_{}{}^{25+}(1s^22p_{1/2})`$ (cf. section III A).
Slopes for the relative electric field enhancement according to Eq. (22) were previously determined in measurements by Bartsch and coworkers for the lighter isoelectronic ions Si<sup>11+</sup>, Cl<sup>14+</sup> and Ti<sup>19+</sup> . The present results are compared with these previous data in Fig. 11. It should be noted that the comparison is only semi-quantitative, because the choice of integration ranges for the calculation of the integrated recombination coefficients is somewhat arbitrary and different cut-off quantum numbers $`n_{\mathrm{cut}}`$ exist for the different ions (cf. Table IV). Clearly, the relative enhancement for a given electric field strength is much lower for Ni<sup>25+</sup> than for the lighter ions studied so far, with the reduction for the step from lithium-like titanium ($`Z=22`$) to nickel ($`Z=28`$) apparently much larger than the step from, e. g., chlorine ($`Z=17`$) to $`Z=22`$. As a general trend with increasing $`Z`$, the radiative decay rates $`A_r`$ decrease $`Z`$ for type I transitions (see Eq. (17)) and $`(Z3)^4`$ for type II transitions, while the autoionization rates are rather independent of $`Z`$. This shifts the range where the low-$`\mathrm{}`$ autoionization rates are larger than the radiative stabilization rates (the condition for DR enhancement by $`\mathrm{}`$-mixing to occur ) down to states of lower principal quantum number $`n`$ for increasing $`Z`$. For Ni<sup>25+</sup> the degree of mixing reached by the typical experimental field strengths appears to be much reduced as compared to lighter ions. In addition, a clear dependence of the electric field enhancement on the the magnetic field strength is no longer observed for Ni<sup>25+</sup>. To clarify the reason for the strong reduction of both the electric and the additional magnetic field effect in the heavy system studied here, detailed quantitative calculations are desirable.
External electric and magnetic fields are ubiquitous in astrophysical or fusion plasmas. Therefore, it is of interest to look into the implications of the field enhancement effect for the Ni<sup>25+</sup> $`\mathrm{\Delta }n=0`$ DR rate coefficient in a plasma. As an example we show in Fig. 12 the ratio of rate coefficients derived from two measurements with and without external electric field. As compared to zero electric field, the recombination coefficient at our highest experimental electric field strength ($`E_{}=270`$ V/cm at $`B_z=80`$ mT) is enhanced by up to 11$`\%`$. This value only represents a lower limit for the enhancement at the given field strength, as we observe DR resonances due to Ni$`{}_{}{}^{24+}(1s^22p_jn\mathrm{})`$ intermediate states only up to $`n_{\mathrm{cut}}150`$. It should also be noted that in our experiments we did not reach the electric field strength where the DR rate enhancement saturates.
## IV Summary and conclusions
The recombination of lithium-like Ni<sup>25+</sup> ions has been experimentally studied in detail. Spectroscopic information on individual DR resonances associated with Ni$`{}_{}{}^{24+}(1s^22p_jn\mathrm{})`$ intermediate states, which have been experimentally resolved up to $`n=32`$, has been extracted and the $`\mathrm{\Delta }n=0`$ DR plasma rate coefficient has been derived. Our experimental result is underestimated by up to a factor of 2 by semi-empirically calculated rate coefficients. At plasma temperatures above 10 eV, results of detailed theoretical calculations are available which agree well with the experiment. At lower temperatures, where Ni<sup>25+</sup> ions may exist in photoionized plasmas, no data for the DR rate coefficient have been previously available.
In the presence of external electric fields up to 300 V/cm, the measured DR resonance strength is enhanced by a factor 1.5; a rather weak effect in comparison with previous measurements at $`Z=11`$, 17 and 22. Due to the overall weakness of the field effect for Ni<sup>25+</sup> ions, a marked dependence of the DR rate enhancement on the strength of a crossed magnetic field as observed for Cl<sup>14+</sup> and Ti<sup>19+</sup> ions has not been detectable in the present investigation.
Experimental limitations prevented us from obtaining an energy resolution in our DRF measurements comparable to that achieved in the field-free measurement. Consequently, we could not resolve a Ni$`{}_{}{}^{24+}(1s^22p_jn\mathrm{})`$ DR resonance with $`n`$ high enough to exhibit a field effect. Such an observation would have facilitated a direct comparison with ab initio calculations, which due to the large number of $`n\mathrm{}m`$ states to be considered, presently can only treat a single $`n`$ manifold of Rydberg states in the presence of crossed $`\stackrel{}{E}`$ and $`\stackrel{}{B}`$ fields. Improvements of this situation can be expected in the near future from the steady increase of computing power on the theoretical side, and on the experimental side from a dedicated electron target which is presently being installed at the TSR. With the electron target and the electron cooler operating at the same time we will be able to perform DRF measurements with continuously cooled ion beams, yielding DRF spectra with increased resolution.
###### Acknowledgements.
We gratefully acknowledge support by the German Federal Ministry for Education and Research (BMBF) through contracts no. 06 GI 848 and no. 06 HD 854. R. P. acknowledges support by the Division of Chemical Sciences, U.S. Department of Energy under contract DE-FG03-97ER14787 with the University of Nevada, Reno. |
warning/0002/hep-lat0002025.html | ar5iv | text | # 1 Introduction
## 1 Introduction
Historically, the measurement by the UA1 Collaboration of a large value for the neutral meson mass-difference $`\mathrm{\Delta }m_d`$ was the first indication that top quark mass had to be very heavy . This mass difference is induced by $`B_d^0`$$`\overline{B}_d^0`$ oscillations. Since then, $`B^0`$$`\overline{B}^0`$ mixing became one of the most important ingredients of current analyses of the unitarity triangle and of CP violation in the Standard Model and beyond .
In the Standard Model the mixing, induced by the box diagrams with an internal top quark, is summarized by the next-to-leading order (NLO) formula
$`\mathrm{\Delta }m_q={\displaystyle \frac{G_F^2}{6\pi ^2}}m_W^2\eta _BS_0(x_t)|V_{tq}V_{tb}^{}|^2m_{B_q}f_{B_q}^2\widehat{B}_{B_q}(q=s,d),`$ (1)
where $`S_0(x_t)`$ is the Inami-Lim function , $`x_t=m_t^2/M_W^2`$ ($`m_t`$ is the $`\overline{\mathrm{MS}}`$ top mass $`m_t^{\overline{\mathrm{MS}}}(m_t^{\overline{\mathrm{MS}}})=165(5)`$ GeV) and $`\eta _B=0.55(1)`$ is the perturbative QCD short-distance NLO correction. The remaining factor, $`f_{B_q}^2\widehat{B}_{B_q}`$, encodes the information of non-perturbative QCD and this is what can be computed on the lattice.
Traditionally, the $`B`$-parameter of the renormalized operator is defined as
$`\overline{B}_q|Q_q^{\mathrm{\Delta }B=2}(\mu )|B_q={\displaystyle \frac{8}{3}}m_{B_q}^2f_{B_q}^2B_{B_q}(\mu ),`$ (2)
where $`Q_q^{\mathrm{\Delta }B=2}=(\overline{b}\gamma _\mu (1\gamma _5)q)(\overline{b}\gamma _\mu (1\gamma _5)q)`$ and $`\mu `$ is the renormalization scale. This definition stems from the vacuum saturation approximation (VSA) in which $`B_{B_q}=1`$. The renormalization group invariant, and scheme-independent, $`B`$ parameter $`\widehat{B}_{B_q}`$ of eq. (1) is defined as
$`\widehat{B}_{B_q}=\alpha _s(\mu )^{\gamma _0/2\beta _0}\left\{1+{\displaystyle \frac{\alpha _s(\mu )}{4\pi }}J\right\}B_{B_q}(\mu ),`$ (3)
where $`\gamma _0=4`$ and $`J`$ depends on the scheme used for renormalizing $`Q_q^{\mathrm{\Delta }B=2}(\mu )`$ (see below).
Besides the $`B^0`$$`\overline{B}^0`$ amplitude, an important quantity for phenomenological applications is given by the ratio
$`{\displaystyle \frac{\mathrm{\Delta }m_s}{\mathrm{\Delta }m_d}}={\displaystyle \frac{|V_{ts}|^2}{|V_{td}|^2}}{\displaystyle \frac{m_{B_s}}{m_{B_d}}}\xi ^2,`$ (4)
where $`\xi =f_{B_s}\sqrt{\widehat{B}_{B_s}}/f_{B_d}\sqrt{\widehat{B}_{B_d}}`$.
In this paper, we have computed the hadronic parameters appearing in eqs. (1) and (4), using a non-perturbatively improved action, and with operators renormalized on the lattice with the non-perturbative method of ref. , as implemented in . Our main results are
$`f_{B_d}\sqrt{\widehat{B}_{B_d}}`$ $`=`$ $`206(28)(7)\mathrm{MeV},\xi ={\displaystyle \frac{f_{B_s}\sqrt{\widehat{B}_{B_s}}}{f_{B_d}\sqrt{\widehat{B}_{B_d}}}}=1.16(7),`$
$`f_{B_s}\sqrt{\widehat{B}_{B_s}}`$ $`=`$ $`237(18)(8)\mathrm{MeV},r_{sd}=\xi ^2\left({\displaystyle \frac{m_{B_s}^2}{m_{B_d}^2}}\right)=1.40(18),`$
$`f_{B_d}`$ $`=`$ $`174(22)_{00}^{+7+4}\mathrm{MeV},f_{B_s}=204(15)_{00}^{+7+3}\mathrm{MeV},`$ (5)
$`{\displaystyle \frac{f_{B_s}}{f_{B_d}}}`$ $`=`$ $`1.17(4)_1^{+0},{\displaystyle \frac{f_{B_d}}{f_{D_s}}}=0.74(5).`$
From eq. (1) we write
$`\mathrm{\Delta }m_d[ps^1]=(0.153\pm 0.010)|V_{td}|^2f_{B_d}^2\widehat{B}_{B_d}[\mathrm{MeV}^2].`$ (6)
We stress that what is really relevant for $`B^0`$$`\overline{B}^0`$ mixing, and can be directly computed on the lattice, is the physical amplitude, corresponding to $`f_{B_d}^2\widehat{B}_{B_d}`$, and not the decay constant $`f_{B_d}`$ and $`\widehat{B}_{B_d}`$ separately. The value of $`f_{B_d}\sqrt{\widehat{B}_{B_d}}`$ given in the first of eqs. (5) comes from the calculation of the amplitude, and thus it includes the correlation between the decay constant and the $`B`$ parameter. Its final error is then smaller than that obtained by combining the results and errors obtained from $`f_{B_d}`$ and $`\widehat{B}_{B_d}`$ from different calculations.
By taking $`|V_{td}|A\lambda ^3\sqrt{(1\rho )^2+\eta ^2}=0.0080(5)`$ (where $`A`$, $`\rho `$, $`\eta `$ and $`\lambda `$ are the Wolfenstein parameters) this gives
$`\mathrm{\Delta }m_d=0.42(12)(4)\mathrm{ps}^1,`$ (7)
where the first error comes from the lattice uncertainty on $`f_{B_d}\sqrt{\widehat{B}_{B_d}}`$ and the second from the error on $`m_t`$ and on $`|V_{td}|`$. The most recent world average of experimental results is :
$`\mathrm{\Delta }m_d^{(\mathrm{exp}.)}=0.473(16)ps^1.`$ (8)
To predict $`\mathrm{\Delta }m_s`$ it is convenient to use the above experimental information and eq. (4) written as
$`\mathrm{\Delta }m_s={\displaystyle \frac{|V_{ts}|^2}{|V_{td}|^2}}\xi ^2\left({\displaystyle \frac{m_{B_s}}{m_{B_d}}}\mathrm{\Delta }m_d\right)^{(\mathrm{exp}.)}.`$ (9)
Using our result for $`\xi ^2`$ and $`|V_{ts}|^2/|V_{td}|^21/[\lambda ^2((1\rho )^2+\eta ^2]=24.4(5.0)`$, we get
$`\mathrm{\Delta }m_s=15.8(2.1)(3.3)ps^1,`$ (10)
to be compared with the experimental lower bound
$`\mathrm{\Delta }m_s>12.4ps^1.`$ (11)
In the above calculations, we have assumed the values of the relevant couplings, namely $`|V_{td}|`$ and $`|V_{ts}|`$, from the unitarity relations of the $`V_{\mathrm{CKM}}`$ matrix. Obviously, from the experimental determinations of $`\mathrm{\Delta }m_q`$ and the hadronic matrix elements we can, instead, constrain $`|V_{tq}|`$. This is what is usually done in the unitarity triangle analyses . For the sake of comparison, we also give the values of the decay constants and their ratios from our previous studies of these quantities, obtained using the same improved action, with a comparable statistics but on a larger lattice
$`f_{B_d}`$ $`=`$ $`173(13)_2^{+34}\mathrm{MeV},f_{B_s}=196(11)_0^{+42}\mathrm{MeV},`$ (12)
$`{\displaystyle \frac{f_{B_s}}{f_{B_d}}}`$ $`=`$ $`1.14(2)(1),{\displaystyle \frac{f_{B_d}}{f_{D_s}}}=0.72(2)_0^{+13}.`$
The last error in the figures above represents the uncertainty coming from the extrapolation (linear vs quadratic) in the inverse meson mass to the $`B`$ mesons. This error is not given in eqs. (5), because in the present study only a linear extrapolations have been made for reasons discussed below.
Besides the quantities relevant to $`B^0`$$`\overline{B}^0`$ mixing, we also give the corresponding quantities for $`D^0`$$`\overline{D}^0`$ mixing. Although $`D^0`$$`\overline{D}^0`$ mixing is expected to be well below the experimental limit in the Standard Model, it may be enhanced in its extensions . We have obtained
$`f_D\sqrt{\widehat{B}_D}`$ $`=`$ $`230(14)_8^{+3}\mathrm{MeV},f_D=207(11)_{00}^{+3+3}\mathrm{MeV},`$ (13)
$`f_{D_s}`$ $`=`$ $`234(9)_{00}^{+3+2}\mathrm{MeV},{\displaystyle \frac{f_{D_s}}{f_D}}=1.13(3)_1^{+0}.`$
Reviews of recent results of quantities considered in this paper can be found in refs. .
The remainder of this paper is as follows: in sec. 2 we give the parameters of the lattice simulation, illustrate the calibration of the lattice spacing and of the quark masses and describe the calculation of the heavy meson spectrum and decay constants; in sec. 3 we discuss the renormalization of the relevant operators, the calculation of their matrix elements and the extrapolation to the physical points; in sec. 4 we give the physical results and discuss the statistical and systematic errors; sec. 5 contains a comparison of our results with other calculations of the same quantities as well as our conclusions.
## 2 Lattice Calibration and Decay Constants
In this section we give some details about our lattice simulation and the calculation of the decay constants, which are a byproduct of this study.
### 2.1 Lattice Setup
The results presented in this study have been obtained on a $`24^3\times 48`$ lattice, using the non-perturbatively improved Clover action at $`\beta =6.2`$ with the Clover coefficient $`c_{_{SW}}=1.614`$, as computed in ref. . The statistical sample consists of 200 independent gauge-field configurations. Statistical errors have been estimated by using the familiar jackknife procedure with 40 jacks, each obtained by decimating 5 configurations from the whole ensemble. The following values for the heavy- and light-quark hopping parameters, $`\kappa _Q`$ and $`\kappa _q`$ respectively, have been used:
* 0.1352 ($`\kappa _{q_1}`$); 0.1349 ($`\kappa _{q_2}`$); 0.1344 ($`\kappa _{q_3}`$),
* 0.1250 ($`\kappa _{Q_1}`$); 0.1220 ($`\kappa _{Q_2}`$); 0.1190 ($`\kappa _{Q_3}`$) .
Although we are not discussing light mesons, it is important to mention some results which will be useful in the analysis of heavy-light meson physics, for example for the extrapolation/interpolation in light-quark masses. More details on the procedures used to calibrate the lattice spacing and fix the light-quark masses can be found in previous publications of our group .
* The ratio of the masses of the light pseudoscalar (P) and vector (V) mesons are:
$`{\displaystyle \frac{m_P}{m_V}}=\{0.597(22)_{q_1},0.682(15)_{q_2},0.761(8)_{q_3}\}.`$ (14)
* We use the method of physical lattice planes to fix the value of the inverse lattice spacing <sup>1</sup><sup>1</sup>1 As in previous publications, we use small-case letters for referring to quantities in physical units (e.g. $`m_P`$ in MeV), and capital letters for denoting the same quantities in lattice units (e.g. $`M_P=m_Pa`$).. This method consists in the following. First we fit the vector meson mass to the form
$`M_V(m_q,m_q)=\alpha _0+\alpha _1\left(M_P(m_q,m_q)\right)^2,`$ (15)
obtaining $`\alpha _0=0.286(16)`$, and $`\alpha _1=1.24(13)`$. Then we fix $`m_V/m_P`$ to the physical kaon ratio ($`m_K^{}/m_K`$) and compare $`M_K^{}`$ to the experimentally measured $`m_K^{}=894`$ MeV. In this way we obtain
$`a^1=2.72(13)\mathrm{GeV}.`$ (16)
This is the value of the inverse lattice spacing which has been used throughout this study.
* For a generic physical quantity in the heavy-light meson sector, $`(m_Q,m_q)`$, the interpolation/extrapolation in the light quark to the strange/up-down ($`s/d`$) mass is performed through the fit
$`(m_Q,m_q)=\alpha _0^Q+\alpha _1^Q\left(M_P(m_q,m_q)\right)^2,`$ (17)
where $`\alpha _{0,1}^Q`$ are the fitting parameters. The light-quark mass corresponds either to the $`d`$-quark, when we extrapolate to $`M_P(m_d,m_d)M_\pi =5.1(3)10^2`$ (as obtained from the lattice-plane method from $`m_\rho /m_\pi `$), or to the strange quark, when we extrapolate to $`M_P(m_s,m_s)M_{\eta _{ss}}=0.266(17)`$ (as inferred from eq. (15) by fixing $`M_V(m_s,m_s)`$ to $`m_\varphi =1020`$ MeV).
### 2.2 Heavy-light Decay Constants
We now discuss the heavy-light meson decay constants, which are important ingredients for physical predictions related to $`B^0`$$`\overline{B}^0`$ mixing. Since this has been extensively discussed in the literature (see for example ref. ), we only recall here the essential steps.
* As usual, hadron masses are extracted by fitting two-point correlation functions. For mesons, the standard form is
$`๐_{_{JJ}}(t)={\displaystyle \underset{\stackrel{}{x}}{}}0|J(\stackrel{}{x},t)J^{}(0)|0\stackrel{t0}{}{\displaystyle \frac{๐ต_J}{M_J}}e^{M_JT/2}\mathrm{cosh}\left[M_J\left({\displaystyle \frac{T}{2}}t\right)\right],`$ (18)
where for $`J(x)=P_5(x)=i\overline{Q}(x)\gamma _5q(x)`$, we choose $`t[16,23]`$, whereas for $`J(x)=V_i(x)=\overline{Q}(x)\gamma _iq(x)`$, we choose $`t[19,23]`$. The time intervals are established in the standard way (i.e. after inspection of the corresponding effective masses). The resulting masses of the pseudoscalar and vector heavy-light mesons, with the light quark interpolated/extrapolated to $`s/d`$ (see eq. (17)), are listed in tab. 1. We also checked that on the same interval ($`t[16,23]`$), the masses of pseudoscalar mesons extracted from the correlation function $`๐_{AP}(t)`$ ($`A_0(x)=\overline{Q}(x)\gamma _0\gamma _5q(x)`$), are indeed the same.
In physical units, the masses directly accessed from our simulation, are:
$`m_{P_d}=\{1.75(8)\mathrm{GeV},2.02(9)\mathrm{GeV},2.26(11)\mathrm{GeV}\},`$ (19)
(20)
$`m_{P_s}=\{1.85(7)\mathrm{GeV},2.11(9)\mathrm{GeV},2.38(10)\mathrm{GeV}\}.`$ (21)
* The pseudoscalar meson (P) decay constant is defined as
$`0|\widehat{A}_0|P(\stackrel{}{p}=0)=i\widehat{F}_PM_P,`$ (22)
where the โhatโ symbol means that the appropriate renormalization constant $`Z_A`$ is taken into account. To determine the decay constant, one computes the ratio
$`{\displaystyle \frac{\underset{\stackrel{}{x}}{}\widehat{A}_0(\stackrel{}{x},t)P_5(0)}{_\stackrel{}{x}P_5(\stackrel{}{x},t)P_5(0)}}`$ $``$ $`\widehat{F}_P{\displaystyle \frac{M_P}{\sqrt{๐ต}_P}}\mathrm{tanh}\left(M_P({\displaystyle \frac{T}{2}}t)\right).`$ (23)
When working with improved Wilson fermions, the axial current is improved by adding to $`A_\mu `$ the operator $`_\mu P`$ with a suitable coefficient
$`0|\widehat{A}_0|P=Z_A\left(0|A_0|P+c_A0|a_0P_5|P\right)=iM_P(\widehat{F}_P^{(0)}+c_Aa\widehat{F}_P^{(1)}),`$ (24)
where the constant, $`c_A=0.037`$, was computed in ref. . It is actually easy to see that
$$aF_P^{(1)}=\frac{\sqrt{๐ต_P}}{M_P}\mathrm{sinh}(M_P).$$
(25)
The size of the correcting term $`c_Aa\widehat{F}_P^{(1)}/\widehat{F}_P^{(0)}`$ for our data is in the range $`(5รท7)\%`$.
* As far as the normalization constant is concerned, its improved value is given by (choice 1.)
$`Z_A(m_Q,m_q)=Z_A(0)(1+b_Aa\overline{m}),`$ (26)
where $`\overline{m}=(m_Q+m_q)/2`$. In the above equation, we have taken
$`am_{Q,q}={\displaystyle \frac{1}{2}}\left({\displaystyle \frac{1}{\kappa _{Q,q}}}{\displaystyle \frac{1}{\kappa _c}}\right).`$ (27)
The value of the critical hopping parameter, $`\kappa _c=0.13580(1)`$, is obtained from the linear fit $`M_P^2m_q0`$, whereas the value of improvement coefficient $`b_A=1.24`$ in (28) is taken from boosted perturbation theory <sup>2</sup><sup>2</sup>2 By working with three light quark species only, we were unable to make a quadratic fit in the quark masses as done in ref. . For a recent non-perturbative determination of $`b_A`$, see ref. . .
In our numerical estimates we have also used
$`Z_A(m_Q,m_q)=Z_A(0)\left[{\displaystyle \frac{\sqrt{1+am_Q}\sqrt{1+am_q}}{1+a\overline{m}}}\right]\left(1+b_Aa\overline{m}\right),`$ (28)
This equation (choice 2.) differs from the choice 1. in that it contains higher order tree-level mass corrections through the so-called KLM factor . As for the value of $`Z_AZ_A(0)`$ in the chiral limit, we take the nonperturbative determination $`Z_A(0)=0.80`$ from ref . This value has been also non-perturbatively computed in refs. ($`Z_A=0.82`$) and ($`Z_A(0)=0.81`$), using different approaches. The values of the lattice decay constants $`\widehat{F}_P`$ for different heavy-quark masses are listed in the third column of tab. 1.
To estimate the values of the physical quantities related to $`B`$-mesons ($`m_{B_d}=5.28`$ GeV, and $`m_{B_s}=5.38`$ GeV), the extrapolations to the physical heavy-quark masses necessarily rely on the heavy-quark symmetry. Accordingly, the decay constants scale as
$`f_P={\displaystyle \frac{\mathrm{\Phi }(m_P)}{\sqrt{m_P}}}\left(1+{\displaystyle \frac{\mathrm{\Phi }^{}(m_P)}{\mathrm{\Phi }(m_P)}}{\displaystyle \frac{1}{m_P}}+\mathrm{}\right).`$ (29)
Apart from an overall (mild) logarithmic correction, the coefficients $`\mathrm{\Phi }(m_P)`$ and $`\mathrm{\Phi }^{}(m_P)`$, are non-perturbative quantities which we extract from a fit to the lattice data. Since we work with three heavy quarks (three values of $`\kappa _Q`$), we have only included the leading $`1/m_P`$ correction by fitting the decay constant to the expression
$`\widehat{F}_P\sqrt{M_P}=\mathrm{\Phi }_0\left(1+{\displaystyle \frac{\mathrm{\Phi }_1}{M_P}}\right).`$ (30)
In tab. 2, we give the physical results, and the values of parameters $`\mathrm{\Phi }_{0,1}`$, in physical units.
The statistical errors are given in parentheses. The systematic errors require some explanation: the first one is obtained from the differences between results obtained by using $`Z_A`$ with choice 1. or 2; the second is obtained by comparing the central values of the procedure described above to those extracted from the study of the ratio $`\widehat{F}_P/M_V`$.
We have not included in the systematic errors: $`i)`$ the uncertainty due to quenching. Pioneering attempts for estimating quenching errors show an increase of about $`20รท30`$ MeV of the value of the decay constants. $`ii)`$ the error due to the truncation of the $`1/m_P`$ expansion in eq. (30). The latter is usually estimated by considering also a quadratic term ($`1/m_P^2`$). With only three heavy-quark masses, we have not done such a fit. From our previous experience , we know that this error has a negligible effect on the value of the $`D`$-meson decay constants, whereas for the $`B`$ meson it represents a major source of systematics ($`30`$ MeV).
Finally, we present an estimate of $`f_{B_d}`$ obtained using the experimental measurement of $`f_{D_s}=241(32)`$ MeV . With the ratio $`f_B/f_{D_s}`$ obtained in our simulation, table 2, we get
$`\left({\displaystyle \frac{f_{B_d}}{f_{D_s}}}\right)^{(\mathrm{latt}.)}=0.74(4)_{00}^{+2+2}=0.74\pm 0.05`$ (31)
$``$ $`f_B=184\pm 24(\mathrm{exp}.)\pm 12(\mathrm{theo}.)\mathrm{MeV}`$ (33)
## 3 Matrix Elements of the Renormalized Operators
In this section, we discuss the renormalization of the relevant four-fermion operators and describe the calculation of the matrix elements introduced in eq. (2). Since we are interested to the $`B^0`$$`\overline{B}^0`$ mixing amplitude, only the parity-even part of the operator, $`Q=(\overline{b}\gamma _\mu q)(\overline{b}\gamma _\mu q)+(\overline{b}\gamma _\mu \gamma _5q)(\overline{b}\gamma _\mu \gamma _5q)`$, will be discussed in the following.
### 3.1 Renormalization of the Four-fermion Operators
On the lattice, the renormalized operator $`Q(\mu )`$ takes the form
$`Q(\mu )=Z(\mu ,g_0^2)\left(O_1+{\displaystyle \underset{i=2}{\overset{5}{}}}\mathrm{\Delta }_i(g_0^2)O_i\right),`$ (34)
where the $`O_i`$ denotes a bare operator and we choose to work in the following parity-even basis
$`O_1`$ $`=`$ $`\overline{b}\gamma _\mu q\overline{b}\gamma _\mu q+\overline{b}\gamma _\mu \gamma _5q\overline{b}\gamma _\mu \gamma _5q`$
$`O_2`$ $`=`$ $`\overline{b}\gamma _\mu q\overline{b}\gamma _\mu q\overline{b}\gamma _\mu \gamma _5q\overline{b}\gamma _\mu \gamma _5q`$
$`O_3`$ $`=`$ $`\overline{b}q\overline{b}q\overline{b}\gamma _5q\overline{b}\gamma _5q`$ (35)
$`O_4`$ $`=`$ $`\overline{b}q\overline{b}q+\overline{b}\gamma _5q\overline{b}\gamma _5q`$
$`O_5`$ $`=`$ $`\overline{b}\sigma _{\mu \nu }q\overline{b}\sigma _{\mu \nu }q.`$
To obtain the renormalized operator, two steps are then necessary. First, one has to subtract operators of the same dimension, which (on the lattice) mix with $`O_1`$ (this is the consequence of the presence of the Wilson term in the fermion action, which explicitly breaks chiral symmetry). This means that one has to compute the subtraction constants $`\mathrm{\Delta }_i(g_0^2)`$. After an appropriate subtraction of the operators $`O_i`$, with $`i1`$, the matrix elements still need (to be finite) an overall renormalization constant provided by $`Z(\mu ,g_0^2)`$. A detailed discussion of the mixing matrix can be found in ref. . In this work, we use the method of refs. to calculate the matching and the renormalization constants non-perturbatively. The computation, at three different values of the scale ($`\mu a=0.71`$, $`1.00`$, $`1.42`$), is performed in the RI-MOM renormalization scheme, in Landau gauge. The results are listed in table 3.
For sufficiently large $`\mu `$, we can use the available NLO anomalous dimension of the operator $`Q(\mu )^{\mathrm{RI}\mathrm{MOM}}`$ , to define the renormalization group invariant (RGI) operator:
$`\widehat{Q}=w^1(\mu )Q(\mu )\alpha _s(\mu )^{\gamma _0/2\beta _0}\left\{1+{\displaystyle \frac{\alpha _s(\mu )}{4\pi }}J\right\}Q(\mu ).`$ (36)
where
$`\gamma _0=4;`$ $`J_{\mathrm{RI}\mathrm{MOM}}={\displaystyle \frac{\mathrm{\hspace{0.33em}17397}2070n_f+104n_f^2}{6(332n_f)^2}}+8\mathrm{log}2.`$ (37)
We also give the NLO part relevant in the $`\overline{\mathrm{MS}}`$ scheme
$`J_{\overline{\mathrm{MS}}}={\displaystyle \frac{\mathrm{\hspace{0.33em}13095}1626n_f+8n_f^2}{6(332n_f)^2}}.`$ (38)
By using $`\alpha _s(m_Z)=0.118`$, and setting $`n_f=4`$, we obtain
$`\mu `$ $`=`$ $`\{1.9\mathrm{GeV},\mathrm{\hspace{0.33em}2.73}\mathrm{GeV},\mathrm{\hspace{0.33em}3.85}\mathrm{GeV}\},`$ (39)
$`w(\mu )_{\mathrm{RI}\mathrm{MOM}}^1`$ $`=`$ $`\{\mathrm{\hspace{0.33em}\hspace{0.33em}1.408},1.453,1.489\}.`$ (41)
In quenched calculations, there is always the embarrassing question whether to use the physical value of $`\alpha _s`$, with the complete formulae for the anomalous dimensions and the $`\beta `$-function, or a โquenched valueโ of the coupling constant (which suffers from intrinsic ambiguities), together with anomalous dimensions and the $`\beta `$-function computed for $`n_f=0`$. The present case is particularly lucky, since by computing $`w(\mu )`$ in the unquenched ($`n_f=4`$) and quenched cases (with $`\mathrm{\Lambda }_{\mathrm{QCD}}^{n_f=0}=250`$ MeV), one finds values of $`\widehat{B}_{B_d}`$ which differ by less than $`2\%`$.
We give in passing, $`w(m_b)_{\overline{\mathrm{MS}}}^1=1.477`$, which will be used later on for a comparison with the results of other authors. In the perturbative calculation of $`w(m_b)_{\overline{\mathrm{MS}}}^1`$, we have used $`m_b=4.25`$ GeV .
### 3.2 Matrix elements
We have computed the relevant three-point correlation functions:
$`๐_Q^{(3)}(\stackrel{}{p},\stackrel{}{q};t_{P_1},t_{P_2})={\displaystyle ๐\stackrel{}{x}๐\stackrel{}{y}e^{i(\stackrel{}{p}\stackrel{}{y}\stackrel{}{q}\stackrel{}{x})}0|P_5(\stackrel{}{x},t_{P_2})Q(\stackrel{}{0},0)P_5^{}(\stackrel{}{y},t_{P_1})|0}.`$ (42)
When the sources are sufficiently separated and far away from the origin, the lowest pseudoscalar mesons are isolated and one has
$`๐_Q^{(3)}(\stackrel{}{p},\stackrel{}{q};t_{P_1},t_{P_2})\stackrel{t_{P_1},t_{P_2}0}{}{\displaystyle \frac{\sqrt{๐ต}_P}{2E_P(\stackrel{}{q})}}e^{E_P(\stackrel{}{q})t_{P_1}}\overline{P}(\stackrel{}{q})|Q|P(\stackrel{}{p}){\displaystyle \frac{\sqrt{๐ต}_P}{2E_P(\stackrel{}{p})}}e^{E_P(\stackrel{}{p})t_{P_2}},`$ (43)
where $`๐ต_P0|P_5|P`$. In order to eliminate irrelevant factors and to extract the $`B`$ parameter we form the following ratio
$`(\stackrel{}{p},\stackrel{}{q};\mu )={\displaystyle \frac{3}{8}}{\displaystyle \frac{๐_Q^{(3)}(\stackrel{}{p},\stackrel{}{q};t_{P_1},t_{P_2};\mu )}{๐_{AP}^{(2)}(\stackrel{}{p};t_{P_2})๐_{AP}^{(2)}(\stackrel{}{q};t_{P_1})}},`$ (44)
where $`\mu `$ indicates the scale dependence. The two point correlation functions are both multiplied by the axial-current renormalization constant in the chiral limit $`Z_A(0)`$. We have not used the improved renormalization constant $`Z_A(m_Q,m_q)`$ since we have not attempted to improve the operator $`Q`$. This is beyond the scope of the present study and would demand the inclusion of many operators of dimension seven, whose coefficients are presently unknown. Our hope is that a part of the uncertainties of $`๐ช(\overline{m}a)`$ cancel in the ratio (44). This has been computed at $`|\stackrel{}{p}|=|\stackrel{}{q}|0`$, corresponding to
$`(\mu ){\displaystyle \frac{3}{8}}{\displaystyle \frac{\overline{P}|Q(\mu )|P}{|0|\widehat{A}_0|P|^2}}B_{P_q}(\mu ).`$ (45)
In this way we obtain the main results of this work. We also considered reasonably low momentum injections to the sources but the additional noise makes these data not useful in practice. For this reason these data will be ignored in the following.
In order to show the quality of the signal, in fig. 1 we show the ratio $`(\mu )`$ defined in eq. (44) at $`\mu =3.8`$ GeV ($`\mu a=1.42`$), for a specific combination of the heavy and light hopping parameters ($`\kappa _Q=\kappa _{Q_2}=0.1220`$ and $`\kappa _q=\kappa _{q_2}=0.1349`$). The ratio is plotted as a function of the time-distance of one of the two sources. The other has been fixed to $`t_{P_1}=16`$. We also performed the simulation by fixing $`t_{P_1}=12`$ and $`t_{P_1}=20`$. From the study of the two-point correlation functions, one isolates the lightest pseudoscalar state around $`t=16`$. That makes $`t_{P_1}=12`$ rather small for our purposes. For $`t_{P_1}=20`$, the sources are not far enough from each other because of the periodic boundary conditions, and this may spoil the signal.
After inspection of the ratios $`(\mu )`$ (for each combination of the heavy and light quarks), a stability plateau is found in the range $`28t33`$. From this plateau we extract the values of $`(\mu )`$, that is the $`B`$ parameters, for each combination of quark masses. At this point, we extrapolate the light quark to the $`s`$\- and $`d`$-quark mass by fitting the $`B`$ parameter to eq. (17). This is illustrated in fig. 2 for $`\kappa _{Q_2}=0.1220`$ at $`\mu =3.8`$ GeV. In tab. 4, we give a detailed list of results for the RI-MOM $`B_{P_q}(\mu )`$, as well as for the RGI $`\widehat{B}_{P_q}`$ (see eq. (36)).
It is interesting to check whether the scale dependence of the extracted $`B_{P_q}(\mu )`$ is well described by the NLO anomalous dimension (eq.( 36)). In fig. 3, we plot the evolution of $`B_{P_q}^{\mathrm{RI}\mathrm{MOM}}(\mu )`$. The three curves (dashed, full and dotted) have been normalized to the $`B`$ parameter at the scale $`\mu a=0.71`$, $`1.01`$ and $`1.42`$ respectively, as given in table 4. We note that the point at $`\mu a=1.01`$ is sensibly lower than what expected on the basis of the perturbative evolution. This effect is a consequence of a fluctuation of the value of $`Z(\mu )`$ at this value of the scale as shown in table 3. We checked that the fluctuation is enhanced by the extrapolation of $`Z(\mu )`$ to the chiral limit and consider it as part of the uncertainty in the determination of the matrix element.
The central values of $`\widehat{B}_{P_q}`$ quoted below correspond to the result obtained converting the RI-MOM $`B`$ parameters at $`\mu a=1.42`$ to the RGI one. The difference with the results obtained by using the other two scales, namely $`\mu a=1`$ and $`\mu a=0.7`$, is included in the estimate of the systematic uncertainty.
## 4 Physical results
In this section, we derive the physical amplitudes from the $`B`$ parameters computed at fixed $`m_Q`$. Before discussing the extrapolation in the heavy-quark mass, let us summarize our results. We have computed the matrix elements of the operator $`Q(\mu )`$ in the RI-MOM scheme at three different scales $`\mu `$. The dependence on light quark is rather smooth, and we have easily performed the extrapolation/interpolation to the $`d`$ and $`s`$ quarks. By using the NLO evolution, we have checked the stability of our results. For each value of the heavy-quark mass, from the $`B`$ parameter in the RI-MOM scheme we have computed $`\widehat{B}_{P_q}`$ by using the result at $`\mu 3.8`$ GeV. This is the only step where perturbation theory comes in the game. At this point we dispose of the following values:
$`\kappa _Q`$ $`=`$ $`\{0.1250,0.1220,0.1190\}`$ (46)
$`\widehat{B}_{P_s}`$ $`=`$ $`\{1.24(2)_{0.09}^{+0.00},\mathrm{\hspace{0.33em}\hspace{0.33em}1.26}(2)_{0.09}^{+0.00},\mathrm{\hspace{0.33em}\hspace{0.33em}1.27}(2)_{0.09}^{+0.00}\}`$ (48)
$`\widehat{B}_{P_d}`$ $`=`$ $`\{1.22(3)_{0.09}^{+0.00},\mathrm{\hspace{0.33em}\hspace{0.33em}1.26}(4)_{0.09}^{+0.00},\mathrm{\hspace{0.33em}\hspace{0.33em}1.27}(4)_{0.09}^{+0.00}\}`$ (50)
In order to extrapolate the above results ito the $`b`$ quark, we rely on the heavy-quark symmetry as in the case of decay constants and fit the $`\widehat{B}`$ parameters to the expression
$`\widehat{B}_{P_q}=c_0^{(q)}\left(1+{\displaystyle \frac{c_1^{(q)}}{M_P}}\right).`$ (51)
The result is shown in fig. 4, and the main numbers are listed in tab. 5. We stress again that we are not able to include terms of order $`1/m_P^2`$ in our fit, because we work with three values of the heavy-quark mass only. Note, however, that our results for the slope $`c_1^{(q)}`$ are much smaller than the predictions of ref. , where the heavy quark was treated non-relativistically. We will comment on this in the next section.
These results, when combined with those in tab. 2, give the following final numbers
$`f_{B_d}\sqrt{\widehat{B}_{B_d}}=206(28)(7)\mathrm{MeV},f_{B_s}\sqrt{\widehat{B}_{B_s}}=237(18)(8)\mathrm{MeV};`$ (52)
(53)
$`\xi {\displaystyle \frac{f_{B_s}}{f_{B_d}}}\sqrt{{\displaystyle \frac{\widehat{B}_{B_s}}{\widehat{B}_{B_d}}}}=1.16(7),r_{sd}=\xi ^2\left({\displaystyle \frac{m_{B_s}^2}{m_{B_d}^2}}\right)^{(\mathrm{exp}.)}=1.40(18).`$ (54)
In practice, for the reasons discussed in the introduction, we have extracted the value of $`f_{B_d}\sqrt{\widehat{B}_{B_d}}`$, as well as the other quantities appearing in the above equation, directly from the calculation of the mixing amplitude.
For completeness, and for comparison with other calculations, we also give the $`B`$ parameter in the $`\overline{\mathrm{MS}}`$ scheme
$`\widehat{B}_{B_d}^{\overline{\mathrm{MS}}}(m_b)`$ $`=`$ $`0.93(8)_{0.6}^{+0.0},`$
$`\widehat{B}_{B_s}^{\overline{\mathrm{MS}}}(m_b)`$ $`=`$ $`0.92(3)_{0.6}^{+0.0}.`$ (55)
## 5 Comparison with Other Calculations and Conclusions
In this section we compare our results with those obtained with different lattice calculations or with other methods.
On the lattice, two approaches have been used so far to compute $`B^0`$$`\overline{B}^0`$ mixing: the direct one (as done in this work) which extrapolates in the heavy-quark mass and the one where the heavy quark is treated with an effective theory. In the latter, either the Heavy Quark Effective Theory (HQET) or Non-Relativistic QCD (NRQCD) were used. There is a large number of lattice studies in which the HQET was invoked. Unfortunately, most of them were plagued by a mistake in the calculation of the renormalization constant. Only recently, this has been corrected in ref. , and the data of refs. have been reanalyzed <sup>3</sup><sup>3</sup>3This is the reason why we quote results from these three papers only.. On the other hand, very recently, the NRQCD treatment of the heavy quark has been employed to compute $`B_{B_q}`$ . The authors also calculated the relevant $`๐ช(1/m_P)`$ corrections. The slope in $`1/m_P`$ that they observe for $`B_{B_d}^{\overline{\mathrm{MS}}}(m_b)`$ is roughly a factor of $`3`$ larger than ours (or the one by ref. ). We stress that there is a general argument which shows that the uncertainty on the $`1/m_P`$ corrections can be easily as large as the corrections themselves unless perturbation theory on the leading term is not pushed to very high orders . This argument found confirmation by the explicit calculation of the $`๐ช(\alpha _s^2)`$ perturbative corrections to the heavy-quark mass in the HQET . Moreover, even the lowest order result of ref. in $`1/m_P`$ is suspicious, since they combine the renormalization constants computed in the HQET with the computation of the matrix elements of NRQCD. Until the renormalization of the relevant operators will not be completed in a consistent theoretical framework, the results of ref. should not be used for comparisons with other approaches.
In the direct approaches, a full relativistic treatment is given to both heavy and light quarks . In this case a major source of systematic error comes from the extrapolation of the results to the physical $`B`$ mesons. Recent results for $`B^0`$$`\overline{B}^0`$ mixing can be found in refs. and . In the second paper, an extrapolation to the continuum ($`a0`$) obtained using non-improved Wilson fermions, has been attempted. They obtain a large central value for the ratio $`r_{sd}`$ ($`r_{sd}=1.76(10)_{42}^{+57}`$), albeit with large errors. The central value is difficult to reconcile with the value of $`\xi `$ and of the decay constants produced by the same authors with the same set of data. In the present work and the one by UKQCD , improved fermions are used but without extrapolation to the continuum. The differences between the present study and that of is that in our case the action is improved non-perturbatively whereas UKQCD uses the mean-field improved action, and that we renormalize non-perturbatively the relevant four-fermion operators. The overall agreement is excellent. This is true also at the values of the heavy-quark masses were lattice measurements have been actually made as shown in fig. 5. Their extrapolated value has an error smaller than ours, probably because they also use a point corresponding to a small value of $`m_P`$, for which the heavy-quark expansion may be questionable.
In table 6, for comparison, we list several recent results of the direct and the HQET approaches.
For completeness, let us mention that also QCD sum rules were employed to compute $`B_{B_d}`$ . Radiative corrections were also included in ref. , where the value for $`B_{B_d}^{\overline{\mathrm{MS}}}(m_b)`$ was found to be compatible with 1 (VSA value). The most recent QCD sum rule estimate has been given by Chernyak , who quotes a value lower than ours and closer to the results obtained from the HQET at lowest order in $`1/m_P`$: $`B_{B_d}^{\overline{\mathrm{MS}}}(5.04\mathrm{GeV})=0.82`$, corresponding to $`\widehat{B}_{B_d}=1.26`$. The effect of SU(3) breaking in $`B_{B_s}/B_{B_d}`$, was studied in the framework of the HQET, combined with chiral perturbation theory . Their result suggests that SU(3) breaking for the $`B_B`$ parameter is practically negligible as confirmed by many lattice calculations. The approach has been extended to quenched chiral perturbation theory , where a pessimistic estimate of quenching errors ($`10`$$`20\%`$ for the matrix element) was quoted.
## Acknowledgements
We warmly thank E. Franco and D. Lin for informative discussions on the subject of this paper. V.L. and G.M. acknowledge MURST for partial support. D.B. thanks INFN for support. V. G. has been supported by CICYT under Grant AEN-96-1718, by DGESIC under Grant PB97-1261 and by the Generalitat Valenciana under Grant GV98-01-80. L. G. has been supported in part under DOE Grant DE-FG02-91ER40676. |
warning/0002/cond-mat0002158.html | ar5iv | text | # Crossover from classical to random-field critical exponents in As-doped TbVO4
## Abstract
Using birefringence techniques we have measured the critical exponents $`\beta `$, $`\gamma `$, and $`\delta `$ in As-doped TbVO<sub>4</sub>, a structural realization of the random-field Ising model where random strain fields are introduced by V-As size mismatch. For pure TbVO<sub>4</sub> we observe the expected classical critical exponents, while for a mixed sample with 15% As concentration our results are $`\beta =0.31\pm 0.03`$, $`\gamma =1.22\pm 0.07`$ and $`\delta =4.2\pm 0.7`$. These values are consistent with the critical exponents for the short range pure Ising model in three dimensions in agreement with a prediction by Toh. The susceptibility data showed a crossover with temperature from classical to random field critical behaviour.
Determination of the critical properties of the random field Ising model (RFIM) has shown encouraging progress recently, both experimentally and theoretically, after many years of uncertainties. Measurements on dilute antiferromagnets in a field, the most frequently studied realization of the RFIM, have been difficult to obtain because dilution of the magnetic species inhibits equilibration close to the critical temperature. However recent experiments by Slanic et al. on samples with only 7% dilution have shown equilibrium behaviour through the transition temperature, thus allowing confident determination of the critical exponents. In the last few years, theoretical investigations making use of a variety of analytical and computational techniques have led to predictions of the critical exponents of the RFIM that show reasonable consistency with each other. A recent analysis by Fortin and Holdsworth supports earlier suggestions that the critical exponents for the RFIM in three dimensions ($`d=3`$) are those of the Ising model for reduced dimension $`d^{}=1.5`$. Nevertheless dimensional reduction has not been rigorously proved, and even if the effective dimension $`d^{}=1.5`$ is correct the uncertainty in some of the calculated critical exponents is quite large. Likewise the experimental situation is still not satisfactory, since the critical exponents measured by Slanic et al. have substantial uncertainties and are only partially consistent with theory. In addition, measurements of specific heat critical exponents in dilute antiferromagnets also appear to disagree with theory.
For another realization of the RFIM, where random strain fields are generated by substitutional impurities in crystals undergoing structural Ising (Jahn-Teller) transitions, the results to date have also not been conclusive. Random fields due to As/V substitutions in DyVO<sub>4</sub>, which has $`d=3`$ Ising exponents, appeared to increase the susceptibility exponent $`\gamma `$ as expected but had no effect on the order parameter exponent $`\beta `$. The interpretation of the effects of random fields in the As-doped DyVO<sub>4</sub> system is complicated by the fact that the true critical behaviour of pure DyVO<sub>4</sub> should be classical due to the long range strain coupling. However because of the relative weakness of the long range to short range interactions, classical exponents are not observable at accessible temperatures $`|t|10^2`$, where $`t=(TT_\mathrm{D})/T_\mathrm{D}`$ is the reduced temperature, leading to uncertainty on what the effects of the random fields would be. We have therefore extended these experiments to the related TbVO<sub>4</sub>/TbAsO<sub>4</sub> system where the critical behaviour of the pure compounds is unequivocally classical, and searched for changes in critical behaviour in mixed crystals due to random fields. Since this system starts from a different universality class, the results cannot be compared directly with results from dilute antiferromagnets, but it is an important system that can independently test theoretical models and predictions of random field effects. In contrast with the large number of theoretical investigations of the short-range RFIM, the literature on the random-strain version of the RFIM is very limited, consisting primarily of a paper by Toh. In this paper, Toh compares the random-strain RFIM with long range forces to that of the short-range RFIM under a renormalization group analysis. The main result is that the critical exponents should change from classical values to values that are close to those of the $`d=3`$ pure short-range Ising model.
The lowest $`4f`$ electronic levels in TbVO<sub>4</sub> consists of 2 singlet states $`18`$ cm<sup>-1</sup> apart and a non-Kramersโ doublet approximately halfway between. Coupling between the doublets and lattice distortions leads to a Tb ion-ion interaction of the Ising form and a tetragonal-orthorhombic phase transition at temperature T<sub>D</sub>. Since the Tb ions are coupled predominantly to $`k0`$ acoustic phonons and to bulk strains, the ion-ion interaction is very long range. The order parameter is the macroscopic strain $`ab`$ where $`a`$ and $`b`$ are basal plane unit cell parameters in the orthorhombic phase. The orthorhombic distortion gives rise to birefringence, $`\mathrm{\Delta }n`$, which is proportional to $`ab`$ (at least to a good approximation).
The full Hamiltonian for the coupled electron-phonon TbVO<sub>4</sub> system in an external magnetic field can be written as
$$H=\frac{1}{/}2\underset{ij}{}J_{\mathrm{๐๐}}\sigma _i^z\sigma _j^z\frac{1}{/}2ฯต\underset{i}{}(1+\tau _i^z)\sigma _i^x\text{B}\underset{i}{}m_i^x$$
(1)
where $`m_i^x=\frac{1}{/}4\text{g}\mu _\text{B}(1+\sigma _i^z)\tau _i^x`$ and $`\sigma ^z,\sigma ^x,\tau ^z\text{and }\tau ^x`$ are Pauli type operators. $`J_{ij}`$ describes the ion-ion interactions, $`2ฯต`$ is the high temperature splitting between the outer singlets and B is the magnetic field applied along the $`x`$ (or $`a`$) axis (i.e. along the 110 direction). The field B is able to induce an orthorhombic distortion because of the strongly anisotropic Tb magnetic moment in the orthorhombic phase. In the mean field approximation, a Landau expansion of the free energy, with $`ฯต=0`$, shows that B<sup>2</sup>/T is effectively an ordering field, and this is supported by the experimental data that follows. For B and $`ฯต`$ small, TbVO<sub>4</sub> is well described by an Ising model Hamiltonian. As the mode softening at the transition in this type of system is anisotropic, classical critical exponents are expected, and observed, rather than $`d=3`$ Ising exponents.
In the mixed compound, Tb(As<sub>x</sub>V<sub>1-x</sub>)O<sub>4</sub>, a fraction $`x`$ of the V atoms are replaced by As atoms, generating random, static strain fields, one component of which has the right symmetry to couple to the order parameter. For $`ฯต=\text{B}=0`$ the Hamiltonian has the form of the RFIM,
$$H=\frac{1}{/}2\underset{ij}{}J_{\mathrm{๐๐}}\sigma _i^z\sigma _j^z\underset{i}{}h_i\sigma _i^z$$
(2)
where $`h_i`$ is the random local strain field which is expected to have a Gaussian distribution about $`h=0`$. In his analysis of this type of random field system with anistropic mode softening, Toh predicts changes to the upper critical dimension, thus modifying the values of critical exponents at $`d=3`$.
Crystals of Tb(As<sub>x</sub>V<sub>1-x</sub>)O<sub>4</sub> with impurity concentrations of $`x=0`$ and 0.15 were prepared using the flux growth method at the University of Oxford. The crystals were cut and polished perpendicular to the $`c`$ axis, with thickness of about 1 mm. They were mounted in a strain-free manner with the $`c`$ axis horizontal in a helium optical cryostat. The crystal could be rotated about a vertical axis, allowing alignment of the $`c`$ axis parallel to the laser beam and one of the $`a`$ axes parallel to a horizontal magnetic field. A circular aperture in the sample holder of about 1 mm in diameter limited the sampled area to a small region of the crystal, and thus reduced effects due to any inhomogeneous composition, temperature and ordering field that may be present in the crystal.
The light source for birefringence experiments was a HeNe laser operating at 543.5 nm, a wavelength that should give reasonable birefringence in this crystal. Photoelastic modulation and lock-in detection allowed sensitive measurement of the phase shift $`\varphi `$ due to the orthorhombic distortion. Adjustment of a Babinet compensator ensured that the detector output was proportional to $`\varphi `$. In a typical experiment, the data acquisition system brought the sample to each desired temperature, waited up to 15 minutes for equilibration, and then ramped the magnetic field up and down while recording field, light intensity, and temperature.
Below T<sub>D</sub>, the orthorhombic distortions are equally likely to be in the 110 or 1$`\overline{1}`$0 direction, leading to the formation of twinned orthorhombic domains separated by domain walls. To avoid the cancellation of birefringence due to multiple domains, the crystal is forced into a single domain by applying a magnetic field that favours the distortion axis parallel to the magnetic field. Thus to determine the birefringence, and hence $`\beta `$, below T<sub>D</sub>, we recorded data over a range of magnetic fields and extrapolated the high magnetic field data back to zero field. With the rather small magnetic fields available ($`<0.35`$ T) it was sometimes difficult to achieve a single domain at temperatures close to T<sub>D</sub> , at $`|`$T - T$`{}_{\mathrm{D}}{}^{}|<0.2`$ K for $`x=0.15`$, but this became progressively easier further away from T<sub>D</sub>. For the $`x=0`$ sample where pinning is presumably weaker, single domain structure at $`|`$T - T$`{}_{\mathrm{D}}{}^{}|<0.75`$ K was relatively easy to achieve. For this reason, the critical isotherm exponent, $`\delta `$, obtained from the dependence of the induced birefringence on ordering field at the transition temperature, is more reliable for pure TbVO<sub>4</sub> than that for the mixed sample. Above T<sub>D</sub>, the change in induced birefringence resulting from a change in the ordering field gives the susceptibility and hence $`\gamma `$. Although the exponents for pure TbVO<sub>4</sub> are known to be classical,, their measurement provides a comparison with the mixed sample where modified critical exponents are expected. In experiments, T<sub>D</sub> appeared to vary slightly from run to run because of effects such as mounting strains and temperature gradients. To minimize the effects of a variable T<sub>D</sub> on the results, efforts were made to measure all exponents of one particular sample on the same run. In the power-law fits to the data, T<sub>D</sub> was chosen to optimize both the $`\beta `$ and $`\gamma `$ fits simultaneously. The consistency in the results of repeated experiments and the quality of the fits give confidence in the results.
The birefringence with increasing and decreasing ordering fields for various temperatures below T<sub>D</sub> in the $`x=0.15`$ sample is shown in Fig. 1. Hysteresis observed in the low field region is attributed to pinning of the multidomain structure. The data from the higher field region where hysteresis is absent were extrapolated to determine the zero field birefringence. A log-log plot of both the $`x=0.15`$ and $`x=0`$ data, fitted to a power law of the form $`\mathrm{\Delta }n|t|^\beta `$ is shown in Fig. 2. Using T$`{}_{\mathrm{D}}{}^{}=29.26\pm 0.03`$ K and data in the range $`0.005<|t|<0.027`$, we obtained a value of $`\beta =0.31\pm 0.03`$ for the mixed sample. At larger $`|t|`$ we found no convincing evidence for crossover behaviour towards the classical exponent. For the pure sample, we obtained $`\beta =0.46\pm 0.06`$ for T$`{}_{\mathrm{D}}{}^{}=32.32\pm 0.04`$ K from data in the range $`0.004<|t|<0.037`$.
Figure 3 presents susceptibility data for selected temperatures above T<sub>D</sub>. At temperatures close to T<sub>D</sub>, these slopes deviate from linearity, attributable in part to the presence of non-zero birefringence at zero field, probably caused by internal strains. The susceptibility exponent $`\gamma `$ was found by fitting the various values of $`\chi (T)`$ to the relation $`\chi ^1|t|^\gamma `$. Figure 4 shows a log-log plot of the susceptibility $`\chi `$ versus the reduced temperature $`|t|`$ for $`x=0`$ and $`0.15`$ samples. As can be seen from Fig. 4, the log-log plot for $`x=0.15`$ shows two linear fits with the data closer to T<sub>D</sub> having a larger slope than that for data further away. Data in the range $`0.027>|t|>0.01`$ (closest to T<sub>D</sub>) were optimized first to a linear fit giving T$`{}_{\mathrm{D}}{}^{}=29.26\pm 0.03`$ K and $`\gamma =1.22\pm 0.07`$. A linear fit to data in the range $`0.100>|t|>0.021`$ (further from T<sub>D</sub>) with T<sub>D</sub> unchanged yields a value of $`\gamma =0.89\pm 0.03`$. We did not attempt to fit these data to a crossover function, but they suggest a crossover temperature near $`|t|=0.024`$ ($`29.96`$ K). The log-log plot of the susceptibility versus reduced temperature for pure TbVO<sub>4</sub> in Fig. 4 does not show a crossover effect. A power law fit in the range $`0.058>|t|>0.005`$ gives T$`{}_{\mathrm{D}}{}^{}=32.32\pm 0.04`$ K and $`\gamma =0.92\pm 0.07`$.
Nonlinearity in the dependence of the order parameter on the ordering field at temperatures below T<sub>D</sub> can be attributed to the progressive detwinning of the sample with increasing field (i. e. changing from a multidomain to a single domain structure). The critical isotherm exponents $`\delta `$ were extracted from the field-induced birefringence data at temperatures closest to that of the previously determined transition temperatures for $`x=0.15`$ and $`x=0`$ samples. Only the higher field $`\mathrm{\Delta }`$n data were fitted to the power law (B<sup>2</sup>/T$`{}_{\mathrm{D}}{}^{})^{1/\delta }`$. For the $`x=0.15`$ sample, we obtained a value of $`\delta =4.2\pm 0.7`$ at $`T=29.24\pm 0.03`$ K and for the pure sample $`\delta =2.6\pm 0.4`$ at $`T=32.33\pm 0.03`$ K. The uncertainty in locating T<sub>D</sub> is incorporated in the uncertainty in $`\delta `$.
Our values of the critical exponents for pure TbVO<sub>4</sub> agree satisfactorily with the values $`\beta =0.5,\gamma =1`$, and $`\delta =3`$ expected for a mean field system. For the random-field sample our values $`\beta =0.31\pm 0.03`$, $`\gamma =1.22\pm 0.07`$, and $`\delta =4.2\pm 0.7`$ are in good agreement with the exponents for the $`d=3`$ Ising model, $`\beta =0.33`$, $`\gamma =1.24`$, and $`\delta =4.8`$.
Tohโs predictions for the effects of random fields on the exponents in this type of system are therefore well supported by our results. Toh noted that while the anisotropic strain interactions in TbVO<sub>4</sub> reduce the upper critical dimension $`d^{}`$ from 4 for the standard Ising model to 2, resulting in classical exponents, the introduction of random strain fields raises $`d^{}`$ from 2 to 4 again, leaving the critical exponents the same as for the standard Ising model in $`d=3`$. In view of the focus in recent theoretical analysis on dimensional reduction it is of interest to examine Tohโs analysis and our results in terms of changes in effective dimension $`d^{}`$ instead of in $`d^{}`$. Thus for the Ising model we regard the upper critical dimension to be fixed at 4; in pure TbVO<sub>4</sub>, the anisotropic strain interactions presumably raise the effective dimension from 3 to 5, giving classical exponents as before. If changes in critical properties due to random fields can be described by dimensional reduction, the consensus prediction is that random fields reduce the effective dimensionality by $`2\eta `$, where $`\eta `$ is the exponent for the decay of magnetization correlations. In the present case, since $`\eta =0.032`$ for $`d=3`$, the dimensionality should be reduced by 1.97 or essentially 2, resulting in $`d^{}=3`$. Hence our results, within their accuracy, are consistent with dimensionality reduction by random fields of $`2\eta 2`$ for this system. |
warning/0002/cond-mat0002166.html | ar5iv | text | # Vesicle propulsion in haptotaxis : a local model
## I Introduction
Phospholipidic vesicles constitute a simple model of cytoplasmic membranes of real cells. A simple model due to Helfrich based on curvature energy has accounted for a variety of equilibrium shapes. The model is based on a minimal energy principle. Some of the shapes (the so-called discocytes) bear strong resemblance with that of an erythrocyte (the red blood cell). Additionally, analysis of flickering (temporal small fluctuations around a given shape) of an erythrocyte by Brochard and Lennon has been quite successfully described by the Helfrich model including hydrodynamics dissipation. The vesicle model has seemed then as a natural candidate, at least in a first stage, for dealing with more complex entities such as those encountered in the realm of biology. In that context, however, most of the features are of nonequilibrium dissipative nature. Very recently several theoretical and experimental investigations have been directed along that line.
We are interested here in vesicle migration, a question on which we have given recently a brief account. Despite the very complex biochemical behaviour of a cell, cells may also exhibit behaviours where simple physical concepts may be evoked. It is well documented that, for example, the migration of the pronephric duct cells in salamanders is regulated by haptotaxis. Haptotaxis is a terminology that is used to express the following fact: when adhesive molecules are present in increasing amounts along an extracellular matrix (or simply on a substrate in vitro), a cell that was constantly making and breaking adhesion with such a molecule would migrate from a region of low concentration to an area where that adhesive molecule was more highly concentrated. There are also evidences that cell migration during embryo development may be guided by an adhesion gradient. In other words cell migration is here guided by a purely external physical factor, while the internal structure (the cytoskeletton) is quite unaffected on the time scale of interest. This feature drastically differs from that of a cell belonging to the immune system where the cytoskeletton plays a decisive role. Despite the fact that the cytoskeletton in pure haptotaxis does not undergo a structural change as is the case during cell crawling, the problem remains very much involved since the cell cytoskeletton dissipation may come to the fore as well as an intricate bond breaking and restoring with the substrate. We shall consider here a pure vesicle moving in haptotaxis. Our belief is that advancement in this field can be achieved only by the progressive refinement of concepts.
We consider a vesicle moving along the substrate thanks to an adhesion gradient. As the vesicle moves, it generates hydrodynamics flow both inside and outside. Hydrodynamics induces nonlocal interactions leading to an effective coupling of two distinct regions on the vesicle. In addition, the two monolayers that form the phospholipidic membrane might slide one relative to the other. Finally during motion the vesicle forms new bonds ahead and destroys others behind, and this process of bond breaking and restoring may be so slow that it may dominate dynamics (see later).
This paper should be regarded as using a very simplistic view in the hope of introducing the concepts of migration and to exhibit in a transparent fashion the way the problem is addressed. We shall keep the description as simple as possible. That is to say: (i) we ignore nonlocality due to hydrodynamics โincorporation of hydrodynamics was briefly discussed in and will be the subject of a forthcoming paperโ, (ii) we confine ourselves to a 2D geometry. Some kind of dissipation due to bond breaking and restoring is introduced in our model. The adoption of a local model (no hydrodynamics) allows one to quite easily obtain analytical results and thus to extract some key ingredients about migration โespecially the role of the adhesion area (length in 2D)โ which turns out to be crucial also when hydrodynamics is included.
The scheme of this paper is organized as follows. In Section II we write down the equations of motion and comment them. In section III we present a forward time integration and present the main result. Section IV presents the solution of the stationary system in a form of a nonlinear eigenvalue problem, where the drift velocity is the eigenvalue. In section V we give an analytical solution. A conclusion and a discussion is presented in section VI.
## II Equation of motion
### A Parameterization
We consider an adhering vesicle, deposited on a flat substratum which is oriented by its normal vector $`\widehat{๐ฒ}`$ (Fig.1). The x-axis is along the wall and represents the direction of vesicle motion occurring by convention from left to right. As stated before we confine ourselves to two dimensions. That is to say, the vesicle morphology is invariant in the z-direction, similar to a tubular vesicle.
The interaction between the vesicle and the substrate is taken into account by introducing an adhesion potential. The range of the potential in realistic situations (typically several $`nm`$) is small in comparison to the vesicle size (several $`\mu m`$), so that it is justified in practice to consider a contact potential, unless otherwise indicated (see later). The energy interaction is then zero if $`y>0`$ and is non vanishing only close to contact (if $`y=0`$). At the junction point between the free part of the vesicle (whose length is denoted as $`L^{}`$) and the adhered part (with length $`L_{adh}`$), the potential undergoes an abrupt change. The contact between the membrane and substrate occurs at two well-defined points $`x_1`$ at the left and $`x_2`$ at the right. These parameters are related to the total length of the curve $`L`$ by $`L=L_{adh}+L^{}=(x_2x_1)+L^{}`$. We use an intrinsic representation of the vesicle contour by introducing $`\psi (s)`$, the angle between the outward normal and the y-axis, and $`s`$ the arc length, as shown on Fig. 1. We only need to consider the function $`\psi (s,t)`$ from $`s=0`$ to $`s=L^{}`$ corresponding to the contact points $`x_1`$ and $`x_2`$, respectively. Because of the contact potential character the adhesion length is completely fixed if the two contact points $`x_1`$ and $`x_2`$ are known. Thus the vesicle shape and its dynamical properties (like the propulsion velocity) are known if $`x_1`$, $`x_2`$ and the function $`\psi (s,t)`$ are determined. The demand that the parametrization of the vesicle be compatible with the adhesion on the substrate it fulfilled by the two geometrical constraints : (i) the distance between both contact points, $`x_2x_1`$, must coincide with the adhesion length $`L_{adh}=x_2x_1`$, (ii) their vertical coordinates $`y_1`$ and $`y_2`$ must have the same value, $`y_2y_1=0`$. These two constraints can be expressed in terms of $`\psi (s,t)`$. For that purpose we use the relations
$$\frac{x}{s}=\mathrm{cos}\psi ,\frac{y}{s}=\mathrm{sin}\psi ,$$
(1)
which allow us to write the two constraints in the following form :
$`{\displaystyle _0^L^{}}{\displaystyle \frac{x}{s}}๐s`$ $`=`$ $`{\displaystyle _0^L^{}}\mathrm{cos}\psi (s)๐s=x_2x_1=L_{adh},`$ (2)
$`{\displaystyle _0^L^{}}{\displaystyle \frac{y}{s}}๐s`$ $`=`$ $`{\displaystyle _0^L}\mathrm{sin}\psi (s)ds=y_2y_1=0.`$ (3)
These are the geometrical constraints. In order to describe vesicle dynamics, we need a dynamical equation for the evolution of $`\psi (s,t)`$. A movement of the vesicle (due for example to an adhesion gradient) is limited by dissipation (such as hydrodynamics etcโฆ). The vesicle reacts to any deviation from equilibrium by its internal forces (bending, possible stretching โor resistance to stretchingโ). Let us first discuss these forces.
### B Energy and forces
All the relevant membrane properties are summarized in the following energy, expressed in 2D, with the dimension of an energy per unit length :
$$E=_๐\kappa \frac{(cc_s)^2}{2}๐s_{x_1}^{x_2}w(x)๐x+_๐\zeta (s)๐s+pS.$$
(4)
The first term is the well known Helfrich curvature energy, with the rigidity $`\kappa `$, the curvature $`c=\psi /s`$ and the spontaneous curvature $`c_s`$ . Because of the 2D-specific conservation law $`c๐s=2\pi `$, any curve displacement leaves unchanged the energy terms associated to the spontaneous curvature. We can thus omit the term associated with $`c_s`$. The second term expresses the adhesion energy and is only integrated on the adhered part of the curve. As we are concerned with an inhomogeneous substratum, the contact potential depends on the variable $`x`$ and is denoted by $`w(x)`$ (with $`w>0`$, meaning that adhesion is favorable). Finally the last two terms ensure the length and surface conservation, respectively. The membrane is a two dimensional incompressible fluid. The phospholipid exchanges with the solvent is virtually absent, and the area per molecule on the vesicle remains constant. This leads to the local length conservation (in the 2D language). The variable $`\zeta `$ is a local Lagrange multiplier which enforces the arc length $`ds`$ to a constant value. The enclosed surface $`S`$ conservation is a consequence of the membrane impermeability and fluid incompressibility. It is ensured by the global Lagrange multiplier $`p`$. The interpretation of $`\zeta `$ and $`p`$ as a tension and a pressure will be discussed latter.
The functional derivative of the energy (4) provides us with the various forces acting on the membrane.
(i) Curvature forces
Under the assumption that the membrane is completely flat on the adhered part, we reduce the integration domain of the first energetic term in eq.(4) to $`[0,L^{}]`$. Making use of the relations $`๐ญ=๐ซ^{}`$ and $`๐ง=(1/c)๐ซ^{\prime \prime }`$ (the prime designates derivative with respect to $`s`$) which are the membrane tangential and normal unit vectors, we obtain for the curvature energy $`E_c`$
$$E_c=\frac{\kappa }{2}_0^L^{}\left(\frac{^2๐ซ}{s^2}\right)^2๐s.$$
(5)
When taking the functional derivative of $`E_c`$ with respect to the position, care must be taken. Indeed the arc length element must also undergo a variation. A convenient formulation avoiding confusion rests on the introduction of a general parametrization $`a[0,1]`$, related to $`s`$ by the metric $`g=(ds/da)^2`$, and is time-independent. The energy expression becomes
$$E_c=\frac{\kappa }{2}_0^1\left[\left(\frac{^2๐ซ}{a^2}\right)^2\frac{1}{g}\left(\frac{^2๐ซ}{a^2}\frac{๐ซ}{a}\right)^2\right]g^{3/2}๐a.$$
(6)
The functional derivative, though straightforward, may be too lengthy if one does not take care in regrouping adequately various terms as explained in . The result can be written in a simple form:
$$๐_c=\frac{\delta E_c}{\delta ๐ซ(s)}=\kappa \left(\frac{^2c}{s^2}+\frac{c^3}{2}\right)๐ง.$$
(7)
The curvature force is, as expected, free of any tangential contribution. The first term in eq.(7), involving the second derivative of the curvature, tends to keep curvature repartition as homogeneous as possible. It is also present in 3D under the more complicated form of the Laplace-Beltrami operator. The second term proportional to $`c^3`$ is in contrast 2D-specific. It tends to increase the size of any convex shape. Note that in 3D the curvature energy is scale invariant, which implies a vanishing curvature force of this type on a sphere. The difference between 2D and 3D can be explained in the following way. Let us consider a finite cylinder of length $`H`$ and radius $`RH`$, which constitutes a good approximation for an infinitely long cylinder. In order to make the cylinder โcloserโ to a sphere, which is the corresponding 3D equilibrium shape, the curvature force would tend to increase the radius and decrease the length so as to bring the cylinder shape as close as possible to a sphere. This gives an intuitive picture of the $`c^3`$ term in 2D.
In the discussion above we did omit the boundary contribution when taking the functional derivative. Since this point is a bit subtle we have postponed it to the end of this section.
(ii) Length and surface constraints
On the free part of the curve, the force which is associated to the first Lagrange multiplier $`\zeta `$ is obtained upon functional derivation of $`E_l=\zeta ๐s`$. The result is
$$๐_l=\delta E_l/\delta ๐ซ(s)=c\zeta (s)๐ง+\frac{d\zeta }{ds}๐ญ.$$
(8)
The normal component is easily identified as a Laplace pressure, whereas the tangential one looks like a Marangoni force (which is encountered when surface tension is inhomogeneous). However, $`\zeta `$ is not exactly similar to a surface tension as for an interface between two fluids. The โtensionโ $`\zeta `$ is not an intrinsic property of the membrane. It adapts itself to the other forces in order to maintain the local length fixed. In other words, the problem is implicitly written in a thermodynamical ensemble with fixed length. This differs from the usual problem for fluid or solid surfaces where the surface tension is fixed instead. Thus $`\zeta `$ is a variable that must be determined self consistently as a Lagrange multiplier, by use of the constraint equation (see appendix in Ref.) :
$$0=\frac{(ds)}{t}=\left(\frac{v_t}{s}+cv_n\right)ds.$$
(9)
This relation (9) simply expresses the condition of vanishing velocity divergence on the curved contour of the vesicle, which is precisely the incompressibility condition in the 2D fluid constituting the membrane (written here in one dimension). A more intuitive way of viewing expression (9) is presented on Fig.2. The Marangoni term is the only tangential term among all membrane forces (see eqs. 7,8,10). It is seen from (8) that the Lagrange multiplier must be uniform at equilibrium. For sake of simplicity and in order to get more insight into analytical understanding, a uniform value will be assigned to $`\zeta `$, even out-of-equilibrium. A discussion of this point will be presented in section VI. This assumption implies some consequences on dynamics (and especially on the tangential velocity) which will be presented in section II C.
Finally we have to consider the force associated to $`E_s=pS`$ :
$$๐_s=\delta E_s/\delta ๐ซ(s)=p๐ง.$$
(10)
The Lagrange multiplier $`p`$ depends only on time; it enforces a constant area. Two physical interpretations can be invoked depending on the situation under consideration. Either we consider an impermeable membrane, and $`p`$ would be the hydrostatic pressure difference between outside and inside; or we choose a model of permeable membrane and $`p`$ would play the role of an osmotic pressure. Both models are equivalent as long as we do not consider hydrodynamic flows.
(iii) Adhesion forces and boundary terms
The functional derivative induces boundary terms at each contact point. The additional variation $`\delta E_c^b`$ and $`\delta E_{w,l}^b`$ for the curvature, adhesion and tension energies, associated with a small displacement $`\delta ๐ซ`$ of the contact points is given by (see )
$`\delta E_c^b=\left[\delta \dot{๐ซ}\left({\displaystyle \frac{\kappa c}{\sqrt{g}}}๐ง\right)\right]_0^L^{}+\left[\delta ๐ซ\left(\kappa {\displaystyle \frac{c}{s}}๐ง{\displaystyle \frac{\kappa c^2}{2}}๐ญ\right)\right]_0^L^{},`$ (11)
$`\delta E_{w,l}^b=\left[\text{}\delta ๐ซ\left(\text{}\zeta ๐ญ+(\zeta w(x))\widehat{๐ฑ}\right)\right]_0^L^{}.`$ (12)
Following the definition of these boundary points, they remain on the substrate. Thus, the accessible values for $`\delta ๐ซ`$ is then reduced to $`\delta ๐ซ\widehat{๐ฑ}`$. Additionally, in order to keep the curvature energy finite, we impose a vanishing value for the contact angle $`\varphi `$ between the membrane and the substrate (see Fig.3). Within our formulation, this constraint does not follow from the energy minimization and has thus to be added into the physical model. More precisely, at the discontinuity point (say $`x_2`$) one has to add to the Helfrich energy a term of the form $`\kappa (\mathrm{\Delta }\psi /\mathrm{\Delta }s)^2`$ which informs us on how would the vesicle on the adhered part feels, so to speak, the behavior of the vesicle at the junction point on the right side. Across the contact point of a vanishing extent, $`\mathrm{\Delta }s0`$, while the angle, if it had to have another value than zero, would make a jump leading to an abnormally increasing curvature energy. We must then impose a vanishing contact angle. These various conditions (motion along the wall and a vanishing contact angle) lead to $`๐ง=\widehat{๐ฒ}`$, $`๐ญ=\widehat{๐ฑ}`$ and $`\delta \dot{๐ซ}\widehat{๐ฑ}`$. It follows then that the term proportional to $`\delta \dot{๐ซ}`$ in eq.(11) vanishes automatically. The second term becomes $`\kappa /2\left(c_2^2\delta x_2c_1^2\delta x_1\right)`$ with $`c_1`$ and $`c_2`$ the curvatures at the left and right contact points. These terms are counterbalanced by adhesion and tension terms (eq.(12)) leading to the relation
$$\frac{\delta E}{\delta x_i}=\left(\kappa \frac{c_i^2}{2}w(x_i)\right)\widehat{๐ฑ},$$
(13)
where the $``$ and $`+`$ signs refer to the rear and fore contact points represented by the subscript $`i=1,2`$. At equilibrium, we recover here the relation $`c=\sqrt{2w/\kappa }`$ .
The energy variation given by eq.(13) can not really be identified as a physical force. It corresponds indeed to a geometrical point displacement. The โforceโ orientation is here parallel to the substrate, whereas the real force acting on the contact point, considered as a material point, is expected to be normal to the substrate. As we have seen above the curvature forces are indeed normal when applied to a an adjacent piece of the membrane (see eq. 7). The present โforceโ has the meaning of how much energy would be involved in displacing the contact point from one position to another. That geometrical point is by its very nature sitting on the substrate, so that the โforceโ associated with its displacement is naturally tangential.
We find it worthwhile to make a short digression. Suppose that the angle is not fixed to zero as we did above. More precisely suppose that the rigidity is so small or the adhesion is so large (see below what does this mean) then the vesicle will be so tense that it would look like a droplet outside some length scale $`\mathrm{}`$ to be determined below (of course within that scale, which is sufficiently close to the substrate, the matching must be tangential). If we do not assume a value for the angle (that is no relation between $`๐ง`$ or $`๐ญ`$ with $`\widehat{๐ฑ}`$ and $`\widehat{๐ฒ}`$), and set $`\delta ๐ซ\widehat{๐ฑ}`$ we find from (13) that $`c=0`$ at the contact (which means a straight line at the contact) and that the angle between that line and the substrate obeys
$`w_i`$ $`=`$ $`\left(1cos(\varphi _i)\right)\zeta `$ (14)
$`\varphi _i`$ $``$ $`\sqrt{{\displaystyle \frac{2w_i}{\zeta }}}\text{ ( for small angles )}`$ (15)
which is nothing but the Young condition. We have neglected $`\kappa c/s`$ in comparison to $`w`$. The justification is as follows. $`\kappa c/s\kappa c_0/\mathrm{}`$, where $`c_0`$ is the true contact curvature given by $`\sqrt{2w/\kappa }`$. The approximation is legitimate provided that the length scale $`\mathrm{}\sqrt{\kappa /w}`$. The length $`\sqrt{\kappa /w}`$ is the radius of curvature at contact. If the scale of interest is outside that internal region, then the droplet limit is justified. It must be emphasized however that the effective contact angle is not an intrinsic property of the adhered membrane, as for a droplet, but it is linked to other parameters (rigidity, the vesicle scaleโon which depends $`\zeta `$โ, etcโฆ ). In particular, the tension $`\zeta `$ is fixed by the reduced volume, which is a global property of the vesicle : different vesicles of the same phospholipid composition, but with different sizes, may have different contact angle on the same substrate.
### C Dynamical equation
An important point which must be emphazised when dealing with dynamics is the identification of the dissipation sources. These are the following: (i) the dissipation in the membrane via molecule rotations (very much like liquid crystals where dissipation is characterized by the Leslie coefficient), (ii) hydrodynamics flows inside and outside the vesicle, (iii) friction between the monolayers, and (iv) bond breaking and restoring with the substrate. It is well known that dissipation associated with rotation (internal dissipation) is negligible in practice, and for free vesicles (no substrate) hydrodynamics seems to be the most important dissipation. Hydrodynamics induces nonlocal interactions and this will be dealt with extensively in a forthcoming paper. Our wish in this paper is to present a pedestrian model, namely a local one, which allows for a complete analytical solution that will help to identify some key ingredients in the migration process. A specific dissipation with the substrate will be introduced later. For the moment we confine our description to the free vesicle case. The local model to be presented here is similar to the so-called Rousse model in the community of polymers. Indeed, for a one dimensional contour in a three dimensional space dynamics becomes local even in the presence of hydrodynamics.
The best way to introduce the dynamical law is to consider a dissipation function, proportional at each point to the square of the velocity :
$$F_d=\frac{\eta }{2}|๐ฏ|^2๐s.$$
(16)
The coefficient $`\eta `$ is here an effective viscosity and has the dimension of a viscosity per unit length. Its numerical value is estimated by $`\eta =\eta _{wat}/R10^2kgm^2s^1`$, with $`\eta _{wat}`$ the water viscosity and $`R`$ a typical vesicle size.
Neglecting inertial terms, the Euler-Lagrange equations become then
$$\frac{\delta E}{\delta ๐ซ}=\frac{\delta F_d}{\delta ๐ฏ}\eta ๐ฏ=๐.$$
(17)
As expected, we find a local proportionality between the membrane velocity $`๐ฏ`$ and the membrane force $`๐`$, which is a nonlinear function of position. In the present picture where the effective tension $`\zeta `$ is space-independent no tangential force appears so that physics will only fix the normal velocity, while the tangential velocity has no physical meaning as described below.
(i) Normal velocity
The normal membrane force is (see eq. (7, 8, 10)) :
$$f_n=\kappa \left(\frac{^2c}{s^2}+\frac{c^3}{2}\right)c\zeta p.$$
(18)
From the dynamical law (17) and the membrane forces expression (18) we obtain the normal velocity as a function of $`\psi (s)`$ :
$$v_n(s)=f_n=\frac{\kappa }{\eta }\left[\frac{^3\psi }{s^3}+\frac{1}{2}\left(\frac{\psi }{s}\right)^3\frac{\zeta }{\kappa }\frac{\psi }{s}\frac{p}{\kappa }\right].$$
(19)
It is convenient to write the dynamical equation in terms of the angle $`\psi `$ and not the curvature $`c`$. The reason is that the boundary conditions are written naturally as a function of $`\psi `$ (tangential matching, $`\psi =\pm \pi `$, and contact curvature $`\psi /s=\sqrt{2w/\kappa }`$).
(ii) Tangential velocity
There is only one tangential contribution to the membrane forces, $`\zeta /s`$, which is zero with the assumption of a uniform tension (see eq. 8). This implies that only the total length is conserved, and not the local one. A dilatation of a part of the membrane is then permitted, as long as the remaining part of the vesicle is contracted in order to keep the total length unchanged (see Fig.4). Within this approximation, there is no energy variation associated to tangential motion, and therefore no forces. In other words we consider the vesicle contour as a mathematical curve, loosing the concept of density : only the shape matters, independently of the points distribution on the curve.
We could equivalently assume that only the normal velocity contributes to dissipation. In that case the dissipation function (16) would take the form $`F_{dn}=\eta /2|v_n|^2๐s`$ and the equation of motion (17) becomes $`๐=\eta v_n๐ง`$. The tangential force $`\zeta /s`$ must then vanish and we get automatically that $`\zeta =`$constant. Thus our assumption of a global Lagrange multiplier can also be viewed as the result of a dissipation due uniquely to normal displacements.
As we have already mentioned, tangential displacements do not induce a geometrical change. If the tangential velocity has no physical meaning (as is the case with a constant $`\zeta `$) its choice should not affect the physics. We are thus at liberty to choose one which is convenient (very much like a gauge-field invariant formulation in electrodynamics). The choice of a gauge is interpreted as a reparametrization of the curve. As seen below the tangential velocity is fixed by the normal velocity once the gauge is specified, but there is naturally no feed back of that tangential velocity on the normal one (the physical one). This is the crucial difference between this non physical tangential velocity and a physical one that may arise in the general case (as discussed in a forthcoming paper). It must be noted however, that whether a tangential velocity is of physical nature or not, the knowledge of the normal velocity is sufficient to describe vesicle dynamics. It is thus only via its influence on the normal velocity that a physical tangential velocity would affect the physics (see below). Still in that case we can introduce a second tangential velocity of geometrical nature that corresponds to the displacement of the representative points of the curve and not to the material ones which are affected by the physical tangential velocity.
In the present model the most convenient parametrization requires a homogeneous points distribution along the free part of the curve, which is expressed as $`d/dt(s(a)/L^{})=0`$. This provides the expression for the โnon physical tangential velocityโ (see appendix in Ref.) :
$$v_t(s)=v_t(0)_0^scv_n๐s+\frac{s}{L^{}}\left(_0^L^{}cv_n๐s+v_t(L^{})v_t(0)\right).$$
(20)
If the free length $`L^{}`$ remains constant during the motion, as happens for a stationary regime, eq.(20) fixing the gauge imposes nothing but a constant distance between two consecutive points on the vesicle. The local length conservation (9), which is physical, seems then to be implied by a gauge!. In reality, once we have adopted a contant tension โimplying a vanishing tangential physical velocityโ, any point distribution is of purely geometrical nature, and we could impose another gauge than the above one, without affecting the physics; the above tangential velocity does not act on the normal velocity. Had we considereed $`\zeta `$ to be non contant, we would then have obtained a physical tangential velocity, which would act on the normal one; use of (9) fixes $`\zeta (s)`$ which in turn acts on the value of $`v_n`$, and then on physics. In the simplistic model we adopt, the tangential velocity is determined a posteriori, independently of the normal velocity. That is why it is only a non physical reparametrization, a โgaugeโ. A remark is in order : in this situation the question of a rolling or sliding motion does not make sense, since both motions differ only by a tangential velocity.
The membrane velocity is given as a function of $`\psi (a)`$, $`a`$ being the auxiliary parametrization of the free part of the vesicle, running from one contact point to the other. In order to obtain a closed system we need a relation between the evolution of $`\psi (a)`$ and the velocities. The temporal derivative of $`\psi `$, for a given $`a`$, is presented in the appendix of Ref.
$$\frac{\psi }{t})_a=cv_t\frac{v_n}{s}.$$
(21)
The last step is the determination of the boundary conditions at the contact with the substrate.
(iii) Contact points velocity
The motion of the contact point is governed by a binding/unbinding mechanism, implying a dissipation law that differs from the bulk dissipation. The most natural way for introducing a dissipation law is the following (with $`\mathrm{\Gamma }`$ a phenomenological dissipation coefficient)
$$\mathrm{\Gamma }\frac{dx_i}{dt}=\frac{\delta E}{\delta x_i}.$$
(22)
Using the energy variation (13) we get the following dynamical law, with $`w_i=w(x_i)`$ and $`v_i=dx_i/dt`$
$$c_1=\sqrt{\frac{2w_1+\mathrm{\Gamma }v_1}{\kappa }},c_2=\sqrt{\frac{2w_2\mathrm{\Gamma }v_2}{\kappa }},$$
(23)
which constitute the dynamical boundary conditions. These out of equilibrium values for the curvature are quite intuitive : the unbinding delay at the rear point induces a larger curvature than at equilibrium, whereas the binding delay at the fore point induces a smaller curvature.
## III Transient behavior
The formalism presented in the first part lends itself very well to analytical computation and stationary shape determination, as will appear in the following paragraphs. Nevertheless, having access to the transient process is highly desirable. In particularly, it checks the dynamical stability of an eventual stationary behavior, obtained after a relaxation. The successive vesicle profiles are determined by a direct numerical implementation of the dynamical equations (19), (20) and (21). Unfortunately, numerical instabilities are difficult to avoid around each contact point (due to a contact adhesion potential). A smoother model, without discontinuities, is more convenient for such an approach. For this reason, in this paragraph devoted to transient processes, the adhesion potential will be supposed to be of small, but non vanishing, range. We rapidly summarize below the small technical changes arising from this model modification. The chosen potential profile is
$$w(๐ซ)=w_0(1+u_0x)(\frac{y_0^4}{y^4}\frac{2y_0^2}{y^2}),$$
(24)
with the new length $`y_0`$ fixing the characteristic distance between the substrate and the membrane, and $`\widehat{w}(x)=w_0(1+u_0x)`$ the minimum of the potential interaction, occurring for $`y=y_0`$. It plays the role of the previous $`w(x)`$. It depends linearly on $`x`$ with an adhesion gradient $`u_0`$. The distance $`y_0`$ is chosen of the order of $`10nm`$, which is small enough in regard of the vesicle size to introduce only small variations between both models. In this case the parametrization is performed on the whole closed curve, and the boundary terms (eqs. 23) are no more relevant. The adhesion forces $`๐_w`$ are obtained by functional derivation of $`w(๐ซ)๐s`$ leading to
$$๐_w=(cw+w๐ง)๐ง.$$
(25)
Additionally the gauge condition fixing the tangential velocity eq. (20) is simplified : the velocity $`v_t(0)`$ is supposed to be zero, without loss of generality, so the first term disappears ; the last term is proportional to the length variation of the total parametrized curve, which is zero because we consider the complete profile and no more the free part of the curve. Thus we obtain for the gauge, replacing eq. (20) :
$$v_t=_0^s๐sv_nc.$$
(26)
Using equations (25) and (26) we finally get the dynamical equation for $`๐ซ`$
$$\frac{๐ซ(a,t)}{t}=\left[\kappa \left(\frac{d^2c}{ds^2}+\frac{1}{2}c^3\right)cw(w๐ง)p\zeta c\right]๐ง+v_t๐ญ.$$
(27)
The Lagrangian multipliers are determined from the following constraint equations :
$`{\displaystyle \frac{dL}{dt}}={\displaystyle cv_n๐s}=0,`$ (28)
$`{\displaystyle \frac{dS}{dt}}={\displaystyle v_n๐s}=0.`$ (29)
The normal component of the velocity in eq. (27) will be denoted by convention as $`v_n=v_n^0pc\zeta `$. With this notation the equation (29) appears as a very simple linear equation system in $`\zeta `$ and $`p`$. Its solution provides the pressure and tension values :
$`\zeta ={\displaystyle \frac{cv_n^0cv_n^0}{c^2c^2}},`$ (30)
$`p=\zeta c+v_n^0.`$ (31)
with the average defined by
$$\mathrm{}\frac{1}{L}_0^L๐s\mathrm{}.$$
(32)
We are now in a position to deal with the numerical anlysis. The dynamics is overdamped and for this reason local in time. Starting from an arbitrary profile, forward time integration provides us with the vesicle evolution. We have checked that (i) for a free vesicle (no substrate) the shape (with no external force) tends towards that obtained by direct energy minimization, (ii) we have also checked that for a homogeneous substrate an arbitrary initial shape evolves after some time to the shapes obtained in by direct energy minimization.
Let us now turn to the non-equilibrium situation ensured by an adhesion gradient. Starting from an initial shape, the vesicle acquires a non-symmetric shape and moves in the gradient direction. After transients have decayed the vesicle acquires a permanent regime with a constant velocity. Figure 5 shows the shape evolution.
## IV Stationary motion : direct numerical solution
The formulation of our problem in terms of a direct stationary problem is very convenient both for a systematic study of the velocity evolution as a function of various parameters. It will also allow us to present a simple analytical solution. It is convenient here to come back to the contact potential model. For a vesicle which has attained a stationary shape and velocity $`V`$ the equations become steady with $`V`$ as an unknown parameter.
For a stationary motion along the x-axis, normal and tangential velocities can be written as functions of the angle $`\psi `$ and of the translational velocity $`V`$ :
$`v_n=V\widehat{๐ฑ}\widehat{๐ง}=V\mathrm{sin}\psi ,`$ (33)
$`v_t=V\widehat{๐ฑ}\widehat{๐ญ}=V\mathrm{cos}\psi .`$ (34)
The shape and velocity are entirely determined from the relation between normal velocities and forces. The equation of motion is obtained from eqs. (19,33):
$$V\mathrm{sin}\psi =\frac{\kappa }{\eta }\left[\frac{^3\psi }{s^3}+\frac{1}{2}\left(\frac{\psi }{s}\right)^3\frac{\zeta }{\kappa }\frac{\psi }{s}\frac{p}{\kappa }\right].$$
(35)
Let us present briefly a counting argument showing that the problem is well defined. Equation (35) is a nonlinear third order differential equation for $`\psi `$, with 3 parameters to be determined : $`\zeta `$, $`p`$ and $`V`$. So we need 6 โinformationsโ.
We have the following equations at our disposal :
$``$ 2 geometrical constraints (eqs. (2) and (3)) :
$$_0^L^{}\mathrm{cos}\psi (s)๐s=L_{adh},_0^L\mathrm{sin}\psi (s)ds=0$$
(36)
$``$ 4 boundary equations corresponding to the contact angles and their first derivatives (the dynamical contact curvatures $`c_1`$ and $`c_2`$ given by equation (23))
$$\psi _1\psi (s=0)=\pi ,\psi _2\psi (s=L^{})=\pi ,\frac{\psi }{s})_{s=0}=c_1,\frac{\psi }{s})_{s=L^{}}=c_2$$
(37)
$``$ 1 equation ensuring that the enclosed surface is equal to the prescribed area.
$``$ 1 equation ensuring that the total length of the curve is precisely the prescribed one, $`L`$, which is related to the two other lengths by
$$L=L^{}+L_{adh}.$$
(38)
There are thus 8 conditions, for only 6 informations needed. The system seems then to be overdetermined. This is not the case. Indeed it must be noted that the problem involves additional unknowns which are $`L^{}`$ and $`L_{adh}`$. So in reality we have 8 unknown parameters as well. The problem is thus well defined.
Once the shape is determined we must in principle evaluate the area and change the parameter $`p`$ until the area coincides with the prescribed one. But since the area is a conjugate variable to $`p`$ we can fix $`p`$ โwhich is more convenientโ and this will fix some area that is treated as free (not imposed in advance). Additionally we are at liberty to prescribe $`L`$ (that fixes some length scale). $`L_{adh}`$ can then be determined if $`L^{}`$ is known; $`L_{adh}`$ can thus be removed from the problem upon using eq.(38). The first constraint (36) becomes then
$$_0^L^{}\mathrm{cos}\psi (s)๐s=LL^{}.$$
(39)
In other words prescribing the total length to $`L`$ and the pressure to $`p`$ lowers the number of unknowns by two. This is so because we do not want to have a specific area, and that $`L^{}`$ and $`L_{adh}`$ are not independent if we treat the total length as known. That is to say we have finally 6 fixed boundary conditions or constraints (36-37) and six parameters which are $`L^{}`$, $`V`$ and $`\zeta `$, plus three constants of integration due to the third order differential equation (35).
A convenient way to solve a differential equation of order $`n`$ is to transform it into a set of $`n`$ first order coupled differential equations. For that purpose we set $`f_1=\psi `$, $`\dot{f_1}=f_2`$ and $`\dot{f_2}=f_3`$ where the dot stands for $`/s`$. Equation (35) then provides us with the expression for $`\dot{f_3}`$ $`^3\psi /s^3`$ as a function of $`f_1`$, $`f_2`$ :
$$\dot{f_3}=\eta /\kappa P_2\mathrm{sin}f_1f_2^3/2p/\kappa +f_2P_3/\kappa F.$$
(40)
In order to make visible the quantities which are treated as unknown parameters we shall use the symbols $`P_i`$ (with $`i=1,2\mathrm{}`$). As stated above there are three parameters $`L^{}=P_1`$, $`V=P_2`$ and $`\zeta =P_3`$. Solution of a set of three equations involves three integration factors. This means that we have 6 unknowns, as argued in the last paragraph. Four physical conditions at the two end points (see eq. (37)) are known. Two constraints are imposed (eq. (36)), and this makes the problem well posed. Note that conditions (36) have an integral form. We find it convenient to rewrite them in a differential form. It is easy to realize that by setting
$$f_4(s)_0^s\mathrm{sin}\psi (s^{})๐s^{},f_5(s)_0^s\mathrm{cos}\psi (s^{})๐s^{},$$
(41)
we can write
$$\dot{f_4}=\mathrm{sin}\psi =\mathrm{sin}f_1,\dot{f_5}=\mathrm{cos}\psi =\mathrm{cos}f_1.$$
(42)
These two functions obviously obey $`f_4(0)=0`$, $`f_5(0)=0`$ , whereas at the second boundary we must impose $`f_4(L^{})=0`$ and $`f_5(L^{})=LL^{}`$ in order to fulfill the two constraints (36, 39). This trick is performed at the expense of two additional functions $`f_4`$ and $`f_5`$ (whose determinations involve two integration constants). We have thus augmented our system by 2 differential equations of first order. The two additional integration constants are precisely fixed by the demand $`f_4(0)=0`$, $`f_5(0)=0`$, whereas the conditions $`f_4(L^{})=0`$ and $`f_5(L^{})=LL^{}`$ are substituted to (36). Finally the shooting NAG code used here requires to invoke the boundary conditions for each function $`f_i`$, with $`i5`$. The boundary conditions for each function is invoked above, except for $`f_3`$ which represents the second derivative of $`\psi `$. This quantity is not known at the boundaries and there is no constraint to be imposed on it. Let $`P_4`$ and $`P_5`$ denote the values of $`f_3`$ at the two end points. We can thus invoke the boundary conditions of $`f_3`$ and the boundary values are quantities which are to be determined. That is to say we introduce two conditions with two additional unknown parameters. We have then in total ten unknowns and ten conditions. Cast into this form our formulation can straightforwardly be implemented into a NAG code (code D02HBF).
In summary the problem to be solved can be written in a standard boundary value problem with unknown parameters :
$$\begin{array}{ccc}\dot{f_1}=f_2,\hfill & f_1(0)=\pi ,\hfill & f_1(P_1)=\pi \hfill \\ \dot{f_2}=f_3,\hfill & f_2(0)=c_1,\hfill & f_2(P_1)=c_2\hfill \\ \dot{f_3}=F,\hfill & f_3(0)=P_4,\hfill & f_3(P_1)=P_5\hfill \\ \dot{f_4}=\mathrm{sin}f_1,\hfill & f_4(0)=0,\hfill & f_4(P_1)=0\hfill \\ \dot{f_5}=\mathrm{cos}f_1,\hfill & f_5(0)=0,\hfill & f_5(P_1)=LP_1\hfill \end{array}$$
(43)
Once the problem is solved the vesicle shape is obtained by making use of equations (1).
If $`w_1=w_2`$, the vesicle is at equilibrium on a homogeneous substrate and one obviously expects a vanishing velocity. This comes out automatically from the above formulation. If we were interested from the beginning in an equilibrium problem, we would then not have introduced $`V`$ as an unknown parameter. In that case because the profile is symmetric the second condition (36) is automatically satisfied, and we would then be left with with nine conditions for nine unknowns.
When $`w_1w_2`$, there is no equilibrium solution for the vesicle, which has to move towards the stronger adhesion region. If we impose a vanishing velocity there is no way to fulfill the second condition (36) (a typical profile would be the one shown on Fig.7) where starting from one end we arrive at the other end at a different height. Arriving at the same height can be achieved only for a specific velocity (or at most a discrete set of solutions), the one we are seeking. Thus the second condition of eq.(36) (which is parametrized by the set of $`P_i`$) can be viewed as โquantizationโ condition. This is a nonlinear eigenvalue problem of Barenblat-Zeldovitch type.
The numerical solution reveals an out-of-equilibrium shape which is significantly different from the equilibrium one, as shown on Fig.6. We note that the curvature in front of the vesicle is higher than the one behind. The reason is that the adhesion energy is higher in the front part, so that the curvature/adhesion balance allows a higher curvature (the vesicle looses curvature energy at the expense of a stronger adhesion).
## V Stationary motion : an analytical solution
The advantage offered by the simplistic picture of our model is the possibility to provide analytical results and thus to shed light on the physical processes that are involved in the problem of vesicle propulsion. It turns out that the equation of motion (35), if multiplied by $`^2\psi /s^2`$, possesses practically a first integral
$$V_0^L^{}\frac{^2\psi }{s^2}\mathrm{sin}(\psi )๐s=\frac{\kappa }{\eta }\left([(c/s)^2/2]_2^1+\frac{1}{2}[c^4/4]_2^1\frac{\zeta }{\kappa }[c^2/2]_2^1\frac{p}{\kappa }[c]_2^1\right).$$
(44)
Each r.h.s. term has an explicit form as a function of the contact curvatures (which are known), and of their first derivatives, which have only negligible contribution for swelled vesicle. The l.h.s. term can be evaluated for a vesicle shape close to a circle. The calculation is sent into appendix and leads to
$$\frac{^2\psi }{s^2}sin(\psi )๐s=4\pi ^2\frac{L_{adh}}{L^2}.$$
(45)
Using the dynamical values for the contact curvature given by eq.(23), we obtain an explicit expression for the velocity
$$V=\frac{L^2\kappa }{(2\pi )^2\eta L_{adh}}\left(\frac{1}{2}[c^4/4]_2^1\frac{\zeta }{\kappa }[c^2/2]_2^1\frac{p}{\kappa }[c]_2^1\right).$$
(46)
In the simple case where $`\mathrm{\Gamma }=0`$ (no dissipation associated with the substrate), expression (46) becomes explicit and provides a good agreement with numerical solution (see Fig.8). The analytical expression for the velocity involves only known parameters, except $`L_{adh}`$ and $`\zeta `$. For the comparison between numerical and analytical results, we took their numerical values.
Limit $`\delta ๐ฐ\mathrm{๐}`$
Another interesting limiting case concerns the small adhesion difference. Expansion of the numerator in eq.(46) to leading order in $`\delta w`$ yields
$$V\frac{\delta w}{\eta /A+\mathrm{\Gamma }}A\frac{w}{\kappa }\frac{R^2}{L_{adh}}\left[1\frac{p}{w}\sqrt{\frac{\kappa }{2w}}\frac{\zeta }{w}\right]$$
(47)
where $`R=L/2\pi `$. The influence of the two dissipation coefficients appears then clearly. It depends on the quantity $`A`$, proportional to the ratio $`R/L_{adh}`$. The bulk dissipation increases with the adhesion length, which seems to be a very robust result, as encountered in the model including hydrodynamics dissipation . The local dissipation represented by $`\mathrm{\Gamma }`$ does not depend on the adhesion length. Only the two contact points matter. Note also that the effective dissipation is $`\eta /A+\mathrm{\Gamma }`$. The bulk dissipation $`\eta /A`$ and the contact one $`\mathrm{\Gamma }`$ play a role of resistances (in an electric analogy) which would be mounted in series.
## VI Discussion and conclusion
This paper has given a first extensive presentation of the problem of vesicle migration in haptotaxis. We have reduced as much as possible the complexity of the problem in order to gain some analytical approximate results. For that purpose we have neglected hydrodynamics which induce nonlocal interactions, and adopted a local model of the Rousse type. The full dynamical problem has been solved by adopting a powerful gauge-field invariant formulation. The dynamical code could account for the transient and the evolution towards a steady-state solution. In that context an introduction of an adhesion potential with a finite, albeit small, range has proven to be necessary in order to circumvent numerical instabilities related to the motion of the contact point. This code has the advantage of dealing with various problems not leading necessarily to permanent motions. For a stationary situation we could cast the problem into a standard boundary value one where the migration velocity appeared as an eigenvalue. This problem is akin to the nonlinear eigenvalue problem of Barenblat-Zeldovitch type. A counting argument showed us that the velocity should belong to a discrete set, only one of them has been identified; we speculate that the solution is unique. The problem could be systematically solved in a fully intrinsic representation of the contour. For a rather tense vesicle we have provided an analytical solution which is in a good agreement with the numerical one. We have identified the role played by the adhesion length in selecting the magnitude of the migration velocity even if no dissipation with the substrate is included. We have also shown that the bond breaking/restoring dissipation and the (effective) bulk one are additive in a way which is analogous to the problem of electrical resistances in series. Bulk dissipation dominates when the ratio of the bulk dissipation coefficient to the contact one exceeds a certain limit, which depends in an intricate manner on various parameters. For real situations, vesicles, and cells in general, are suspended in aqueous solutions. It is therefore highly important to include hydrodynamics. Moreover the Lagrange multiplier $`\zeta `$ is a local quantity. We have recently given a brief account on these questions. An extensive discussion will be presented in the near future.
## A Derivation of equation 44
$$D=_0^L^{}\psi ^{\prime \prime }\mathrm{sin}\psi ds=_0^L^{}(\psi ^{})^2\mathrm{cos}\psi ds$$
We write $`\psi =2\pi s/L^{}\pi +ฯต`$, which implies to first order in $`ฯต^{}`$
$`D`$ $`=`$ $`\left({\displaystyle \frac{2\pi }{L^{}}}\right)^2{\displaystyle _0^L^{}}\mathrm{cos}\psi ds{\displaystyle \frac{4\pi }{L^{}}}{\displaystyle _0^L^{}}ฯต^{}\mathrm{cos}(2\pi s/L^{}\pi +ฯต)๐s`$ (A1)
$`=`$ $`\left({\displaystyle \frac{2\pi }{L^{}}}\right)^2L_{adh}+{\displaystyle \frac{4\pi }{L^{}}}{\displaystyle _0^L^{}}ฯต^{}\mathrm{cos}(2\pi s/L^{}+ฯต)๐s`$ (A2)
We then make use of the following relation:
$$\frac{d}{ds}(\mathrm{sin}\psi )=\left(\frac{2\pi }{L^{}}+ฯต^{}\right)\mathrm{cos}(2\pi s/L^{}+ฯต)$$
The integral between $`0`$ and $`L^{}`$ of the l.h.s. trem vanishes. We then obtain
$$_0^L^{}ฯต^{}\mathrm{cos}(2\pi s/L^{}+ฯต)๐s=\frac{2\pi }{L^{}}_0^L^{}\mathrm{cos}\psi ds=\frac{2\pi }{L^{}}L_{adh}$$
The sought after relation has then the form
$$D=\left(\frac{2\pi }{L^{}}\right)^2L_{adh}$$
Fig. 1 Notations used in the text.
Fig. 2 A geometrical explanation of the arc length variation with time.
Fig. 3 Force equilibrium at the fore contact point in the small rigidity limit.
Fig. 4 Translation of a circle obtained with a purely normal motion. The right part is dilated, whereas the left part is contracted.
Fig. 5 Successive vesicle profiles. The first one with open circles is an arbitrarily chosen initial shape. It relaxes to a permanent shape marked by filled circles on a inhomogeneous substrate.
Fig. 6 Out of equilibrium adhering vesicle profiles. $`V`$ is measured in units of $`100\mu m`$ and $`W`$ in units of $`10^4mJ/m^2`$.
Fig. 7 Geometrical constraint on the curve.
Fig. 8 Evolution of the vesicle velocity as a function of the adhesion difference. |
warning/0002/cond-mat0002144.html | ar5iv | text | # Hamiltonian dynamics and geometry of phase transitions in classical XY models
## I Introduction
The present paper deals with the study of the microscopic Hamiltonian dynamical phenomenology associated to thermodynamical phase transitions. This general subject is addressed in the special case of planar, classical Heisenberg XY models in two and three spatial dimensions. A preliminary presentation of some of the results and ideas contained in this paper has been already given in .
There are several reasons to tackle the Hamiltonian dynamical counterpart of phase transitions. On the one hand, we might wonder whether our knowledge of the already wide variety of dynamical properties of Hamiltonian systems can be furtherly enriched by considering the dynamical signatures, if any, of phase transitions. On the other hand, it is a-priori conceivable that also the theoretical investigation of the phase transition phenomena could benefit of a direct investigation of the natural microscopic dynamics. In fact, from a very general point of view, we can argue that in those times where microscopic dynamics was completely unaccessible to any kind of investigation, statistical mechanics has been invented just to replace dynamics. During the last decades, the advent of powerful computers has made possible, to some extent, a direct access to microscopic dynamics through the so called molecular dynamical simulations of the statistical properties of โmacroscopicโ systems.
Molecular dynamics can be either considered as a mere alternative to Monte Carlo methods in practical computations, or it can be also seen as a possible link to concepts and methods (those of nonlinear Hamiltonian dynamics) that could deepen our insight about phase transitions. In fact, by construction, the ergodic invariant measure of the Monte Carlo stochastic dynamics, commonly used in numerical statistical mechanics, is the canonical Gibbs distribution, whereas there is no general result that guarantees the ergodicity and mixing of natural (Hamiltonian) dynamics. Thus, the general interest for any contribution that helps in clarifying under what conditions equilibrium statistical mechanics correctly describes the average properties of a large collection of particles, safely replacing their microscopic dynamical description.
Actually, as it has been already shown and confirmed by the results reported below, there are some intrinsically dynamical observables that clearly signal the existence of a phase transition. Notably, Lyapunov exponents appear as sensitive measurements for phase transitions. They are also probes of a hidden geometry of the dynamics, because Lyapunov exponents depend on the geometry of certain โmechanical manifoldsโ whose geodesic flows coincide with the natural motions. Therefore, a peculiar energy โ or temperature โ dependence of the largest Lyapunov exponent at a phase transition point also reflects some important change in the geometry of the mechanical manifolds.
As we shall discuss throughout the present paper, also the topology of these manifolds has been discovered to play a relevant role in the phase transition phenomena (PTP).
Another strong reason of interest for the Hamiltonian dynamical counterpart of PTP is related to the equivalence problem of statistical ensembles. Hamiltonian dynamics has its most natural and tight relationship with microcanonical ensemble. Now, the well known equivalence among all the statistical ensembles in the thermodynamic limit is valid in general in the absence of thermodynamic singularities, i.e. in the absence of phase transitions. This is not a difficulty for statistical mechanics as it might seem at first sight , rather, this is a very interesting and intriguing point.
The inequivalence of canonical and microcanonical ensembles in presence of a phase transition has been analytically shown for a particular model by Hertel and Thirring , it is mainly revealed by the appearance of negative values of the specific heat and has been discussed by several authors .
The microcanonical description of phase transitions seems also to offer many advantages in tackling first order phase transitions , and seems considerably less affected by finite-size scaling effects with respect to the canonical ensemble description . This non-equivalence problem, together with certain advantages of the microcanonical ensemble, strenghtens the interest for the Hamiltonian dynamical counterpart of PTP. Let us briefly mention the existing contributions in the field.
Butera and Caravati , considering an XY model in two dimensions, found that the temperature dependence of the largest Lyapunov exponent changes just near the critical temperature $`T_c`$ of the Kosterlitz-Thouless phase transition. Other interesting aspects of the Hamiltonian dynamics of the XY model in two dimensions have been extensively considered in , where a very rich phenomenology is reported. Recently, the behaviour of Lyapunov exponents has been studied in Hamiltonian dynamical systems: i) with long-range interactions , ii) describing either clusters of particles or magnetic or gravitational models exhibiting phase transitions, iii) in classical lattice field theories with $`O(1)`$, $`O(2)`$ and $`O(4)`$ global symmetries in two and three space dimensions , iv) in the XY model in two and three space dimensions , v) in the โ$`\mathrm{\Theta }`$ \- transitionโ of homopolymeric chains . The pattern of $`\lambda (T)`$ close to the critical temperature $`T_c`$ is model-dependent. The behaviour of Lyapunov exponents near the transition point has been considered also in the case of first- order phase transitions . It is also worth mentioning the very intriguing result of Ref., where a glassy transition is accompanied by a sharp jump of $`\lambda (T)`$.
$`\lambda (T)`$ always detects a phase transition and, even if its pattern close to the critical temperature $`T_c`$ is model-dependent, it can be used as an order parameter โ of dynamical origin โ also in the absence of a standard order parameter (as in the case of the mentioned โ$`\mathrm{\Theta }`$-transitionโ of homopolymers and of the glassy transition in amorphous materials). This appears of great prospective interest also in the light of recently developed analytical methods to compute Lyapunov exponents (see Section IV).
Among Hamiltonian models with long-range interactions exhibiting phase transitions, the most extensively studied is the mean-field XY model , whose equilibrium statistical mechanics is exactly described, in the thermodynamic limit, by mean-field theory . In this system, the theoretically predicted temperature dependence of the largest Lyapunov exponent $`\lambda `$ displays a non-analytic behavior at the phase transition point.
The aims of the present paper are
* to investigate the dynamical phenomenology of Kosterlitz-Thouless and of second order phase transitions in the $`2d`$ and $`3d`$ classical Heisenberg XY models respectively;
* to highlight the microscopic dynamical counterpart of phase transitions through the temperature dependence of the Lyapunov exponents, also providing some physical interpretation of abstract quantities involved in the geometric theory of chaos (in particular among vorticity, Lyapunov exponents and sectional curvatures of configuration space);
* to discuss the hypothesis that phase transition phenomena could be originated by suitable changes in the topology of the constant energy hypersurfaces of phase space, therefore hinting to a mathematical characterization of phase transitions in the microcanonical ensemble.
The paper is organized as follows: Sections $`II`$ and $`III`$ are devoted to the dynamical investigation of the $`2d`$ and $`3d`$ XY models respectively. In Section $`IV`$ the geometric description of chaos is considered, with the analytic derivation of the temperature dependence of the largest Lyapunov exponent, the geometric signatures of a second-order phase transition and the topological hypothesis. Section $`V`$ contains a presentation of the relationship between the extrinsic geometry and topology of the energy hypersurfaces of phase space and thermodynamics; the results of some numeric computations are also reported. Finally, Section $`VI`$ is devoted to summarize the achievements reported in the present paper and to discuss their meaning.
## II $`2d`$ XY model
We considered a system of planar, classical โspinsโ (in fact rotators) on a square lattice of $`N=n\times n`$ sites, and interacting through the ferromagnetic interaction $`V=_{i,j}J๐_i๐_j`$ (where $`|๐_i|=1)`$. The addition of standard, i.e. quadratic, kinetic energy term leads to the following choice of the Hamiltonian
$$H=\underset{i,j=1}{\overset{n}{}}\left\{\frac{p_{i,j}^2}{2}+J[2\mathrm{cos}(q_{i+1,j}q_{i,j})\mathrm{cos}(q_{i,j+1}q_{i,j})]\right\},$$
(1)
where $`q_{i,j}`$ are the angles with respect to a fixed direction on the reference plane of the system. In the usual definition of the XY model both the kinetic term and the constant term $`2JN`$ are lacking; however, their contribution does not modify the thermodynamic averages (because they usually depend only on the configurational partition function, $`Z_C=_{i=1}^Ndq_i\mathrm{exp}[\beta V(q)]`$, the momenta being trivially integrable when the kinetic energy is quadratic). Thus, as we tackle classical systems, the choice of a quadratic kinetic energy term is natural because it corresponds to $`\frac{1}{2}_{i=1}^N|\dot{๐}_i|^2`$, written in terms of the momenta $`p_{i,j}`$ canonically conjugated to the lagrangian coordinates $`q_{i,j}`$. The constant term $`2JN`$ is introduced to make the low energy expansion of Eq. (1) coincident with the Hamiltonian of a system of weakly coupled harmonic oscillators.
The theory predicts for this model a Kosterlitz-Thouless phase transition occurring at a critical temperature estimated around $`T_cJ`$. Many Monte Carlo simulations of this model have been done in order to check the predictions of the theory. Among them, we quote those of Tobochnik and Chester and of Gupta and Baillie which, on the basis of accurate numerical analysis, confirmed the predictions of the theory and fixed the critical temperature at $`T_c=0.89`$ ($`J=1`$).
The analysis of the present work is based on the numerical integration of the equations of motion derived from Hamiltonian (1). The numerical integration is performed by means of a bilateral, third order, symplectic algorithm , and it is repeated at several values of the energy density $`ฯต=E/N`$ ($`E`$ is the total energy of the system which depends upon the choice of the initial conditions). While the Monte Carlo simulations perform statistical averages in the canonical ensemble, Hamiltonian dynamics has its statistical counterpart in the microcanonical ensemble. Statistical averages are here replaced by time averages of relevant observables. In this perspective, from the microcanonical definition of temperature $`1/T=S/E`$, where $`S`$ is the entropy, two definitions of temperature are available: $`T=\frac{2}{N}K`$ (where $`K`$ is the kinetic energy per degree of freedom), if $`S=\mathrm{log}_{i=1}^Ndq_idp_i\mathrm{\Theta }(H(p,q)E)`$, where $`\mathrm{\Theta }()`$ is the Heaviside step function, and $`\stackrel{~}{T}=\left[\left(\frac{N}{2}1\right)K^1\right]^1`$, if $`S=\mathrm{log}_{i=1}^Ndq_idp_i\delta (H(p,q)E)`$ . $`T`$ (or $`\stackrel{~}{T}`$) are numerically determined by measuring the time average of the kinetic energy $`K`$ per degree of freedom (or its inverse), i.e. $`T=lim_t\mathrm{}\frac{2}{N}\frac{1}{t}_0^t๐\tau K(\tau )`$ (and similarly for $`\stackrel{~}{T}`$). There is no appreciable difference in the outcomes of the computations of temperature according to these two definitions.
### A Dynamical analysis of thermodynamical observables
#### 1 Order parameter
The order parameter for a system of planar โspinsโ whose Hamiltonian is invariant under the action of the group $`O(2)`$, is the bidimensional vector
$$๐=(\underset{i,j=1}{\overset{n}{}}๐_{i,j}^x,\underset{i,j=1}{\overset{n}{}}๐_{i,j}^y)(\underset{i,j=1}{\overset{n}{}}\mathrm{cos}q_{i,j},\underset{i,j=1}{\overset{n}{}}\mathrm{sin}q_{i,j}),$$
(2)
which describes the mean spin orientation field. After the Mermin-Wagner theorem, we know that no symmetry-breaking transition can occur in one and two dimensional systems with a continuous symmetry and nearest-neighbour interactions. This means that, at any non-vanishing temperature, the statistical average of the total magnetization vector is necessarily zero in the thermodynamic limit. However, a vanishing magnetization is not necessarily expected when computed by means of Hamiltonian dynamics at finite $`N`$. In fact, statistical averages are equivalent to averages computed through suitable markovian Monte Carlo dynamics that a-priori can reach any region of phase space, whereas in principle a true ergodicity breaking is possible in the case of differentiable dynamics. Also an โeffectiveโ ergodicity breaking of differentiable dynamics is possible, when the relaxation times โ of time to ensemble averages โ are very fastly increasing with $`N`$ .
This model has two integrable limits: coupled harmonic oscillators and free rotators, at low and high temperatures respectively. Hereafter, $`T`$ is meant in units of the coupling constant $`J`$.
For a lattice of $`N=10\times 10`$ sites, Figure 1 shows that at low temperatures ($`T<0.5`$)โ being the system almost harmonic โ we can observe a persistent memory of the total magnetization associated with the initial condition, which, on the typical time scales of our numeric simulations ($`10^6`$ units of proper time), looks almost frozen.
By raising the temperature above a first threshold $`T_00.6`$, the total magnetization vector โ observed on the same time scale โ starts rotating on the plane where it is confined. A further increase of the temperature induces a faster rotation of the magnetization vector together with a slight reduction of its average modulus.
At temperatures slightly greater than $`1`$, we observe that already at $`N=10\times 10`$ a random variation of the direction and of the modulus of the vector $`๐(t)`$ sets in.
At $`T>1.2`$, we observe a fast relaxation and, at high temperatures ($`T10`$), a sort of saturation of chaos.
At a first glance, the results reported in Fig. 1 could suggest the presence of a phase transition associated with the breaking of the $`O(2)`$ symmetry. In fact, having in mind the Landau theory, the ring-shaped distribution of the instantaneous magnetization shown by Fig. 1 is the typical signature of an $`O(2)`$-broken symmetry phase and the spot-like patterns around zero are proper to the unbroken symmetry phase.
The apparent contradiction of these results with the Mermin-Wagner theorem is resolved by checking whether the observed phenomenology is stable with $`N`$. Thus, some simulations have been performed at larger values of $`N`$. At any temperature, we found that the average modulus $`|๐(t)|_t`$ of the vector $`๐(t)`$, computed along the trajectory, systematically decreases by increasing $`N`$. However, for temperatures lower than $`T_0`$, the $`N`$-dependence of the order parameter is very weak, whereas, for temperatures greater than $`T_0`$, the $`N`$-dependence of the order parameter is rather strong. In Fig. 2 two extreme cases ($`N=10\times 10`$ and $`N=200\times 200`$) are shown for $`T=0.74`$. The systematic trend of $`|๐(t)|`$ toward smaller values at increasing $`N`$ is consistent with its expected vanishing in the limit $`N\mathrm{}`$.
At $`T=1`$, Fig. 3 shows that, when the lattice dimension is greater than $`50\times 50`$, $`๐(t)`$ displays random variations both in direction (in the interval \[0,$`2\pi `$\]) and in modulus (between zero and a value which is smaller at larger $`N`$).
#### 2 Specific heat
By means of the recasting of a standard formula which relates the average fluctuations of a generic observable computed in canonical and microcanonical ensembles , and by specializing it to the kinetic energy fluctuations, one obtains a microcanonical estimate of the canonical specific heat
$$c_V(T)=\frac{C_V}{N}\{\begin{array}{cc}c_V(ฯต)=\frac{k_Bd}{2}\left[1\frac{Nd}{2}\frac{K^2K^2}{K^2}\right]^1,\hfill & \\ T=T(ฯต)\hfill & \end{array}$$
(3)
where $`d`$ is the number of degrees of freedom for each particle. Time averages of the kinetic energy fluctuations computed at any given value of the energy density $`ฯต`$ yield $`C_V(T)`$, according to its parametric definition in Eq.(3).
ยฟFrom the microcanonical definition $`1/C_V=T(E)/E`$ of the constant volume specific heat, a formula can be worked out , which is exact at any value of $`N`$ (at variance with the expression (3)). It reads
$$c_V=\frac{C_V}{N}=[N(N2)KK^1]^1$$
(4)
and it is the natural expression to be used in Hamiltonian dynamical simulations of finite systems.
The numerical simulations of the Hamiltonian dynamics of the $`2d`$ XY model โ computed with both Eqs.(3) and (4) โ yield a cuspy pattern for $`c_V(T)`$ peaked at $`T1`$ (Fig. 4). This is in good agreement with the outcomes of canonical Monte Carlo simulations reported in Ref. , where a pronounced peak of $`c_V(T)`$ was detected at $`T1.02`$.
By varying the lattice dimensions, the peak height remains constant, in agreement with the absence of a symmetry-breaking phase transition.
#### 3 Vorticity
Another thermodynamic observable which can be studied is the vorticity of the system. Let us briefly recall that if the angular differences of nearby โspinsโ are small, we can suppose the existence of a continuum limit function $`\theta (๐ซ)`$ that conveniently fits a given spatial configuration of the system. Spin waves correspond to regular patterns of $`\theta (๐ซ)`$, whereas the appearance of a singularity in $`\theta (๐ซ)`$ corresponds to a topological defect, or a vortex, in the โspinโ configuration. When such a defect is present, along any closed path $`๐`$ that contains the centre of the defect, one has
$$_๐\theta (๐ซ)๐๐ซ=2\pi q,q=0,\pm 1,\pm 2,\mathrm{}$$
(5)
indicating the presence of a vortex whose intensity is $`q`$. For a lattice model with periodic boundary conditions, there is an equal number of vortices and antivortices (i.e. vortices rotating in opposite directions). Thus, the vorticity of our model can be defined as the mean total number of equal sign vortices per unit volume. In order to compute the vorticity $`๐ฑ`$ as a function of temperature, we have averaged the number of positive vortices along the numerical phase space trajectories. On the lattice, $`๐ซ`$ is replaced by the multi-index $`๐ข`$ and $`_\mu \theta _i=q_{๐ข+\mu }q_๐ข`$, then the number of elementary vortices is counted: the discretized version of $`_{\mathrm{}}\theta d๐ซ=1`$ amounts to one elementary vortex on a plaquette. Thus $`๐ฑ`$ is obtained by summing over all the plaquettes.
Our results are in agreement with the values obtained by Tobochnik and Chester by means of Monte Carlo simulations with $`N=60\times 60`$.
As shown in Fig. 5, on the $`10\times 10`$ lattice, the first vortex shows up at $`T0.6`$ and on the $`40\times 40`$ lattice at $`T0.5`$, when the system changes its dynamical behavior, increasing its chaoticity (see next Subsection). At lower temperatures, vortices are less probable, due to the fact that the formation of vortex has a minimum energy cost. Below $`T1`$, the vortex density steeply grows with a power law $`๐ฑ(T)T^{10}`$. The growth of $`๐ฑ`$ then slows down, until the saturation is reached at $`T10`$.
### B Lyapunov exponents and chaoticity
The values of the largest Lyapunov exponent $`\lambda _1`$ have been computed using the standard tangent dynamics equations \[see Eqs. (11) and (60)\], and are reported in Fig. 6.
Below $`T0.6`$, the dynamical behavior is nearly the same as that of harmonic oscillators and the excitations of the system are only โspin-wavesโ.
In the interval $`[0.,0.6]`$, the observed temperature dependence $`\lambda _1(T)T^2`$ is equivalent to the $`\lambda _1(ฯต)ฯต^2`$ dependence (since at low temperature $`T(ฯต)ฯต`$), already found โ analytically and numerically โ in the quasi-harmonic regime of other systems and characteristic of weakly chaotic dynamics .
Above $`T0.6`$, vortices begin to form and correspondingly the largest Lyapunov exponent signals a โqualitativeโ change of the dynamics through a steeper increase vs. $`T`$.
At $`T0.9`$, where the theory predicts a Kosterlitz - Thouless phase transition, $`\lambda _1(T)`$ displays an inflection point.
Finally, at high temperatures, the power law $`\lambda _1(T)T^{1/6}`$ is found.
## III $`3d`$ XY model
In order to extend the dynamical investigation to the case of second-order phase transitions, we have studied a system described by an Hamiltonian having at the same time the main characteristics of the $`2d`$ model and the differences necessary to the appearance of a spontaneous symmetry-breaking below a certain critical temperature. The model we have chosen is such that the spin rotation is constrained on a plane and only the lattice dimension has been increased, in order to elude the โno goโ conditions of the Mermin-Wagner theorem. This is simply achieved by tackling a system defined on a cubic lattice of $`N=n\times n\times n`$ sites and described by the Hamiltonian
$`H`$ $`=`$ $`{\displaystyle \underset{i,j,k=1}{\overset{n}{}}}\{{\displaystyle \frac{p_{i,j,k}^2}{2}}+J[3\mathrm{cos}(q_{i+1,j,k}q_{i,j,k})`$ (6)
$``$ $`\mathrm{cos}(q_{i,j+1,k}q_{i,j,k})\mathrm{cos}(q_{i,j,k+1}q_{i,j,k})]\}.`$ (7)
### A Dynamical analysis of thermodynamical observables
The basic thermodynamical phenomenology of a second-order phase transition is characterized by the existence of equilibrium configurations that make the order parameter bifurcating away from zero at some critical temperature $`T_c`$ and by a divergence of the specific heat $`c_V(T)`$ at the same $`T_c`$. Therefore, this is the obvious starting point for the Hamiltonian dynamical approach.
#### 1 Order parameter
Below a critical value of the temperature, the symmetry-breaking in a system invariant under the action of the $`O(2)`$ group, appears as the selection โ by the average magnetization vector of Eq. (2)โ of a preferred direction among all the possible, energetically equivalent choices. By increasing the lattice dimension, the symmetry breaking is therefore characterized by a sort of simultaneous โfreezingโ of the direction of the order parameter $`๐`$ and of the convergence of its modulus to a non-zero value.
Figure 7 shows that in the $`3d`$ lattice, at $`T<2`$, i.e. in the broken-symmetry phase (as we shall see in the following), the dynamical simulations yield a thinner spread of the longitudinal fluctuations by increasing $`N`$ โ that is, $`|๐|`$ oscillates by exhibiting a trend to converge to a non-zero value โ and that the transverse fluctuations damp, โfixingโ the direction of the oscillations. This direction depends on the initial conditions.
Moreover, the dynamical analysis provides us with a better detail than a simple distinction between regular and chaotic dynamics. In fact, it is possible to distinguish between three different dynamical regimes (Fig. 8).
At low temperatures, up to $`T0.8`$, one observes the persistency of the initial direction and of an equilibrium value of the modulus $`|๐|`$ close to one.
At $`0.8<T<2.2`$, one observes transverse oscillations, whose amplitude increases with temperature.
At $`T>2.2`$, the order parameter exhibits the features typical of an unbroken symmetry phase. In fact, it displays fluctuations peaked at zero, whose dispersion decreases by increasing the temperature (bottom of Fig. 8) and, at a given temperature, by increasing the lattice volume (Fig. 9a,b).
We can give an estimate of the order parameter by evaluating the average of the modulus $`|๐(t)|=\rho (T)`$. At $`T<2.2`$, the $`N`$-dependence is given mainly by the rotation of the vector, while the longitudinal oscillations are moderate, as shown in Fig. 10. At temperatures above $`T2.2`$, we observe the squeezing of $`\rho (T)`$ to a small value.
The existence of a second order phase transition can be recognized by comparing the temperature behavior and the $`N`$-dependence of the thermodynamic observables computed for the $`2d`$ and the $`3d`$ models. Both systems exhibit the rotation of the magnetization vector and small fluctuations of its modulus when they are considered on small lattices. In the $`2d`$ model the average modulus of the order parameter is theoretically expected to vanish logarithmically with $`N`$, what seems qualitatively compatible with the weak $`N`$ dependence shown in Fig. 2, whereas in the $`3d`$ model we observe a stability with $`N`$ of $`|๐|`$, suggesting the convergence to a non-zero value of the order parameter also in the limit $`N\mathrm{}`$, as shown in Fig. 7.
$`T2.2`$ is an approximate value of the critical temperature $`T_c`$ of the second-order phase transition. This value will be refined in the following Subsection. No finite-size scaling analysis has been performed for two different reasons: i) our main concern is a qualitative phenomenological analysis of the Hamiltonian dynamics of phase transitions rather than a very accurate quantitative analysis, ii) finite-size effects are much weaker in the microcanonical ensemble than in the canonical ensemble .
#### 2 Specific heat
As in the $`2d`$ model, numerical simulations of the Hamiltonian dynamics have been performed with both Eqs.(3) and (4). The outcomes show a cusplike pattern of the specific heat, whose peak makes possible a better determination of the critical temperature. By increasing the lattice dimension up to $`N=15\times 15\times 15`$, the cusp becomes more pronounced, at variance with the case of the $`2d`$ model. Fig. 11 shows that this occurs at the temperature $`T_c2.17`$.
#### 3 Vorticity
The definition of the vorticity in the $`3d`$ case is not a simple extension of the $`2d`$ case. Vortices are always defined on a plane and if all the โspinsโ could freely move in the three-dimensional space, the concept of vortices would be meaningless. For the $`3d`$ planar (anisotropic) model considered here, vortices can be defined and studied on two-dimensional subspaces of the lattice. The variables $`q_{i,j,k}`$ do not contain any information about the position of the plane where the reference direction to measure the angles $`q_{i,j,k}`$ is assigned. Dynamics is completely independent of this choice, which has no effect on the Hamiltonian. Moreover, as the Hamiltonian is symmetric with respect to the lattice axes, the three coordinate-planes are equivalent. This equivalence implies that vortices can contemporarily exist on three orthogonal planes. Though the usual pictorial representation of a vortex can hardly be maintained, its mathematical definition is the same as in the $`2d`$ lattice case. Hence three vorticity functions exist and their average values - at a given temperature - should not differ, what is actually confirmed by numerical simulations.
The vorticity function vs. temperature is plotted in Fig. 12. On a lattice of $`10\times 10\times 10`$ spins, the first vortex is observed at $`T0.8`$. The growth of the average density of vortices is very fast up to the critical temperature, above which the saturation is reached.
### B Lyapunov exponents and symmetry-breaking phase transition
A quantitative analysis of the dynamical chaoticity is provided by the temperature dependence of the largest Lyapunov exponent.
Figure 13 shows the results of this computation. At low temperatures, in the limit of quasi-harmonic oscillators, the scaling law is again found to be $`\lambda _1(T)T^2`$ and, at high temperatures, the scaling law is again $`\lambda _1(T)T^{1/6}`$, as in the $`2d`$ case. In the temperature range intermediate between $`T0.8`$ and $`T_c2.17`$, there is a linear growth of $`\lambda _1(T)`$. At the critical temperature, the Lyapunov exponent exhibits an angular point. This makes a remarkable difference between this system undergoing a second order phase transition and its $`2d`$ version, undergoing a Kosterlitz- Thouless transition. In fact, the analysis of the $`2d`$ model has shown a mild transition between the different regimes of $`\lambda _1(T)`$ (inset of Fig. 12), whereas in the $`3d`$ model this transition is sharper (inset of Fig. 13).
We have also computed the temperature dependence of the largest Lyapunov exponent of Markovian random processes which replace the true dynamics on the energy surfaces $`\mathrm{\Sigma }_E`$ (see Appendix). The results are shown in Fig. 14. The dynamics is considered strongly chaotic in the temperature range where the patterns $`\lambda _1(T)`$ are the same for both random and differentiable dynamics, i.e. when differentiable dynamics mimics, to some extent, a random process. The dynamics is considered weakly chaotic when the value $`\lambda _1`$ resulting from random dynamics is larger than the value $`\lambda _1`$ resulting from differentiable dynamics. The transition from weak to strong chaos is quite abrupt. Figure 14 shows that the pattern of the largest Lyapunov exponent computed by means of the random dynamics reproduces that of the true Lyapunov exponent at temperatures $`TT_c`$. This means that the setting in of strong thermodynamical disorder corresponds to the setting in of strong dynamical chaos. The โwindowโ of strong chaoticity starts at $`T_c`$ and ends at $`T10`$. The existence of a second transition from strong to weak chaos is due to the existence, for $`T\mathrm{}`$, of the second integrable limit (of free rotators), whence chaos cannot remain strong at any $`T>T_c`$.
## IV Geometry of dynamics and phase transitions
Let us briefly recall that the geometrization of the dynamics of $`N`$-degrees-of-freedom systems defined by a Lagrangian $`=KV`$, in which the kinetic energy is quadratic in the velocities: $`K=\frac{1}{2}a_{ij}\dot{q}^i\dot{q}^j`$, stems from the fact that the natural motions are the extrema of the Hamiltonian action functional $`๐ฎ_H=๐t`$, or of the Maupertuisโ action $`๐ฎ_M=2K๐t`$. In fact, also the geodesics of Riemannian and pseudo-Riemannian manifolds are the extrema of a functional, the arc-length $`\mathrm{}=๐s`$, with $`ds^2=g_{ij}dq^idq^j`$. Hence, a suitable choice of the metric tensor allows for the identification of the arc-length with either $`๐ฎ_H`$ or $`๐ฎ_M`$, and of the geodesics with the natural motions of the dynamical system. Starting from $`๐ฎ_M`$, the โmechanical manifoldโ is the accessible configuration space endowed with the Jacobi metric
$$(g_J)_{ij}=[EV(q)]a_{ij},$$
(8)
where $`V(q)`$ is the potential energy and $`E`$ is the total energy. A description of the extrema of Hamiltonโs action $`๐ฎ_H`$ as geodesics of a โmechanical manifoldโ can be obtained using Eisenhartโs metric on an enlarged configuration spacetime ($`\{q^0t,q^1,\mathrm{},q^N\}`$ plus one real coordinate $`q^{N+1}`$), whose arc-length is
$$ds^2=2V(\{q\})(dq^0)^2+a_{ij}dq^idq^j+2dq^0dq^{N+1}.$$
(9)
The manifold has a Lorentzian structure and the dynamical trajectories are those geodesics satisfying the condition $`ds^2=Cdt^2`$, where $`C`$ is a positive constant. In the geometrical framework, the (in)stability of the trajectories is the (in)stability of the geodesics, and it is completely determined by the curvature properties of the underlying manifold according to the Jacobi equation
$$\frac{^2\xi ^i}{ds^2}+R_{jkm}^i\frac{dq^j}{ds}\xi ^k\frac{dq^m}{ds}=0,$$
(10)
whose solution $`\xi `$, usually called Jacobi or geodesic variation field, locally measures the distance between nearby geodesics; $`/ds`$ stands for the covariant derivative along a geodesic and $`R_{jkm}^i`$ are the components of the Riemann curvature tensor. Using the Eisenhart metric (9), the relevant part of the Jacobi equation (10) is
$$\frac{d^2\xi ^i}{dt^2}+R_{0k0}^i\xi ^k=0,i=1,\mathrm{},N$$
(11)
where the only non-vanishing components of the curvature tensor are $`R_{0i0j}=^2V/q_iq_j`$. Equation (11) is the tangent dynamics equation, which is commonly used to measure Lyapunov exponents in standard Hamiltonian systems. Having recognized its geometric origin, it has been devised in Ref. a geometric reasoning to derive from Eq.(11) an effective scalar stability equation that, independently of the knowledge of dynamical trajectories, provides an average measure of their degree of instability. An intermediate step in this derivation yields
$$\frac{d^2\xi ^j}{dt^2}+k_R(t)\xi ^j+\delta K^{(2)}(t)\xi ^j=0,$$
(12)
where $`k_R=K_R/N`$ is the Ricci curvature along a geodesic defined as $`K_R=\frac{1}{v^2}R_{ij}\dot{q}^i\dot{q}^j`$, with $`v^2=\dot{q}^i\dot{q}_i`$ and $`R_{ij}=R_{ikj}^k`$, and $`\delta K^{(2)}`$ is the local deviation of sectional curvature from its average value . The sectional curvature is defined as $`K^{(2)}=R_{ijkl}\xi ^i\dot{q}^j\xi ^k\dot{q}^l/\xi ^2\dot{q}^2`$.
Two simplifying assumptions are made: $`(i)`$ the ambient manifold is almost isotropic, i.e. the components of the curvature tensor โ that for an isotropic manifold (i.e. of constant curvature) are $`R_{ijkm}=k_0(g_{ik}g_{jm}g_{im}g_{jk})`$, $`k_0=const`$ โ can be approximated by $`R_{ijkm}k(t)(g_{ik}g_{jm}g_{im}g_{jk})`$ along a generic geodesic $`\gamma (t)`$; $`(ii)`$ in the large $`N`$ limit, the โeffective curvatureโ $`k(t)`$ can be modeled by a gaussian and $`\delta `$-correlated stochastic process. Hence, one derives an effective stability equation, independent of the dynamics and in the form of a stochastic oscillator equation ,
$$\frac{d^2\psi }{dt^2}+[k_0+\sigma _k\eta (t)]\psi =0,$$
(13)
where $`\psi ^2|\xi |^2`$. The mean $`k_0`$ and variance $`\sigma _k`$ of $`k(t)`$ are given by $`k_0=K_R/N`$ and $`\sigma _k^2=(K_RK_R)^2/N`$, respectively, and the averages $``$ are geometric averages, i.e. integrals computed on the mechanical manifold. These averages are directly related with microcanonical averages, as it will be seen at the end of Section V. $`\eta (t)`$ is a gaussian $`\delta `$-correlated random process of zero mean and unit variance.
The main source of instability of the solutions of Eq.(13), and therefore the main source of Hamiltonian chaos, is parametric resonance, which is activated by the variations of the Ricci curvature along the geodesics and which takes place also on positively curved manifolds . The dynamical instability can be enhanced if the geodesics encounter regions of negative sectional curvatures, such that $`k_R+\delta K^{(2)}<0`$, as it is evident from Eq. (12).
In the case of Eisenhart metric, it is $`K_R\mathrm{\Delta }V=_{i=1}^N(^2V/q_i^2)`$ and $`K^{(2)}=R_{0i0j}\xi ^i\xi ^j/\xi ^2(^2V/q^iq^j)\xi ^i\xi ^j/\xi ^2`$. The exponential growth rate $`\lambda `$ of the quantity $`\psi ^2+\dot{\psi }^2`$ of the solutions of Eq. (13), is therefore an estimate of the largest Lyapunov exponent that can be analytically computed. The final result reads
$$\lambda =\frac{\mathrm{\Lambda }}{2}\frac{2k_0}{3\mathrm{\Lambda }},\mathrm{\Lambda }=\left(2\sigma _k^2\tau +\sqrt{\frac{64k_0^3}{27}+4\sigma _k^4\tau ^2}\right)^{\frac{1}{3}},$$
(14)
where $`\tau =\pi \sqrt{k_0}/(2\sqrt{k_0(k_0+\sigma _k)}+\pi \sigma _k)`$; in the limit $`\sigma _k/k_01`$ one finds $`\lambda \sigma _k^2`$.
### A Signatures of phase transitions from geometrization of dynamics
In the geometric picture, chaos is mainly originated by the parametric instability activated by the fluctuating curvature felt by geodesics, i.e. the fluctuations of the (effective) curvature are the source of the instability of the dynamics. On the other hand, as it is witnessed by the derivation of Eq. (13) and by the equation itself, a statistical-mechanical-like treatment of the average degree of chaoticity is made possible by the geometrization of the dynamics. The relevant curvature properties of the mechanical manifolds are computed, at the formal level, as statistical averages, like other thermodynamic observables. Thus, we can expect that some precise relationship may exist between geometric, dynamic and thermodynamic quantities. Moreover, this implies that phase transitions should correspond to peculiar effects in the geometric observables.
In the particular case of the $`2d`$ XY model, the microcanonical average kinetic energy $`K`$ and the average Ricci curvature $`K_R`$ computed with the Eisenhart metric are linked by the equation
$$K_R=\underset{i,j=1}{\overset{N}{}}\frac{^2V}{^2q_{i,j}}=2J\underset{i,j=1}{\overset{N}{}}\mathrm{cos}(q_{i+1,j}q_{i,j})+\mathrm{cos}(q_{i,j+1}q_{i,j})=2(JV),$$
(15)
so that
$$H=Nฯต=K+V\frac{K}{N}=ฯต2J+\frac{1}{2}\frac{K_R}{N}.$$
(16)
Being the temperature defined as $`T=2K/N`$ (with $`k_B=1`$) and being $`d=1`$ (because each spin has only one rotational degree of freedom), from Eq.(3) it follows that
$$c_V=\frac{1}{2}\left(1\frac{1}{2}\frac{\sigma _k^2/N}{T^2}\right)^1.$$
(17)
In the special case of these XY systems, it is possible to link the specific heat and the Ricci curvature by inserting Eq.(16) into the usual expression for the specific heat at constant volume. Thus, one obtains the equation
$$c_V=\frac{1}{2N}\frac{K_R(T)}{T}.$$
(18)
The appearance of a peak in the specific heat function at the critical temperature has to correspond to a suitable temperature dependence of the Ricci curvature.
In the $`3d`$ model, the potential energy and the Ricci curvature are proportional, according to: $`\frac{1}{N}V=3\frac{1}{2N}K_R`$.
Another interesting point is the relation between a geometric observable and the vorticity function in both models. As already seen in previous sections, the vorticity function is a useful signature of the dynamical chaoticity of the system. From the geometrical point of view, the enhancement of the instability of the dynamics with respect to the parametric instability due to curvature fluctuations, is linked to the probability of obtaining negative sectional curvatures along the geodesics (as discussed for $`1d`$ XY model in Ref.). In fact, when vortices are present in the system, there will surely be two neighbouring spins with an orientation difference greater than $`\pi /2`$, such that, if $`i,j`$ and $`i+1,j`$ are their coordinates on the lattice, it follows that
$$q_{i+1,j}q_{i,j}>\frac{\pi }{2}\mathrm{cos}(q_{i+1,j}q_{i,j})<0.$$
(19)
The sectional curvature relative to the plane defined by the velocity $`๐ฏ`$ along a geodesic and a generic vector $`\xi ๐ฏ`$ is
$$K^{(2)}=\underset{i,j,k,l=1}{\overset{N}{}}\frac{^2V}{q_{i,j}q_{k,l}}\frac{\xi ^{i,j}\xi ^{k,l}}{\xi ^2}.$$
(20)
For the $`2d`$ XY model, it is
$$K^{(2)}=\frac{J}{\xi ^2}\underset{i,j=1}{\overset{N}{}}\{\mathrm{cos}(q_{i+1,j}q_{i,j})[\xi ^{i+1,j}\xi ^{i,j}]^2+\mathrm{cos}(q_{i,j+1}q_{i,j})[\xi ^{i,j+1}\xi ^{i,j}]^2\}.$$
(21)
Thus, a large probability of having a negative value of the cosine of the difference among the directions of two close spins corresponds to a larger probability of obtaining negative values of the sectional curvatures along the geodesics; here for $`\xi `$ the geodesic separation vector of Eq.(11) is chosen.
In the $`3d`$ model, the sectional curvature relative to the plane defined by the velocity $`๐ฏ`$ and a generic vector $`\xi ๐ฏ`$ is
$`K^{(2)}`$ $`=`$ $`{\displaystyle \frac{J}{\xi ^2}}{\displaystyle \underset{i,j,k=1}{\overset{N}{}}}\{\mathrm{cos}(q_{i+1,j,k}q_{i,j,k})[\xi ^{i+1,j,k}\xi ^{i,j,k}]^2+`$ (22)
$`+`$ $`\mathrm{cos}(q_{i,j+1,k}q_{i,j,k})[\xi ^{i,j+1,k}\xi ^{i,j,k}]^2+\mathrm{cos}(q_{i,j,k+1}q_{i,j,k})[\xi ^{i,j,k+1}\xi ^{i,j,k}]^2\}`$ (23)
and again the probability of finding negative values of $`K^{(2)}`$ along a trajectory is limited to the probability of finding vortices.
The mean values of the geometric quantities entering Eq.(13) can be numerically computed by means of Monte Carlo simulations or by means of time averages along the dynamical trajectories. In fact, due to the lack of an explicit expression for the canonical partition function of the system, these averages are not analytically computable. For sufficiently high temperatures, the potential energy becomes negligible with respect to the kinetic energy, and each spin is free to move independently from the others. Thus, in the limit of high temperatures, one can estimate the configurational partition function $`Z_C=_\pi ^\pi _๐ขdq_๐ขe^{\beta V(q)}`$ by means of the expression
$`Z_C`$ $`=`$ $`e^{2\beta JN}{\displaystyle _\pi ^\pi }{\displaystyle \underset{i,j=1}{\overset{N}{}}}dq_{i,j}\mathrm{exp}\{\beta J{\displaystyle \underset{i,j=1}{\overset{N}{}}}[\mathrm{cos}(q_{i+1,j}q_{i,j})+\mathrm{cos}(q_{i,j+1}q_{i,j})]\}`$ (24)
$``$ $`e^{2\beta JN}{\displaystyle _\pi ^\pi }{\displaystyle \underset{i,j=1}{\overset{N}{}}}du_{i,j}dv_{i,j}\mathrm{exp}\{\beta J{\displaystyle \underset{i,j=1}{\overset{N}{}}}[\mathrm{cos}(u_{i,j})+\mathrm{cos}(v_{i,j})]\}`$ (25)
after the introduction of $`u_{i,j}=q_{i+1,j}q_{i,j}`$ and $`v_{i,j}=q_{i,j+1}q_{i,j}`$ as independent variables. In this way, some analytical estimates of the average Ricci curvature $`k_0(T)`$ and of its r.m.s. fluctuations $`\sigma _k^2(T)`$ have been obtained for the $`2d`$ model (Fig. 15). For temperatures above the temperature of the Kosterlitz-Thouless transition, these estimates are in agreement with the numerical computations on a $`N=10\times 10`$ lattice. It is confirmed that Hamiltonian dynamical simulations, already on rather small lattices, are useful to predict, with a good approximation, the thermodynamic limit behavior of relevant observables. Moreover, the good quality of the high temperature estimate gives a further information: at the transition temperature, the correlations among the different degrees of freedom are destroyed, confirming the strong chaoticity of the dynamics.
The same high temperature estimates of $`k_0(T)`$ and $`\sigma _k^2(T)`$ have been performed for the $`3d`$ system. In Fig. 16, the numerical determination of $`\sigma _k^2(T)`$ shows the appearance of a very pronounced peak at the phase transition point which is not predicted by the analytic estimate, whereas the average Ricci curvature $`k_0(T)`$ is in agreement with the analytic values of the high temperature estimate, computed by spin decoupling, above the critical temperature, as in the $`2d`$ model.
### B Geometric observables and Lyapunov exponents
We have seen that the largest Lyapunov exponent is sensitive to the phase transition and at the same time we know that it is also related to the average curvature properties of the โmechanical manifoldsโ. Thus, the geometric observables $`k_0(T)`$ and $`\sigma _k^2(T)`$ above considered can be used to estimate the Lyapunov exponents, as well as to detect the phase transition.
In principle, by means of Eq.(14), one can evaluate the largest Lyapunov exponent without any need of dynamics, but simply using global geometric quantities of the manifold associated to the physical system. For $`2d`$ and $`3d`$ XY models, fully analytic computations are possible only in the limiting cases of high and low temperatures. Microcanonical averages of $`k_0`$ and $`\sigma _k^2`$ at arbitrary $`T`$ have been numerically computed through time averages. We can call this hybrid method semi-analytic.
In Fig. 17, the results of the semi-analytic prediction of the Lyapunov exponents for the $`2d`$ model are plotted vs. temperature and compared with the numerical outcomes of the tangent dynamics. As one can see, the prediction formulated on the basis of Eq.(14) underestimates the numerical values given by the tangent dynamics. The semi-analytic prediction can be improved by observing that the replacement of the sectional curvature fluctuation $`\delta K^{(2)}`$ in Eq.(12) with a fraction of the Ricci curvature \[which underlies the derivation of Eq.(13)\] underestimates the frequency of occurrence of negative sectional curvatures, which was already the case of the $`1d`$ XY model . The correction procedure can be implemented by evaluating the probability $`P(T)`$ of obtaining a negative value of the sectional curvature along a generic trajectory and then by operating the substitution
$$K_R(T)\frac{K_R(T)}{1+P(T)\alpha }.$$
(26)
The parameter $`\alpha `$ is a free parameter to be empirically estimated. Its value ranges from $`100`$ to $`200`$, without appreciable differences in the final result. It resumes the non trivial information about the more pronounced tendency of the trajectories towards negative sectional curvatures with respect to the predictions of the geometric model describing the chaoticity of the dynamics.
The probability $`P(T)`$ is estimated through the occurrence along a trajectory of negative values of the sum of the coefficients that appear in the definition of $`K^{(2)}`$ \[Eqs.(21) and (23)\]
$$P(T)\frac{_\pi ^\pi \mathrm{\Theta }(\mathrm{cos}(q_{k+1,l}q_{k,l})\mathrm{cos}(q_{k,l+1}q_{k,l}))\mathrm{exp}[\beta V(๐ช)]_{k,l=1}^Ndq_{k,l}}{_\pi ^\pi \mathrm{exp}[\beta V(๐ช)]_{k,l=1}^Ndq_{k,l}},$$
(27)
averaged over all the sites $`k,l(1,\mathrm{},N)`$; $`\mathrm{\Theta }`$ is the step function.
Alternatively, owing to the already remarked relation between vorticity and sectional curvature $`K^{(2)}`$, $`P(T)`$ can be replaced by the average density of vortices
$$K_R(T)\frac{K_R(T)}{1+\overline{\alpha }๐ฑ(T)},$$
(28)
where $`\overline{\alpha }`$ a free parameter. Actually, in the $`2d`$ model, the two corrections, one given by Eq.(26) with $`P(T)`$ of Eq. (27), the other given by Eq.(28) with the vorticity function in place of $`P(T)`$, convey the same information. The semi-analytic predictions of $`\lambda _1(T)`$ with correction are reported in Fig. 17.
In the limits of high and low temperatures, $`\lambda _1(T)`$ can be given the analytic forms $`\lambda _1(T)T^{1/6}`$ at high temperature, and $`\lambda _1(T)T^2`$ at low temperature. In the former case, the high temperature approximation (25) is used, and in the latter case the quasi-harmonic oscillators approximation is done. The deviation of $`\lambda _1(T)`$ from the quasi-harmonic scaling, starting at $`T0.6`$ and already observed to correspond to the appearance of vortices, finds here a simple explanation through the geometry of dynamics: vortices are associated with negative sectional curvatures, enhancing chaos.
By increasing the spatial dimension of the system, it becomes more and more difficult to accurately estimate the probability of obtaining negative sectional curvatures. The assumption that the occurrence of negative values of the cosine of the difference between the directions of two nearby spins is nearly equal to $`P(T)`$, is less effective in the $`3d`$ model than in the $`2d`$ one. Again, the vorticity function can be assumed as an estimate of $`P(T)`$ \[Eq. (28)\]. The quality of the results has a weak dependence upon the parameter $`\alpha `$. The correction remains good, with $`\alpha `$ belonging to a broad interval of values ($`100รท200`$). In the limits of high and low temperatures, the model predicts correctly the same scaling laws of the $`2d`$ system.
In Fig. 18 the semi-analytic predictions for the Lyapunov exponents, with and without correction, are plotted vs. temperature together with the numerical results of the tangent dynamics. It is noticeable that the prediction of Eq. (14) is able to give the correct asymptotic behavior of the Lyapunov exponents also at low temperatures, the most difficult part to obtain by means of dynamical simulations.
### C A topological hypothesis
We have seen in Fig. 16 that a sharp peak of the Ricci-curvature fluctuations $`\sigma _\kappa ^2(T)`$ is found for the $`3d`$ model in correspondence of the second order phase transition, whereas, for the $`2d`$ model, $`\sigma _\kappa ^2(T)`$ appears regular and in agreement with the theoretically predicted smooth pattern. On the basis of heuristic arguments, in Refs. we suggested that the peak of $`\sigma _\kappa ^2`$ observed for the $`3d`$ XY model, as well as for $`2d`$ and $`3d`$ scalar and vector lattice $`\phi ^4`$ models, might originate in some change of the topology of the mechanical manifolds. In fact, in abstract mathematical models, consisting of families of surfaces undergoing a topology change โ i.e. a loss of diffeomorphicity among them โ at some critical value of a parameter labelling the members of the family, we have actually observed the appearance of cusps of $`\sigma _K^2`$ at the transition point between two subfamilies of surfaces of different topology, $`K`$ being the Gauss curvature.
Actually, for the mean-field XY model, where both $`\sigma _\kappa ^2(T)`$ and $`\lambda _1(T)`$ have theoretically been shown to loose analyticity at the phase transition point, a direct evidence of a โspecialโ change of the topology of equipotential hypersurfaces of configuration space has been given . Other indirect and direct evidences of the actual involvement of topology in the deep origin of phase transitions have been recently given for the lattice $`\phi ^4`$ model.
In the following Section we consider the extension of this topological point of view about phase transitions from equipotential hypersurfaces of configuration space to constant energy hypersurfaces of phase space.
## V Phase space geometry and thermodynamics.
In the preceding Section we have used some elements of intrinsic differential geometry of submanifolds of configuration space to describe the average degree of dynamical instability (measured by the largest Lyapunov exponent). In the present Section we are interested in the relationship between the extrinsic geometry of the constant energy hypersurfaces $`\mathrm{\Sigma }_E`$ and thermodynamics.
Hereafter, phase space is considered as an even-dimensional subset $`\mathrm{\Gamma }`$ of $`^{2N}`$ and the hypersurfaces $`\mathrm{\Sigma }_E=\{(p_1,\mathrm{},p_N,q_1,\mathrm{},q_N)|H(p_1,\mathrm{},p_N,q_1,\mathrm{},q_N)=E\}`$ are manifolds that can be equipped with the standard Riemannian metric induced from $`^{2N}`$. If, for example, a surface is parametrically defined through the equations $`x^i=x^i(z^1,\mathrm{},z^k)`$, $`i=1,\mathrm{},2N`$, then the metric $`g_{ij}`$ induced on the surface is given by $`g_{ij}(z^1,\mathrm{},z^k)=_{n=1}^{2N}\frac{x^n}{z^i}\frac{x^n}{z^j}`$. The geodesic flow associated with the metric induced on $`\mathrm{\Sigma }_E`$ from $`^{2N}`$ has nothing to do with the Hamiltonian flow that belongs to $`\mathrm{\Sigma }_E`$. Nevertheless, it exists an intrinsic Riemannian metric $`g_S`$ of phase space $`\mathrm{\Gamma }`$ such that the geodesic flow of $`g_S`$, restricted to $`\mathrm{\Sigma }_E`$, coincides with the Hamiltonian flow ($`g_S`$ is the so called Sasaki lift to the tangent bundle of configuration space of the Jacobi metric $`g_J`$ that we mentioned in a preceding Section).
The link between extrinsic geometry of the $`\mathrm{\Sigma }_E`$ and thermodynamics is estabilished through the microcanonical definition of entropy
$$S=k_B\mathrm{log}_{\mathrm{\Sigma }_E}\frac{d\sigma }{H},$$
(29)
where $`d\sigma =\sqrt{det(g)}dx_1\mathrm{}dx_{2N1}`$ is the invariant volume element of $`\mathrm{\Sigma }_E^{2N}`$, $`g`$ is the metric induced from $`^{2N}`$ and $`x_1\mathrm{}x_{2N1}`$ are the coordinates on $`\mathrm{\Sigma }_E`$.
Let us briefly recall some necessary definitions and concepts that are needed in the study of hypersurfaces of euclidean spaces.
A standard way to investigate the geometry of an hypersurface $`\mathrm{\Sigma }^m`$ is to study the way in which it curves around in $`^{m+1}`$: this is measured by the way the normal direction changes as we move from point to point on the surface. The rate of change of the normal direction $`๐`$ at a point $`x\mathrm{\Sigma }`$ is described by the shape operator $`L_x(๐ฏ)=_๐ฏ๐=(N_1๐ฏ,\mathrm{},N_{m+1}๐ฏ)`$, where $`๐ฏ`$ is a tangent vector at $`x`$ and $`_๐ฏ`$ is the directional derivative of the unit normal $`๐`$. As $`L_x`$ is an operator of the tangent space at $`x`$ into itself, there are $`m`$ independent eigenvalues $`\kappa _1(x),\mathrm{},\kappa _m(x)`$, which are called the principal curvatures of $`\mathrm{\Sigma }`$ at $`x`$. Their product is the Gauss-Kronecker curvature: $`K_G(x)=_{i=1}^m\kappa _i(x)=\mathrm{det}(L_x)`$, and their sum is the so-called mean curvature: $`M_1(x)=\frac{1}{m}_{i=1}^m\kappa _i(x)`$. The quadratic form $`L_x(๐ฏ)๐ฏ`$, associated with the shape operator at a point $`x`$, is called the second fundamental form of $`\mathrm{\Sigma }`$ at $`x`$.
It can be shown that the mean curvature of the energy hypersurfaces is given by
$$M_1(x)=\frac{1}{2N1}\left(\frac{H(x)}{H(x)}\right),$$
(30)
where $`H(x)/H(x)`$ is the unit normal to $`\mathrm{\Sigma }_E`$ at a given point $`x=(p_1,\mathrm{},p_N,q_1,\mathrm{},q_N)`$, and $`=(/p_1,\mathrm{},/q_N)`$, whence the explicit expression
$`(2N1)M_1=`$ $``$ $`{\displaystyle \frac{1}{H}}\left[N+{\displaystyle \underset{๐ข}{}}\left({\displaystyle \frac{^2V}{q_๐ข^2}}\right)\right]`$ (31)
$`+`$ $`{\displaystyle \frac{1}{H^3}}\left[{\displaystyle \underset{๐ข}{}}p_๐ข^2+{\displaystyle \underset{๐ข,๐ฃ}{}}\left({\displaystyle \frac{^2V}{q_๐ขq_๐ฃ}}\right)\left({\displaystyle \frac{V}{q_๐ข}}\right)\left({\displaystyle \frac{V}{q_๐ฃ}}\right)\right],`$ (32)
where $`๐ข,๐ฃ`$ are multi-indices according to the number of spatial dimensions.
The link between geometry and physics stems from the microcanonical definition of the temperature
$$\frac{1}{T}=\frac{S}{E}=\frac{1}{\mathrm{\Omega }_\nu (E)}\frac{d\mathrm{\Omega }_\nu (E)}{dE},$$
(33)
where we used Eq.(29) with $`k_B=1`$, $`\nu =2N1`$, and $`\mathrm{\Omega }_\nu (E)=_{\mathrm{\Sigma }_E}๐\sigma /H`$. ยฟFrom the formula
$$\frac{d^k}{dE^k}\left(_{\mathrm{\Sigma }_E}\alpha ๐\sigma \right)(E^{})=_{\mathrm{\Sigma }_E^{}}A^k(\alpha )๐\sigma ,$$
(34)
where $`\alpha `$ is an integrable function and $`A`$ is the operator $`A(\alpha )=\frac{}{H}\left(\alpha \frac{H}{H}\right),`$ it is possible to work out the result
$$\frac{1}{T}=\frac{1}{\mathrm{\Omega }_\nu (E)}\frac{d\mathrm{\Omega }_\nu (E)}{dE}=\frac{1}{\mathrm{\Omega }_\nu }_{\mathrm{\Sigma }_E}\frac{d\sigma }{H}\left[2\frac{M_1^{}}{H}\frac{\mathrm{}H}{H^2}\right]\frac{1}{\mathrm{\Omega }_\nu }_{\mathrm{\Sigma }_E}\frac{d\sigma }{H}\frac{M_1^{}}{H},$$
(35)
where $`M_1^{}=(H/H)`$ is directly proportional to the mean curvature (30). In the last term of Eq.(35) we have neglected a contribution which vanishes as $`๐ช(1/N)`$. Eq. (35) provides the fundamental link between extrinsic geometry and thermodynamics . In fact, the microcanonical average of $`M_1^{}/H`$, which is a quantity tightly related with the mean curvature of $`\mathrm{\Sigma }_E`$, gives the inverse of the temperature, whence other important thermodynamic observables can be derived. For example, the constant volume specific heat
$$\frac{1}{C_V}=\frac{T(E)}{E},$$
(36)
using Eq.(33), yields
$$C_V=\left(\frac{S}{E}\right)^2\left(\frac{^2S}{E^2}\right)^1,$$
(37)
becoming at large $`N`$
$$C_V=\frac{M_1^{}}{H}_{mc}^2\left[\frac{1}{\mathrm{\Omega }_\nu }\frac{d}{dE}_{\mathrm{\Sigma }_E}\frac{d\sigma }{H}\left(\frac{M_1^{}}{H}+R(E)\right)\frac{M_1^{}}{H}_{mc}^2\right]^1,$$
(38)
where the subscript $`mc`$ stands for microcanonical average, and $`R(E)`$ stands for the quantities of order $`๐ช(1/N)`$ neglected in the last term of Eq.(33) (a-priori, its derivative can be non negligible and has to be taken into account). Eq. (38) highlights a more elaborated link between geometry and thermodynamics: the specific heat depends upon the microcanonical average of $`M_1^{}/H`$ and upon the energy variation rate of the surface integral of this quantity.
Remarkably, the relationship between curvature properties of the constant energy surfaces $`\mathrm{\Sigma }_E`$ and thermodynamic observables given by Eqs.(33) and (38) can be extended to embrace also a deeper and very interesting relationship between thermodynamics and topology of the constant energy surfaces. Such a relationship can be discovered through a reasoning which, though approximate, is highly non-trivial, for it makes use of a deep theorem due to Chern and Lashof . As $`H=\{_ip_i^2+[_iV(q)]^2\}^{1/2}`$ is a positive quantity increasing with the energy, we can write
$$\frac{1}{T}=\frac{1}{\mathrm{\Omega }_\nu }\frac{d\mathrm{\Omega }_\nu }{dE}\frac{1}{\mathrm{\Omega }_\nu }_{\mathrm{\Sigma }_E}\frac{d\sigma }{H}\frac{M_1^{}}{H}=D(E)\frac{1}{\mathrm{\Omega }_\nu }_{\mathrm{\Sigma }_E}๐\sigma M_1,$$
(39)
where we have introduced the factor function $`D(E)`$ in order to extract the total mean curvature $`_{\mathrm{\Sigma }_E}๐\sigma M_1`$; $`D(E)`$ has been numerically found to be smooth and very close to $`1/H^2_{mc}`$ (see Section V A and Fig. 19). Then, recalling the expression of a multinomial expansion
$$(x_1+\mathrm{}+x_\nu )^\nu =\underset{_{\{n_i\},{\scriptscriptstyle n_k}=\nu }}{}\frac{\nu !}{n_1!\mathrm{}n_\nu !}x_1^{n_1}\mathrm{}x_\nu ^{n_\nu },$$
(40)
and identifying the $`x_i`$ with the principal curvatures $`k_i`$, one obtains
$$M_1^\nu =\nu !\underset{i=1}{\overset{\nu }{}}k_i+R=\nu !K+R,$$
(41)
where $`K=_ik_i`$ is the Gauss-Kronecker curvature, and $`R`$ is the sum (40) without the term with the largest coefficient ($`n_k=1,k`$). Using $`\nu !\nu ^\nu e^\nu \sqrt{2\pi \nu }`$,
$$M_1^\nu \nu ^\nu e^\nu \sqrt{4\pi N}K+R$$
(42)
is obtained. The above mentioned theorem of Chern and Lashof states that
$$_{\mathrm{\Sigma }_E}|K|๐\sigma Vol[๐_1^\nu ]\underset{i=0}{\overset{\nu }{}}b_i(\mathrm{\Sigma }_E),$$
(43)
i.e. the total absolute Gauss-Kronecker curvature of a hypersurface is related with the sum of all its Betti numbers $`b_i(\mathrm{\Sigma }_E)`$. The Betti numbers are diffeomorphism invariants of fundamental topological meaning , therefore their sum is also a topologic invariant. $`๐_1^\nu `$ is a hypersphere of unit radius. Combining Eqs. (42) and (43) and integrating on $`\mathrm{\Sigma }_E`$, we obtain
$$_{\mathrm{\Sigma }_E}|M_1^\nu |๐\sigma \nu ^\nu e^\nu \sqrt{2\pi \nu }_{\mathrm{\Sigma }_E}|K|๐\sigma +_{\mathrm{\Sigma }_E}|R|๐\sigma ๐\underset{i=0}{\overset{\nu }{}}b_i(\mathrm{\Sigma }_E)+(E),$$
(44)
with the shorthands $`๐=\nu ^\nu e^\nu Vol(S_1^\nu )`$ and $`=_{\mathrm{\Sigma }_E}|R|๐\sigma `$.
Now, with the aid of the inequality $`f^{1/n}๐\mu f๐\mu ^{1/n}`$, we can write
$$_{\mathrm{\Sigma }_E}|M_1|๐\sigma =_{\mathrm{\Sigma }_E}|M_1^\nu |^{1/\nu }๐\sigma \left|_{\mathrm{\Sigma }_E}M_1^\nu ๐\sigma \right|^{1/\nu }.$$
(45)
If $`M_10`$ everywhere on $`\mathrm{\Sigma }_E`$, then $`\left|_{\mathrm{\Sigma }_E}M_1^\nu ๐\sigma \right|^{1/\nu }=\left(_{\mathrm{\Sigma }_E}|M_1^\nu |๐\sigma \right)^{1/\nu }`$, whence, in the hypothesis that $`M_10`$ is largely prevailing , $`\left|_{\mathrm{\Sigma }_E}M_1^\nu ๐\sigma \right|^{1/\nu }\left(_{\mathrm{\Sigma }_E}|M_1^\nu |๐\sigma \right)^{1/\nu }`$. Under the same assumption, $`_{\mathrm{\Sigma }_E}M_1๐\sigma _{\mathrm{\Sigma }_E}|M_1|๐\sigma `$ and therefore
$$_{\mathrm{\Sigma }_E}M_1๐\sigma _{\mathrm{\Sigma }_E}|M_1^\nu |^{1/\nu }๐\sigma \left|_{\mathrm{\Sigma }_E}M_1^\nu ๐\sigma \right|^{1/\nu }\left(_{\mathrm{\Sigma }_E}|M_1^\nu |๐\sigma \right)^{1/\nu }\left[๐\underset{i=0}{\overset{\nu }{}}b_i(\mathrm{\Sigma }_E)+(E)\right]^{1/\nu }.$$
(46)
Finally,
$`{\displaystyle \frac{1}{T(E)}}={\displaystyle \frac{1}{\mathrm{\Omega }_\nu }}{\displaystyle \frac{d\mathrm{\Omega }_\nu }{dE}}`$ $``$ $`{\displaystyle \frac{M_1^{}}{H}}_{mc}={\displaystyle \frac{1}{\mathrm{\Omega }_\nu }}{\displaystyle _{\mathrm{\Sigma }_E}}{\displaystyle \frac{d\sigma }{H}}{\displaystyle \frac{M_1^{}}{H}}=D(E){\displaystyle \frac{1}{\mathrm{\Omega }_\nu }}{\displaystyle _{\mathrm{\Sigma }_E}}๐\sigma M_1`$ (47)
$``$ $`{\displaystyle \frac{D(E)}{\mathrm{\Omega }_\nu }}\left[๐{\displaystyle \underset{i=0}{\overset{\nu }{}}}b_i(\mathrm{\Sigma }_E)+(E)\right]^{1/\nu }.`$ (48)
Equation (48) has the remarkable property of relating the microcanonical definition of temperature of Eq.(39) with a topologic invariant of $`\mathrm{\Sigma }_E`$. The Betti numbers can be exponentially large with $`N`$ \[for example, in the case of $`N`$-tori $`๐^N`$, they are $`b_k=\left(\genfrac{}{}{0pt}{}{N}{k}\right)`$\], so that the quantity $`(b_k)^{1/N}`$ can converge, at arbitrarily large $`N`$, to a non-trivial limit (i.e. different from one). Thus, even though the energy dependence of $``$ is unknown, the energy variation of $`b_i(\mathrm{\Sigma }_E)`$ must be mirrored โ at any arbitrary $`N`$ โ by the energy variation of the temperature. By considering Eq.(38) in the light of Eq.(48), we can expect that some suitably abrupt and major change in the topology of the $`\mathrm{\Sigma }_E`$ can reflect into the appearance of a peak of the specific heat, as a consequence of the associated energy dependence of $`b_k(\mathrm{\Sigma }_E)`$ and of its derivative with respect to $`E`$. In other words, we see that a link must exist between thermodynamical phase transitions and suitable topology changes of the constant energy submanifolds of the phase space of microscopic variables. The arguments given above, though in a still rough formulation, provide a first attempt to make a connection between the topological aspects of the microcanonical description of phase transitions and the already proposed topological hypothesis about topology changes in configuration space and phase transitions .
Direct support to the topological hypothesis has been given by the analytic study of a mean-field XY model and by the numerical computation of the Euler characteristic $`\chi `$ of the equipotential hypersurfaces $`\mathrm{\Sigma }_v`$ of the configuration space in a $`2d`$ lattice $`\phi ^4`$ model . The Euler characteristic is the alternate sum of all the Betti numbers of a manifold, so it is another topological invariant, but it identically vanishes for odd dimensional manifolds, like the $`\mathrm{\Sigma }_E`$. In Ref., $`\chi (\mathrm{\Sigma }_v)`$ neatly reveals the symmetry-breaking phase transition through a sudden change of its variation rate with the potential energy density $`v`$. A sudden โsecond orderโ variation of the topology of the $`\mathrm{\Sigma }_v`$ appears in both Refs. as the requisite for the appearance of a phase transition. These results strenghten the arguments given in the present Section about the role of the topology of the constant energy hypersurfaces. In fact, the larger is $`N`$, the smaller are the relative fluctuations $`\delta ^2V^{1/2}/V`$ and $`\delta ^2K^{1/2}/K`$ of the potential and kinetic energies respectively. At very large $`N`$, the product manifold $`\mathrm{\Sigma }_v^{N1}\times ๐_t^{N1}`$, with $`vV`$ and $`tK`$, $`v+t=E`$, is a good model manifold to represent the part of $`\mathrm{\Sigma }_E`$ that is overwhelmingly sampled by the dynamics and that therefore constitutes the effective support of the microcanonical measure on $`\mathrm{\Sigma }_E`$. The kinetic energy submanifolds $`๐_t^{N1}=\{(p_1,\mathrm{},p_N)^N|_{i=1}^N\frac{1}{2}p_i^2=t\}`$ are hyperspheres.
In other words, at very large $`N`$ the microcanonical measure mathematically extends over a whole energy surface but, as far as physics is concerned, a non-negligible contribution to the microcanonical measure is in practice given only by a small subset of an energy surface. This subset can be reasonably modeled by the product manifold $`\mathrm{\Sigma }_v^{N1}\times ๐_t^{N1}`$, because the total kinetic and total potential energies - having arbitrarily small fluctuations, provided that $`N`$ is large enough - can be considered almost constant. Thus, since $`๐_t^{N1}`$ at any $`t`$ is always an hypersphere, a change in the topology of $`\mathrm{\Sigma }_v^{N1}`$ directly entails a change of the topology of $`\mathrm{\Sigma }_v^{N1}\times ๐_t^{N1}`$, that is of the effective model-manifold for the subset of $`\mathrm{\Sigma }_E`$ where the dynamics mainly โlivesโ at a given energy $`E`$.
At small $`N`$, the model with a single product manifold is no longer good and should be replaced by the non-countable union $`_v\mathrm{\Sigma }_v^{N1}\times ๐_{Ev}^{N1}`$, with $`v`$ assuming all the possible values in a real interval $``$. From this fact the smoothing of the energy dependence of thermodynamic variables follows. Nevertheless, the geometric and topologic signals of the phase transition can remain much sharper than the thermodynamic signals also at small $`N`$ $`(<100)`$, as it is witnessed by the $`2d`$ lattice $`\phi ^4`$ model .
Finally, let us comment about the relationship between intrinsic geometry, in terms of which we discussed the geometrization of the dynamics, and extrinsic geometry, dealt with in the present Section.
The most direct and intriguing link is estabilished by the expression for microcanonical averages of generic observables of the kind $`A(q)`$, with $`q=(q_1,\mathrm{},q_N)`$,
$$A_{mc}=\frac{1}{\mathrm{\Omega }_{2N}(E)}_{H(p,q)E}d^Npd^NqA(q)=\frac{1}{Vol(M_E)}_{V(q)E}d^Nq[EV(q)]^{N/2}A(q),$$
(49)
where $`M_E=\{q^N|V(q)E\}`$. Eq. (49) is obtained by means of a Laplace-transform method ; it is remarkable that $`[EV(q)]^Ndet(g_J)`$, where $`g_J`$ is the Jacobi metric whose geodesic flow coincides with newtonian dynamics (see Section $`V`$), therefore $`d^Nq[EV(q)]^{N/2}d^Nq\sqrt{det(g_J)}`$ is the invariant Riemannian volume element of $`(M_E,g_J)`$. Thus,
$$\frac{1}{Vol(M_E)}_{V(q)E}d^Nq[EV(q)]^{N/2}A(q)\frac{1}{Vol(M_E)}_{M_E}d^Nq\sqrt{det(g_J)}A(q),$$
(50)
which means that the microcanonical averages $`A(q)_{mc}`$ can be expressed as Riemannian integrals on the mechanical manifold $`(M_E,g_J)`$.
In particular, this also applies to the microcanonical definition of entropy
$$S=k_B\mathrm{log}_{H(p,q)E}d^Npd^Nq=k_B\mathrm{log}_0^E๐E^{}_{\mathrm{\Sigma }_E^{}}\frac{d\sigma }{H},$$
(51)
which is alternative to that given in Eq.(29), though equivalent to it in the large $`N`$ limit. We have
$`S`$ $`=`$ $`k_B\mathrm{log}\left[{\displaystyle \frac{1}{C\mathrm{\Gamma }(N/2+1)}}{\displaystyle _{V(q)E}}d^Nq[EV(q)]^{N/2}\right]`$ (52)
$``$ $`k_B\mathrm{log}{\displaystyle _{M_E}}d^Nq\sqrt{det(g_J)}+const.,`$ (53)
where the last term gives the entropy as the logarithm of the Riemannian volume of the manifold.
The topology changes of the surfaces $`\mathrm{\Sigma }_v^{N1}`$, that are to be associated with phase transitions, will deeply affect also the geometry of the mechanical manifolds $`(M_E,g_J)`$ and $`(M\times ^2,g_E)`$ and, consequently, they will affect the average instability properties of their geodesic flows. In fact, Eq.(14) links some curvature averages of these manifolds with the numeric value of the largest Lyapunov exponent. Loosely speaking, major topology changes of $`\mathrm{\Sigma }_v^{N1}`$ will affect microcanonical averages of geometric quantities computed through Eq.(49), likewise entropy, computed through Eq.(53).
Thus, the peculiar temperature patterns displayed by the largest Lyapunov exponent at a second-order phase transition point โ in the present paper reported for the $`3d`$ $`XY`$ model, in Ref. reported for lattice $`\phi ^4`$ models โ appear as reasonable consequences of the deep variations of the topology of the equipotential hypersurfaces of configuration space.
We notice that topology seems to provide a common ground to the roots of microscopic dynamics and of thermodynamics and, notably, it can account for major qualitative changes simultaneously occurring in both dynamics and thermodynamics when a phase transition is present.
### A Some preliminary numerical computations
Let us briefly report on some preliminary numerical computations concerning the extrinsic geometry of the hypersurfaces $`\mathrm{\Sigma }_E`$ in the case of the $`3d`$ XY model.
The first point about extrinsic geometry that we numerically addressed was to check whether the inverse of the temperature, that appears in Eq.(39), can be reasonably factorized into the product of a smooth โdeformation factorโ $`D(E)`$ and of the total mean curvature $`_{\mathrm{\Sigma }_E}M_1๐\sigma `$. To this purpose, the two independently computed quantities $`1/H^2_{mc}`$ and $`D(E)=[_{\mathrm{\Sigma }_E}(d\sigma /H)(M_1^{}/H)]/[_{\mathrm{\Sigma }_E}๐\sigma M_1]`$ are compared in Fig. 19, showing that actually $`_{\mathrm{\Sigma }_E}(d\sigma /H)(M_1^{}/H)1/H^2_{mc}_{\mathrm{\Sigma }_E}๐\sigma M_1`$. In other words, $`D(E)1/H^2_{mc}`$ and no โsingularโ feature in its energy pattern seems to exist, what suggests that $`_{\mathrm{\Sigma }_E}๐\sigma M_1`$ has to convey all the information relevant to the detection of the phase transition. There is no reason to think that the validity of the factorization given in Eq.(39) is limited to the special case of the XY model.
The other point that we tackled concerns an indirect quantification of how a phase space trajectory curves around and knots on the $`\mathrm{\Sigma }_E`$ to which it belongs. We can expect that the way in which an hypersurface $`\mathrm{\Sigma }_E`$ is โfilledโ by a phase space trajectory living on it will be affected by the geometry and the topology of the $`\mathrm{\Sigma }_E`$. In particular, we computed the normalized autocorrelation function of the time series $`M_1[x(t)]`$ of the mean curvature at the points of $`\mathrm{\Sigma }_E`$ visited by the phase space trajectory, that is, the quantity
$$\mathrm{\Gamma }(\tau )=\delta M_1(t+\tau )\delta M_1(t)_t,$$
(54)
where $`\delta M_1(t)=M_1(t)M_1(t^{})_t^{}`$ is the fluctuation with respect to the average (the โprocessโ $`M_1(t)`$ is supposed stationary). Our aim was to highlight the extrinsic geometric-dynamical counterpart of a symmetry-breaking phase transition.
The practical computation of $`\mathrm{\Gamma }(\tau )`$ proceeds by working out the Fourier power spectrum $`|\stackrel{~}{M}_1(\omega )|^2`$ of $`M_1[x(t)]`$, obtained by averaging $`15`$ spectra computed by an FFT algorithm with a mesh of $`2^{15}`$ points and a sampling time $`\mathrm{\Delta }t=0.1`$. Some typical results for $`\mathrm{\Gamma }(\tau )`$, obtained at different temperatures, are reported in Fig.20. The patterns $`\mathrm{\Gamma }(\tau )`$ display a first regime of very fast decay, which is not surprising because of the chaoticity of the trajectories at any energy, followed by a longer tail of slower decay. An autocorrelation time $`\tau _{corr}`$ can be defined through the first intercept of $`\mathrm{\Gamma }(\tau )`$ with an almost-zero level ($`\mathrm{\Gamma }=0.01`$). In Fig.21 we report the values of $`\tau _{corr}`$ so defined vs. temperature. In correspondence of the phase transition (whose critical temperature is marked by a vertical dotted line), $`\tau _{corr}`$ changes its temperature dependence: by lowering the temperature, below the transition $`\tau _{corr}(T)`$ rapidly increases, whereas it mildly decreases above the transition. Below $`T0.9`$, where the vortices disappear, the autocorrelation functions of $`M_1`$ look quite different and it seems no longer possible to coherently define a correlation time. This result has an intuitive meaning and confirms that the phase transition corresponds to a change in the microscopic dynamics, as already signaled by the largest Lyapunov exponent; however, notice that the correlation times $`\tau _{corr}(T)`$ are much longer than the inverse values of the corresponding $`\lambda _1(T)`$. Qualitatively, $`\lambda _1(T)`$ and $`\tau _{corr}^1(T)`$ look similar, however the two functions are not simply related.
## VI Discussion and perspectives
The microscopic Hamiltonian dynamics of the classical Heisenberg XY model in two and three spatial dimensions has been numerically investigated. This has been possible after the addition to the Heisenberg potentials of a standard (quadratic) kinetic energy term. Special emphasis has been given to the study of the dynamical counterpart of phase transitions, detected through the time averages of conventional thermodynamic observables, and to the new mathematical concepts that are brought about by Hamiltonian dynamics.
The motivations of the present study are given in the Introduction. Let us now summarize what are the outcomes of our investigations and comment about their meaning. There are three main topics, tightly related one to the other:
* the phenomenological description of phase transitions through the natural, microscopic dynamics in place of the usual Monte Carlo stochastic dynamics;
* the investigation, in presence of phase transitions, of certain aspects of the (intrinsic) geometry of the mechanical manifolds where the natural dynamics is represented as a geodesic flow;
* the discussion of the relationship between the (extrinsic) geometry of constant energy hypersurfaces of phase space and thermodynamics.
About the first point, we have found that microscopic Hamiltonian dynamics very clearly evidences the presence of a second order phase transition through the time averages of conventional thermodynamic observables. Moreover, the familiar sharpening effects, at increasing $`N`$, of the specific heat peak and of the order parameter bifurcation are observed. The evolution of the order parameter with respect to the physical time (instead of the fictitious Monte Carlo time) is also accessible, showing the appearance of Goldstone modes and that, in presence of a second order phase transition, there is a clear tendency to the freezing of transverse fluctuations of the order parameter when $`N`$ is increased. The โfreezingโ is observed together with a reduction of the longitudinal fluctuations, i.e. the rotation of the magnetization vector slows down, preparing the breaking of the $`O(2)`$ symmetry at $`N\mathrm{}`$. At variance, when a Kosterlitz-Thouless transition is present, at increasing $`N`$ the magnetization vector has a faster rotation and a smaller norm, preparing the absence of symmetry-breaking in the $`N\mathrm{}`$ limit as expected.
Remarkably, to detect phase transitions, microscopic Hamiltonian dynamics provides us with additional observables of purely dynamical nature, i.e. without statistical counterpart: Lyapunov exponents. Similarly to what we and other authors already reported for other models (see Introduction), also in the case of the $`3d`$ XY model a peculiar temperature pattern of the largest Lyapunov exponent shows up in presence of the second order phase transition, signaled by a โcuspyโ point. By comparing the patterns $`\lambda _1(T)`$ given by Hamiltonian dynamics and by a suitably defined random dynamics respectively, we suggest that the transition between thermodynamically ordered and disordered phases has its microscopic dynamical counterpart in a transition between weak and strong chaos. Though $`aposteriori`$ physically reasonable, this result is far from obvious, because the largest Lyapunov exponent measures the average local instability of the dynamics, which $`apriori`$ has little to do with a collective, and therefore global, phenomenon such as a phase transition. The effort to understand the reason of such a sensitivity of $`\lambda _1`$ to a second order phase transition and to other kinds of transitions, as mentioned in the Introduction, is far reaching.
Here we arrive to the second point listed above. In the framework of a Riemannian geometrization of Hamiltonian dynamics, the largest Lyapunov exponent is related to the curvature properties of suitable submanifolds of configuration space whose geodesics coincide with the natural motions. In the mathematical light of this geometrization of the dynamics, and after the numerical evidence of a sharp peak of curvature fluctuations at the phase transition point, the peculiar pattern of $`\lambda _1(T)`$ is due to some major change occurring to the geometry of mechanical manifolds at the phase transition. Elsewhere, we have conjectured that indeed some major change in the topology of configuration space submanifolds should be the very source of the mentioned major change of geometry.
Thus, we have made a first attempt to provide an analytic argument supporting this topological hypothesis (third point of the above list). This is based on the appearance of a non trivial relationship between the geometry of constant energy hypersurfaces of phase space with their topology and with the microcanonical definition of thermodynamics. Even still in a preliminary formulation, our reasoning already seems to indicate the topology of energy hypersurfaces as the best candidate to explain the deep origin of the dynamical signature of phase transitions detected through $`\lambda _1(T)`$.
The circumstance, mentioned in the preceding Section, of the persistence at small $`N`$ of geometric and topologic signals of the phase transition that are much sharper than the thermodynamic signals is of prospective interest for the study of phase transition phenomena in finite, small systems, a topic of growing interest thanks to the modern developments \- mainly experimental - about the physics of nuclear, atomic and molecular clusters, of conformational phase transitions in homopolymers and proteins, of mesoscopic systems, of soft-matter systems of biological interest. In fact, some unambiguous information for small systems - even about the existence itself of a phase transition - could be better obtained by means of concepts and mathematical tools outlined here and in the quoted papers. Here we also join the very interesting line of thought of Gross and collaborators about the microcanonical description of phase transitions in finite systems.
Let us conclude with a speculative comment about another possible direction of investigation related with this signature of phase transitions through Lyapunov exponents. In a field-theoretic framework, based on a path-integral formulation of classical mechanics , Lyapunov exponents are defined through the expectation values of suitable operators. In the field-theoretic framework, ergodicity breaking appears related to a supersymmetry breaking , and Lyapunov exponents are related to mathematical objects that have many analogies with topological concepts .
The new mathematical concepts and methods, that the Hamiltonian dynamical approach brings about, could hopefully be useful also in the study of more โexoticโ transition phenomena than those tackled in the present work. Besides the above mentioned soft-matter systems, this could be the case of transition phenomena occurring in amorphous and disordered materials.
## VII acknowledgments
We warmly thank L. Casetti, E.G.D. Cohen,R. Franzosi and L. Spinelli for many helpful discussions. During the last year C.C. has been supported by the NSF (Grant # 96-03839) and by the La Jolla Interfaces in Science program (sponsored by the Burroughs Wellcome Fund). This work has been partially supported by I.N.F.M., under the PAIS Equilibrium and non-equilibrium dynamics in condensed matter systems, which is hereby gratefully acknowledged.
## VIII Appendix A
Let us briefly explain how a random markovian dynamics is constructed on a given constant energy hypersurface of phase space. The goal is to compare the energy dependence of the largest Lyapunov exponent computed for the Hamiltonian flow and for a suitable random walk respectively. One has to devise an algorithm to generate a random walk on a given energy hypersurface such that, once the time interval $`\mathrm{\Delta }t`$ separating two successive steps is assigned, the average increments of the coordinates are equal to the average increments of the same coordinates for the differentiable dynamics integrated with a time step $`\mathrm{\Delta }t`$. In other words, the random walk has to roughly mimick the differentiable dynamics with the exception of its possible time-correlations.
One starts with a random initial configuration of the coordinates $`q_i,i=1,2,\mathrm{},N`$, uniformly distributed in the interval $`[0,2\pi ]`$, and with a random gaussian-distributed choice of the coordinates $`p_i`$. The random pseudo-trajectory is generated according to the simple scheme
$`(q_i)_{(k+1)\mathrm{\Delta }t}`$ $``$ $`(q_i)_{k\mathrm{\Delta }t}+\alpha _qG_{i,k}\mathrm{\Delta }t`$ (55)
$`(p_i)_{(k+1)\mathrm{\Delta }t}`$ $``$ $`(p_i)_{k\mathrm{\Delta }t}+\alpha _pG_{i,k}\mathrm{\Delta }t,`$ (56)
where $`\mathrm{\Delta }t`$ is the time interval associated to one step $`kk+1`$ in the markovian chain, $`G_{i,k}`$ are gaussian distributed random numbers with zero expectation value and unit variance; the parameters $`\alpha _q`$ and $`\alpha _p`$ are the variances of the processes $`(q_i)_k`$ and $`(p_i)_k`$. These variances are functions of the energy per degree of freedom $`\epsilon `$. They have to be set equal to the numerically computed average increments of the coordinates obtained along the differentiable trajectories integrated with the same time step $`\mathrm{\Delta }t`$, that is
$`\alpha _q(ฯต)`$ $`=`$ $`\left[{\displaystyle \frac{1}{N}}{\displaystyle \underset{i=1}{\overset{N}{}}}{\displaystyle \frac{(q_i(t+\mathrm{\Delta }t)q_i(t))^2}{\mathrm{\Delta }t}}\right]^{1/2}_t\left[{\displaystyle \frac{1}{N}}{\displaystyle \underset{i=1}{\overset{N}{}}}p_i^2\right]^{1/2}_t\sqrt{T}`$ (57)
$`\alpha _p(ฯต)`$ $`=`$ $`\left[{\displaystyle \frac{1}{N}}{\displaystyle \underset{i=1}{\overset{N}{}}}{\displaystyle \frac{(p_i(t+\mathrm{\Delta }t)p_i(t))^2}{\mathrm{\Delta }t}}\right]^{1/2}_t\left[{\displaystyle \frac{1}{N}}{\displaystyle \underset{i=1}{\overset{N}{}}}\dot{p}_i^2\right]^{1/2}_t,`$ (58)
where $`T`$ is the temperature. Then, in order to make minimum the energy fluctuations around any given value of the total energy, a criterium to accept or reject a new step along the markovian chain has to be assigned. A similar problem has been considered by Creutz, who developed a Monte Carlo microcanonical algorithm , where a โMaxwellian demonโ gives a part of its energy to the system to let it move to a new configuration, or gains energy from the system, if the new proposed configuration produces an energy lowering. If the demon does not have enough energy to allow an energy increasing update of the coordinates, no coordinate change is performed. In this way, the total energy remains almost constant with only small fluctuations. As usual in Monte Carlo simulations, it is appropriate to fix the parameters so as the acceptance rate of the proposed updates of the configurations is in the range $`30\%`$$`60\%`$.
A reliability check of the so defined random walk, and of the adequacy of the phase space sampling through the number of steps adopted in each run, is obtained by computing the averages of typical thermodynamic observables of known temperature dependences.
An improvement to the above described โdemonโ algorithm has been obtained through a simple reprojection on $`\mathrm{\Sigma }_E`$ of the updated configurations ; the coordinates generated by means of (56) are corrected with the formulae
$$q_i(k\mathrm{\Delta }t)q_i(k\mathrm{\Delta }t)+\left[\frac{(\frac{H}{q_i})\mathrm{\Delta }E}{_{i=1}^N(p_j^2+(\frac{H}{q_j})^2)}\right]_{x_R(k\mathrm{\Delta }t)}$$
(59)
$`p_i(k\mathrm{\Delta }t)p_i(k\mathrm{\Delta }t)\left[{\displaystyle \frac{p_i\mathrm{\Delta }E}{_{j=1}^N(p_j^2+(\frac{H}{q_j})^2)}}\right]_{x_R(k\mathrm{\Delta }t)},`$
where $`\mathrm{\Delta }E`$ is the difference between the energy of the new configuration and the reference energy, and $`x_R(k\mathrm{\Delta }t)`$ denotes the random phase space trajectory. At each assigned energy, the computation of the largest Lyapunov exponent $`\lambda _1^R`$ of this random trajectory is obtained by means of the standard definition
$$\lambda _1^R=\underset{n\mathrm{}}{lim}\frac{1}{n\mathrm{\Delta }t}\underset{k=1}{\overset{n}{}}\mathrm{log}\frac{\zeta ((k+1)\mathrm{\Delta }t)}{\zeta (k\mathrm{\Delta }t)},$$
(60)
where $`\zeta (t)(\xi (t),\dot{\xi }(t))`$ is given by the discretized version of the tangent dynamics
$$\frac{\xi _i((k+1)\mathrm{\Delta }t)2\xi _i(k\mathrm{\Delta }t)+\xi _i((k1)\mathrm{\Delta }t)}{\mathrm{\Delta }t^2}+\left(\frac{^2V}{q_iq_j}\right)_{x_R(k\mathrm{\Delta }t)}\xi _j(k\mathrm{\Delta }t)=0.$$
(61)
For wide variations of the parameters ($`\mathrm{\Delta }t`$ and acceptance rate), the resulting values of $`\lambda _1^R`$ are in very good agreement. Moreover, the algorithm is sufficiently stable and the final value of $`\lambda _1^R`$ is independent of the choice of the initial condition.
A more refined algorithm could be implemented by constructing a random markovian process $`q(t_k)[q_1(t_k),\mathrm{},q_N(t_k)]`$ performing an importance sampling of the measure $`d\mu =[EV(q)]^{N/21}dq`$ in configuration space. In fact, similarly to what is reported in Eq.(49), one has $`_{H(p,q)=E}d^Npd^Nq=const_{V(q)E}d^Nq[EV(q)]^{N/21}`$. A random process obtained by sampling such a measure โ with the additional property of a relation between the average increment and the physical time step $`\mathrm{\Delta }t`$ as discussed above โ would enter into Eq.(61) to yield $`\lambda _1^R`$. However, this would result in much heavier numerical computations (with some additional technical difficulty at large $`N`$) which was not worth in view of the principal aims of the present work. |
warning/0002/hep-ex0002063.html | ar5iv | text | # 1 Introduction
## 1 Introduction
The detectors having low background and low threshold are required for investigation of rare processes involving low energy neutrino and weak interacting particles. Such detectors can be used effectively for search for the dark matter, for measurement of neutrino magnetic moment and coherent neutrino scattering by nuclei and for investigation of solar neutrino problem.
For above investigations one needs low background detector of mass several $`kg`$ and with threshold less than 1000 $`eV`$. The cryogenic and germanium detectors partly correspond to these requirements. The drawback of the cryogenic detectors is the complexity of their production and use. The drawback of germanium detectors is a rather high threshold $`2รท10`$ $`KeV`$ which is due to a leakage current and electronic and microphonic noises. It would be very attractive to provide effective decreasing of the detector threshold using the internal proportional amplification of signal. Internal proportional amplification in the semiconductor detectors is realized now in the silicon avalanche photodiodes (APD) , where a gain of about $`10^2รท10^4`$ is achieved by avalanche multiplication of electrons at electric field $`(5รท6)\times 10^5`$ $`V/cm`$ in narrow $`pn`$ junction with sensitive volume several $`mm^3`$. Below we shall demonstrate the possibility of realizing the germanium detector with internal amplification (GDA) and present its design.
## 2 Principles of GDA and its design.
In semiconductor detectors ( as in a gas proportional counter-PC, or multiwire proportional chamber-MWPC) the conditions for internal proportional amplification of electrons can be fulfilled. It is known that in APD critical electric field $`E_{cr}`$, that provides multiplication of electrons at room temperature, is equal $`(5รท6)\times 10^5`$ $`V/cm`$ . The $`E_{cr}`$ for germanium at liquid nitrogen temperature can be defined from the dependence of electron drift velocity on electric field and energy of production of electron-hole pairs and photons . For germanium at 77K the magnitude of $`E_{cr}`$ derived in this fashion is equal $`9\times 10^4`$ $`V/cm`$. In APD and gas PC the critical electric field is produced by a different way. In the first case $`E_{cr}`$ is achieved by high concentration of impurities in narrow junction. As a result the sensitive volume of APD is several $`mm^3`$. In gas PC the $`E_{cr}`$ can be achieved by special configuration of the electric field, due to large difference of cathode and anode sizes. In high purity germanium (HPGe) with a sensitive volume near 100 $`cm^3`$ the $`E_{cr}`$ can be obtained by the same manner.
Electric field in the gas cylindrical PC is of the form:
$$E(r)=\frac{V}{r\mathrm{ln}\left(R_2/R_1\right)}$$
(1)
where $`V`$ is applied voltage, $`R_1`$ and $`R_2`$ are the radii of cathode and anode correspondingly, $`r`$ is a distance from anode.
One can see from (1) that at $`V=10^3`$ $`V`$, $`R_1=0.001`$ $`cm`$ and $`R_2=1`$ $`cm`$ $`E(r)`$ is about $`10^5`$ $`V/cm`$ near the anode. Unlike gas PC the electric field in the coaxial detector from high purity germanium is defined not only by $`V`$, $`R_1`$ and $`R_2`$ but the concentration of donor (n-type) or acceptor (p-type) impurities as well. The magnitude of the volume charge in the sensitive volume of crystal depends on these impurities. The electric field in coaxial detector from HPGe with regards to the impurities will be :
$$E(r)=\frac{Ne}{2ฯต}r\frac{\left[V+(Ne/4ฯต)\left(R_2^2R_1^2\right)\right]}{r\mathrm{ln}\left(R_2/R_1\right)},$$
(2)
where $`N`$ is impurity concentration, $`e`$ is the electron charge,$`ฯต`$ is the dielectric constant of germanium.
The electric field (2) can be expressed with a depletion voltage $`V_d`$, which is a minimum applied voltage necessary to neutralize the volume charge and to provide the sensitive region in the whole crystal volume. The $`V_d`$ for the coaxial detector, taking into account that $`R_2R_1`$, may be given by :
$$V_d\frac{Nq}{4ฯต}R_2^2.$$
(3)
and the equation (2) may be rewritten for n-type germanium in the final form :
$$E(r)=\frac{2V_d}{R_2^2}r\frac{VV_d}{r\mathrm{ln}\left(R_2/R_1\right)}.$$
(4)
The dependence $`E(r)`$ on $`r`$ is shown on fig.1 for coaxial detector from HPGe. Electric field near the anode reaches $`E_{cr}`$ \- magnitude required for avalanche multiplication of electrons. The coaxial germanium detector with internal amplification is the more appropriate for the low background spectrometers but the possibility of fabrication of inner electrode of 20 micron radius is highly conjectural presently. So we shall consider the more realistic problem - fabrication of planar germanium detector with internal amplification, multistrip planar germanium detector, similar in design to MWPC. The electric field in MWPC is of the form for one dimension case (the coordinates x and y relate to an centered on the wire, x is parrelel to the wire plane, y is perpendicular) :
$$E(0,y)=\frac{\pi V}{s\left[\frac{\pi L}{s}\mathrm{ln}\frac{\pi d}{s}\right]}\mathrm{coth}\frac{\pi y}{s}$$
(5)
where $`V`$ is applied voltage, $`s`$ is wire spacing , $`d`$ is a diameter of the wire and $`L`$ is the thickness of the planar detector. As in the case of coaxial germanium detector one must take into account for multistrip germanium detector the depletion voltage $`V_d`$. In the case of planar germanium detector $`V_d=\frac{Ne}{2ฯต}L^2`$. The electric field for multistrip germanium detector is of the form:
$$E(0,y)=\frac{2V_d}{L^2}y\frac{\pi \left(VV_d\right)}{s\left[\frac{\pi L}{s}\mathrm{ln}\frac{\pi d}{s}\right]}\mathrm{coth}\frac{\pi y}{s}$$
(6)
where $`d`$ is the strip width. The dependence of $`E(0,y)`$ on $`y`$ for multistrip germanium detector is shown on fig.2. In the cases being considered the electric field near the anode is sufficient for avalanche multiplication of electrons ($`E>10^5`$ $`V/cm`$). The amplification factor can be estimated as $`K=2^{h/l}`$ where $`l`$ is a free electron path for inelastic scattering and $`h`$ is a length of avalanche region where $`E>E_{cr}`$. The $`l`$ in germanium at 77$`K`$ is equal 0.5 micron and for $`L=3`$ $`cm`$ $`h`$ is equal 5 micron (see fig.2) so for this case it is possible to achieve $`K=10^3`$. If one does not need high amplification factor it is possible to decrease $`V`$ or to increase the strip width $`s`$.
The GDA is assumed to use for investigation of rare processes so the spectrometer including GDA must have large mass of detector, it can be fabricated from separate modules of mass about 0.7 $`kg`$ each. One module represents the multistrip planar germanium detector from HPGe having impurity concentration less than $`10^{10}`$ $`cm^3`$, measures $`70\times 70\times 30`$ $`mm^3`$ (see fig.3). The 12 anode strips of 20 micron width and of 65 $`mm`$ length are fabricated by photomask method . The cathode area is equal $`65\times 65`$ $`mm^2`$ and the fiducial volume is equal to 130 $`cm^3`$. There are the guard electrodes in the anode and cathode planes . The anode strips can be connected together however it is more convenient to take signals from separate strips to suppress the background Compton gamma-quanta.
For the fabrication of GDA it is necessary to use the germanium crystals of uniform distribution of impurities to provide homogenous electric field near the anode. Second in importance it is the providing small depth and width of junction layer under the strips so the electric field near the strips is defined by junction dimensions. The design of GDA must gurantee the reliable cooling of crystal since the critical electric field and amplification factor depend on free path of charge carriers which in its turn depends on temperature.
## 3 The energy resolution and threshold of GDA.
The energy resolution of semiconductor detector will be given by
$$\mathrm{\Delta }E=\sqrt{(\mathrm{\Delta }E_{int})^2+(\mathrm{\Delta }E_{el})^2}$$
(7)
where $`\mathrm{\Delta }E_{int}`$ is the intrinsic energy resolution of detector which is defined by statistical fluctuation in the number of charge carriers created in detector sensitive volume, $`\mathrm{\Delta }E_{el}`$ is the energy resolution which is defined by associated electronics. In the case of GDA these two terms will be of form
$$\mathrm{\Delta }E_{int}=2.34\sqrt{\epsilon E(F+f)K^2}$$
(8)
and
$$\mathrm{\Delta }E_{el}=\frac{4.52\epsilon }{e}\sqrt{\frac{0.6kT}{\tau S}C^2+kT\tau \left[\frac{1}{R_\mathrm{\Sigma }}+\frac{e}{2kT}(I_s+I_bfK^2)\right]}$$
(9)
where $`\epsilon `$ is the energy to create one pair of charge carriers, $`E`$ is the energy deposited in the detector, $`F`$ is the Fano factor, $`f`$ is the excess noise factor due to the fluctuation of multiplication, $`K`$ is the amplification factor, $`e`$ is the electron charge, $`T`$ is the absolute temperature of the resistors, $`C`$ is the total capacitance presented at the input of preamplifier, $`\tau `$ is the time constant of the RC circuits of the preamplifier, $`S`$ is the steepness of the field effect transistor characteristic, $`R_\mathrm{\Sigma }`$ is the resistance at the input of the preamplifier, $`I_s`$ is the detector leakage surface current and $`I_b`$ is the detector leakage balk current due to the thermal generation of charge carriers. According to the calculation for GDA with K$`>`$10 one must take into account in the formula for $`\mathrm{\Delta }E_{el}`$ only the last term due to the bulk leakage current and the formula (7) for GDA can be rewritten as
$$\mathrm{\Delta }E2.36K\sqrt{\epsilon E(F+f)+10^4I_b\tau f}$$
(10)
where $`\epsilon `$ and $`E`$ is in eV , $`I_b`$ is in nA and $`\tau `$ is in $`\mu sec`$. The GDA energy threshold is defined by $`I_b`$ or more exactly by the last term in (10):
$$E_{th}2.36\sqrt{10^4I_b\tau f}$$
(11)
One can estimate the $`E_{th}`$ for microstrip planar detector from HPGe of volume 100 $`cm^3`$ with internal amplification: at N=$`10^{10}`$$`cm^3`$, $`I_b`$=0.01 nA per one strip, $`\tau `$=0.5 $`\mu sec`$ and $`f`$=0.5 : $`E_{th}`$ 12 eV.
The dependance of relative energy resolution $`\mathrm{\Delta }`$E/E on energy in the more interesting energy range $`50รท5000`$ $`eV`$ for GDA is shown in table.
The internal amplification of GDA degrades the performance somewhat, but such energy resolution of GDA is adequate to the investigation of above problems. It is interesting to note that common planar detector from HPGe of volume 100 $`cm^3`$ produced by โCanberraโ (type GL3825R) has relative energy resolution about 8$`\%`$ at E=5900 eV.
## 4 The prospects of GDA use.
Presently germanium detectors are in considerable use in low background measurements for high purity of germanium crystals: their content of radioactive impurity does not exceed $`10^{13}g/g`$ . If the internal amplification is realized in germanium detectors and their threshold is lowered to several $`eV`$ the possibility for their use will considerably increase. Firstly, the lowering of the detector threshold makes it possible to extend the kinematical domain of investigations, secondly, it enables in some cases significantly to decrease the level of background. Following are brief discussions of prospects of GDA using for neutrino magnetic moment ( $`\mu _\nu `$ ) measurement, for investigation of neutrino coherent scattering by nuclei and for search for โlightโ WIMPs.
Neutrino magnetic moment(NMM) measurement.
A laboratory measurement of the NMM is based on its contribution to the (anti)neutrino - electron scattering. For a non-zero NMM the differential over the kinetic energy $`T`$ of the recoil electron cross section $`d\sigma /dT`$ is given by the sum of the weak interaction cross section and the electromagnetic one. At a small recoil energy the weak part practically constant, while the electromagnetic one grows as $`1/T`$ towards low energies. For improving the sensitivity to $`\mu _\nu `$ it is necessary to lower the threshold for detecting the recoil electrons as far as the background conditions allow. Now there is number of the projects dedicated to the NMM measurements with detectors of mass $`2รท1000`$ kg and with thresholds of $`3รท500`$ KeV. Authors of proposal are going to use germanium detector with mass of 2 kg and threshold of 3 KeV to achieve for two years reactor measurement the limit on $`\mu _\nu 3\times 10^{11}\mu _B`$, with $`\mu _B=e/(2m_e)`$ being the Bhorh magneton. The use of GDA with internal amplification in this experiment would provide possibility to achieve limit on $`\mu _\nu `$ about 2$`\times 10^{11}\mu _B`$. However, one can not significantly low the GDA threshold in the reactor experiment for two reasons. Firstly, one must take into account atomic binding effect for electrons. The total $`\stackrel{~}{\nu }_e`$ cross section will be 30$`\%`$ less due to the binding effects ( in the energy range $`200รท3000`$ eV ). The secondly at T$`<`$500 eV the effect neutrino coherent scattering (NCS) by germanium nuclei will surpass the effect of $`\stackrel{~}{\nu }e`$ scatttering and NCS will present physical background relative to the $`\stackrel{~}{\nu }e`$ scattering. Nevertheless the low GDA threshold can be used in full measure if one uses artifical neutrino source (ANS) instead of reactor. Nowadays there are several proposals on development ANS with activity $`5รท40`$ MCi . Although the development and construction of the ANS is rather expensive and challenging problem the use of ANS offers some advantages over the reactor based experiments:
1) ANS can provide significantly higher neutrino flux density (up to $`10^{15}`$ 1/$`cm^2s`$),
2) ANS can be used in deep underground laboratory , where level of background is significanly lower than in shallow box near the reactor,
3) one can use the optimum ratio of effect and background measurement times, that is impossible in the reactor experiment,
4) the uncertainities of the neutrino spectrum and flux density are small.
Presently the more appropriate ANS for NMM measurement is a tritium source proposed in for measurement of NMM by means of low threshold cryogenic silicon detector. The end-point energy of tritium beta spectra is 18,6 KeV and for maximum recoil electron energy a kinematical limitation gives:
$$T_{max}=\frac{2E_\nu ^2}{2E_\nu +m_ec^2}=1260eV.$$
(12)
If one uses tritium ANS with activity 40 MCi then the neutrino density flux will be 6$`\times 10^{14}`$ 1/$`cm^2s`$ in the detector position. The limit on NMM about 3$`\times 10^{12}\mu _B`$ can be achieved during one year measurement time with silicon detector mass of 3 kg in the energy range $`1รท300`$ eV. The very same limit on NMM can be achieved by use the GDA with mass of 3 kg during one year of measurement time in the energy range $`30รท500`$ eV. In the GDA case the count rate due to electromagnetic neutrino interaction will be 0.24 event/d (at $`\mu _\nu =3\times 10^{12}\mu _B`$), the count rate due to weak interaction will be one order of magnitude less. Now let us consider the background count rate. If one takes into account that $`\sigma _{EM}ln(T_{max}/T_{min})`$ then the value of $`\sigma _{EM}`$ in the energy range $`3รท50`$ KeV($`\mathrm{\Delta }E_1`$) and $`30รท500`$ eV($`\mathrm{\Delta }E_2`$) will be the very same. The best level of background for germanium detector was achieved in and it is equal 0.1 event/keV$``$kg$``$d. If one assume that level of background does not depend on energy in the energy range of interest then the level of background in the second energy interval will be two order of magnitude less since $`(\mathrm{\Delta }E_2/\mathrm{\Delta }E_1)10^2`$ and will be 0.15 event/d. In addition in the case of tritium ANS there is no physical background due to the NCS since in this case the maximum recoil energy of germanium nuclei $`T_{max}^n`$ is equal 0.01 eV.
Neutrino coherent scattering by nuclei.
The use of GDA can open up the possibility to investigate the coherent neutrino scattering by nuclei. This interaction is of great importance into interstellar processes and up to now did not observed in laboratory because the very low kinetic energy transferred to nucleus in the process of neutrino-nucleus scattering. For germanium nucleus and reactor antineutrino spectrum the maximum kinetic recoil energy is about $`T_{max}^n`$=2500 eV and taking into account quenching effect only one third of this energy are imparted on ionization . So one needs germanium detector with threshold significantly lower 800 eV for investigation of NCS by germanium nuclei. Such low threshold can be provided by GDA having internal amplification about 10<sup>2</sup>. The differential cross section of NCS by nuclei can be expressed by :
$$\frac{d\sigma _W^c}{dT}\frac{G_F^2}{4\pi }N^2M\left(1\frac{MT^n}{2E_\nu ^2}\right)$$
(13)
$$\frac{d\sigma _{EM}^c}{dT}\left(\frac{\mu _\nu }{\mu _B}\right)^2\frac{\pi \alpha ^2}{m_e^2}Z^2\left(\frac{1}{T^n}\right)$$
(14)
where M,N and Z are the mass, neutron and charge numbers of nuclei correspondingly and T<sup>n</sup> \- nucleus recoil energy. If one uses GDA with mass of 3 kg and threshold 30 eV in the reactor experiment the count rate due NCS in the energy range $`30รท300`$ eV will be near 100 event/d at antineutrino intensity $`2\times 10^{13}`$ $`\nu /cm^2s`$ due to the weak interaction. The count rate due to the electromagnetic interaction (at $`\mu _\nu =2\times 10^{11}\mu _B`$) will be 0.45 event/d. The level of detector background will be only 0.1 event/d if one uses the above estimations of detector background, taking into account that the energy range is equal 0.3 KeV in this case. So one can see that use of GDA for measurement NCS makes it possible :
* to improve the weak interaction constants,
* to investigate the neutrino oscillations by alternative way,
* to make more precise the quenching factor for germanium nuclei at low energy transfer which is of interest for dark matter experiments.
The GDA use in the Dark Matter Experiment.
Presently in Dark Matter Experiments the main efforts are directed at decreasing of the background and lowering of the energy threshold of the detector, since for more appropriate dark matter candidate- WIMPs or weakly interacting massive particles - the expected count rate for WIMPs-nuclear scattering is in the range $`0.001รท1.0`$ event/kg$``$d and expected recoil energy $`T^n`$ lies in the wide energy range beginning of some tens $`eV`$ and above. Significant success was achieved by CDMS (Cryogenic Dark Matter Search) collaboration in radiation background decreasing . The rejection of 99$`\%`$ of photon background was achieved use detectors which simultaneously measure the recoil energy in both- photon and charge mediated signals. However, this method is effective at detector threshold above 15 KeV. So experiments for search for WIMPs with mass lower than 10 Gev/$`c^2`$ (at very low T) are carried out now by CRESST (Cryogenic Rare Event Search with Superconducting Thermometers) collaboration which plans to use cryogenic detector including four 250 g sapphire detectors with threshold 500 eV. The using of several GDA modules of mass 1 kg each with threshold about 30 eV would be very effective in this investigations.
The investigation of Solar Neutrino Problem with use GDA can provide possibility to detect simultaneously the whole neutrino spectrum but one needs of course to use large mass detector in this case. It is beleived that GDA being developed will find also use in the applied fields.
The authors would like to thank V.S.Kaftanov for his participation in this work and Yu.Kamishkov for his interest to this work and for enlightening discussions. |
warning/0002/gr-qc0002003.html | ar5iv | text | # Maxwell equation, Shroedinger equation, Dirac equation, Einstein equation defined on the multifractal sets of the time and the space
## 1 abstract
What forms will have an equations of modern physics if the dimensions of our time and space are fractional? The generalized equations enumerated by title are presented by help the generalized fractional derivatives of Riemann-Liouville.
## 2 Introduction
In the articles - the generalized fractional Riemann-Liouville derivatives (GFD) are determined and the fractal theory of time and space (and some others physical questions) basing on the using GFD for functions defined on a multifractal sets are presented. The multifractal time and space sets are characterized by fractal dimensions $`d_t(๐ซ(t),t)`$ and $`d_๐ซ(t(๐ซ),๐ซ)`$. In this paper the generalization of main equations of the modern physics are presented for the multifractal time and space in the frame of multifractal model of time and space presented in -. These equations gives in a little corrections for the known equations for the case when the fractal dimensions (FD) of time $`d_t`$ and space $`d_๐ซ`$ are $`d_t=1+\epsilon (๐ซ(t),t)`$ (and so on $`d_๐ซ`$) and the FD are slightly differs from unity, i.e. $`|\epsilon |<<1`$, that is valid for small densities of Lagrangians in points $`t`$, $`๐ซ`$, i.e. for weak forces in the domain of space and time near $`๐ซ,t`$. All the equations may be received by means of the principle of minimum of fractal dimensions functional (see ) and from this principle the generalized Eulerโs equations may be write down. We use more simple method in this article, consisting in the replacing the ordinary derivatives in the known equations by GFD (it may be ground by comparison with the generalized Euler equations). Before receiving the equations we remind the main definitions and designations of the theory -:
Generalized fractional derivatives (GFD):
We begin from remembering of the fractional Riemann-Liouville derivatives definitions :
$$D_{+,t}^df(t)=\left(\frac{d}{dt}\right)^n\underset{a}{\overset{t}{}}๐t^{}\frac{f(t^{})}{\mathrm{\Gamma }(nd)(tt^{})^{dn+1}}$$
(1)
$$D_{,t}^df(t)=(1)^n\left(\frac{d}{dt}\right)^n\underset{t}{\overset{b}{}}๐t^{}\frac{f(t^{})}{\mathrm{\Gamma }(nd)(t^{}t)^{dn+1}}$$
Let a function $`f(t)`$ of variable $`t`$ is defined on multifractal set $`S_t`$ which consist from multifractal subsets $`s_i(t_i)`$. We shall see subsets $`s_i(t_i)`$ as the โpointsโ $`t_i`$ (with a continuous distribution for different multifractal subsets $`s_i(t_i)`$ of multifractal set $`S_t`$ ordered by values of $`t`$. Let the function $`d(t_i)=d(t)`$ is continuous and describes their fractional dimensions (in some cases coinciding with local fractal dimensions of set $`S_t`$ as function $`t`$. For the elementary generalization the definitions (1)-(2) are used physical reasons and variable $`t`$ is interpreted as a time. For a continuous functions $`f(t)`$ (the generalized functions defined on the class of the finitary functions (see), the fractional derivatives of the Riemann - Liouville are continuous also. So for infinitesimal intervals of time the functionals (1)-(2) will vary on an infinitesimal quantity. For the continuous function $`d(t)`$ the changes it thus also will be infinitesimal. It allows, as the elementary generalization (1) that is suitable for describe the changes the function $`f(t)`$ defined on multifractal subsets $`s(t)`$ (as well as in the (1)-(2)), to take into account the summary influence of a kernel of integral $`(tt^{})^{d(t)n+1}\mathrm{\Gamma }^1(nd(t))`$, depending from $`d(t)`$, on the $`f(t)`$ in all points of integration and, instead of (1)-(2) to write the integral which takes into account all this influences. Thus, we enter the following definitions (generalized fractional derivatives and integrals (GFD)), that account also dependence $`d(t)`$ from time and vector parameter $`๐ซ(t)`$ (i.e. $`d_td_t(๐ซ,t)`$):
$$D_{+,t}^{d_t}f(t)=\left(\frac{d}{d_t}\right)^n\underset{a}{\overset{t}{}}๐t^{}\frac{f(t^{})}{\mathrm{\Gamma }(nd_t(t^{}))(tt^{})^{d_t(t^{})n+1}}$$
(2)
$`D_{,t}^{d_t}f(t)=(1)^n\times `$ (3)
$`\times `$ $`\left({\displaystyle \frac{d}{dt}}\right)^n{\displaystyle \underset{t}{\overset{b}{}}}๐t^{}{\displaystyle \frac{f(t^{})}{\mathrm{\Gamma }(nd_t(t^{}))(t^{}t)^{d_t(t^{})n+1}}}`$
In (2)-(3), as well as in (1)-(2), $`a`$ and $`b`$ stationary values defined on an infinite axis (from $`\mathrm{}`$ to $`\mathrm{}`$), $`a<b`$ , $`n1d_t<n`$, $`n=\{d_t\}+1`$, $`\{d_t\}`$\- the integer part of $`d_t0`$, $`n=0`$ for $`d_t<0`$. The only difference the (2)-(3) from the (1)-(2) is: $`d_t=d_t(๐ซ(t),t)`$\- fractional dimensions (further will be used for it terms โ fractal dimensions โ (FD) or โ the global fractal dimension (FD)โ of subset $`s_t`$) is the function of time and coordinates, instead of stationary values in the (1)-(2). Similar to (1)-(2), it is possible to define the GFD, (that coincides for integer values of fractional dimensions $`d_๐ซ(๐ซ,t)`$ with derivatives respect to vector variable $`๐ซ`$) $`D_{+,๐ซ}^{d_๐ซ}f(๐ซ,t)`$ respect to vector $`๐ซ(t)`$ variables (spatial coordinates). We pay attention, that definitions (1)-(2) are a special case of Hadamard derivatives .
2. The connection between the fractional dimensions (FD)of time and space with Lagrangian functions of energy densities read:
$$d_t=1+\underset{i,\alpha }{}\beta _{i,\alpha }L_{i,\alpha }(t,๐ซ,\mathrm{\Phi }_i,\psi _i)$$
(4)
In (4) $`\alpha `$ takes value: $`\alpha =t,๐ซ`$. More complicated dependencies of $`d_\alpha `$ at $`L_{\alpha ,i}`$ are considered in . Note that relation (4) (and similar expression for $`d_๐ซ`$ does not contain any limitations on the value of $`\beta _iL_{i,\alpha }(t,๐ซ,\mathrm{\Phi }_i,\psi _i)`$ unless such limitations are imposed on the corresponding Lagrangians, and therefore $`d_t`$ can reach any whatever high or small value.
3. Letโs write now the equations enumerated in title in fractal time and space by using GFD:
a) Maxwell equations:
$`{\displaystyle \underset{i=1}{\overset{3}{}}}D_{,i,r}^{d_r}D_{+,i,r}^{d_r}A_j(x){\displaystyle \frac{1}{c^2}}D_{,t}^{d_t}D_{+,t}^{d_t}A_j(x)+m^2A_j(x)={\displaystyle \frac{4\pi }{c}}j_j(x),`$ (5)
$`j_j=eD_{+,j,t}^{d_j}r_i`$
$$D_{+,j,r}^{d_j}A_j(x)=0$$
(6)
In (5)-(6) FD $`d_j`$ is equal to $`d_๐ซ`$ for $`j=1,2,3`$ and $`d_t`$ for $`j=0`$ and introduced the mass of foton for providing existence of GFD on infinity (then it must be select equal zero).
b) Shreodinger equation
$$i\mathrm{}D_{+,t}^{d_t}\psi (๐ซ,t)=\frac{\mathrm{}^2}{2m}D_{,r}^{d_๐ซ}D_{+,r}^{d_r}\psi (๐ซ,t)e^2(๐ซ,t)\psi (๐ซ,t)$$
(7)
where in $`D_{,r}^{d_r},D_{+,r}^{d_r}`$ operators $``$ replaced by operators $`ie/\mathrm{}c๐`$)
c)Dirac equation
$$[i\gamma _i(D_{+,i}^{d_i}ieA_i(x))m]\psi (x)=0$$
(8)
where $`\gamma _i`$ are Dirac matrices. It is necessary to make the difference for GFD $`D_i^{d_i}`$ with respect to $`t`$ or with respect to $`๐ซ`$ :
$$D_i^{d_i}=D_t^{d_t},D_๐ซ^{d_r}$$
(9)
For atomic electrons the main role plays the electric fields of nucleus. So the density of Lagrangians energy that defined the FG $`d_t`$ may be selected as
$$d_t=1+\beta \mathrm{\Phi }(๐ซ,t)1+\frac{e^2}{rM_0c^2}$$
(10)
where $`M_0`$ is the mass of electrical charge body that originate electrical field. It is easy to demonstrate that on the distances of the first Bohrโs radius in atoms the fractional corrections to time dimensions (difference the $`d_t`$ from unity) have values $`10^8`$, so the fractal corrections to electron energy $`E`$ in atoms will be have values $`10^8E`$. It lay out (or in limits domain) of experimental possibilities of the modern experiment.
d) Einstein equation
It is possible to receive the generalization of general relativity equation by using two ways. In the first way it is necessary to introduce a parallel displacement in the Riemann space with fractional dimensions that may be done without difficulties for weak fields (may be it is possible to determine the parallel displacement and for strong fields by the same relations). In that case the carrier of a measure is the Riemann space and we obtain the determination for covariant derivatives in Riemann space with fractional dimensions
$$D_{\pm ,\alpha }^{d_i}t^{\mu \nu }=D_{\pm ,\alpha }^{d_i}t^{\mu \nu }+\gamma _{\alpha \beta }^\nu t^{\mu \beta }i=t,r$$
(11)
where $`t^{\mu \nu }`$ tensor and $`\gamma ^{\mu \nu }`$ metric tensor the Riemann โfour-dimension space with fractional dimensionsโ, $`D_{\pm ,\alpha }^{d_i}`$ is GFD, $`\gamma _{\alpha \beta }^\nu `$ are Christoffelโs symbols
$$\gamma _{\alpha \beta }^\nu =\frac{1}{2}\gamma ^{\nu \sigma }(D_{\pm ,\alpha }^{d_i}\gamma _{\beta \sigma }+D_{\pm ,\beta }^{d_i}\gamma _{\alpha \sigma }+D_{\pm ,\sigma }^{d_i}\gamma _{\alpha \beta }$$
(12)
The equations for gravitation field tensor $`\stackrel{~}{\mathrm{\Phi }}^{\mu \nu }`$=$`\sqrt{\gamma }\mathrm{\Phi }^{\mu \nu }`$, $`\gamma =det(\gamma _{\mu \nu })`$, $`\stackrel{~}{t}^{\mu \nu }`$= $`\sqrt{\gamma }t^{\mu \nu }`$, $`L`$ \- is a scalar density of matter (see in details (5)-(6)) than read
$$\gamma ^{\alpha \beta }D_{,\alpha }^{\nu ,d_i}D_{+,\beta }^{\nu ,d_i}\mathrm{\Phi }^{\mu \nu }+b^2\stackrel{~}{\mathrm{\Phi }}^{\mu \nu }=\lambda \frac{\delta L}{\delta \gamma ^{\mu \nu }}=\lambda \stackrel{~}{t}^{\mu \nu }(\gamma ^{\mu \nu },\mathrm{\Phi }_A)$$
(13)
where $`b`$ is a constant value that necessary to introduce for using more broad sets of functions with GFD and it after calculations may be put zero. The equation for curvature tensor (with GFD ) have an usual form
$$R^{\mu \nu }\frac{1}{2}\gamma ^{\mu \nu }R=\frac{8\pi }{\sqrt{g}}T^{\mu \nu }$$
(14)
$$D_{\pm ,\mu }^{d_i}\stackrel{~}{g}^{\mu \nu }=0$$
(15)
The equation (15) describes the boundary conditions for $`g^{\mu \nu }`$ on the Universe surface. Stress, that equations (11)-(13) describe gravitation fields in the Riemann space with fractional dimensions, i.e. the carrier of measure is the Riemann space. For the case of weak fields the generalized covariant derivatives may be represented as (see )
$$D_{\pm ,\alpha }^{d_i}t^{\mu \nu }^{}D_{\pm ,\alpha }^{d_i}t^{\mu \nu }+^{\prime \prime }D_{\pm ,\alpha }^{d_i}t^{\mu \nu }$$
(16)
The $`{}_{}{}^{}D_{\pm ,\alpha }^{d_i}`$ in (16) describes the contribution from FD of time and space, the member $`{}_{}{}^{\prime \prime }D_{\pm ,\alpha }^{d_i}`$ describes the contribution from Riemann space with integer dimensions.
The second way for describing the gravitation fields in the fractal time and space ( by GFD using ) consists in the other the measure carrier selection. It is more simplest way to select the measure carrier as the flat four-dimensions pseudo Euclidean Minkowski space. In that case may be used (as a base) the system of reference which coincide for FD equal to unit with Cartesian system of reference (we remember that in the fractal theory of time and space there are an absolute systems of reference). The equations of the gravitation in that case will be analogies the equation of the theory in which all derivatives replaced on GFD and metric tensor $`\gamma ^{\mu \nu }`$ are consist the functions (functionals) originated by fractional dimensions (i.e. it must be the function of $`L`$ \- Lagrangian energy densities of gravitation fields). These equations have the form
$$\gamma ^{\alpha \beta }D_{,\alpha }^{\nu ,d_i}D_{+,\beta }^{\nu ,d_i}\mathrm{\Phi }^{\mu \nu }+b^2\stackrel{~}{\mathrm{\Phi }}^{\mu \nu }=\lambda \frac{\delta L}{\delta \gamma ^{\mu \nu }}=\lambda \stackrel{~}{t}^{\mu \nu }(\gamma ^{\mu \nu },\mathrm{\Phi }_A)$$
(17)
The equation (17) differs from equation ( 13) by three aspects:
a) the metric tensor $`\gamma ^{^{\mu \nu }}`$ now determined on the Minkowski space with fractional dimensions;
b) it differs by dependencies of metrics tensor $`\gamma ^{\mu \nu }`$ from $`L`$ ( because there are no dependencies in $`\gamma ^{\mu \nu }`$ of $`L`$ originating through the Riemann metric tensor ), there are only dependencies originating through FD;
c) the reason of appearance of the dependencies the $`\gamma ^{\mu \nu }`$ at $`L`$ lay in the originate it by the fractal dimensions of time and space. If FG are integer the (16) coincide with equation of the theory . For weak fields GFD may be represented only by FD covariant derivatives (only one member in right part of (16)) and in that case may be represented by metric tensor $`g^{\mu \nu }`$ of an โeffectiveโ Riemann space with integer dimensions (see ). We pay attention that the corresponding results of the theory for connections between metric tensor $`\gamma ^{\mu \nu }`$ of Minkowski space with โeffectiveโ metric tensor $`g^{\mu \nu }`$ of Riemann space and gravitation tensor are the special case of our theory. In general case the metric tensor of Minkowski space are complicated function of gravitation field tensor.
We leave for readers an interesting task to generalize the equations of quantum gravitation for multifractal time and space.
## 3 Relations between GFD and ordinary derivatives for $`d_\alpha `$ near integer values
If $`d_\alpha `$$`n`$ where $`n`$ is an integer number, for example $`d_\alpha `$=$`1+\epsilon (๐ซ(t),t)`$, $`\alpha =๐ซ,t`$, in that case it is possible represent GFD by approximate relations (see )
$$D_{+,x_\alpha }^{1+ฯต}f(๐ซ(t),t)=\frac{}{x_\alpha }f(๐ซ(t),t)+\frac{}{x_\alpha }[\epsilon (๐ซ(t),t)f(๐ซ(t)),t)]$$
(18)
The replacement in generalized Maxwell equations (5)-(6) , Shroedinger equation (7), Dirac equation (8), Einstein equation (12) (and so on ) the GFD defined by (2) by approximation of GFD defined in the (18) gives a possibility to solve numerous tasks in fractal space and time and calculate the corrections from fractional dimensions for these tasks.
## 4 Conclusion
1. In this paper were presented the main equation of modern physics defined on multifractal sets of time and space. In case integer dimensions of time and space all of them coincide with the known equation. The value of correction to integer dimensions of time and space in conditions of Earth are very small.So, the correction to dimension of time that gives the gravitational field of Earth on the ground of Earth is equal $`10^9`$. The corrections from electric field nucleus on atomic distances from nucleus are $`10^8`$. So, it may be neglected by these corrections, but only in the cases of weak fields. In the case of strong fields all generalized equations becomes in the integral fractional equations. The last donโt have singularities and only this fact originate the interest to these equations. We pay attention that in the fractal theory of time and space our Universe is the open system (the statistical theory of open system are in ,
2.Can these equation be used for strong fields (i.e. for FD that differs a lot of from integer values). This question now is open. We may only say that answer on this question concerns with answer on other question: is it possible to use the method of ordinary Lagrangians of quantum and classical theories for describing the strong fields? If answer on last question is positive, there are a hope that the positive answer for the first question will be correct. |
warning/0002/cond-mat0002346.html | ar5iv | text | # Conductivity Due to Classical Phase Fluctuations in a Model For High-Tc Superconductors
## I Acknowledgments
This work was supported by the National Science Foundation, Grant DMR97-31511. We are grateful for valuable conversations with J. Orenstein, D. J. Bergman and T. R. Lemberger. |
warning/0002/cond-mat0002174.html | ar5iv | text | # Foundations of Dissipative Particle Dynamics
## I Introduction
The non-equilibrium behavior of fluids continues to present a major challenge for both theory and numerical simulation. In recent times, there has been growing interest in the study of so-called โmesoscaleโ modeling and simulation methods, particularly for the description of the complex dynamical behavior of many kinds of soft condensed matter, whose properties have thwarted more conventional approaches. As an example, consider the case of complex fluids with many coexisting length and time scales, for which hydrodynamic descriptions are unknown and may not even exist. These kinds of fluids include multi-phase flows, particulate and colloidal suspensions, polymers, and amphiphilic fluids, including emulsions and microemulsions. Fluctuations and Brownian motion are often key features controlling their behavior.
From the standpoint of traditional fluid dynamics, a general problem in describing such fluids is the lack of adequate continuum models. Such descriptions, which are usually based on simple conservation laws, approach the physical description from the macroscopic side, that is in a โtop downโ manner, and have certainly proved successful for simple Newtonian fluids . For complex fluids, however, equivalent phenomenological representations are usually unavailable and instead it is necessary to base the modeling approach on a microscopic (that is on a particulate) description of the system, thus working from the bottom upwards, along the general lines of the program for statistical mechanics pioneered by Boltzmann . Molecular dynamics (MD) presents itself as the most accurate and fundamental method but it is far too computationally intensive to provide a practical option for most hydrodynamic problems involving complex fluids. Over the last decade several alternative โbottom upโ strategies have therefore been introduced. Hydrodynamic lattice gases , which model the fluid as a discrete set of particles, represent a computationally efficient spatial and temporal discretization of the more conventional molecular dynamics. The lattice-Boltzmann method , originally derived from the lattice-gas paradigm by invoking Boltzmannโs Stosszahlansatz, represents an intermediate (fluctuationless) approach between the top-down (continuum) and bottom-up (particulate) strategies, insofar as the basic entity in such models is a single particle distribution function; but for interacting systems even these lattice-Boltzmann methods can be subdivided into bottom-up and top-down models .
A recent contribution to the family of bottom-up approaches is the dissipative particle dynamics (DPD) method introduced by Hoogerbrugge and Koelman in 1992 . Although in the original formulation of DPD time was discrete and space continuous, a more recent re-interpretation asserts that this model is in fact a finite-difference approximation to the โtrueโ DPD, which is defined by a set of continuous time Langevin equations with momentum conservation between the dissipative particles . Successful applications of the technique have been made to colloidal suspensions , polymer solutions and binary immiscible fluids . For specific applications where comparison is possible, this algorithm is orders of magnitude faster than MD . The basic elements of the DPD scheme are particles that represent rather ill-defined โmesoscopicโ quantities of the underlying molecular fluid. These dissipative particles are stipulated to evolve in the same way that MD particles do, but with different inter-particle forces: since the DPD particles are pictured to have internal degrees of freedom, the forces between them have both a fluctuating and a dissipative component in addition to the conservative forces that are present at the MD level. Newtonโs third law is still satisfied, however, and consequently momentum conservation together with mass conservation produce hydrodynamic behavior at the macroscopic level.
Dissipative particle dynamics has been shown to produce the correct macroscopic (continuum) theory; that is, for a one-component DPD fluid, the Navier-Stokes equations emerge in the large scale limit, and the fluid viscosity can be computed . However, even though dissipative particles have generally been viewed as clusters of molecules, no attempt has been made to link DPD to the underlying microscopic dynamics, and DPD thus remains a foundationless algorithm, as is that of the hydrodynamic lattice gas and a fortiori the lattice-Boltzmann method. It is the principal purpose of the present paper to provide an atomistic foundation for dissipative particle dynamics. Among the numerous benefits gained by achieving this, we are then able to provide a precise definition of the term โmesoscaleโ, to relate the hitherto purely phenomenological parameters in the algorithm to underlying molecular interactions, and thereby to formulate DPD simulations for specific physicochemical systems, defined in terms of their molecular constituents. The DPD that we derive is a representation of the underlying MD. Consequently, to the extent that the approximations made are valid, the DPD and MD will have the same hydrodynamic descriptions, and no separate kinetic theory for, say, the DPD viscosity will be needed once it is known for the MD system. Since the MD degrees of freedom will be integrated out in our approach the MD viscosity will appear in the DPD model as a parameter that may be tuned freely.
In our approach, the โdissipative particlesโ (DP) are defined in terms of appropriate weight functions that sample portions of the underlying conservative MD particles, and the forces between the dissipative particles are obtained from the hydrodynamic description of the MD system: the microscopic conservation laws carry over directly to the DPD, and the hydrodynamic behavior of MD is thus reproduced by the DPD, albeit at a coarser scale. The mesoscopic (coarse-grained) scale of the DPD can be precisely specified in terms of the MD interactions. The size of the dissipative particles, as specified by the number of MD particles within them, furnishes the meaning of the term โmesoscopicโ in the present context. Since this size is a freely tunable parameter of the model, the resulting DPD introduces a general procedure for simulating microscopic systems at any convenient scale of coarse graining, provided that the forces between the dissipative particles are known. When a hydrodynamic description of the underlying particles can be found, these forces follow directly; in cases where this is not possible, the forces between dissipative particles must be supplemented with the additional components of the physical description that enter on the mesoscopic level.
The DPD model which we derive from molecular dynamics is formally similar to conventional, albeit foundationless, DPD . The interactions are pairwise and conserve mass and momentum, as well as energy . Just as the forces conventionally used to define DPD have conservative, dissipative and fluctuating components, so too do the forces in the present case. In the present model, the role of the conservative force is played by the pressure forces. However, while conventional dissipative particles possess spherical symmetry and experience interactions mediated by purely central forces, our dissipative particles are defined as space-filling cells on a Voronoi lattice whose forces have both central and tangential components. These features are shared with a model studied by Espaรฑol . This model links DPD to smoothed particle hydrodynamics and defines the DPD forces by hydrodynamic considerations in a way analogous to earlier DPD models. Espaรฑol et al. have also carried out MD simulations with a superposed Voronoi mesh in order to measure the coarse grained inter-DP forces.
While conventional DPD defines dissipative particle masses to be constant, this feature is not preserved in our new model. In our first publication on this theory , we stated that, while the dissipative particle masses fluctuate due to the motion of MD particles across their boundaries, the average masses should be constant. In fact, the DP-masses vary due to distortions of the Voronoi cells, and this feature is now properly incorporated in the model.
We follow two distinct routes to obtain the fluctuation-dissipation relations that give the magnitude of the thermal forces. The first route follows the conventional path which makes use of a Fokker-Planck equation . We show that the DPD system is described in an approximate sense by the isothermal-isobaric ensemble. The second route is based on the theory of fluctuating hydrodynamics and it is argued that this approach corresponds to the statistical mechanics of the grand canonical ensemble. Both routes lead to the same result for the fluctuating forces and simulations confirm that, with the use of these forces, the measured DP temperature is equal to the MD temperature which is provided as input. This is an important finding in the present context as the most significant approximations we have made underlie the derivation of the thermal forces.
## II Coarse-graining molecular dynamics: from micro to mesoscale
The essential idea motivating our definition of mesoscopic dissipative particles is to specify them as clusters of MD particles in such a way that the MD particles themselves remain unaffected while all being represented by the dissipative particles. The independence of the molecular dynamics from the superimposed coarse-grained dissipative particle dynamics implies that the MD particles are able to move between the dissipative particles. The stipulation that all MD particles must be fully represented by the DPโs implies that while the mass, momentum and energy of a single MD particle may be shared between DPโs, the sum of the shared components must always equal the mass and momentum of the MD particle.
### A Definitions
Full representation of all the MD particles can be achieved in a general way by introducing a sampling function
$$f_k(๐ฑ)=\frac{s(๐ฑ๐ซ_k)}{_ls(๐ฑ๐ซ_l)}.$$
(1)
where the positions $`๐ซ_k`$ and $`๐ซ_l`$ define the DP centers, $`๐ฑ`$ is an arbitrary position and $`s(๐ฑ)`$ is some localized function. It will prove convenient to choose it as a Gaussian
$$s(๐ฑ)=\mathrm{exp}(x^2/a^2)$$
(2)
where the distance $`a`$ sets the scale of the sampling function, although this choice is not necessary. The mass, momentum and internal energy $`E`$ of the $`k`$th DP are then defined as
$`M_k`$ $`=`$ $`{\displaystyle \underset{i}{}}f_k(๐ฑ_i)m,`$ (3)
$`๐_k`$ $`=`$ $`{\displaystyle \underset{i}{}}f_k(๐ฑ_i)m๐ฏ_i,`$ (4)
$`{\displaystyle \frac{1}{2}}M_kU_k^2+E_k`$ $`=`$ $`{\displaystyle \underset{i}{}}f_k(๐ฑ_i)\left({\displaystyle \frac{1}{2}}mv_i^2+{\displaystyle \frac{1}{2}}{\displaystyle \underset{ji}{}}V_{MD}(r_{ij})\right)`$ (5)
$``$ $`{\displaystyle \underset{i}{}}f_k(๐ฑ_i)ฯต_i,`$ (6)
where $`๐ฑ_i`$ and $`๐ฏ_i`$ are the position and velocity of the $`i`$th MD particle, which are all assumed to have identical masses $`m`$, $`๐_k`$ is the momentum of the $`k`$th DP and $`V_{MD}(r_{ij})`$ is the potential energy of the MD particle pair $`ij`$, separated a distance $`r_{ij}`$. The particle energy $`ฯต_i`$ thus contains both the kinetic and a potential term. The kinematic condition
$$\dot{๐ซ}_k=๐_k๐_k/M_k$$
(7)
completes the definition of our dissipative particle dynamics.
It is generally true that mass and momentum conservation suffice to produce hydrodynamic behavior. However, the equations expressing these conservation laws contain the fluid pressure. In order to get the fluid pressure a thermodynamic description of the system is needed. This produces an equation of state, which closes the system of hydrodynamic equations. Any thermodynamic potential may be used to obtain the equation of state. In the present case we shall take this potential to be the internal energy $`E_k`$ of the dissipative particles, and we shall obtain the equations of motion for the DP mass, momentum and energy. Note that the internal energy would also have to be computed if a free energy had been chosen for the thermodynamic description. For this reason it is not possible to complete the hydrodynamic description without taking the energy flow into account. As a byproduct of this the present DPD also contains a description of the heat flow and corresponds to the recently introduced DPD with energy conservation . Espaรฑol previously introduced an angular momentum variable describing the dynamics of extended particles : this is needed when forces are non-central in order to avoid dissipation of energy in a rigid rotation of the fluid. Angular momentum could be included on the same footing as momentum in the following developments. However for reasons both of space and conceptual economy we shall omit it in the present context, even though it is probably important in applications where hydrodynamic precision is important. In the following sections, we shall use the notation $`๐ซ`$, $`M`$, $`๐`$ and $`E`$ with the indices $`k,l,m`$ and $`n`$ to denote DPโs while we shall use $`๐ฑ`$, $`m`$, $`๐ฏ`$ and $`ฯต`$ with the indices $`i`$ and $`j`$ to denote MD particles.
### B Equations of motion for the dissipative particles based on a microscopic description
The fact that all the MD particles are represented at all instants in the coarse-grained scheme is guaranteed by the normalization condition $`_kf_k(๐ฑ)=1`$. This implies directly that
$`{\displaystyle \underset{k}{}}M_k`$ $`=`$ $`{\displaystyle \underset{i}{}}m`$ (8)
$`{\displaystyle \underset{k}{}}๐_k`$ $`=`$ $`{\displaystyle \underset{i}{}}m๐ฏ_i`$ (9)
$`{\displaystyle \underset{k}{}}E_k^{\text{tot}}`$ $`=`$ $`{\displaystyle \underset{k}{}}\left({\displaystyle \frac{1}{2}}M_k๐_k^2+E_k\right)={\displaystyle \underset{i}{}}ฯต_i;`$ (10)
thus with mass, momentum and energy conserved at the MD level, these quantities are also conserved at the DP level. In order to derive the equations of motion for dissipative particle dynamics we now take the time derivatives of Eqs. (6). This gives
$`{\displaystyle \frac{\text{d}M_k}{\text{d}t}}`$ $`=`$ $`{\displaystyle \underset{i}{}}\dot{f}_k(๐ฑ_i)m`$ (11)
$`{\displaystyle \frac{\text{d}๐_k}{\text{d}t}}`$ $`=`$ $`{\displaystyle \underset{i}{}}\left(\dot{f}_k(๐ฑ_i)m๐ฏ_i+f_k(๐ฑ_i)๐
_i\right)`$ (12)
$`{\displaystyle \frac{\text{d}E_k^{\text{tot}}}{\text{d}t}}`$ $`=`$ $`{\displaystyle \underset{i}{}}\left(\dot{f}_k(๐ฑ_i)ฯต_i+f_k(๐ฑ_i)\dot{ฯต}_i\right)`$ (13)
where $`\text{d}/\text{d}t`$ is the substantial derivative and $`๐
_i=m\dot{๐ฏ}_i`$ is the force on particle $`i`$.
The Gaussian form of $`s`$ implies that
$`\dot{s}(๐ฑ)=(2/a^2)\dot{๐ฑ}๐ฑs(๐ฑ)`$. This makes it possible to write
$$\dot{f}_k(๐ฑ_i)=f_{kl}(๐ฑ_i)(๐ฏ_i^{}๐ซ_{kl}+๐ฑ_i^{}๐_{kl})$$
(14)
where the overlap function $`f_{kl}`$ is defined as $`f_{kl}(๐ฑ)(2/a^2)f_k(๐ฑ)f_l(๐ฑ)`$, $`๐ซ_{kl}(๐ซ_k๐ซ_l)`$ and $`๐_{kl}(๐_k๐_l)`$, and we have rearranged terms so as to get them in terms of the centered variables
$`๐ฏ_i^{}`$ $`=`$ $`๐ฏ_i{\displaystyle \frac{(๐_k+๐_l)}{2}}`$ (15)
$`๐ฑ_i^{}`$ $`=`$ $`๐ฑ_i{\displaystyle \frac{(๐ซ_k+๐ซ_l)}{2}}.`$ (16)
Before we proceed with the derivation of the equations of motion it is instructive to work out the actual forms of $`f_k(๐ฑ)`$ and $`f_{kl}(๐ฑ)`$ in the case of only two particles $`k`$ and $`l`$. Using the Gaussian choice of $`s`$ we immediately get
$$f_k(๐ฑ)=\frac{1}{1+\left[\mathrm{exp}((๐ฑ(๐ซ_k+๐ซ_l)/2)(๐ซ_{kl})/(a^2))\right]^2}.$$
(17)
The overlap function similarly follows:
$$f_{kl}(๐ฑ)=\frac{1}{2a^2}\mathrm{cosh}^2\left(\left(๐ฑ\frac{๐ซ_k+๐ซ_l}{2}\right)\left(\frac{๐ซ_{kl}}{a^2}\right)\right).$$
(18)
These two functions are shown in Fig.1. Note that the scale of the overlap region is not $`a`$ but $`a^2/|๐ซ_k๐ซ_l|`$. Dissipative particle interactions only take place where the overlap function is non-zero. This happens along the dividing line which is equally far from the two particles. The contours of non-zero $`f_{kl}`$ thus define a Voronoi lattice with lattice segments of length $`l_{kl}`$. This Voronoi construction is shown in Fig. 2 in which MD particles in the overlap region defined by $`f_{kl}>0.1`$, are shown, though presently not actually simulated as dynamic entities. The volume of the Voronoi cells will in general vary under the dynamics. However, even with arbitrary dissipative particle motion the cell volumes will approach zero only exceptionally, and even then the identities of the DP particles will be preserved so that they subsequently re-emerge.
#### 1 Mass equation
The mass equation (11) takes the form
$$\frac{\text{d}M_k}{\text{d}t}\underset{l}{}\dot{M}_{kl}$$
(19)
where
$$\dot{M}_{kl}=\underset{i}{}f_{kl}(๐ฑ_i)m(๐ฏ_i^{}๐ซ_{kl}+๐ฑ_i^{}๐_{kl}).$$
(20)
The $`๐ฏ_i^{}`$ term will be shown to be negligible within our approximations. The $`๐ฑ_i^{}๐_{kl}`$-term however describes the geometric effect that the Voronoi cells do not conserve their volume: The relative motion of the DP centers causes the cell boundaries to change their orientation. We will return to give this โboundary twistingโ term a quantitative content when the equations of motion are averagedโan effect which was overlooked in our first publication of this theory where it was stated that $`\dot{M}_{kl}=0`$.
#### 2 Momentum equation
The momentum equation (12) takes the form
$`{\displaystyle \frac{\text{d}๐_k}{\text{d}t}}`$ $`=`$ $`{\displaystyle \underset{li}{}}f_{kl}(๐ฑ_i)m๐ฏ_i(๐ฏ_i^{}๐ซ_{kl}+๐ฑ_i^{}๐_{kl})`$ (21)
$`+`$ $`{\displaystyle \underset{li}{}}f_k(๐ฑ_i)๐
_i`$ (22)
We can write the force as $`๐
_i=m๐ +_j๐
_{ij}`$ where the first term is an external force and the second term is the internal force caused by all the other particles. Newtonโs third law then takes the form $`๐
_{ij}=๐
_{ji}`$. The last term in Eq. (22) may then be rewritten as
$$\underset{i}{}f_k(๐ฑ_i)๐
_i=M_k๐ +\underset{ij}{}f_k(๐ฑ_i)๐
_{ij}$$
(23)
where
$`{\displaystyle \underset{ij}{}}f_k(๐ฑ_i)๐
_{ij}`$ $`=`$ $`{\displaystyle \underset{ij}{}}f_k(๐ฑ_i)๐
_{ji}`$ (24)
$`=`$ $`{\displaystyle \underset{ij}{}}f_k(๐ฑ_j+\mathrm{\Delta }๐ฑ_{ij})๐
_{ji}`$ (25)
$``$ $`{\displaystyle \underset{ij}{}}f_k(๐ฑ_j)๐
_{ji}{\displaystyle \underset{ij}{}}\left(\mathrm{\Delta }๐ฑ_{ij}f_k(๐ฑ_i)\right)๐
_{ji}`$ (26)
$`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{ij}{}}\left(\mathrm{\Delta }๐ฑ_{ij}f_k(๐ฑ_i)\right)๐
_{ji}`$ (27)
$`=`$ $`{\displaystyle \underset{l}{}}\left\{{\displaystyle \underset{ij}{}}{\displaystyle \frac{1}{2}}f_{kl}(๐ฑ_i)๐
_{ij}\mathrm{\Delta }๐ฑ_{ij}\right\}๐ซ_{kl}`$ (28)
where $`\mathrm{\Delta }๐ฑ_{ij}=๐ฑ_i๐ฑ_j`$, we have Taylor expanded $`f_k(๐ฑ)`$ around $`๐ฑ_j`$ and used a result similar to Eq. (14) to evaluate $`f_k(๐ฑ)`$. In passing from the third to the fourth line in the above equations we have moved the first term on the right hand side to the left hand side and divided by two. Now, if we group the last term above with the $`๐ซ_{kl}`$ term in Eq. (22), make use of Eq. (16), and do some rearranging of terms we get
$`{\displaystyle \frac{\text{d}๐_k}{\text{d}t}}`$ $`=`$ $`M_k๐ +{\displaystyle \underset{l}{}}\dot{M}_{kl}{\displaystyle \frac{๐_k+๐_l}{2}}`$ (29)
$`+`$ $`{\displaystyle \underset{li}{}}f_{kl}(๐ฑ_i)๐ท_i^{}๐ซ_{kl}`$ (30)
$`+`$ $`{\displaystyle \underset{li}{}}f_{kl}(๐ฑ_i)m๐ฏ_i^{}๐ฑ_i^{}๐_{kl}`$ (31)
where we have used the relation $`\dot{M}_k=_l\dot{M}_{kl}`$ and defined the general momentum-flux tensor
$$๐ท_i=m๐ฏ_i๐ฏ_i+\frac{1}{2}\underset{j}{}๐
_{ij}\mathrm{\Delta }๐ฑ_{ij}.$$
(32)
This tensor is the momentum analogue of the mass-flux vector $`m๐ฏ_i`$. The prime indicates that the velocities on the right hand side are those defined in Eq. (16). The tensor $`๐ท_i`$ describes both the momentum that the particle carries around through its own motion and the momentum exchanged by inter-particle forces. It may be arrived at by considering the momentum transport across imaginary cross sections of the volume in which the particle is located.
#### 3 Energy equation
In order to get the microscopic energy equation of motion we proceed as with the mass and momentum equations and the two terms that appear on the right hand side of Eq. (13).
Taking $`V_{MD}`$ to be a central potential and using the relations $`V_{MD}(r_{ij})=V_{MD}^{}(r_{ij})๐_{ij}=๐
_{ij}`$ and $`\dot{V}_{MD}(r_{ij})=V_{MD}^{}(r_{ij})๐_{ij}๐ฏ_{ij}=๐
_{ij}๐ฏ_{ij}`$ where $`๐ฏ_{ij}=๐ฏ_i๐ฏ_j`$ we get the time rate of change of the particle energy
$$\dot{ฯต_i}=m๐ ๐ฏ_i+\frac{1}{2}\underset{ji}{}๐
_{ij}(๐ฏ_i+๐ฏ_j).$$
(33)
This gives the first term of Eq. (13) in the form
$$\underset{i}{}f_k(๐ฑ_i)\dot{ฯต}=๐_k๐ +\frac{1}{2}\underset{ij}{}f_k(๐ฑ_i)๐
_{ij}(๐ฏ_i+๐ฏ_j).$$
(34)
The last term of this equation is odd under the exchange $`ij`$ and exactly the same manipulations as in Eq. (28) may be used to give
$`{\displaystyle \underset{i}{}}f_k(๐ฑ_i)\dot{ฯต}`$ $`=`$ $`๐_k๐ `$ (35)
$`+`$ $`{\displaystyle \underset{l,ij}{}}f_{kl}(๐ฑ_i){\displaystyle \frac{1}{4}}๐
_{ij}(๐ฏ_i+๐ฏ_j)\mathrm{\Delta }๐ฑ_{ij}๐ซ_{kl}`$ (36)
$`=`$ $`๐_k๐ +{\displaystyle \underset{l,ij}{}}f_{kl}(๐ฑ_i)({\displaystyle \frac{1}{4}}๐
_{ij}(๐ฏ_i^{}+๐ฏ_j^{})`$ (37)
$`+`$ $`{\displaystyle \frac{1}{2}}๐
_{ij}{\displaystyle \frac{๐_k+๐_l}{2}})\mathrm{\Delta }๐ฑ_{ij}๐ซ_{kl}`$ (38)
where for later purposes we have used Eqs. (16) to get the last equation. The last term of Eq. (13) is easily written down using Eq. (14). This gives
$$\underset{i}{}\dot{f}_k(๐ฑ_i)ฯต_i=\underset{li}{}f_{kl}(๐ฑ_i)(๐ฏ_i^{}๐ซ_{kl}+๐ฑ_i^{}๐_{kl})ฯต_i.$$
(39)
As previously we write the particle velocities in terms of $`๐ฏ_i^{}`$. The corresponding expression for the particle energy is $`ฯต_i=ฯต_i^{}+m๐ฏ_i^{}(๐_k+๐_l)/2+(1/2)m((๐_k+๐_l)/2)^2`$ where the prime in $`ฯต_i^{}`$ denotes that the particle velocity is $`๐ฏ_i^{}`$ rather than $`๐ฏ_i`$. Equation (39) may then be written
$`{\displaystyle \underset{i}{}}\dot{f}_k(๐ฑ_i)ฯต_i`$ $`=`$ $`{\displaystyle \underset{l}{}}{\displaystyle \frac{1}{2}}\dot{M}_{kl}\left({\displaystyle \frac{๐_k+๐_l}{2}}\right)^2`$ (40)
$`+`$ $`{\displaystyle \underset{li}{}}f_{kl}(๐ฑ_i)\left(ฯต_i^{}๐ฏ_i^{}+m๐ฏ_i^{}๐ฏ_i^{}{\displaystyle \frac{๐_k+๐_l}{2}}\right)๐ซ_{kl}`$ (41)
$`+`$ $`{\displaystyle \underset{li}{}}f_{kl}(๐ฑ_i)ฯต_i๐ฑ_i^{}๐_{kl}.`$ (42)
Combining this equation with Eq. (38) we obtain
$`\dot{E}_k^{\text{tot}}`$ $`=`$ $`{\displaystyle \underset{li}{}}f_{kl}(๐ฑ_i)\left(๐_{ฯตi}^{}+\mathrm{\Pi }_i^{}{\displaystyle \frac{๐_k+๐_l}{2}}\right)๐ซ_{kl}`$ (43)
$`+`$ $`M_k๐_k๐ +{\displaystyle \underset{l}{}}{\displaystyle \frac{1}{2}}\dot{M}_{kl}\left({\displaystyle \frac{๐_k+๐_l}{2}}\right)^2`$ (44)
$`+`$ $`{\displaystyle \underset{li}{}}f_{kl}(๐ฑ_i)\left(ฯต_i^{}+m๐ฏ_i^{}\left({\displaystyle \frac{๐_k+๐_l}{2}}\right)\right)๐ฑ_i^{}๐_{kl}.`$ (45)
where the momentum-flux tensor is defined in Eq. (32) and we have identified the energy-flux vector associated with a particle $`i`$
$$๐_{ฯตi}=ฯต_i๐ฏ_i+\frac{1}{4}\underset{ij}{}๐
_{ij}(๐ฏ_i+๐ฏ_j)\mathrm{\Delta }๐ฑ_{ij}.$$
(46)
Again the prime denotes that the velocities are $`๐ฏ_i^{}`$ rather than $`๐ฏ_i`$. To get the internal energy $`\dot{E}_k`$ instead of $`\dot{E}_k^{\text{tot}}`$ we note that $`\text{d}(๐_k^2/2M_k)/\text{d}t=๐_k\dot{๐}_k(1/2)\dot{M}_k๐_k^2`$. Using this relation, the momentum equation Eq. (31), as well as the substitution $`(๐_k+๐_l)/2=๐_k๐_{kl}/2`$ in Eq. (45), followed by some rearrangement of the $`\dot{M}_{kl}`$ terms we find that
$`\dot{E}_k^{\text{tot}}`$ $`=`$ $`{\displaystyle \frac{\text{d}}{\text{d}t}}\left({\displaystyle \frac{1}{2}}M_k๐_k^2\right)`$ (47)
$`+`$ $`{\displaystyle \underset{l}{}}{\displaystyle \frac{1}{2}}\dot{M}_{kl}\left({\displaystyle \frac{๐_{kl}}{2}}\right)^2+{\displaystyle \underset{li}{}}f_{kl}(๐ฑ_i)\left(๐_{ฯตi}^{}\mathrm{\Pi }_i^{}{\displaystyle \frac{๐_{kl}}{2}}\right)๐ซ_{kl}`$ (48)
$`+`$ $`{\displaystyle \underset{li}{}}f_{kl}(๐ฑ_i)\left(ฯต_i^{}m๐ฏ_i^{}{\displaystyle \frac{๐_{kl}}{2}}\right)๐ฑ_i^{}๐_{kl}.`$ (49)
This equation has a natural physical interpretation. The first term represents the translational kinetic energy of the DP as a whole. The remaining terms represent the internal energy $`E_k`$. This is a purely thermodynamic quantity which cannot depend on the overall velocity of the DP, i.e. it must be Galilean invariant. This is easily checked as the relevant terms all depend on velocity differences only.
The $`\dot{M}_{kl}`$ term represents the kinetic energy received through mass exchange with neighboring DPs. As will become evident when we turn to the averaged description, the term involving the momentum and energy fluxes represents the work done on the DP by its neighbors and the heat conducted from them. The $`ฯต_i^{}`$-term represents the energy received by the DP due to the same โboundary twistingโ effect that was found in the mass equation. Upon averaging, the last term proportional to $`๐ฏ_i^{}`$ will be shown to be relatively small since $`๐ฏ_i^{}=0`$ in our approximations. This is true also in the mass and momentum equations. Equations (20), (31) and (49) have the coarse grained form that will remain in the final DPD equations. Note, however, that they retain the full microscopic information about the MD system, and for that reason they are time-reversible. Equation (31) for instance contains only terms of even order in the velocity. In the next section terms of odd order will appear when this equation is averaged.
It can be seen that the rate of change of momentum in Eq. (31) is given as a sum of separate pairwise contributions from the other particles, and that these terms are all odd under the exchange $`lk`$. Thus the particles interact in a pairwise fashion and individually fulfill Newtonโs third law; in other words, momentum conservation is again explicitly upheld. The same symmetries hold for the mass conservation equation (20) and energy equation (45).
## III Derivation of dissipative particle dynamics: average and fluctuating forces
We can now investigate the average and fluctuating parts of Eqs. (49), (31) and (20). In so doing we shall need to draw on a hydrodynamic description of the underlying molecular dynamics and construct a statistical mechanical description of our dissipative particle dynamics. For concreteness we shall take the hydrodynamic description of the MD system in question to be that of a simple Newtonian fluid . This is known to be a good description for MD fluids based on Lennard-Jones or hard sphere potentials, particularly in three dimensions . Here we shall carry out the analysis for systems in two spatial dimensions; the generalization to three dimensions is straight forward, the main difference being of a practical nature as the Voronoi construction becomes more involved.
We shall begin by specifying a scale separation between the dissipative particles and the molecular dynamics particles by assuming that
$$|๐ฑ_i๐ฑ_j|<<|๐ซ_k๐ซ_l|,$$
(50)
where $`๐ฑ_i`$ and $`๐ฑ_j`$ denote the positions of neighbouring MD particles. Such a scale separation is in general necessary in order for the coarse-graining procedure to be physically meaningful. Although for the most part in this paper we are thinking of the molecular interactions as being mediated by short-range forces such as those of Lennard-Jones type, a local description of the interactions will still be valid for the case of long-range Coulomb interactions in an electrostatically neutral system, provided that the screening length is shorter than the width of the overlap region between the dissipative particles. Indeed, as we shall show here, the result of doing a local averaging is that the original Newtonian equations of motion for the MD system become a set of Langevin equations for the dissipative particle dynamics. These Langevin equations admit an associated Fokker-Planck equation. An associated fluctuation-dissipation relation relates the amplitude of the Langevin force to the temperature and damping in the system.
### A Definition of ensemble averages
With the mesoscopic variables now available, we need to define the correct average corresponding to a dynamical state of the system. Many MD configurations are consistent with a given value of the set $`\{๐ซ_k,M_k,๐_k,E_k\}`$, and averages are computed by means of an ensemble of systems with common instantaneous values of the set $`\{๐ซ_k,M_k,๐_k,E_k\}`$. This means that only the time derivatives of the set $`\{๐ซ_k,M_k,๐_k,E_k\}`$, i.e. the forces, have a fluctuating part. In the end of our development approximate distributions for $`๐_k`$โs and $`E_k`$โs will follow from the derived Fokker-Planck equations. These distributions refer to the larger equilibrium ensemble that contains all fluctuations in $`\{๐ซ_k,M_k,๐_k,E_k\}`$.
It is necessary, to compute the average MD particle velocity $`๐ฏ`$ between dissipative particle centers, given $`\{๐ซ_k,M_k,๐_k,E_k\}`$. This velocity depends on all neighboring dissipative particle velocities. However, for simplicity we shall only employ a โnearest neighborโ approximation, which consists in assuming that $`๐ฏ`$ interpolates linearly between the two nearest dissipative particles. This approximation is of the same nature as the approximation used in the Newtonian fluid stressโstrain relation which is linear in the velocity gradient. This implies that in the overlap region between dissipative particles $`k`$ and $`l`$
$$๐ฏ^{}=๐ฏ^{}(๐ฑ)=\frac{๐ฑ^{}๐ซ_{kl}}{r_{kl}^2}๐_{kl},$$
(51)
where the primes are defined in Eqs. (16) and $`r_{kl}=|๐ซ_k๐ซ_l|`$.
A preliminary mathematical observation is useful in splitting the equations of motion into average and fluctuating parts. Let $`r(๐ฑ)`$ be an arbitrary, slowly varying function on the $`a^2/r_{kl}`$ scale. Then we shall employ the approximation corresponding to a linear interpolation between DP centers, that $`r(๐ฑ)=(1/2)(r_k+r_l)`$ where $`๐ฑ`$ is a position in the overlap region between DP k and l and $`r_k`$ and $`r_l`$ are values of the function $`r`$ associated with the DP centers k and l respectively.
Then
$`{\displaystyle \underset{i}{}}`$ $`f_{kl}(๐ฑ_i)r(๐ฑ){\displaystyle ๐x๐y\frac{\rho _k+\rho _l}{2}f_{kl}(๐ฑ)\frac{r_k+r_l}{2}}`$ (52)
$``$ $`{\displaystyle \frac{l_{kl}}{2a^2}}{\displaystyle \frac{\rho _k+\rho _l}{2}}{\displaystyle \frac{r_k+r_l}{2}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}๐x^{}\mathrm{cosh}^2(x^{}r_{kl}/a^2)`$ (53)
$`=`$ $`{\displaystyle \frac{l_{kl}}{r_{kl}}}{\displaystyle \frac{\rho _k+\rho _l}{2}}{\displaystyle \frac{r_k+r_l}{2}},`$ (54)
where $`\frac{\rho _k+\rho _l}{2}`$ is the MD particle number density and we have used the identity $`\mathrm{tanh}^{}(x)=\mathrm{cosh}^2(x)`$. We will also need the first moment in $`๐ฑ^{}`$
$`{\displaystyle \underset{i}{}}`$ $`f_{kl}(๐ฑ_i)`$ $`๐ฑ_i^{}r(๐ฑ_i){\displaystyle ๐x๐y\frac{\rho _k+\rho _l}{2}f_{kl}(๐ฑ)๐ฑ^{}\frac{r_k+r_l}{2}}`$ (55)
$``$ $`{\displaystyle \frac{1}{2a^2}}{\displaystyle \frac{\rho _k+\rho _l}{2}}{\displaystyle \frac{r_k+r_l}{2}}{\displaystyle ๐x๐y\mathrm{cosh}^2\left(\frac{xr_{kl}}{a^2}\right)y๐ข_{kl}}`$ (56)
$`=`$ $`{\displaystyle \frac{l_{kl}}{2r_{kl}}}L_{kl}{\displaystyle \frac{\rho _k+\rho _l}{2}}{\displaystyle \frac{r_k+r_l}{2}}๐ข_{kl}`$ (57)
where the unit vectors $`๐_{kl}=๐ซ_{kl}/r_{kl}`$ and $`๐ข_{kl}`$ are shown in Fig. 3, we have used the fact that the integral over $`x๐_{kl}\mathrm{cosh}^2\mathrm{}`$ vanishes since the integrand is odd, and the last equation follows by the substitution $`x(a^2/r_{kl})x`$. In contrast to the vector $`๐_{kl}`$ the vector $`๐ข_{kl}`$ is even under the exchange $`kl`$, as is $`L_{kl}`$. This is a matter of definition only as it would be equally permissible to let $`๐ข_{kl}`$ and $`L_{kl}`$ be odd under this exchange. However, it is important for the symmetry properties of the fluxes that $`๐ข_{kl}`$ and $`L_{kl}`$ have the same symmetry under $`kl`$.
### B The mass conservation equation
Taking the average of Eq. (20), we observe that the first term vanishes if Eq. (51) is used, and the second term follows directly from Eq. (57). We thus obtain
$$\dot{M}_k=\underset{l}{}(\dot{M}_{kl}+\dot{\stackrel{~}{M}}_{kl})$$
(58)
where
$$\dot{M}_{kl}=\underset{li}{}f_{kl}m(๐ฑ_i)๐ฑ_i^{}๐_{kl}=\frac{l_{kl}}{2r_{kl}}L_{kl}\frac{\rho _k+\rho _l}{2}๐ข_{kl}๐_{kl},$$
(59)
and $`\dot{\stackrel{~}{M}}_{kl}=\dot{M}_{kl}\dot{M}_{kl}`$. The finite value of $`\dot{M}_{kl}`$ is caused by the relative DP motion perpendicular to $`๐_{kl}`$. This is a geometric effect intrinsic to the Voronoi lattice. When particles move the Voronoi boundaries change their orientation, and this boundary twisting causes mass to be transferred between DPโs. This mass variation will be visible in the energy flux, though not in the momentum flux. It will later be shown that the effect of mass fluctuations in the momentum and energy equations may be absorbed in the force and heat flux fluctuations.
### C The momentum conservation equation
Using Eq. (59) we may split Eq. (31) into average and fluctuating parts to get
$`{\displaystyle \frac{\text{d}๐_k}{\text{d}t}}`$ $`=`$ $`M_k๐ `$ (60)
$`+`$ $`{\displaystyle \underset{l}{}}\dot{M}_{kl}{\displaystyle \frac{๐_k+๐_l}{2}}+{\displaystyle \underset{li}{}}f_{kl}(๐ฑ_i)๐ท_i๐ซ_{kl}`$ (61)
$`+`$ $`{\displaystyle \underset{i}{}}f_{kl}(๐ฑ_i)m๐ฏ_i^{}๐ฑ_i^{}๐_{kl}+{\displaystyle \underset{l}{}}\stackrel{~}{๐
}_{kl},`$ (62)
where the fluctuating force or, equivalently, the momentum flux is
$`\stackrel{~}{๐
}_{kl}`$ $`=`$ $`{\displaystyle \underset{i}{}}f_{kl}(๐ฑ_i)[(๐ท_i๐ท_i)๐ซ_{kl}`$ (63)
$`+`$ $`m(๐ฏ_i^{}๐ฑ_i^{}๐ฏ_i^{}๐ฑ_i^{})๐_{kl}]`$ (64)
$`+`$ $`\dot{\stackrel{~}{M}}_{kl}{\displaystyle \frac{๐_k+๐_l}{2}}.`$ (65)
Note that by definition $`\stackrel{~}{๐
}_{lk}=\stackrel{~}{๐
}_{kl}`$. The fact that we have absorbed mass fluctuations with the fluctuations in $`\stackrel{~}{๐
}_{kl}`$ deserves a comment. In general force fluctuations will cause mass fluctuations, which in turn will couple back to cause momentum fluctuations. The time scale over which this will happen is $`t_\eta =r_{kl}^2/\eta `$, where $`\eta `$ is the dynamic viscosity of the MD system. This is the time it takes for a velocity perturbation to decay over a distance of $`r_{kl}`$. Perturbations mediated by the pressure, i.e. sound waves, will have a shorter time. In the sequel we shall need to make the assumption that the forces are Markovian, and it is clear that this assumption may only be valid on time scales larger than $`t_\eta `$. Since the time scale of a hydrodynamic perturbation of size $`l`$, say, is also given as $`l^2/\eta `$ this restriction implies the scale separation requirement $`r_{kl}^2<<l^2`$, consistent with the scale $`r_{kl}`$ being mesoscopic.
Since $`๐ท_i`$ is in general dissipative in nature, Eq. (62) and its mass- and energy analogue will be referred to as DPD1. It is at the point of taking the average in Eq. (62) that time reversibility is lost. Note, however, that we do not claim to treat the introduction of irreversibility into the problem in a mathematically rigorous way. This is a very difficult problem in general which so far has only been realized by rigorous methods in the case of some very simple dynamical systems with well defined ergodic properties . We shall instead use the constitutive relation for a Newtonian fluid which, as noted earlier, is an emergent property of Lennard-Jones and hard sphere MD systems, to give Eq. (62) a concrete content. The momentum-flux tensor then has the following simple form
$$\rho ๐ท_i=m\rho \mathrm{๐ฏ๐ฏ}+๐p\eta (๐ฏ+(๐ฏ)^T)$$
(66)
where $`p`$ is the pressure and $`๐ฏ`$ the average velocity of the MD fluid, <sup>T</sup> denotes the transpose and $`๐`$ is the identity tensor . In the above equation we have for simplicity assumed that the bulk viscosity $`\zeta =(2/d)\eta `$ where $`d`$ is the space dimension 2. The modifications to include an independent $`\zeta `$ are completely straight forward.
Using the assumption of linear interpolation (Eq. (51)), the advective term $`\rho \mathrm{๐ฏ๐ฏ}`$ vanishes in the frame of reference of the overlap region since there $`๐ฏ^{}0`$. The velocity gradients in Eq. (66) may be evaluated using Eq. (51); the result is
$$๐ฏ+(๐ฏ)^T=\frac{1}{r_{kl}}\left(๐_{kl}๐_{kl}+๐_{kl}๐_{kl}\right).$$
(67)
Note further that $`_ll_{kl}`$ is in fact a surface integral over the DP surface. Consequently
$$\underset{l}{}l_{kl}๐_{kl}g_k=0$$
(68)
for any function $`g_k`$ that does not depend on $`l`$. In particular we have $`_ll_{kl}๐_{kl}(p_k+p_l)/2=_ll_{kl}๐_{kl}p_{kl}/2`$, where $`p_{kl}=p_kp_l`$. Combining Eqs. (66), (54) and (67), Eq. (62) then takes the form
$`{\displaystyle \frac{\text{d}๐_k}{\text{d}t}}=M_k๐ +{\displaystyle \underset{l}{}}\dot{M}_{kl}{\displaystyle \frac{๐_k+๐_l}{2}}`$ (69)
$``$ $`{\displaystyle \underset{l}{}}l_{kl}\left({\displaystyle \frac{p_{kl}}{2}}๐_{kl}+{\displaystyle \frac{\eta }{r_{kl}}}\left(๐_{kl}+(๐_{kl}๐_{kl})๐_{kl}\right)\right)`$ (70)
$`+`$ $`{\displaystyle \underset{l}{}}\stackrel{~}{๐
}_{kl},`$ (71)
where we have assumed that the pressure $`p`$, as well as the average velocity, interpolates linearly between DP centers, and we have omitted the $`๐ฏ_i^{}๐ฑ_i^{}0`$ term. Note that all terms except the gravity term on the right hand side of Eq. (71) are odd when $`kl`$. This shows that Newtonโs third law is unaffected by the approximations made and that momentum conservation holds exactly. The same statements can be made for the mass equation and the energy equation. The pressure will eventually follow from an equation of state of the form $`p_k=p(E_k,V_k,M_k)`$ where $`V_k`$ is the volume and $`M_k`$ is the mass of DP $`k`$.
### D The energy conservation equation
Splitting Eq. (49) into an average and a fluctuating part gives
$`\dot{E}_k`$ $`=`$ $`{\displaystyle \underset{li}{}}f_{kl}(๐ฑ_i)\left(๐_{ฯตi}^{}\mathrm{\Pi }_i^{}{\displaystyle \frac{๐_{kl}}{2}}\right)๐ซ_{kl}`$ (72)
$`+`$ $`{\displaystyle \underset{li}{}}f_{kl}(๐ฑ_i)ฯต_i^{}๐ฑ_i^{}๐_{kl}`$ (73)
$`+`$ $`{\displaystyle \underset{l}{}}{\displaystyle \frac{1}{2}}\dot{M}_{kl}\left({\displaystyle \frac{๐_{kl}}{2}}\right)^2`$ (74)
$``$ $`{\displaystyle \underset{l}{}}\stackrel{~}{๐
}_{kl}{\displaystyle \frac{๐_{kl}}{2}}+\stackrel{~}{q}_{kl}.`$ (75)
where we have defined
$`\stackrel{~}{q}_{kl}`$ $`=`$ $`{\displaystyle \underset{i}{}}f_{kl}(๐ฑ_i)(๐_{ฯตi}^{}๐_{ฯตi}^{})๐ซ_{kl}+{\displaystyle \frac{\dot{\stackrel{~}{M}}_{kl}}{2}}\left({\displaystyle \frac{๐_{kl}}{2}}\right)^2`$ (76)
$`+`$ $`{\displaystyle \underset{i}{}}f_{kl}(๐ฑ_i)[(ฯต_i^{}๐ฑ_i^{}ฯต_i^{}๐ฑ_i^{})`$ (77)
$``$ $`m{\displaystyle \frac{๐_{kl}}{2}}๐ฏ_i^{}๐ฑ_i^{}]๐_{kl}`$ (78)
i.e. the fluctuations in the heat flux also contains the energy fluctuations caused by mass fluctuations. This is like the momentum case.
Note that in taking the average in Eq. (75) the $`๐ท๐_{kl}`$ product presents no problem as $`๐_{kl}`$ is kept fixed under this average. If we had averaged over different values of $`๐_{kl}`$ the product of velocities in $`๐ท๐_{kl}`$ would have caused difficulties. Equation (75) is the third component in the description at the DPD1 level.
The average of the energy flux vector $`๐_ฯต`$ is taken to have the general form
$$\rho ๐_ฯต=ฯต๐ฏ+\sigma ๐ฏ\lambda T$$
(79)
where $`\sigma =๐ท\rho \mathrm{๐ฏ๐ฏ}`$ is the stress tensor, and $`\lambda `$ the thermal conductivity and $`T`$ the local temperature. Note that in Eq. (49) only $`๐_ฯต^{}`$ appears. Since $`๐ฏ^{}\mathrm{๐}`$ we have $`๐_ฯต^{}=\lambda T`$. Averaging of Eq. (75) gives
$`\dot{E}_k={\displaystyle \underset{l}{}}l_{lk}\lambda {\displaystyle \frac{T_{kl}}{r_{kl}}}`$ (80)
$``$ $`{\displaystyle \underset{l}{}}l_{lk}\left({\displaystyle \frac{p_k+p_l}{2}}๐_{kl}{\displaystyle \frac{\eta }{r_{kl}}}(๐_{kl}+(๐_{kl}๐_{kl})๐_{kl})\right){\displaystyle \frac{๐_{kl}}{2}}`$ (81)
$`+`$ $`{\displaystyle \underset{l}{}}{\displaystyle \frac{1}{2}}\dot{M}_{kl}\left({\displaystyle \frac{๐_{kl}}{2}}\right)^2+{\displaystyle \frac{l_{kl}}{4r_{kl}}}L_{kl}๐ข_{kl}๐_{kl}\left({\displaystyle \frac{E_k}{V_k}}+{\displaystyle \frac{E_l}{V_l}}\right)`$ (82)
$``$ $`{\displaystyle \underset{l}{}}\stackrel{~}{๐
}_{kl}{\displaystyle \frac{๐_{kl}}{2}}+\stackrel{~}{q}_{kl}.`$ (83)
where $`T_{kl}=T_kT_l`$ is the temperature difference between DPโs $`k`$ and $`l`$, and we have used linear interpolation to write $`ฯต_1^{}=(1/2)(E_k/V_k+E_l/V_l)`$. The first term above describes the heat flux according to Fourierโs law. The next non-fluctuating terms, which are multiplied by $`๐_{kl}/2`$ represent the (rate of) work done by the interparticle forces, and the $`\stackrel{~}{๐
}_{kl}`$ term represents the work done by the fluctuating force.
As has been pointed out by Avalos et al and Espanol the work done by $`\stackrel{~}{๐
}_{kl}`$ has the effect that it increases the thermal motion of the DPโs at the expense of a reduction in $`E_k`$. This is the case here as well since the above $`\stackrel{~}{๐
}_{kl}๐_{kl}`$ term always has a positive average due to the positive correlation between the force and the velocity increments.
Equation (83) is identical in form to the energy equation postulated by Avalos and Mackie , save for the fact that here the conservative force $`((p_k+p_l)/2)๐_{kl}๐_{kl}/2`$ (which sums to zero under $`_k`$) is present. The pressure forces in the present case correspond to the conservative forces in conventional DPDโit will be observed that they are both derived from a potential. However, while the conservative force in conventional DPD must be thought to be carried by some field external to the particles, the pressure force in our model has its origin within the particles themselves. There is also a small difference between the present form of Fourierโs law and the description of thermal conduction employed by Avalos and Mackie. While the heat flux here is taken to be linear in differences in $`T`$, Avalos and Mackie use a flux linear in differences in $`(1/T)`$. As both transport laws are approximations valid to lowest order in differences in $`T`$, they should be considered equivalent.
With the internal energy variable at hand it is possible to update the pressure and temperature $`T`$ of the DPโs provided an equation of state for the underlying MD system is assumed, and written in the form $`P=P(E,V,m)`$ and $`T=T(E,V,m)`$. For an ideal gas these are the well known relations $`PV=(2/d)E`$ and $`k_BT=(2/d)mE`$.
Note that we only need the average evolution of the pressure and temperature. The fluctuations of $`p`$ are already contained in $`\stackrel{~}{๐
}_{kl}`$ and the effect of temperature fluctuations is contained within $`\stackrel{~}{q}_{kl}`$.
At this point we may compare the forces arising in the present model to those used in conventional DPD. In conventional DPD the forces are pairwise and act in a direction parallel to $`๐_{kl}`$, with a conservative part that depends only on $`r_{kl}`$ and a dissipative part proportional to $`(๐_{kl}๐_{kl})๐_{kl}`$ . The forces in our new version of DPD are pairwise too. The analog of the conservative force, $`l_{kl}(p_{kl}/2)๐_{kl}`$, is central and its $`๐ซ`$ dependence is given by the Voronoi lattice. When there is no overlap $`l_{kl}`$ between dissipative particles their forces vanish. (A cutโoff distance, beyond which no physical interactions are permitted, was also present in the earlier versions of DPDโsee, for example, Ref. โwhere it was introduced to simplify the numerical treatment.) Due to the existence of an overlap region in our model, the dissipative force has both a component parallel to $`๐_{kl}`$ and a component parallel to the relative velocity $`๐_{kl}`$. However, due to the linear nature of the stressโstrain relation in the Newtonian MD fluid studied here, this force has the same simple linear velocity dependence that has been postulated in the literature.
The friction coefficient is simply the viscosity $`\eta `$ of the underlying fluid times the geometric ratio $`l_{kl}/r_{kl}`$. As has been pointed out both in the context of DPD and elsewhere, the viscosity is generally not proportional to a friction coefficient between the particles. After all, conservative systems like MD are also described by a viscosity. Generally the viscosity will be caused by the combined effect of particle interaction (dissipation, if any) and the momentum transfer caused by particle motion. The latter contribution is proportional to the mean free path. The fact that the MD viscosity $`\eta `$, the DPD viscosity and the friction coefficient are one and the same therefore implies that the mean free path effectively vanishes. This is consistent with the space filling nature of the particles. See Sec. VI B for a further discussion of the zero viscosity limit.
Note that constitutive relations like Eqs. (66) and (79) are usually regarded as components of a top-down or macroscopic description of a fluid. However, any bottom-up mesoscopic description necessarily relies on the use of some kind of averaging procedure; in the present context, these relations represent a natural and convenient although by no means a necessary choice of average. The derivation of emergent constitutive relations is itself part of the programme of non-equilibrium statistical mechanics (kinetic theory), which provides a link between the microscopic and the macroscopic levels. However, as noted above, no general and rigorous procedure for deriving such relations has hitherto been realised; in the present theoretical treatment, such assumed constitutive relations are therefore a necessary input in the linking of the microscopic and mesoscopic levels.
## IV Statistical mechanics of dissipative particle dynamics
In this section we discuss the statistical properties of the DPโs with the particular aim of obtaining the magnitudes of $`\stackrel{~}{๐
}_{kl}`$ and $`\stackrel{~}{q}_{kl}`$. We shall follow two distinct routes that lead to the same result for these quantities, one based on the conventional Fokker-Planck description of DPD, and one based on Landauโs and Lifshitzโs fluctuating hydrodynamics .
It is not straightforward to obtain a general statistical mechanical description of the DP-system. The reason is that the DPโs, which exchange mass, momentum, energy and volume, are not captured by any standard statistical ensemble. For the grand canonical ensemble, the system in question is defined as the matter within a fixed volume, and in the case of a the isobaric ensemble the particle number is fixed. Neither of these requirements hold for a DP in general.
A system which exchanges mass, momentum, energy and volume without any further restrictions will generally be ill-defined as it will lose its identity in the course of time. The DPโs of course remain well-defined by virtue of the coupling between the momentum and volume variables: The DP volumes are defined by the positions of the DP-centers and the DP-momenta govern the motion of the DP-centers. Hence the quantities that are exchanged with the surroundings are not independent and the ensemble must be constructed accordingly.
However, for present purposes we shall leave aside the interesting challenge of designing the statistical mechanical properties of such an ensemble, and derive the magnitude of $`\stackrel{~}{๐
}_{kl}`$ and $`\stackrel{~}{q}_{kl}`$ from two different approximations. The approximations are both justifiable from the assumption that $`\stackrel{~}{๐
}_{kl}`$ and $`\stackrel{~}{q}_{kl}`$ have a negligible correlation time. It follows that their properties may be obtained from the DP behavior on such short time scales that the DP-centers may be assumed fixed in space. As a result, we may take either the DP volume or the system of MD-particles fixed for the relevant duration of time. Hence for the purpose of getting $`\stackrel{~}{๐
}_{kl}`$ and $`\stackrel{~}{q}_{kl}`$ we may use either the isobaric ensemble, applied to the DP system, or the grand canonical ensemble, applied to the MD system. We shall find the same results from either route. The analysis of the DP system using the isobaric ensemble follows the standard procedure using the Fokker-Planck equation, and the result for the equilibrium distribution is only valid in the short time limit. The analysis of the MD system corresponding to the grand canonical ensemble could be conducted along the similar lines. However, it is also possible to obtain the magnitude of $`\stackrel{~}{๐
}_{kl}`$ and $`\stackrel{~}{q}_{kl}`$ directly from the theory of fluctuating hydrodynamics since this theory is derived from coarse-graining the fluid onto a grid. The pertinent fluid velocity and stress fields thus result from averages over fixed volumes associated with the grid points: Since mass flows freely between these volumes the appropriate ensemble is thus the grand canonical one.
### A The isobaric ensemble
We consider the system of $`N_k1`$ MD particles inside a given DP<sub>k</sub> at a given time, say all the MD particles with positions that satisfy $`f_k(๐ฑ_i)>1/2`$ at time $`t_0`$. At later times it will be possible to associate a certain volume per particle with these particles, and by definition the system they form will exchange volume and energy but not mass. We consider all the remaining DPโs as a thermodynamic bath with which DP<sub>k</sub> is in equilibrium. The system defined in this way will be described by the Gibbs free energy and the isobaric ensemble. Due to the diffusive spreading of MD-particles, this system will only initially coincide with the DP; during this transient time interval, however, we may treat the DPโs as systems of fixed mass and describe them by the approximation $`\dot{M}_{kl}=0`$. The magnitudes of $`\stackrel{~}{q}`$ and $`\stackrel{~}{๐
}`$ follow in the form of fluctuation-dissipation relations from the Fokker-Planck equivalent of our Langevin equations. The mathematics involved in obtaining fluctuation-dissipation relations is essentially well-known from the literature , and our analysis parallels that of Avalos and Mackie . However, the fact that the conservative part of the conventional DP forces is here replaced by the pressure and that the present DPโs have a variable volume makes a separate treatment enlightening.
The probability $`\rho (V_k,๐_k,E_k)`$ of finding DP<sub>k</sub> with a volume $`V_k`$, momentum $`๐_k`$ and internal energy $`E_k`$ is then proportional to $`\mathrm{exp}(S_T/k_B)`$ where $`S_T`$ is the entropy of all DPโs given that the values $`(V_k,๐_k,E_k)`$ are known for DP<sub>k</sub>. If $`S^{}`$ denotes the entropy of the bath we can write $`S_T`$ as
$`S_T`$ $`=`$ $`S^{}(V_TV_k,๐_T๐_k,E_T{\displaystyle \frac{P_k^2}{2M_k}}E_k)+S_k`$ (84)
$``$ $`S^{}(V_T,๐_T,E_T){\displaystyle \frac{S^{}}{E}}\left(E_k+{\displaystyle \frac{P_k^2}{2M_k}}\right){\displaystyle \frac{S^{}}{V}}V_k`$ (85)
$``$ $`{\displaystyle \frac{S^{}}{๐}}๐_k+S_k`$ (86)
where the derivatives are evaluated at $`(V_T,๐_T,E_T)`$ and thus characterize the bath only. Assuming that $`๐_T`$ vanishes there is nothing in the system to give the vector $`S^{}/๐`$ a direction, and it must therefore vanish as well . The other derivatives give the pressure $`p_0`$ and temperature $`T_0`$ of the bath and we obtain
$$S_T=S^{}(V_T,๐_T,E_T)\frac{1}{T_0}\left(G_k+\frac{P_k^2}{2M_k}\right)$$
(87)
where the Gibbs free energy has the standard form $`G_k=E_k+p_0V_kT_0S_k`$. Since there is nothing special about DP<sub>k</sub> it immediately follows that the the full equilibrium distribution has the form
$$\rho ^{\text{eq}}=Z^1(T_0,p_0)\mathrm{exp}\left(\beta _0\underset{k}{}\frac{P_k^2}{2M_k}+G_k\right),$$
(88)
where $`\beta _0=1/(k_BT_0)`$. The temperature $`T_k=(S_k/E_k)^1`$ and pressure $`p_k=T_k(S_k/V_k)`$ will fluctuate around the equilibrium values $`T_0`$ and $`p_0`$. The above distribution is analyzed by Landau and Lifshitz who show that the fluctuations have the magnitude
$$\mathrm{\Delta }P_k^2=\frac{k_BT_0}{V_k\kappa _S},\mathrm{\Delta }T_k^2=\frac{k_BT_0^2}{Vc_v}$$
(89)
where the isentropic compressibility $`\kappa _S=(1/V)(V/P)_S`$ and the specific heat capacity $`c_v`$ are both intensive quantities. Comparing our expression with the distribution postulated by Avalos and Mackie, we have replaced the Helmholtz by the Gibbs free energy in Eq. (88). This is due to the fact that our DPโs exchange volume as well as energy.
We write the fluctuating force as
$$\stackrel{~}{๐
}_{kl}=๐_{kl}W_{kl}+๐_{kl}W_{kl}$$
(90)
where, for reasons soon to become apparent, we have chosen to decompose $`\stackrel{~}{๐
}_{kl}`$ into components parallel and perpendicular to $`๐_{kl}`$. The $`W`$โs are defined as Gaussian random variables with the correlation function
$`W_{kl\alpha }(t)W_{nm\beta }(t^{})`$ $`=`$ $`\delta _{\alpha \beta }\delta (tt^{})(\delta _{kn}\delta _{lm}+\delta _{km}\delta _{ln})`$ (91)
where $`\alpha `$ and $`\beta `$ denote either $``$ or $``$. The product of $`\delta `$ factors ensures that only equal vectorial components of the forces between a pair of DPโs are correlated, while Newtonโs third law guarantees that $`๐_{kl}=๐_{lk}`$. Likewise the fluctuating heat flux takes the form
$$\stackrel{~}{q}_{kl}=\mathrm{\Lambda }_{kl}W_{kl}$$
(92)
where $`W_{kl}`$ satisfies Eq. (91) without the $`\delta _{\alpha \beta }`$ factor and energy conservation implies $`\mathrm{\Lambda }_{kl}=\mathrm{\Lambda }_{lk}`$.
The force correlation function then takes the form
$`\stackrel{~}{๐
}_{kn}(t)\stackrel{~}{๐
}_{lm}(t^{})`$ $`=`$ $`(๐_{kn}๐_{lm}+๐_{kn}๐_{lm})`$ (94)
$`(\delta _{kl}\delta _{nm}+\delta _{km}\delta _{ln})\delta (tt^{})`$
$``$ $`๐_{klnm}(\delta _{kl}\delta _{nm}+\delta _{km}\delta _{ln})\delta (tt^{})`$ (95)
where we have introduced the second order tensor $`๐_{knlm}`$.
It is a standard result in non-equilibrium statistical mechanics that a Langevin description of a dynamical variable $`๐ฒ`$
$$\dot{๐ฒ}=๐(๐ฒ)+\stackrel{~}{๐}$$
(96)
where $`\stackrel{~}{๐}`$ is a delta-correlated force has an equivalent probabilistic representation in terms of the Fokker-Planck equation
$$\frac{\rho (๐ฒ,t)}{t}=(๐(๐ฒ)\rho (๐ฒ))+\frac{1}{2}:(๐(๐ฒ)\rho (๐ฒ))$$
(97)
where $``$ denotes derivatives with respect to $`๐ฒ`$ and $`\rho (๐ฒ,t)`$ is the probability distribution for the variable $`๐ฒ`$ at time $`t`$, $`\stackrel{~}{๐}(๐ฒ,t)\stackrel{~}{๐}(๐ฒ,t^{})=๐\delta (tt^{})`$ and $`๐`$ is a symmetric tensor of rank two .
In the preceding paragraph, $`๐`$ denotes all the fluctuating terms in Eqs. (71) and (83). Using the above definitions and $`\dot{M}_{kl}=0`$ it is a standard matter to obtain the Fokker-Planck equation
$$\frac{\rho }{t}=(L_0+L_{\text{DIS}}+L_{\text{DIF}}),\rho $$
(98)
where
$`L_0`$ $`=`$ $`{\displaystyle \underset{k}{}}{\displaystyle \frac{}{๐ซ_k}}๐_k+{\displaystyle \underset{kl}{}}l_{kl}({\displaystyle \frac{}{๐_k}}๐_{kl}{\displaystyle \frac{p_{kl}}{2}}`$ (99)
$`+`$ $`{\displaystyle \frac{}{E_k}}๐_{kl}๐_{kl}{\displaystyle \frac{p_k+p_l}{4}})`$ (100)
$`L_{\text{DIS}}`$ $`=`$ $`{\displaystyle \underset{kl}{}}l_{kl}\left({\displaystyle \frac{}{๐_k}}๐
_{kl}^D{\displaystyle \frac{}{E_k}}\left({\displaystyle \frac{๐_{kl}}{2}}๐
_{kl}^D\lambda {\displaystyle \frac{T_{kl}}{r_{kl}}}\right)\right)`$ (101)
$`L_{\text{DIF}}`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{kl}{}}(๐_{klkl}{\displaystyle \frac{}{๐_k}}_{kl}{\displaystyle \frac{}{E_k}}(๐_{klkl}{\displaystyle \frac{๐_{kl}}{2}}_{kl}`$ (102)
$``$ $`\mathrm{\Lambda }_{kl}^2({\displaystyle \frac{}{E_k}}{\displaystyle \frac{}{E_l}}))),`$ (103)
$`๐
_{kl}^D=(\eta /r_{kl})(๐_{kl}+(๐_{kl}๐_{kl})๐_{kl})`$, and the sum $`_{kl}`$ runs over both $`k`$ and $`l`$. The operator $`_{kl}`$ is defined as in Ref. :
$$_{kl}=\left(\frac{}{๐_k}\frac{}{๐_l}\right)\frac{๐_{kl}}{2}\left(\frac{}{E_k}\frac{}{E_l}\right).$$
(104)
The steady-state solution of Eq. (98) is already given by Eq. (88); following conventional procedures we can obtain the fluctuation-dissipation relations for $`๐`$ and $`\mathrm{\Lambda }`$ by inserting $`\rho ^{\text{eq}}`$ in Eq. (98).
Apart from the tensorial nature of $`๐_{klkl}`$ the operators $`L_{\text{DIS}}`$ and $`L_{\text{DIF}}`$ are essentially identical to those published earlier in conventional DPD . However, the โLiouvilleโ operator $`L_0`$ plays a somewhat different role as it contains the $`/E_k`$ term, corresponding to the fact that the pressure forces do work on the DPโs to change their internal energy.
While $`L_0\rho ^{\text{eq}}`$ conventionally vanishes exactly by construction of the inter-DP forces, here it vanishes only to order $`1/N_k`$. In order to evaluate $`L_0\rho ^{\text{eq}}`$ we need the following relationship
$$\frac{}{๐ซ_k}=\frac{1}{2}\underset{kl}{}l_{kl}๐_{kl}\left(\frac{}{V_l}\frac{}{V_k}\right),$$
(105)
which is derived by direct geometrical consideration of the Voronoi construction. By repeated use of Eq. (68) it is then a straightforward algebraic task to obtain
$$L_0\rho ^{\text{eq}}=\frac{\rho ^{\text{eq}}}{4}\underset{kl}{}l_{kl}๐_{kl}๐_k\left[\frac{p_l}{E_l}\frac{p_{kl}T_{kl}}{k_BT_kT_l}\right],$$
(106)
which does not vanish identically. However, note that if we estimate $`E_lN_lk_BT`$ we obtain $`p_l/E_l(1/N_k)(p_l/k_BT)`$. Similarly we may estimate $`p_{kl}`$ and $`T_{kl}`$ from Eq. (89) to obtain
$$\frac{p_{kl}T_{kl}}{k_BT_kT_l}\frac{\sqrt{\mathrm{\Delta }P^2\mathrm{\Delta }T^2}}{k_BT_kT_l}=\frac{1}{N_k}\sqrt{\frac{N_k/V_k}{\kappa _Sc_vT_0^2}}.$$
(107)
The last square root is an intensive quantity of the order $`p_0/(k_BT_0)`$, as may be easily demonstrated for the case of an ideal gas. Since each separate quantity that is contained in the differences in the square brackets of Eq. (106) is of the order $`p_0/T_0`$ we have shown that they cancel up to relative order $`1/N_k1`$. In fact, it is not surprising that Langevin equations which approximate local gradients to first order only in the corresponding differences, like $`T_{kl}`$, give rise to a Fokker-Planck description that contains higher order correction terms.
Having shown that $`L_0\rho ^{\text{eq}}`$ vanishes to a good approximation we may proceed to obtain the fluctuation-dissipation relations from the equation $`(L_{\text{DIS}}+L_{\text{DIF}})\rho ^{\text{eq}}=0`$. It may be noted from Eq. (103) that this equation is satisfied if
$`(l_{kl}๐
_{kl}^D+{\displaystyle \frac{1}{2}}๐_{klkl}_{kl})\rho ^{\text{eq}}`$ $`=`$ $`0`$ (108)
$`\left(l_{kl}\lambda {\displaystyle \frac{T_{kl}}{r_{kl}}}+{\displaystyle \frac{1}{2}}\mathrm{\Lambda }_{kl}^2\left({\displaystyle \frac{}{E_k}}{\displaystyle \frac{}{E_l}}\right)\right)\rho ^{\text{eq}}`$ $`=`$ $`0.`$ (109)
Using the identity
$$๐_{kl}๐_{kl}+๐ข_{kl}๐ข_{kl}=๐$$
(110)
where $`๐ข_{kl}`$ a vector normal to $`๐_{kl}`$, we may show that Eq. (109) implies that
$`\omega _{kl}^2`$ $`=`$ $`2\omega _{kl}^2=4\eta k_B\mathrm{\Theta }_{kl}{\displaystyle \frac{l_{kl}}{r_{kl}}}`$ (111)
$`\mathrm{\Lambda }_{kl}^2`$ $`=`$ $`2k_BT_kT_l\lambda {\displaystyle \frac{l_{kl}}{r_{kl}}},`$ (112)
where $`\mathrm{\Theta }_{kl}^1=(1/2)(T_k^1+T_l^1)`$.
### B $`\stackrel{~}{๐
}`$ from fluctuating hydrodynamics
Having derived the fluctuation-dissipation relations from the approximation of the isobaric ensemble we now derive the same result from fluctuating hydrodynamics, which corresponds to the grand canonical ensemble. We shall only derive the magnitude of $`\stackrel{~}{๐
}_{kl}`$ since $`\stackrel{~}{q}`$ follows on the basis of the same reasoning.
Fluctuating hydrodynamics is based on the conservation equations for mass, momentum and energy with the modification that the momentum and energy fluxes contain an additional fluctuating term. Specifically, the momentum flux tensor takes the form $`P+\rho \mathrm{๐ฏ๐ฏ}+\sigma ^{}`$, where $`P`$ is the pressure, $`๐ฏ`$ is the velocity field and the viscous stress tensor is given as
$$\sigma ^{}=\eta \left(๐ฏ+๐ฏ^T\frac{2}{d}๐ฏ\right)+\zeta ๐ฏ+๐ฌ,$$
(113)
where $`๐ฌ`$ is the fluctuating component of the momentum flux. From the same approximations as we used in deriving Eq. (112), i.e. a negligible correlation time for the fluctuating forces, Landau and Lifshitz derive
$`๐ฌ(๐ฑ,t)\mathrm{๐ง๐ฌ}(๐ฑ^{},0)๐ง=2k_BT\left(\eta (1+\mathrm{๐ง๐ง})+(\zeta {\displaystyle \frac{2}{d}}\eta )\mathrm{๐ง๐ง}\right)`$ (114)
$`\delta (t)\{\begin{array}{cc}\frac{1}{\mathrm{\Delta }V_n}& \text{if }๐ฑ,๐ฑ^{}\epsilon \mathrm{\Delta }V_n\\ 0& \text{otherwise}\end{array}`$ (117)
where $`๐ง`$ is an arbitrary unit vector and $`n`$ labels the volume element $`\mathrm{\Delta }V_n`$. By following the derivations presented by Landau and Lifshitz, it may be noted that nowhere is it assumed that the $`\mathrm{\Delta }V_n`$โs are cubic or stationary.
By making the identifications $`\zeta =(2/d)\eta `$, $`๐ง๐_{kl}`$ $`\stackrel{~}{๐
}_{kl}=l_{kl}๐ฌ๐_{kl}`$, $`\mathrm{\Delta }V_nV_{kl}`$, (shown in Fig. 4), and $`T=\mathrm{\Theta }_{kl}`$ we may immediately write down
$`\stackrel{~}{๐
}_{kl}(t)\stackrel{~}{๐
}_{nm}(0)`$ $`=`$ $`2{\displaystyle \frac{k_B\mathrm{\Theta }_{kl}l_{kl}^2}{V_{kl}}}\eta (1+๐_{kl}๐_{nm})\delta (t)`$ (118)
$`(`$ $`\delta _{kn}\delta _{lm}+\delta _{km}\delta _{ln})`$ (119)
where again the last sum of $`\delta `$-factors ensures that $`kl`$ and $`nm`$ denote the same DP pair. Observing from Fig. 4 that $`V_{kl}=l_{kl}r_{kl}`$, it now follows directly from Eq. (119) that
$`\stackrel{~}{๐
}_{kl}(t)๐_{kl}\stackrel{~}{๐
}_{nm}(0)๐_{nm}`$ $`=`$ $`2\stackrel{~}{๐
}_{kl}(t)๐ข_{kl}\stackrel{~}{๐
}_{nm}(0)๐ข_{nm}`$ (120)
$`=`$ $`4k_B\mathrm{\Theta }_{kl}{\displaystyle \frac{l_{kl}}{r_{kl}}}\eta \delta (t)`$ (121)
$`(`$ $`\delta _{kn}\delta _{lm}+\delta _{km}\delta _{ln})`$ (122)
which is nothing but the momentum part of Eq. (112). That the fluctuating heat flux $`\stackrel{~}{q}`$ produces the form of fluctuation-dissipation relations given in Eq. (112) follows from a similar analysis. Thus the approximation of fixed DP volume $`V_k`$ produces the same result as the approximation of fixed number of MD particles $`N_k`$. This is due to the fact that both approximations are based on the assumption that the DPโs are only considered within a time interval which is longer than the correlation time of the fluctuations but shorter than the time needed for the DPโs to move significantly.
The result given in Eq. (122) was derived from the somewhat arbitrary choice of discretizaton volume $`V_{kl}`$; this is the volume which corresponds to the segment $`l_{kl}`$ over which all forces have been taken as constant. It is thus the smallest discretization volume we may consistently choose. It is reassuring that Eq. (122) also follows from different choices of $`\mathrm{\Delta }V_n`$. For example, one may readily check that Eq. (122) is obtained if we split $`V_{kl}`$ in two along $`r_{kl}`$ and consider $`\stackrel{~}{๐
}_{kl}`$ to be the sum of two independent forces acting on the two parts of $`l_{kl}`$.
We are now in a position to quantify the average component $`\dot{\stackrel{~}{E}}_k_{lk}\stackrel{~}{๐
}_{kl}๐_{kl}/2`$ of the fluctuations in the internal energy given in Eq. (83). Writing the velocity in response to $`\stackrel{~}{๐
}_{kl}`$ as $`\stackrel{~}{๐}_k=_{lk}_{\mathrm{}}^t\text{d}t^{}\stackrel{~}{F}_{kl}(t^{})/M_k`$, we get that $`\dot{\stackrel{~}{E}}_k=_{\mathrm{}}^t\text{d}t^{}\stackrel{~}{F}_{kl}(t^{})\stackrel{~}{F}_{kl}(t)`$ which by Eqs. (112) and (95) becomes $`\dot{\stackrel{~}{E}}_k=(1/M_k)3l_{kl}\eta k_B\mathrm{\Theta }_{kl}/r_{kl}`$. This result is the same as one would have obtained applying the rules of Itรด calculus to $`\stackrel{~}{๐}_k^2/(2M_k)`$. It yields the modified, though equivalent, energy equation
$`\dot{E}_k={\displaystyle \underset{l}{}}l_{lk}\lambda {\displaystyle \frac{T_{kl}}{r_{kl}}}`$ (123)
$``$ $`{\displaystyle \underset{l}{}}l_{lk}\left({\displaystyle \frac{p_k+p_l}{2}}๐_{kl}{\displaystyle \frac{\eta }{r_{kl}}}(๐_{kl}+(๐_{kl}๐_{kl})๐_{kl})\right){\displaystyle \frac{๐_{kl}}{2}}`$ (124)
$``$ $`{\displaystyle \underset{l}{}}\stackrel{~}{๐
^{}}_{kl}{\displaystyle \frac{๐_{kl}}{2}}3{\displaystyle \frac{l_{kl}}{r_{kl}}}\eta k_B\mathrm{\Theta }_{kl}+\stackrel{~}{q}_{kl}.`$ (125)
where we have written $`\stackrel{~}{๐
^{}}_{kl}`$ with a prime to denote that it is uncorrelated with $`๐_{kl}`$. In a numerical implementation this implies that $`\stackrel{~}{๐
^{}}_{kl}`$ must be generated from a different random variable than $`\stackrel{~}{๐
}_{kl}`$, which was used to update $`๐_{kl}`$.
The fluctuation-dissipation relations Eqs. (112) complete our theoretical description of dissipative particle dynamics, which has been derived by a coarse-graining of molecular dynamics. All the parameters and properties of this new version of DPD are related directly to the underlying molecular dynamics, and properties such as the viscosity which are emergent from it.
## V Simulations
While the present paper primarily deals with theoretical developments we have carried out simulations to test the equilibrium behavior of the model in the case of the isothermal model. This is a crucial test as the derivation of the fluctuating forces relies on the most significant approximations. The simulations are carried out using a periodic Voronoi tesselation described in detail elsewhere .
Figure 5 shows the relaxation process towards equilibrium of an initially motionless system. The DP temperature is measured as $`๐_k^2/(2M_k)`$ for a system of DPs with internal energy equal to unity. The simulations were run for 4000 iterations of 5000 dissipative particles and a timestep $`\text{d}t=0.0005`$ using an initial molecular density $`\rho =5`$ for each DP. The molecular mass was taken to be $`m=1`$, the viscosity was set at $`\eta =1`$, the expected mean free path is 0.79, and the Reynolds number (See Sec. VI B) is Re=2.23. It is seen that the convergence of the DP system towards the MD temperature is good, a result that provides strong support for the fluctuation-dissipation relations of Eq. (112).
## VI Possible applications
### A Multiscale phenomena
For most practical applications involving complex fluids, additional interactions and boundary conditions need to be specified. These too must be deduced from the microscopic dynamics, just as we have done for the interparticle forces. This may be achieved by considering a particulate description of the boundary itself and including molecular interactions between the fluid MD particles and other objects, such as particles or walls. Appropriate modifications can then be made on the basis of the momentum-flux tensor of Eq. (32), which is generally valid.
Consider for example the case of a colloidal suspension, which is shown in Fig. 6. Beginning with the hydrodynamic momentum-flux tensor Eq. (32) and Eq. (71), it is evident that we also need to define an interaction region where the DPโcolloid forces act: the DPโcolloid interaction may be obtained in the same form as the DPโDP interaction of Eq. (71) by making the replacement $`l_{kl}L_{kI}`$, where $`L_{kI}`$ is the length (or area in 3D) of the arc segment where the dissipative particle meets the colloid (see Fig. 6) and the velocity gradient $`r_{kl}^1((๐_{kl}๐_{kl})๐_{kl}+๐_{kl})`$ is that between the dissipative particle and the colloid surface. The latter may be computed using $`๐_k`$ and the velocity of the colloid surface together with a no-slip boundary condition on this surface. In Eq. (112) the replacement $`l_{kl}L_{KI}`$ must also be made.
Although previous DPD simulations of colloidal fluids have proved rather successful at low to intermediate solids volume fractions, they break down for dense systems whose solids volume fraction exceeds a value of about 40% because the existing method is unable to handle multiple lengthscale phenomena. However, our new version of the algorithm provides the freedom to define dissipative particle sizes according to the local resolution requirements as illustrated in Fig. 6. In order to increase the spatial resolution where colloidal particles are within close proximity it is necessary and perfectly admissible to introduce a higher density of dissipative particles there; this ensures that fluid lubrication and hydrodynamic effects are properly maintained. After these dissipative particles have moved it may be necessary to re-tile the DP system; this is easily achieved by distributing the mass and momentum of the old dissipative particles on the new ones according to their area (or volume in 3D). Considerations of space prevent us from discussing this problem further in the present paper, but we plan to report in detail on such dense colloidal particle simulations using our method in future publications. We note in passing that a wide variety of other complex systems exist where modeling and simulation are challenged by the presence of several simultaneous length scales, for example in polymeric and amphiphilic fluids, particularly in confined geometries such as porous media .
### B The low viscosity limit and high Reynolds numbers
In the kinetic theory derived by Marsh, Backx and Ernst the viscosity is explicitly shown to have a kinetic contribution $`\eta _K=\rho D/2`$ where $`D`$ is the DP self diffusion coefficient and $`\rho `$ the mass density. The kinetic contribution to the viscosity was measured by Masters and Warren within the context of an improved theory. How then can the viscosity $`\eta `$ used in our model be decreased to zero while kinetic theory puts the lower limit $`\eta _K`$ to it?
To answer this question we must define a physical way of decreasing the MD viscosity while keeping other quantities fixed, or, alternatively rescale the system in a way that has the equivalent effect. The latter method is preferable as it allows the underlying microscopic system to remain fixed. In order to do this we non-dimensionalize the DP momentum equation Eq. (71).
For this purpose we introduce the characteristic equilibrium velocity, $`U_0=\sqrt{k_BT/M}`$, the characteristic distance $`r_0`$ as the typical DP size. Then the characteristic time $`t^{}=r_0/U_0`$ follows.
Neglecting gravity for the time being Eq. (71) takes the form
$`{\displaystyle \frac{\text{d}๐_k^{}}{\text{d}t^{}}}={\displaystyle \underset{l}{}}l_{kl}^{}\left({\displaystyle \frac{p_{kl}^{}}{2}}๐_{kl}+{\displaystyle \frac{1}{\text{Re}}}\left(๐_{kl}^{}+(๐_{kl}^{}๐_{kl})๐_{kl}\right)\right)`$ (126)
$`+`$ $`{\displaystyle \underset{l}{}}{\displaystyle \frac{l_{kl}^{}L_{kl}^{}}{2r_{kl}^{}}}{\displaystyle \frac{\rho _k^{}+\rho _l^{}}{2}}๐ข_{kl}๐_{kl}^{}{\displaystyle \frac{๐_k^{}+๐_l^{}}{2}}+{\displaystyle \underset{l}{}}\stackrel{~}{๐
}_{kl}^{},`$ (127)
where $`๐_k^{}=๐_k/(MU_0)`$, $`p_{kl}^{}=p_{kl}r_0^2/(MU_0^2)`$, $`M=\rho r_0^2`$ in 2d, the Reynolds number $`\text{Re}=U_0r_0\rho /\eta `$ and $`\stackrel{~}{๐
}_{kl}^{}=(r_0/MU_0^2)\stackrel{~}{๐
}_{kl}`$ where $`\stackrel{~}{๐
}_{kl}`$ is given by Eqs. (90) and (112). A small calculations then shows that if $`\stackrel{~}{๐
}_{kl}^{}`$ is related to $`\omega _{kl}^{}`$ and $`t^{}`$ like $`\stackrel{~}{๐
}_{kl}`$ related to $`\omega _{kl}`$ and $`t`$, then
$$\omega _{kl}^{}_{}{}^{}2\frac{1}{\text{Re}}\frac{k_BT}{MU_0^2}\frac{1}{\text{Re}}$$
(128)
where we have neglected dimensionless geometric prefactors like $`l_{kl}/r_{kl}`$ and used the fact that the ratio of the thermal to kinetic energy by definition of $`U_0`$ is one.
The above results imply that when the DPD system is measured in non-dimensionalized units everything is determined by the value of the mesoscopic Reynolds number Re. There is thus no observable difference in this system between increasing $`r_0`$ and decreasing $`\eta `$.
Returning to dimensional units again the DP diffusivity may be obtained from the Stokes-Einstein relation as
$$D=\frac{k_BT}{ar_0\eta }$$
(129)
where $`a`$ is some geometric factor ($`a=6\pi `$ for a sphere) and all quantities on the right hand side except $`r_0`$ refer directly to the underlying MD. As we are keeping the MD system fixed and increasing Re by increasing $`r_0`$, it is seen that $`D`$ and hence $`\eta _K`$ vanish in the process.
We note in passing that if $`D`$ is written in terms of the mean free path $`\lambda `$: $`D=\lambda \sqrt{k_BT/(\rho r_0^2)}`$ and this result is compared with Eq. (129) we get $`\lambda ^{}=\lambda /r_01/r_0`$ in 2d, i.e. the mean free path, measured in units of the particle size decreases as the inverse particle size. This is consistent with the decay of $`\eta _K`$. The above argument shows that decreasing $`\eta `$ is equivalent to keeping the microscopic MD system fixed while increasing the DP size, in which case the mean free path effects on viscosity is decreased to zero as the DP size is increased to infinity. It is in this limit that high Re values may be achieved.
Note that in this limit the thermal forces $`\stackrel{~}{๐
}_{kl}\text{Re}^{1/2}`$ will vanish, and that we are effectively left with a macroscopic, fluctuationless description. This is no problem when using the present Voronoi construction. However, the effectively spherical particles of conventional DPD will freeze into a colloidal crystal, i.e. into a lattice configuration in this limit. Also while conventional DPD has usually required calibration simulations to determine the viscosity, due to discrepancies between theory and measurements, the viscosity in this new form of DPD is simply an input parameter. However, there may still be discrepancies due to the approximations made in going from MD to DPD. These approximations include the linearization of the inter-DP velocity fields, the Markovian assumption in the force correlations and the neglect of a DP angular momentum variable.
None of the conclusions from the above arguments would change if we had worked in three dimensions in stead of two.
## VII Conclusions
We have introduced a systematic procedure for deriving the mesoscopic modeling and simulation method known as dissipative particle dynamics from the underlying description in terms of molecular dynamics.
Figure 7 illustrates the structure of the theoretical development of DPD equations from MD as presented in this paper. The initial coarse graining leads to equations of essentially the same structure as the final DPD equations. However, they are still invariant under time- reversal. The label DPD1 refers to Eqs. (58), (62) and (75), whereas the DPD2 equations have been supplemented with specific constitutive relations both for the non-equilibrium fluxes (momentum and heat) and an equilibrium description of the thermodynamics. These equations are Eqs. (71) and (83) along with Eqs. (112). The development we have made which is shown in Fig. 7 does not claim to derive the irreversible DPD equations from the reversible ones of molecular dynamics in a rigorous manner, although it does illustrate where the transition takes place with the introduction of molecular averages. The kinetic equations of this new DPD satisfy an $`H`$-theorem, guaranteeing an irreversible approach to the equilibrium state. Note that in passing to the time-asymmetric description by the introduction of the averaged description of Eq. (66), a time asymmetric non-equilibrium ensemble is required .
This is the first time that any of the various existing mesoscale methods have been put on a firm โbottom upโ theoretical foundation, a development which brings with it numerous new insights as well as practical advantages. One of the main virtues of this procedure is the capability it provides to choose one or more coarse-graining lengthscales to suit the particular modeling problem at hand. The relative scale between molecular dynamics and the chosen dissipative particle dynamics, which may be defined as the ratio of their number densities $`\rho _{\text{DPD}}/\rho _{\text{MD}}`$, is a free parameter within the theory. Indeed, this rescaling may be viewed as a renormalisation group procedure under which the fluid viscosity remains constant: since the conservation laws hold exactly at every level of coarse graining, the result of doing two rescalings, say from MD to DPD$`\alpha `$ and from DPD$`\alpha `$ to DPD$`\beta `$, is the same as doing just one with a larger ratio, i.e. $`\rho _{\text{DPD}\beta }/\rho _{\text{MD}}=(\rho _{\text{DPD}\beta }/\rho _{\text{DPD}\alpha })(\rho _{\text{DPD}\alpha }/\rho _{\text{MD}})`$.
The present coarse graining scheme is not limited to hydrodynamics. It could in principle be used to rescale the local description of any quantity of interest. However, only for locally conserved quantities will the DP particle interactions take the form of surface terms as here, and so it is unlikely that the scheme will produce a useful description of non-conserved quantities.
In this context, we note that the bottom-up approach to fluid mechanics presented here may throw new light on aspects of the problem of homogeneous and inhomogeneous turbulence. Top-down multiscale methods and, to a more limited extent, ideas taken from renormalisation group theory have been applied quite widely in recent years to provide insight into the nature of turbulence ; one might expect an alternative perspective to emerge from a fluid dynamical theory originating at the microscopic level, in which the central relationship between conservative and dissipative processes is specified in a more fundamental manner. From a practical point of view it is noted that, since the DPD viscosity is the same as the viscosity emergent from the underlying MD level, it may be treated as a free parameter in the DPD model, and thus high Reynolds numbers may be reached. In the $`\eta 0`$ limit the model thus represents a potential tool for hydrodynamic simulations of turbulence. However, we have not investigated the potential numerical complications of this limit.
The dissipative particle dynamics which we have derived is formally similar to the conventional version, incorporating as it does conservative, dissipative and fluctuating forces. The interactions are pairwise, and conserve mass and momentum as well as energy. However, now all these forces have been derived from the underlying molecular dynamics. The conservative and dissipative forces arise directly from the hydrodynamic description of the molecular dynamics and the properties of the fluctuating forces are determined via a fluctuationโdissipation relation.
The simple hydrodynamic description of the molecules chosen here is not a necessary requirement. Other choices for the average of the general momentum and energy flux tensors Eqs. (46) and (32) may be made and we hope these will be explored in future work. More significant is the fact that our analysis permits the introduction of specific physicochemical interactions at the mesoscopic level, together with a well-defined scale for this mesoscopic description.
While the Gaussian basis we used for the sampling functions is an arbitrary albeit convenient choice, the Voronoi geometry itself emerged naturally from the requirement that all the MD particles be fully accounted for. Well defined procedures already exist in the literature for the computation of Voronoi tesselations and so algorithms based on our model are not computationally difficult to implement. Nevertheless, it should be appreciated that the Voronoi construction represents a significant computational overhead. This overhead is of order $`N\mathrm{log}N`$, a factor $`\mathrm{log}N`$ larger than the most efficient multipole methods in principle available for handling the particle interactions in molecular dynamics. However, the prefactors are likely to be much larger in the particle interaction case.
Finally we note the formal similarity of the present particulate description to existing continuum fluid dynamics methods incorporating adaptive meshes, which start out from a top-down or macroscopic description. These top-down approaches include in particular smoothed particle hydrodynamics and finite-element simulations. In these descriptions too the computational method is based on tracing the motion of elements of the fluid on the basis of the forces acting between them . However, while such top-down computational strategies depend on a macroscopic and purely phenomenological fluid description, the present approach rests on a molecular basis.
###### Acknowledgements.
It is a pleasure to thank Frank Alexander, Bruce Boghosian and Jens Feder for many helpful and stimulating discussions. We are grateful to the Department of Physics at the University of Oslo and Schlumberger Cambridge Research for financial support which enabled PVC to make several visits to Norway in the course of 1998; and to NATO and the Centre for Computational Science at Queen Mary and Westfield College for funding visits by EGF to London in 1999 and 2000. |
warning/0002/astro-ph0002460.html | ar5iv | text | # The hard X-ray emission of luminous infrared galaxies
## 1 Introduction.
The nature and the energy source of the Luminous Infrared Galaxies (LIGs) has been a matter of debate since their discovery by IRAS, more than 10 years ago. This class of sources is characterized by a very high luminosity emitted mostly in the farโinfrared ($`\mathrm{L}_{\mathrm{IR}}>10^{11}\mathrm{L}_{}`$). Such a luminosity is higher than in normal galaxies ($`\mathrm{L}_{\mathrm{IR}}10^{10}`$ erg s<sup>-1</sup> for a large spiral like the Milky Way or M31) and therefore an additional source of energy is required. The infrared emission of LIGs is due to the presence of a large amount of dust that reprocesses optical and UV radiation to IR wavelengths. Since the primary emission is often not directly observable, its origin remains controversial.
Two main mechanisms have been invoked in this context: an intense starburst activity and/or an Active Galactic Nucleus (AGN). It is now widely accepted that in general both sources can be present in the LIGs, but the relative contribution of each is still unclear.
Recent studies with ISO of a sample of ULIGs (Ultra Luminous Infrared Galaxies, $`\mathrm{L}_{\mathrm{IR}}>10^{12}\mathrm{L}_{}`$) indicate that the bulk of the IR emission is due to starburst activity (Genzel et al. 1998, Lutz et al. 1998). The diagnostic used in these works is mainly based on mid-infrared spectroscopy: in particular, the policyclic aromatic hydrocarbon features at 7.7 $`\mu `$m are claimed to be better than optical lines at discriminating between starburst and AGN, because of the lower extinction.
On the other side, numerous X-ray observations performed in the last few years, which will be in part reviewed in this paper, suggest that an AGN contribution to the observed IR luminosity is common.
Resolving this issue would be important not only for the understanding of the LIG phenomenon, but also to assess the contribution of the LIGs to the diffuse X-ray background (XRB). The knowledge of the fraction of the bolometric luminosity which emerges in the hard X-ray domain would allow the inclusion of the LIGs in the models of the XRB. These models are by and large successful in synthesizing the XRB from the contributions of individual AGNs; they must however include a dominant contribution from absorbed, type 2 AGNs, which up to now have been observed directly only at low luminosities and low redshifts. The zeroth-order extrapolation, which assumes that type 2s and type 1s evolve in the same way (โunified modelโ), has to face several discrepancies, and could be cured only by adding extra type 2 sources at intermediate or high redshifts (Gilli et al. 1999). Some of the LIGs could very well be the required additional sources, but their mean X-ray properties are still poorly known.
The expected X-ray spectrum, and the relative contribution of the X-ray emission to the bolometric luminosity are very different in the cases of AGN and starburst dominance, therefore hard (2-10 keV) X-ray data can be a powerful diagnostic.
The typical X-ray starburst spectrum is a thermal continuum with temperatures ranging from a fraction of a keV to several keV (in units of kT), plus several emission lines at low energies, and a K<sub>ฮฑ</sub> feature of highly ionized iron at the energy of E$`6.7`$ keV (Cappi et al. 1999, Persic et al. 1998). The 2-10 keV luminosity is less than one part in 10<sup>3</sup> of the total emission.
The AGN spectrum between 2 and 10 keV is strongly dependent on the amount of obscuration suffered by the primary continuum component: if the absorbing column density $`\mathrm{N}_\mathrm{H}`$ is lower than $`10^{24}`$ cm<sup>-2</sup>, the direct continuum spectrum is visible at energies higher than the photoelectric cutoff E<sub>C</sub>, whose value depends on N<sub>H</sub>. The mean direct continuum, as deduced from the spectra of unabsorbed AGNs, is well represented by a powerlaw with photon index $`\mathrm{\Gamma }1.7`$. In addition, a K<sub>ฮฑ</sub> iron line is observed at energies ranging from 6.4 keV (neutral iron) to 6.95 keV (hydrogen-like iron). The current view is that the iron line originates both from the accretion disc and from a reprocessing medium located farther away (1โ10 pc): the first component can be broadened by relativistic effects and is suppressed if the absorbing column density is higher than several times 10<sup>23</sup> cm<sup>-2</sup>, the second component is narrow, and its energy depends on the ionization degree of the reprocessing medium. The equivalent width (EW) of the line is EW $``$ 100-200 eV with respect to the intrinsic continuum, and can be much higher (up to some keV) if the continuum is heavily absorbed at the line energy (while the line, produced in a different region, is still visible).
If the column density is higher than $`10^{24}`$ cm<sup>-2</sup> (i.e. the medium is Compton thick), the intrinsic emission is completely absorbed and only a reflected component survives. The X-ray spectrum of Compton thick sources is generally characterized by an Fe line with large equivalent width ($``$ 1 keV) and a reflection dominated continuum flatter than the intrinsic spectrum (see Maiolino et al. 1998 for further details). While the observed 2โ10 keV contribution to the bolometric luminosity can be higher than 10% in unabsorbed AGNs, in Compton thick Seyfert 2s it can be lower than 0.1%.
Summarizing, the X-ray continuum of a starburst is expected to have a thermal high energy cutoff and a soft spectrum with emission lines, while an AGN is characterized by a hard power law with a possible low energy cutoff due to absorption or, in the extremely absorbed cases (Compton thick), a flat reflection spectrum. Cold (E$``$ 6.4 keV) and broad iron lines, or a large relative X-ray luminosity are clear signatures of an AGN, whereas narrow lines of ionized iron or a relatively weak X-ray emission are possible in both scenarios.
In this paper we study the sample of all the luminous infrared galaxies observed so far in hard X rays (2-10 keV). We combine the X-ray information with the IRAS photometry, and find interesting correlations between X-ray properties and IR colours. We propose a simple model that matches the observed correlations and is compatible with the available spectroscopic and spectropolarimetric data at near-IR and optical wavelengths.
In Sect. 2 we describe in detail our sample, in Sect. 3 we analyze the X-ray properties of the sources and the correlation between the X and IR emission. Finally, in Sect. 4 we discuss our results in the framework of a geometrical two-parameter model. Conclusions and future work are summarized in Sect. 5. In the Appendix we briefly report the results of our own spectral analysis of those sources which have not been published elsewhere. Throughout this paper we use the cosmological parameters H$`{}_{0}{}^{}=75`$ km s<sup>-1</sup> Mpc<sup>-1</sup> and q$`{}_{0}{}^{}=0.5`$.
## 2 The sample
The selection of our sample is primarily based on the existence of data in the 2-10 keV energy range. We collected all the galaxies with published 2-10 keV data, and those with unpublished observations in the ASCA and BeppoSAX archives. Within the X-ray sample we selected all the sources detected in at least three IRAS bands with L$`{}_{IR}{}^{}>10^{11}`$ L. The final sample consists of 78 objects. The X-ray data were obtained from the literature for 63 sources, from the ASCA Public Archive for 10, and from BeppoSAX observations for the remaining 5. The IR data were obtained from the NASA Extragalactic Database.
We briefly discuss the 15 unpublished sources in the Appendix.
Table 1 gives all the relevant information about the sample. The infrared (8-1000 $`\mu `$m) luminosity has been obtained from the IRAS fluxes with the following prescription (see, for instance, Sanders & Mirabel 1996):
$$F_{IR}=\mathrm{1.8\; 10}^{11}(13.48f_{12}+5.16f_{25}+2.58f_{60}+f_{100})$$
(1)
were $`f_{12},f_{25},f_{60},f_{100}`$ are the flux densities in the four IRAS filters expressed in Jansky and $`F_{IR}`$ is expressed in units of erg cm<sup>-2</sup> s<sup>-1</sup>. When we have only an upper limit for one of the four IRAS points we use one half of the measured upper limit.
The sample is very heterogeneous, and to assess the validity of any conclusion one must consider the selection criteria of the many subsamples which were merged.
We note from Table 1 that most objects are optically classified as AGNs and some of them are among the best known and widely studied type 1 and type 2 AGNs. In fact, 46 objects out of 78 were selected by the original authors because of their well known AGN activity, irrespective of their IR properties. The remaining 32 were instead selected for observation because of their high IR luminosity (many of them are ULIGs, i.e. sources with $`\mathrm{L}_{\mathrm{IR}}>10^{12}`$ L, see Fig. 1). Nevertheless, also this second half is strongly biased toward AGN-dominated sources: many targets were selected because of the presence of some AGN indicator, such as broad lines in the polarized optical spectrum, or compact radio cores with high brightness temperature. A few objects had no AGN indicator of any sort, and were observed to assess the energy source in a โgenericโ LIG; they were chosen among bright IR sources without previous hard X-ray data, so they are underluminous in X rays with respect to the average LIG, and in Fig. 1 they are adiacent to well known starbursts.
In Fig. 1 we show the X-ray to IR flux ratio versus the infrared luminosity. The different codes correspond to the selection criteria which were the drivers of the X-ray observation. We note that the much discussed correlation between the IR luminosity and the occurrence of AGN activity (Sanders & Mirabel 1996, Genzel et al. 1998, Lutz et al. 1998) does not appear in Fig. 1, just because of the strong bias in favor of AGN-dominated sources discussed above. Well known AGNs and well known starbursts<sup>1</sup><sup>1</sup>1 The classification of ARP 220 is ambiguous (see Smith et al. 1998, Genzel et al. 1998, Kim, Veilleux & Sanders 1998); here we have adopted the starburst classification, but moving the source to the Seyfert 2 class would not change any of our conclusions. cluster very clearly in the upper and lower half of the diagram, respectively, confirming the diagnostic value of the X-ray to IR luminosity ratio. We further note that all flavours of sources are represented in the sample, although their relative numbers are altered with respect to a proper, unbiased sample, and they are spread over a large region of the parameter space: thus the true distributions of the various observables among LIGs cannot be derived from our data set, but any correlation among these observables is likely to be real, as discussed in the following.
Finally, we note that some of the AGNs in the sample are radio loud. However, with the exception of a blazar (3C273), their radio โloudnessโ does not appear to introduce any anomalous behaviour, in terms of their X-ray or IR properties, with respect to the radio quiet sources (which are the great majority); this will be discussed in Sect. 3.2.
## 3 X and IR properties of the sample
### 3.1 X rays
1) X-ray absorption in AGNs The X-ray properties of our sample are extremely varied. Most of the objects optically classified as type 1 Seyferts of QSOs are bright X-ray emitters (relative to their IR luminosity), and have no signs of large absorption in their spectra (with the exception of the three Broad Absorption Line QSOs, which will be discussed later). A significant fraction of the narrow line AGNs are also relatively bright in the X rays and have spectra characterized by Compton thin absorbtion ($`\mathrm{N}_\mathrm{H}<10^{24}\mathrm{cm}^2`$). In other narrow line AGNs the 2โ10 keV spectrum is very flat and/or the Fe line at 6.4 keV has a large equivalent width ($``$ 1 keV); as discussed in the Introduction both these features identify reflection dominated AGNs whose direct X-ray radiation is obscured along our line of sight by Compton thick gas ($`\mathrm{N}_\mathrm{H}>10^{24}\mathrm{cm}^2`$). For 11 narrow line AGNs, most of which are new sources whose analysis is presented in the Appendix, there is no spectral information in the hard X-rays, and only a weak X-ray detection or an upper limit are available. In these cases either the AGN is heavily obscured (Compton thick), or it is intrinsically weak (and the IR luminosity dominated by the starburst), or a combination of the two. If the flux of the (reddening corrected) \[OIII\]$`\lambda `$5007 narrow line is assumed to be a fair, isotropic indicator of the intrinsic luminosity, then the ratio between the flux of this line and the observed X-ray flux provides information about the absorption affecting the nuclear X-ray source. In particular, as discussed in Bassani et al. (1999) and Maiolino et al. (1998), a flux ratio X/\[OIII\] lower than unity is typical of Compton thick sources. Of the 11 X-ray faint AGNs, 7 have measured \[OIII\] and Balmer decrement; their X/\[OIII\] ratio is always $`<1`$, and indicates that these must be Compton thick, reflection dominated sources. For the remaining 4 AGNs, no conclusion can be drawn on the X-ray absorption and luminosity.
2) X-ray spectral indices: We have studied the (instrinsic) X-ray spectral indices of a sub-sample of objects with an X/IR flux ratio higher than 10<sup>-2</sup>. As discussed in the following, this condition selects mostly broad line AGNs with no X-ray absorption or narrow line AGNs with mild X-ray absorption, so that the intrinsic spectral index can be reliably derived by correcting for the observed photoelectric cutoff. Also, we only considered sources with X-ray observations down to 0.5 keV, in order to detect possible soft thermal components, and dropped all the Narrow Line Seyfert 1s, which are known to have an intrinsic spectrum different (steeper) from โnormalโ Seyferts. This sub-sample was divided in two groups according to the IR colour<sup>2</sup><sup>2</sup>2In this work we use an infrared colour defined as the ratio between the fluxes in the IRAS 25 $`\mu `$m and 60$`\mu `$m bands, that, according to Eq. 1, is given by C$`=\frac{2f_{25}}{f_{60}}`$., C, being larger or smaller than 1.1. As shown in Fig. 2, objects with cold IR colour have an X-ray spectrum softer than usually measured in Seyfert 1s and quasars: the photon index of 6 sources out of 8 is $`1.98<\mathrm{\Gamma }<2.5`$. As for the group with warm IR colour, only 2 sources out of 14 have a photon index $`\mathrm{\Gamma }>2`$, while the remaining 12 have $`\mathrm{\Gamma }<2`$ (Fig. 2).
We must caution that the data are not homogeneous: for many objects we have the result of the spectral fitting with a single powerlaw, while for others the fit includes a separate thermal component to account for the soft excess. In the latter cases we re-fitted the composite model spectra with a single powerlaw, to make the results comparable with each other. Of course our simple procedure is far from precise, and a more detailed analysis would require a new fit to the original data.
This finding on the spectral indices is relevant to the explanation we propose for the correlations of Fig. 3, and will be discussed in the next Section.
3) The BeppoSAX subsample: In the Appendix we report the analysis of 5 BeppoSAX observations of LIGs. Three of them are ULIGs with compact radio emission and brightness temperature T$`{}_{B}{}^{}>10^7`$ K, one (NGC 1073) is a Seyfert 2 with a strong \[O III\] emission; the last one (IRAS 00198-7926) is a Seyfert 2 with warm IR colour. In addition to these, our sample contains 8 more type 2 sources observed with BeppoSAX: IRAS 11058-1131, IRAS 22017+0319, IRAS 20210+1121, IRAS 14454-4343, studied by Ueno et al. (2000); NGC 1365 (Risaliti et al. 2000), IRAS 09104+4109 (Franceschini et al. 2000), NGC 1068 (Matt et al. 1997) and NGC 6240 (Vignati et al. 1999). The 13 sources listed above constitute a first small sample observed up to 200 keV by means of the PDS instrument. Among these sources only three, IRAS 22017+0319, NGC 1365 and NGC 6240, clearly show a direct emission, with a typical Seyfert 2 spectrum that extends from 1 to 200 keV and an absorbing column density of 4$`\times 10^{22}`$ cm<sup>-2</sup>, 4$`\times 10^{23}`$ cm<sup>-2</sup> and 2$`\times 10^{24}`$ cm<sup>-2</sup>, respectively. The remaining 10 sources are very weak in X rays (with the exception of NGC 1068, that has a good signal-to-noise spectrum, even if it is completely Compton-thick up to 200 keV, Matt et al. 1997), and have a flat 2-10 keV spectral index (when at all constrained), all hints of heavy obscuration. For 6 out of these 10 sources we have a detection in the PDS (marginal in all objects but NGC 1068) that is consistent with the extrapolation towards higher energies of the 2-10 keV best fit. The most straightforward interpretation of these objects is that they are Compton thick and reflectionโdominated even in the 10โ200 keV band, which implies a column density $`>10^{25}`$ cm<sup>-2</sup>. Alternatively, some of these AGNs might be fading and their cold reflection dominated spectrum could be the echo of their past activity, as observed in NGC 2992 (Weaver et al. 1996, Gilli et al. 2000), Mkn 3 (Iwasawa et al. 1994) and NGC 4051 (Guainazzi et al. 1998). In particular, the latter might be the case for Mkn 273. For one object, IRAS 09104+4109, a marginal detection of an excess was obtained, implying $`\mathrm{N}_\mathrm{H}5\times 10^{24}`$ cm<sup>-2</sup>. For the remaining three sources we can only put a lower limit to $`\mathrm{N}_\mathrm{H}`$ of $`10^{24}`$ cm<sup>-2</sup>, either because the PDS upper limit is inconclusive, or because there is a confusing source in the PDS field of view.
We will discuss in Sect. 4.1 the implication of such a large number of reflection dominated sources.
### 3.2 X-IR colour-magnitude diagram
An infrared indicator of the presence of an AGN is the ratio between the fluxes at 25 and 60 $`\mu `$m (see Note 1 for its definition). The 60 $`\mu `$m emission is mainly due to reprocessing of the UV-optical radiation by dust heated at intermediate temperatures ($``$50K), that is present both in starbursts and in the circumnuclear torus of AGNs. The 25 $`\mu `$m emission is due to warmer dust (T$``$ 100 K) that is more abundant in the central regions of AGNs. We note that this simple scheme is altered by the obscuration thought to be present in type 2 AGNs, since the 25 $`\mu `$m emission can be absorbed as well (assuming a Galactic dust to gas ratio, the optical depth at 25 $`\mu `$m is $`\tau >1`$ for $`\mathrm{N}_\mathrm{H}>10^{23}`$ cm<sup>-2</sup>, Draine 1989). This effect will be analyzed in detail in the following Section.
In Fig. 3 we plot the X/IR flux ratio versus the 25-60 $`\mu `$m colour<sup>3</sup><sup>3</sup>3We excluded from the plot in Fig. 3 four of the objects listed in Table 1: 3C 273, whose infrared emission is known to be dominated by a blazar component, at least in outburst; and IRAS 15307+3252, PG 1148+549 and PG 1634+706, that are the only three objects with redshift as high as $``$1, where the K-correction could be important.. The different symbols are related to the optical classification (Seyfert 1/QSO, Seyfert 2, starburst) and to the Compton thin / Compton thick classification derived from the X-ray analysis discussed in Sect. 3.1.
A clear correlation is apparent in Fig. 3: type 1 AGNs have preferentially high X/IR ratios and warm infrared colours. Moving towards lower 25/60$`\mu `$m ratios we find lower X/IR ratios and an increasing fraction of obscured AGNs at first, and of starbursts afterwards. Together with the correlation we note also a large scatter, which will be discussed in detail in the next Section.
The result presented in Fig. 3 is in qualitative agreement with general expectations. In addition, it poses interesting quantitative constraints on the physics of LIGs. Since the boundaries depicted in Fig. 3 are crucial for our subsequent analysis, we must: (1) check if the lack of objects with low X-ray flux and warm IR colour and/or with high X-ray flux and cold IR colour can be ascribed to the selection biases illustrated in the previous Section; (2) compare the LIGs with non-LIG sources, in order to see if the correlation of Fig. 3 is a characteristic of the former or is present in the generic hard X-ray extragalactic population.
(1) The shortage of objects in the top-left region of Fig. 3 is very likely real, since Xโray loud objects should have been detected easily. For instance, the same region is populated in Fig. 4. Concerning the lack of sources in the bottom-right part of the Figure, we have studied a sample of IR-warm LIGs not observed in X rays: from the IR colour and the IR flux we have estimated the minimum X-ray flux that a source should have in order to follow the correlation of Fig. 3. If the minimum flux were higher than -say- $`10^{11}`$ erg cm<sup>-2</sup> s<sup>-1</sup> for a significant number of sources, then we should conclude that probably these objects do not match the correlation, because it would be very improbable that such X-ray loud sources escaped detection in all previous X-ray surveys. We used to this purpose the catalog of de Grijp (1985), which is a selection of the IR-warmest sources in the IRAS Point Source Catalog. We found that for all the LIGs of this catalog not included in our sample the minimum required X-ray flux is lower than about 10<sup>-12</sup> erg cm<sup>-2</sup> s<sup>-1</sup>. Summarizing, there are no obvious bias effects which could ascribe the observed correlation to selection problems.
(2) The second point has been investigated by checking if a similar correlation exists also for sources with L $`<10^{11}`$L. Our control sample of non-LIGs is composed of $`50`$ AGNs with hard X-ray data available in the literature. The sample is not complete neither representative of a particular class of sources, but it shows nonetheless (see Fig. 4) that the absence of strong X-ray emitters with cold IR colours is a feature typical of the LIGs. On the contrary, the lack of IR warm and X-ray weak objects could be a general property of hard X-ray sources.
Finally, we note that the radio loud (broad line) AGNs, marked with a horizontal segment in Fig. 3, do not have an anomalous behaviour with respect to type 1 radio quiet objects. Instead, the three broad absorption line (BAL) QSOs (identified with a filled circle in Fig. 3) are characterized by a low X-ray luminosity compared to the other broad line AGNs. This is a well known property of BAL QSOs and it is ascribed to absorption associated to the outflowing medium responsible for the BALs (Brandt et al. 1999, Crenshaw et al. 1999).
## 4 Discussion, and a basic scenario
Our approach is to compare the observed distribution of Fig. 3 with various combinations of starburst and AGN spectra, the latter absorbed according to different prescriptions.
The horizontal line in the upper part of Fig. 5 represents a model with only the AGN component: the rightmost point, at coordinates 2.1 (colour) and 0.1 (flux ratio), is a typical location for bright Seyfert 1s. This โstarting pointโ cannot be the average of the Seyfert 1s in our sample, because some of them are contaminated by starburst activity. On the other hand, mean values of lower luminosity objects could contain a contribution from the host. We then adopted as a template the bright Seyfert 1 IC4329A, that is not a LIG, but has an infrared luminosity ($`\mathrm{L}_{\mathrm{IR}}6\times `$ 10<sup>10</sup>L) significantly higher than normal galactic values, no indication of starburst activity and a well studied, warm IR spectrum. The X-ray source is modelled as a powerlaw with photon index $`\mathrm{\Gamma }=1.7`$, and the absorbing material โcovering the X-ray source in the same way as the infrared sourceโ is assumed to have a Galactic dustโtoโgas ratio and extinction curve. In the infrared, only the 25$`\mu `$m flux is affected by the absorption, while the 60$`\mu `$m flux and the bolometric flux remain constant.
The model is obviously incorrect, and does not fit the observed trend even for those many objects that are known to be dominated by the AGN component. The prediction is a fast decrease of the colour parameter at an almost constant X/IR flux ratio, whereas the data show that the two quantities decline together. To decrease the X-ray flux by an amount comparable to the 25$`\mu `$m flux, the absorbing material must have properties different from the Galactic gas, or it must be distributed differently against the X-ray and the IR sources.
The warm dust emitting the bulk of the 25 $`\mu `$m emission is located at several pc from the nuclear source and therefore suffers a much lower absorption than the nuclear X-ray source. This might explain the discrepancy between model and data. In a scenario like that of Fig. 6, if the medium-IR source is at a fixed height above the plane, one could crudely assume that it sees a fixed small fraction of the absorption seen by the X-ray source, for a substantial range of inclination angles.
$$\mathrm{N}_H(\mathrm{IR})=k\times \mathrm{N}_H(\mathrm{X})$$
(2)
where k is a constant. We have adopted a Galactic gasโtoโdust ratio and changed the IR absorption with respect to the X-ray one, but we could have taken an identical $`\mathrm{N}_\mathrm{H}`$ toward both sources and a lower $`A_V/\mathrm{N}_H`$. Then the constant k in the equation above, which in our model is a geometrical factor, would represent the correction to the dustโtoโgas ratio.
The oblique rightmost line in Fig. 5 represents the double absorber, AGN dominated model described in Fig. 6. The N<sub>H</sub> values indicated by the labels refer to the X-ray absorption, while the parameter k was adjusted (k $`70`$) to fit the boundary of the populated region. In order to add a starburst component, we took the average colour and X/IR flux ratio of the six starburst-dominated sources in our sample, since in these cases the possible contamination from additional components (such as a very obscured AGN) should be negligible. The dotted curves give different degrees of mixing, starting from the right with objects 100% dominated by a more or less absorbed AGN.
In summary, the position of an object along the main extension ($``$ vertical) of the plot is a measure of the amount of absorption incurred by the AGN component, while the position across the plot, along the horizontal direction, is a measure of the amount of mixing with a starburst component.
This interpretation is supported by the optical identifications. However, there are some further points which support the proposed picture: 1) According to our scheme, the type 2 sources having the warmest IR colour allowed by their X/IR flux ratio are dominated by the AGN component, with low absorption on large scales. Therefore we expect to find broad polarized or near-IR broad lines more often in these objects than in colder ones (see also Heisler, Lumsden & Bailey 1997). In objects with colder IR colors the starburst is more and more important, and the AGN lines could be diluted by the radiation contributed by the starburst, or suppressed by the large scale absorption associated with it. There are only a few type 2 sources with the required data (spectropolarimetry and/or near-IR spectra), but those with broad polarized or near-IR broad lines all lie along the AGN locus in Fig. 5, as expected.
2) The X-ray signature of the starburst contribution is a thermal component with typical kT values of 0.1-3 keV. Therefore, in this case, the overall X-ray spectrum should be steeper, and the fit with a single powerlaw should give a photon index larger than in typical AGNs, say $`\mathrm{\Gamma }>2`$. The evidence presented in Fig. 2 goes exactly in this direction: the large majority of the IR-cold sources have steep (starburst-like) X-ray indices while IR-warm sources have flatter (AGN-like) indices.
Our model seems to provide a qualitative explanation for several observed trends. However, the Sy1s located above and to the right of the โstarting pointโ of Fig. 5 give an idea of the intrinsic spread of these two quantities; the spread is comparable to the horizontal width of the plot, and indicates that care must be exercised in drawing strong conclusions from the IR colour.
### 4.1 Non-detections at E $`>`$ 20 keV
In Sect. 3.1 we briefly presented a subsample of 13 LIGs with observations up to 200 keV, provided by the PDS onboard BeppoSAX. Observations at several tens of keV allow the exploration of the $`\mathrm{N}_\mathrm{H}`$ range 10<sup>24</sup> cm$`{}_{}{}^{2}<\mathrm{N}_\mathrm{H}<10^{25}`$ cm<sup>-2</sup>, since these column densities suppress completely the direct radiation up to about 10 keV, but become translucent at higher energies. Well known examples are NGC 4945 (Done et al. 1996) and the Circinus galaxy (Matt et al. 1999). A significant lack of objects with 10<sup>24</sup> cm$`{}_{}{}^{2}<\mathrm{N}_\mathrm{H}<10^{25}`$ cm<sup>-2</sup> was already pointed out in a sample of optically selected Seyfert 2s (Risaliti et al. 1999). Here we seem to find the same result for the LIGs: among all sources obscured between 2 and 10 keV, and having PDS observations, 6 are reflection dominated up to 100 keV, implying an $`\mathrm{N}_\mathrm{H}>10^{25}`$ cm<sup>-2</sup> (or, alternatively, a fading nucleus); only one has 10<sup>24</sup> cm$`{}_{}{}^{2}<\mathrm{N}_\mathrm{H}<10^{25}`$ cm<sup>-2</sup>, while for the others no firm conclusion can be drawn.
This finding, if confirmed for a larger and well selected sample, would have important consequences in the synthesis models of the X-ray background. If a large fraction of sources are reflection dominated up to 100 keV, the integrated flux of the absorbed AGNs is dimmed and they are on average less detectable at a given flux. Therefore, the ratio between absorbed and unabsorbed AGNs should be increased with respect to the local value to account for the hard XRB intensity and hard X-ray source counts (Gilli et al. 1999).
### 4.2 An unbiased sample of LIGs
As we already discussed in Sect. 2, the sample studied in this paper is heterogeneous and incomplete. However, since it contains all the X-ray data currently available, we can extract from it the least biased subsample that it is possible to assemble at present. To do this, we ordered the LIGs of the Bright Galaxy Sample (BGS, Soifer et al. 1987, Sanders et al. 1995) according to their total infrared flux, and selected the first $`n`$ of them so as to maximize the fraction of objects with hard X-ray data. It is worth noting that by selecting in total IR flux we remove a bias that affects the BGS catalogue (and any randomly selected subsample). In fact the BGS is flux-limited at 60$`\mu `$m and therefore the effective cut in terms of total IR flux is higher for the AGNs (which on average have a lower fraction of the IR luminosity emitted at 60$`\mu `$m) than for the starbursts (which have the peak of their emission at 60 $`\mu `$m). The magnitude of the effect can be estimated by using the sources listed in Table 1 and the equation quoted in Sect. 2: we find that the fraction of the IR flux emitted at 60$`\mu `$m is 50-55% for the starbursts and 15-20% for the IR-warmest AGNs. So, in order to eliminate the bias, one has to remain a factor of at least $`3`$ above the limit of the parent catalog.
The number of objects included in the subsample is 23, which corresponds to a limiting total IR flux of 1.6 10<sup>-9</sup> erg cm<sup>-2</sup> s<sup>-1</sup>. In this way the fraction of X-ray observed sources is highest (15 out of 23, i.e. 65%) and the lowest value of the 60$`\mu `$m flux density is 15.5 Jy, indeed a factor $`3`$ higher than the limiting flux of the BGS (5.24 Jy). The sources are listed in Table 2.
Although a complete statistical study of their properties is at present impossible because of the lack of X-ray data for 8 sources, we can nonetheless draw some preliminary conclusions. According to the optical classification, 5 of the 15 X-ray observed sources are starbursts, while the remaining 10 are AGNs. Among the latter, 6 are Comptonโthick and only one (NGC 1365) is Comptonโthin with a column density $`\mathrm{N}_\mathrm{H}10^{23}`$ cm<sup>-2</sup>. In three sources, the AGN seen in the optical is completely invisible in the X rays. Finally, the 25/60 $`\mu `$m flux ratio is generally low, and indicates that the emission of the warm dust associated with the AGN is either low compared with the starburst 60$`\mu `$m emission, or absorbed.
The statement we can make is that the AGNs are rather common among the luminous infrared galaxies (13 out of 23 in the whole sample), but their contribution to the bolometric luminosity is not easy to assess. The AGN could provide the higher fraction of the energy only if its direct emission were entirely reprocessed; alternatively, the contribution of the starburst could be important or even dominant.
## 5 Conclusions
This paper deals with the sample of luminous infrared galaxies (LIGs) observed to date in hard X rays (2-10 keV). The sample is affected by a selection bias in favor of AGN-dominated sources, nevertheless it covers a region of the parameter space large enough to allow some general conclusions. For a significant fraction of the sample (15/78) hard X-ray data are analyzed and published for the first time.
The main results of our work can be summarized as follows:
1) The X-ray properties of the LIGs are very different from source to source: the X-ray brightest objects are Seyfert 1 and lowโabsorption Seyfert 2 galaxies, with a X/IR flux ratio as high as 0.1-0.2. At the opposite extreme are the starburst dominated objects, for which the X/IR flux ratio is $`10^4`$ or lower. Interestingly, we note that a significant fraction of the sources optically classified as AGNs do not show any indication of nuclear activity in X rays. Thus, either the AGN contribution is negligible, or the direct emission must be absorbed by a column density $`\mathrm{N}_\mathrm{H}>10^{24}`$ cm<sup>-2</sup>.
2) A subsample of 13 LIGs was observed up to 200 keV. Six out of 11 sources that are Comptonโthick in the 2-10 keV range are reflection dominated also from 10 to 200 keV, thus implying a column density $`\mathrm{N}_\mathrm{H}>10^{25}`$ cm<sup>-2</sup> (or, alternatively, a fading nucleus). Only two have an excess in the 15-100 keV range, while for the last three the hard 15-100 keV X-ray emission remains unconstrained. This result is somewhat puzzling, for the 10$`{}_{}{}^{24}10^{25}`$ cm<sup>-2</sup> interval of column density appears to be underpopulated both for the LIGs and the non-LIG local Seyfert 2 galaxies.
3) The shape of the correlation between the X/IR flux ratio and the 25/60 $`\mu `$m IR colour suggests that the 25$`\mu `$m emission is absorbed by a lower column density than the X-ray emission or, alternatively, that the absorbing material has a dustโtoโgas ratio lower than Galactic.
The X-IR correlation is well reproduced by a model in which both AGNs and starbursts contribute to the total emission. The differences in the IR colours are mainly due to the different contribution of the starburst component, while the X/IR flux ratio is mainly determined by the amount of absorption affecting the AGN.
Additional evidence in support of the model is the detection of broad optical polarized or infrared lines only in the IR-warmest sources (where we see directly the inner region of the torus), and the steepness of the X-ray spectra of the IR-colder objects (where the starburst contributes strongly).
4) Finally, we assembled an unbiased sample of LIGs from the BGS catalog, in a way which maximizes the fraction of sources with hard X-ray data available (15 out of 23, i.e. 65%). The sample is limited in total IRAS flux (i.e. from 12$`\mu `$m to 100$`\mu `$m). This selection criterion removes a bias against AGNs present in most other LIG samples that, instead, are limited in the 60$`\mu `$m flux. From the optical classification and from the IR and X-ray data we find that AGNs are common among the LIGs (13 out of 23 host an active nucleus) but they are either weak or heavily obscured.
###### Acknowledgements.
Three of us (RG, RM and MS) were supported in part by the Italian Ministry for University and Research (MURST) through grant Cofin 98-02-32. The detailed comments of the referee, K. Iwasawa, have greatly improved the presentation.
## Appendix A Unpublished sources
We present here a brief discussion of the X-ray data of 15 sources of our sample for which no information is available in the literature. Among the 10 sources observed by ASCA, only 3 (TOL 1351-375, IRAS 05189-2524 and MKN 1048) have a good signal-to-noise. For them a detailed study is possible, and the main spectral parameters can be determined. The remaining 7 sources (3 classified as starbursts and 4 as AGNs) have detections at a few sigmas and therefore we can only estimate a flux and sometimes a photon index. The 5 sources observed by BeppoSAX are all very faint in the X-rays, and for 4 of them the X-ray observation does not provide any firm indication of the presence of an AGN. For the fifth object, MKN 266, the detection of a cold iron line and the flatness of the spectrum suggest the presence of a heavily obscured AGN. All the errors quoted are at the 90% level of confidence.
a) Sources from the ASCA public archive: We have reduced and analyzed the ASCA GIS observations of 11 sources. All the data have been retrieved from the ASCA public archive. IRAS 03158+4227: The GIS data of IRAS 03158+4227 give a detection at a low signalโtoโnoise level ($`6\sigma `$). The spectrum is well fitted by a powerlaw with $`\mathrm{\Gamma }2.8`$. The 2-10 keV flux is F=5.4$`\pm 1\times 10^{13}`$ erg cm<sup>-2</sup> s<sup>-1</sup>. There is no evidence of the iron line, and the steepness of the spectrum suggests a thermal origin of the observed emission. IRAS 03158+4227 is classified optically as Seyfert 2, therefore it should host an AGN. From the low X-ray luminosity, relative to the infrared emission, and from the absence of AGN indicators in the X rays we deduce that the AGN is obscured by a column density $`\mathrm{N}_\mathrm{H}>10^{24}`$ cm<sup>-2</sup>. IRAS 07598+6508: The source IRAS 07598+6508 has not been detected in the GIS observation with an exposure time of 83500 sec. We can only derive an upper limit to the 2-10 keV flux, F<sup>MAX</sup>=8 10<sup>-14</sup> erg cm<sup>-2</sup> s<sup>-1</sup> (at a 90% confidence level). IRAS 17208-0014: The source IRAS 17208-0014 is marginally detected in the GIS. The flux we derive is F=3$`\pm 1\times 10^{13}`$ erg cm<sup>-2</sup> s<sup>-1</sup>. TOL 1351-375: TOL 1351-375 is a Seyfert 1.9 galaxy observed by ASCA in 1997. The spectrum obtained after the standard data reduction has a good signal-to-noise and has been fitted with a multi-component model consisting of an absorbed powerlaw, an iron K<sub>ฮฑ</sub> line and a thermal component with kT$``$1.2 keV. The best fit gives the following values: photon index $`\mathrm{\Gamma }=1.94\pm 0.13`$; absorbing column density $`\mathrm{N}_\mathrm{H}=1.6\pm 3\times 10^{22}`$ cm<sup>-2</sup>; rest frame line energy E=6.38$`{}_{0.3}{}^{}{}_{}{}^{+0.2}`$ keV; line equivalent width EW=175$`\pm `$120 eV. The measured 2-10 keV flux is 3.8 $`10^{12}`$ erg cm<sup>-2</sup> s<sup>-1</sup>. IRAS 18293-3413: This source has been detected at a 5$`\sigma `$ level. There is no evidence of an AGN contribution. Fitting the data with a blackbody we obtain kT=0.85 keV. The 2-10 keV flux is 9 $`10^{14}`$ erg cm<sup>-2</sup> s<sup>-1</sup>. NGC 1614: The spectrum of NGC 1614 is well fitted by a powerlaw with $`\mathrm{\Gamma }=1.55\pm 0.4`$, while a termal component is rejected. There is no evidence of an iron line. The measured 2-10 keV flux is 5.6 $`10^{13}`$ erg cm<sup>-2</sup> s<sup>-1</sup>. IRAS 05189-2524: This source was observed by ASCA for $``$ 40000 seconds. The spectrum is well fitted by a two-component model, consisting of a powerlaw with photon index $`\mathrm{\Gamma }=1.89_{0.34}^{+0.35}`$ absorbed by a column density of $`\mathrm{N}_\mathrm{H}=4.7_{1.1}^{+1.4}\times 10^{22}`$ cm<sup>-2</sup>, and a thermal component with kT=0.88$`{}_{0.35}{}^{}{}_{}{}^{+0.89}`$. The iron line is not detected. From the non-detection we estimate an upper limit to the equivalent width of 235 eV. The measured 2-10 keV flux is 5.3$`\times 10^{12}`$ erg cm<sup>-2</sup> s<sup>-1</sup>. In summary, the E$`>2`$ keV spectrum is typical of a Comptonโthin type 2 AGN, while at lower energies an extra thermal component is present, that can be either associated to the AGN itself or due to starburst activity. NGC 7212: The Seyfert 2 galaxy NGC 7212 has a flat 0.5-10 keV spectrum, well fitted by a powerlaw with photon index $`\mathrm{\Gamma }=0.75_{0.55}^{+0.50}`$. An emission line with peak energy E=6.06$`{}_{0.40}{}^{}{}_{}{}^{+0.35}`$ (compatible with a K<sub>ฮฑ</sub> iron line at a confidence level of 90%) is detected at a 2$`\sigma `$ level, with a best fit equivalent width EW=1.0$`{}_{0.9}{}^{}{}_{}{}^{+1.3}`$ keV. From these spectral features we deduce that the source is Compton-thick. The measured 2-10 keV flux is F=9.8$`\times 10^{13}`$ erg cm<sup>-2</sup> s<sup>-1</sup>. MKN 1048: The GIS spectrum of MKN 1048 is well fitted by a simple powerlaw with photon index $`\mathrm{\Gamma }=1.60_{0.01}^{+0.02}`$, a value significantly lower than the average photon index found in the X-ray spectra of quasars. The K<sub>ฮฑ</sub> iron line is marginally detected, at a 2$`\sigma `$ level (best fit equivalent width EW=96$`{}_{}{}^{+86}{}_{83}{}^{}`$ eV. The measured 2-10 keV flux is F=9.9$`\times 10^{12}`$ erg cm<sup>-2</sup> s<sup>-1</sup>. WAS 49b The Seyfert 2 galaxy WAS 49b presents a very unusual X-ray spectrum. A good analytical fit is provided by a single powerlaw with photon index $`\mathrm{\Gamma }=0.5`$, a very low value typical of reflection-dominated spectra. In contrast with this interpretation, the iron line is only marginally detected with an equivalent width EW $`250\pm 200`$ eV and the 2-10 keV flux is relatively high with respect to the infrared emission (F$`{}_{210}{}^{}=1.2\times 10^{12}`$ erg cm<sup>-2</sup> s<sup>-1</sup>), thus indicating that the source is Compton - thin. A similar good fit is obtained by a model composed by a powerlaw with $`\mathrm{\Gamma }=1.7`$ absorbed by a column density $`\mathrm{N}_\mathrm{H}10^{23}`$ cm<sup>-2</sup> and a flat ($`\mathrm{\Gamma }=0.6`$) powerlaw with a 1 keV normalization equal to 1/6 that of the first component. Since WAS 49b is part of a triple sistem, we suggest that this flat contribution is due to the diffuse emission of the hot intra-cluster gas.
b) Sources observed with BeppoSAX: MKN 273: MKN 273 was observed with BeppoSAX in 1998 and was detected by the MECS (1.65-10 keV) at a 2-10 keV flux of 3.5$`\times 10^{13}`$ erg cm<sup>-2</sup> s<sup>-1</sup>, significantly lower than the flux measured by ASCA 2 years before (F$`{}_{210}{}^{}=7\times 10^{13}`$ erg cm<sup>-2</sup> s<sup>-1</sup>, Iwasawa 1998). Fitting the data with a powerlaw and an iron line we obtain a spectral index ($`\mathrm{\Gamma }=1.1_{0.6}^{+0.9}`$, poorly constrained due to the weakness of the continuum, and a marginal detection of the iron line, at a 3$`\sigma `$ level (best fit equivalent width EW=1.2$`{}_{}{}^{+2}{}_{1}{}^{}`$ keV). The detection in the PDS is also weak. We conclude that either the AGN at the centre of this source is absorbed by a column density $`\mathrm{N}_\mathrm{H}>10^{25}`$ cm<sup>-2</sup> or, alternatively, that the central source has faded and only the reflected component is now visible (this would also explain the flux decrease between the ASCA and SAX observations). MKN 266: MKN 266 was detected by the MECS and marginally (at a 2$`\sigma `$ level) by the PDS. The 2-10 keV spectrum is remarkably flat (photon index $`\mathrm{\Gamma }=0.7_{0.3}^{+0.4}`$) suggesting the presence of a heavily obscured AGN. The iron line is detected at a 90% confidence level, with a best fit equivalent width of 575 eV. The absence of any excess in the PDS is an indication of a column density higher than 10<sup>25</sup> cm<sup>-2</sup>. The measured 2-10 keV flux is 5.6 $`10^{13}`$ erg cm<sup>-2</sup> s<sup>-1</sup>. NGC 2623: NGC 2623 was detected by the MECS at 3.5$`\sigma `$, thus we can only estimate the 2-10 keV flux, F$`{}_{210}{}^{}8\times 10^{14}`$ erg cm<sup>-2</sup> s<sup>-1</sup>. The detection in the PDS is weak and cannot be associated with certainty to NGC 2623, because of the presence in the PDS field ($`100`$) of another source that could emit in the hard X rays. MKN 1073: MKN 1073 was not detected by the MECS and therefore we can only estimate an upper limit to the 2-10 keV flux, F$`{}_{210}{}^{}>710^{14}`$ erg cm<sup>-2</sup> s<sup>-1</sup>. The signal in the PDS is very high, with an exceptionally steep spectrum ($`\mathrm{\Gamma }3`$), but it is probably emitted by an unidentified source in a nearby cluster, that has a strong emission in the 2-10 keV band as well. IRAS 00198-7926 The source IRAS 00198-7926 was observed by the BeppoSAX instruments for $``$ 20000 seconds and was not detected. From this non-detection we estimate an upper limit to the 2-10 keV flux, F$`{}_{210}{}^{}>1\times 10^{13}`$ erg cm<sup>-2</sup> s<sup>-1</sup>. |
warning/0002/astro-ph0002416.html | ar5iv | text | # A Lens Mapping Algorithm for Weak Lensing
## 1 Introduction
The mapping of the distribution of matter in extended regions around rich clusters of galaxies, from the systematic distortions of background galaxies due to gravitational lensing, has increasingly become feasible and popular over the last decade. This *shear* is measured from the quadrupole moments of the images of the background galaxies in terms of the local โpolarizationโ of an image compared to its assumed intrinsic circular form. In order to remove the effect of the source ellipticity distribution, the parameters are averaged over a large number of source galaxies. The two-dimensional mass distribution of the lens is reconstructed from the shear map away from the critical lines (โweak reconstructionโ). Following the pioneering work of Tyson et al. (1990), the idea was quantitatively developed by Kaiser and Squires (1993, also Kaiser 1995 and Squires and Kaiser 1996) and Seitz and Schneider (1995, 1996). Since then, several interesting variants on this theme have appeared in the literature, e.g. using maximum likelihood cluster reconstruction (Bartelmann et al. 1996, Seitz et al. 1998), or using methods based on the variational principle (Lombardi and Bertin 1999) or maximum entropy (Bridle et al. 1998).
Since the shear data contain complete information about the mass distribution in two independent fields $`\gamma _1/(1-\kappa )`$ and $`\gamma _2/(1-\kappa )`$, it is possible to obtain several algorithms to estimate the mass distribution from the measured shear. The observed shear is available only on a finite grid. Due to the unknown intrinsic distribution of the ellipticities of the source galaxies, and the effect of the distortion of the PSF due to observing conditions (seeing, tracking etc.), the measured shear is noisy. It is therefore desirable to develop new algorithms in the hope that they might be able to deal with the noise better than other methods.
In this paper we develop an algorithm (LM: Lens Mapping algorithm) for the lens mass reconstruction from the measured reduced shear, $`g_i=\gamma _i/(1-\kappa )`$. The method involves two steps. First, the lens mapping is reconstructed from the reduced shear, which can be done uniquely with the assumption that the lens mapping goes to identity far away from the lens. As the second step, we show in ยง3 that for a sub-critical lens, the surface mass density can be reconstructed completely from the derived lens mapping. In reality, however, the measured shear is available only in a finite region of the lens plane. In this case the lens mapping cannot be uniquely obtained and hence the LM algorithm exhibits a mass-sheet degeneracy, which is also present in other methods of mass reconstruction. In ยง3.1, we characterize this degeneracy in the context of the LM algorithm. The performance of this method with discretely sampled, noisy data is dealt with in ยง4, where we demonstrate the various features of the LM algorithm, by reconstructing the mass distribution for an analytically given shear field, and separately showing the effects of discrete sampling and of uniform noise (due to measurement error) on the input data.
## 2 Weak Lensing and the Mass-Sheet Degeneracy
For most practical applications of gravitational lensing the lens can be considered to be thin. Under the small angle approximation the lens equation is given by
$$๐ฒ=๐ฑ-\mathbf{}\psi (๐ฑ),$$
(1)
where the source angular position is denoted as $`๐ฒ`$, the image angular position as $`๐ฑ`$ and the relativistic potential $`\psi `$ satisfies the equation $`\kappa (๐ฑ)\mathrm{\Sigma }/\mathrm{\Sigma }_{\mathrm{crit}}=\frac{1}{2}\mathbf{}^2\psi `$, where $`\mathrm{\Sigma }_{\mathrm{crit}}=(c^2/4\pi G)(d_\mathrm{s}/d_\mathrm{l}d_{\mathrm{ls}})`$. Here $`d_\mathrm{s},d_\mathrm{l}`$ and $`d_{\mathrm{ls}}`$ are respectively the angular diameter distance to the source, to the lens and between the lens and the source. The coordinates for source and image are small compared to unity and we can treat their components as Cartesian. It is convenient to define two planes: that containing the source is called the source plane and that containing the cluster, the lens plane (or image plane).
The coordinate differentials in the source plane and the corresponding differentials in the lens plane are related through $`dy_i=M_{ij}^1dx_j`$ and the inverse of the magnification matrix $`M`$ is given by
$$M^1=(1-\kappa )\left(\begin{array}{cc}1-g_1& g_2\\ g_2& 1+g_1\end{array}\right),$$
(2)
where $`๐ g_1+ig_2=(\gamma _1+i\gamma _2)/(1-\kappa )`$, $`\gamma _1=\frac{1}{2}(\psi _{,11}-\psi _{,22})`$ and $`\gamma _2=\psi _{,12}`$ and the subscripts denote differentiation with respect to the two components of the image coordinates. From the quadrupole moments of the surface brightness of the images of the background galaxies the reduced shear can be measured unambiguously in the regions where $`\kappa <1`$ (for observational details see Kaiser 1999 or Bartelmann & Schneider 2000).
The continuity of $`๐ฒ(x_1,x_2)`$ implies $`\mathrm{}_{jk}^2y_i=\mathrm{}_{kj}^2y_i`$, which, along with, eq. (2) gives
$$\mathrm{K}_{,l}\mathrm{G}_{ij}^1-\mathrm{K}_{,j}\mathrm{G}_{il,j}^1=\mathrm{G}_{il,j}^1-\mathrm{G}_{ij,l}^1,$$
(3)
where $`\mathrm{K}\mathrm{ln}(1-\kappa )`$ and $`G^1M^1/(1-\kappa )`$. On multiplying (3) by the inverse of $`\mathrm{G}^1`$, and taking the trace of the resulting equation gives us
$$_l\mathrm{ln}(1-\kappa )=\underset{i=1}{\overset{2}{}}\underset{j=1}{\overset{2}{}}\left[\mathrm{G}_{ij}(\mathrm{G}_{il,j}^1-\mathrm{G}_{ij,l}^1)\right].$$
(4)
This equation was first derived by Kaiser (1995). It is clear that replacing $`1-\kappa `$ on the left hand side with $`\lambda (1-\kappa )`$, where $`\lambda `$ is a constant, does not affect the equation. Therefore any particular solution of this equation can be used to obtain a one-parameter degenerate family of functions, all of which satisfy (4). This is known as the mass-sheet degeneracy.
## 3 The Lens Mapping Algorithm
In general the lens equation (1) is a many-to-one mapping, and consequently the inverse of the lens equation $`๐ฑ=๐ฑ(๐ฒ)`$ has several branches, with no single branch completely specifying the lens mapping. However, if the lens is sub-critical, then the lens mapping is one-to-one everywhere. In this case both the mapping and its inverse are uniquely defined and completely specify the mass distribution. Here we show that in this case it is possible to obtain the lens mapping completely from the reduced shear, with the additional assumption that it goes to identity sufficiently far away from the lens.
For a non-critical lens, open curves in the source plane are mapped to open curves in the lens plane. In particular, any infinite straight line is mapped to an infinite open curve in the lens plane. In Fig. 1a, we schematically show the images of the two perpendicular straight lines passing through the source at $`(y_1,y_2)`$. The point of intersection of the two curves gives the position of the image ($`๐_1`$, $`๐_2`$).
From (2) we can obtain the equation which the image of an arbitrary infinite straight line in the source plane satisfies in the lens plane. In particular, we obtain the equations which map the coordinate grid lines of a Cartesian coordinate system in the source plane to the lens plane
$`dy_1=0{\displaystyle \frac{dx_1}{dx_2}}={\displaystyle \frac{g_2}{1-g_1}};dy_2=0{\displaystyle \frac{dx_2}{dx_1}}={\displaystyle \frac{g_2}{1+g_1}}.`$ (5)
These are first-order ordinary differential equations and can be uniquely integrated through any point in the lens plane. If we consider the source position as the intersection of the lines $`y_1=\mathrm{constant}`$ and $`y_2=\mathrm{constant}`$ in the source plane, then the image of this point will be the intersection of the images of these two curves in the lens plane. In this manner we obtain a one-to-one mapping from the lens to the source plane. Away from the lens singularities we can integrate the equations (5) to obtain the numerical solutions
$`x_1=X(x_2;๐);x_2=Y(x_1;๐)`$ (6)
where $`X(x_2;๐)`$ is the mapped curve for $`dy_1=0`$, passing through the point $`๐`$ in the lens plane, and $`Y(x_1;๐)`$ is the mapped curve for $`dy_2=0`$, passing through the point $`๐`$ in the lens plane (see Fig. 1a). The auxiliary variable $`๐`$ in the above equations also explicitly represents the family of curves which the equations (5) generate when integrated with the initial conditions $`X(\chi _2)=\chi _1`$ and $`Y(\chi _1)=\chi _2`$, respectively, for the two equations.
In Fig. 1b we show the image of a Cartesian grid that is equally spaced in the source plane, lensed by a two-component model, where both clusters are represented by Pseudo Isothermal Circular Mass Distributions (PICMD) with the same core radius $`r_c`$, separated by $`10r_c`$ along the $`x_1`$ axis. The central values of the dimensionless surface density for the two components are $`\kappa _0=0.8`$ and $`\kappa _0=0.5`$ respectively. Fig. 1b illustrates the integral curves of (5).
To obtain the source position corresponding to the image position $`๐`$ we use the assumption that far away from the lens the lens mapping is $`y_ix_i`$ (see Fig. 1a). Therefore the integrated curves should go to the original unperturbed source grid lines at large $`|๐ฑ|`$. This is ensured by the following asymptotic conditions
$`\underset{x_2\pm \mathrm{}}{lim}X(x_2;๐)=y_1;\underset{x_1\pm \mathrm{}}{lim}Y(x_1;๐)=y_2.`$ (7)
The lens equation can now be formally written as
$`y_1=X(\mathrm{};๐ฑ);y_2=Y(\mathrm{};๐ฑ),`$ (8)
where the auxiliary variables $`๐`$ are now interpreted as the image coordinates $`๐ฑ`$. The trace of the magnification matrix which is given in terms of $`\kappa `$ as
$$\frac{\mathrm{}X(\mathrm{};๐ฑ)}{\mathrm{}x_1}+\frac{\mathrm{}Y(\mathrm{};๐ฑ)}{\mathrm{}x_2}=2\left[1-\kappa (๐ฑ)\right],$$
(9)
provides an estimate of $`\kappa `$ at any point. Fig. 3 shows the lens mass distribution thus reconstructed from the shear represented by Fig. 1b, assuming that the shear is noise-free and is known everywhere.
It might appear surprising that we have been able to obtain the lens mapping solely from shear. To understand how, let us recall that our assumption that the lens mapping becomes identity far away from the lens removes the mass-sheet degeneracy We also note that it is essentially equivalent to the assumption in Kaiserโs method that the mass density vanishes *f*ar away from the center, though quantitatively the notion of โfarโ is different in the two methods. For a finite field we will see that the degeneracy reappears.
The above discussion applies to points away from the critical lines of the lens mapping. At these singularities, the eqs. (5) become non-integrable. This becomes apparent by writing the eqs. (5) in terms of $`\gamma _i`$s and $`\kappa `$,
$`{\displaystyle \frac{dx_1}{dx_2}}={\displaystyle \frac{\gamma _2}{1-\kappa -\gamma _1}};{\displaystyle \frac{dx_2}{dx_1}}={\displaystyle \frac{\gamma _2}{1-\kappa +\gamma _1}}.`$ (10)
For a sufficiently smooth mass distribution, the potential and consequently $`\gamma _2`$ is non-singular everywhere. Since the Jacobian of the lens mapping, $`(1-\kappa )^2-\gamma ^2`$, is positive far away from the lens, we conclude that $`(1-\kappa )^2-\gamma _1^2>\gamma _2^2`$. From this inequality it is clear that $`(1-\kappa )^2-\gamma _1^2`$ can become zero only when the Jacobian $`(1-\kappa )^2-\gamma ^2`$ also becomes zero. Since $`(1-\kappa )^2-\gamma _1^2=0`$ implies that either $`1-\kappa -\gamma _1=0`$ or $`1-\kappa +\gamma _1=0`$, we conclude that both the equations (5) cannot be integrated simultaneously, which implies that a correspondence between the lens plane and source plane cannot be obtained. Thus the LM algorithm cannot be applied in regions very close to the critical lines of the Lens mapping.
### 3.1 Finite Field, Critical Lens
The algorithm as described in the previous section is applicable to cases where the reduced shear is known everywhere. In practice, this information is available only in a limited region of the lens plane. We show here that the LM algorithm can be generalized to this case. For a finite field, eqs. (7) can be applied only up to the edge of region in which data is available, and the correct limit cannot be evaluated. It is easy to see that due to the unknown magnification factor at the edge, these equations provide an incorrect measure of the derivatives in eq. (9).
Let us consider the differentials of $`X`$ and $`Y`$ (as defined in eq. 6) at the edge $`๐ฑ^\mathrm{e}`$,
$`\mathrm{\Delta }X(x_2^\mathrm{e},๐)={\displaystyle \frac{\mathrm{\Delta }X}{\mathrm{\Delta }y_1}}\mathrm{\Delta }y_1;\mathrm{\Delta }Y(x_1^\mathrm{e},๐)={\displaystyle \frac{\mathrm{\Delta }Y}{\mathrm{\Delta }y_2}}\mathrm{\Delta }y_2.`$ (11)
Since at the edge of the image the derivatives of $`X`$ and $`Y`$ are $`>1`$ we obtain slightly higher values of $`\mathrm{\Delta }y_i`$. Substituting the value of the derivatives in eq. (11) we obtain
$`\mathrm{\Delta }X(x_2^\mathrm{e},๐)`$ $`=`$ $`[1-\kappa (๐ฑ^\mathrm{e})][1-g_1(๐ฑ^\mathrm{e})]\mathrm{\Delta }y_1,`$ (12)
$`\mathrm{\Delta }Y(x_1^\mathrm{e},๐)`$ $`=`$ $`[1-\kappa (๐ฑ^\mathrm{e})][1+g_1(๐ฑ^\mathrm{e})]\mathrm{\Delta }y_2.`$ (13)
This gives us a way of obtaining the exact value of the derivatives of the lens mapping, and thus completely removing the mass-sheet degeneracy. However, in reality we only know the reduced shear, and consequently the correction factors at the edges, only up to a factor of $`1-\kappa `$. For finite fields we can impose the boundary condition $`\kappa =0`$ to obtain the mass distribution. This corrects the generalized degeneracy up to a factor of $`(1-\kappa )`$.
This algorithm can be applied to a critical lens as well. Since the algorithm depends on our ability to integrate equations (5) without hitting a singularity, we see that the method can be used for the region which lie outside the critical lines. In Fig. 3, where the model has $`\kappa _0=2`$, we see that our algorithm works well outside the critical region.
## 4 Noisy and Discretely Sampled Data
We now present a few examples to illustrate the application of the LM algorithm to noisy data. Since we wish to compare the original mass distribution with that reconstructed by the algorithm, we correct for the mass-sheet degeneracy in all the examples by the prescription given in ยง 3.1. We consider a single component lens modeled as PICMD. We simulate the signal by analytically calculating the reduced shear from a model. We illustrate the effect of finite sampling of data and the presence of noise with separate examples.
We first consider the case where the reduced shear is known at all the points from an analytical model. We choose a lens with the central, dimensionless surface mass density given by $`\kappa _0=0.7`$, and reconstruct the mass distribution by applying the LM algorithm on a $`20\times 20`$ grid. The noise in the reconstruction, which is very small everywhere ($`\mathrm{}<1\%`$), can be entirely accounted for as arising from the inaccuracies in the evaluation of various integrals and derivatives required by the algorithm.
To illustrate the effect of finite sampling of shear in the image plane, we consider a model with a lens velocity dispersion of $`\sigma _v=1100`$ km/s and with a core radius $`R_c=50`$ kpc. The source galaxies are assumed to be at the redshift $`z_s=1`$ and the lens is at $`z_l=0.2`$. We sample the reduced shear in pixels of size $`20^{\prime \prime }\times 20^{\prime \prime }`$. Using a bicubic spline interpolation routine to evaluate the reduced shear at all the points, the mass distribution is reconstructed using the LM algorithm on a $`20\times 20`$ grid as before. The resulting percentage error in the reconstructed mass distribution (given by $`100\times (\kappa _{\mathrm{rec}}-\kappa _{\mathrm{true}})/\kappa _{\mathrm{true}}`$) is shown in Fig. 5. We note that the excess noise near the center of the distribution is mainly due to the fact that the reduced shear is most rapidly varying close to the center and therefore the interpolated shear values are noisier at those points. This noise will vanish if the grid on which the shear is evaluated were made finer, but then in the real world, one is limited by the number of galaxies over which the value of the reduced shear has to be averaged, and we cannot use a substantially finer grid without increasing shot noise.
As our last example we consider the same case as the last one, but add a uniform noise with a constant signal-to-noise ratio of $`10`$ to the given shear field, before the application of interpolation. The noise in the reconstructed mass distribution (Fig. 5) can be seen to be around twice that of the noise in the shear field.
## 5 Conclusions
We have described a simple algorithm for mass reconstruction of a cluster lens by directly obtaining the lens mapping from the measured reduced shear and taking its derivatives to compute the surface mass density $`\kappa `$.
The method works best for sub-critical lenses, where it can give the mass distribution at all points, but it can work for critical lenses as well in limited regions of the lens plane, away from the critical lines. The algorithm is shown to have a mass-sheet degeneracy of the same type as exists in the other methods of reconstruction if the field in which the reduced shear is available is finite. We have tested the algorithm on discretely sampled noisy (simulated) data and have found that it reproduces the mass distribution within acceptable limits.
TDS thanks the University Grants Commission (India) and IUCAA for providing support for this work. We thank Yuri A. Shchekinov and Prof. J. Ehlers for useful discussions, and an anonymous referee for helpful suggestions. |
warning/0002/hep-ph0002197.html | ar5iv | text | # 1 Introduction
## 1 Introduction
An outstanding problem in astrophysics concerns the origin and nature of magnetic fields in the galaxies and in the clusters of galaxies . The uniformity of the magnetic fields strength in the several galaxies and the recent observation of magnetic fields with the same intensity in high red-shift protogalactic clouds suggest that galactic and intergalactic magnetic fields may have a primordial origin. Hopefully, a confirmation to this intriguing hypothesis will come from the forthcoming balloon and satellite missions looking at the anisotropies of the Cosmic Microwave Background Radiation. In fact, among other very important cosmological parameters, the observations performed by these surveyors may be able to detect the imprint of primordial magnetic fields on the temperature and polarizations acoustic peaks .
Besides observations, a considerable amount of theoretical work, based on the particle physics standard model as well as on its extensions, has been done which support the hypothesis that the production of magnetic fields may actually be occurred during the big-bang .
Quite independently from the astrophysical considerations, several authors paid some effort to investigate the possible effects that cosmic magnetic fields may have for some relevant physical processes which occurred in the early Universe like the big-bang nucleosynthesis and the electroweak baryogenesis (EWB). The latter is the main subject of this contribution.
Since, before the electroweak phase transition (EWPT) to fix the unitary gauge is a meaningless operation, the electromagnetic field is undefined above the weak scale and we can only speaks in terms of hyperelectric and hypermagnetic fields. The importance of a possible primordial hypercharge magnetic fields for the electroweak baryogenesis scenario is three-fold. In fact, we will show that hypercharge magnetic fields can affect the dynamics of the EWPT, they change the rate of the baryon number violating anomalous processes and, finally, hypermagnetic fields may themselves carry baryon number if they possess a non trivial topology. These effects will be shortly reviewed respectively in the section 2,3, and 4 of this contribution.
## 2 The effect of a strong hypermagnetic fields on the EWPT
As it is well known, the properties of the EWPT are determined by the Higgs field effective potential. In the framework of the minimal standard model (MSM), by accounting for radiative corrections from all the known particles and for finite temperature effects, one obtains that
$$V_{\mathrm{eff}}(\varphi ,T)\frac{1}{2}(\mu ^2\alpha T^2)\varphi ^2T\delta \varphi ^3+\frac{1}{4}(\lambda \delta \lambda _T)\varphi ^4.$$
(1)
where $`\varphi `$ is the radial component of the Higgs field and $`T`$ is the temperature (for the definitions of the coefficients see e.g. Ref.).
A strong hypermagnetic field can produce corrections to the effective potential as it affects the charge particles propagators. It was shown, however, that these kind of corrections are not the most relevant effect that strong hypermagnetic fields may produce on the EWPT. In fact, it was recently showed by Giovannini and Shaposhnikov and by Elmfors, Enqvist and Kainulainen that hypermagnetic fields can affect directly the Gibbs free energy (in practice the pressure) difference between the broken and the unbroken phase, hence the strength of the transition. The effect can be understood in analogy with the Meissner effect, i.e. the expulsion of the magnetic field from superconductors as consequence of photon getting an effective mass inside the specimen. In our case, it is the $`Z`$โcomponent of the hypercharge $`U(1)_Y`$ magnetic field which is expelled from the broken phase just because $`Z`$โbosons are massive in that phase. Such a process has a cost in terms of free energy. Since in the broken phase the hypercharge field decompose into
$$A_\mu ^Y=\mathrm{cos}\theta _wA_\mu \mathrm{sin}\theta _wZ_\mu ,$$
(2)
we see that the Gibbs free energy in the broken and unbroken phases are
$`G_b`$ $`=`$ $`V(\varphi ){\displaystyle \frac{1}{2}}\mathrm{cos}^2\theta _w(B_Y^{ext})^2,`$ (3)
$`G_u`$ $`=`$ $`V(0){\displaystyle \frac{1}{2}}(B_Y^{ext})^2.`$ (4)
where $`B_Y^{ext}`$ is the external hypermagnetic field. In other words, compared to the case in which no magnetic field is present, the energy barrier between unbroken and broken phase, hence the strength of the transition, is enhanced by the quantity $`{\displaystyle \frac{1}{2}}\mathrm{sin}^2\theta _w(B_Y^{ext})^2`$. According to the authors of refs. this effect can have important consequence for the electroweak baryogenesis scenario.
In any scenario of baryogenesis it is crucial to know at which epoch do the sphaleronic transitions, which violate the sum ($`B+L`$) of the baryon and lepton numbers, fall out of thermal equilibrium. Generally this happens at temperatures below $`\overline{T}`$ such that
$$\frac{E(\overline{T})}{\overline{T}}A,$$
(5)
where $`E(T)`$ is the sphaleron energy at the temperature $`T`$ and $`A3545`$, depending on the poorly known prefactor of the sphaleron rate. In the case of baryogenesis at the electroweak scale one requires the sphalerons to drop out of thermal equilibrium soon after the electroweak phase transition. It follows that the requirement $`\overline{T}=T_c`$, where $`T_c`$ is the critical temperature, turns eq. (5) into a lower bound on the higgs vacuum expectation value (VEV),
v(Tc)Tc
>
1.
>
๐ฃsubscript๐๐subscript๐๐1\frac{v(T_{c})}{T_{c}}\mathrel{\vbox{\hbox{$>$}\nointerlineskip\hbox{$\sim$}}}1\,. (6)
As a result of intense research in the recent years , it is by now agreed that the standard model (SM) does not have a phase transition strong enough as to fulfill eq. (6), whereas there is still some room left in the parameter space of the minimal supersymmetric standard model (MSSM).
The interesting observation made in Refs. is that a magnetic field for the hypercharge $`U(1)_Y`$ present for $`T>T_c`$ may help to fulfill Eq.(6). In fact, it follows from the Eqs.(3), that in presence of the magnetic field the critical temperature is defined by the expression
$$V(0,T_c)V(\varphi ,T_c)=\frac{1}{2}\mathrm{sin}^2\theta _w(B_Y^{ext}(T_c))^2.$$
(7)
This expression implies a smaller value of $`T_c`$ than that it would take in the absence of the magnetic field, hence a larger value of the ratio (6). On the basis of the previous considerations and several numerical computations, the authors of Refs. concluded that for some reasonable values of the magnetic field strength the EW baryogenesis can be revived even in the standard model. In the next section we shall reconsider critically this conclusion.
## 3 The sphaleron in a magnetic field
The sphaleron , is a static and unstable solution of the field equations of the electroweak model, corresponding to the top of the energy barrier between two topologically distinct vacua. In the limit of vanishing Weinberg angle, $`\theta _w0`$, the sphaleron is a spherically symmetric, hedgehog-like configuration of $`SU(2)`$ gauge and Higgs fields. No magnetic moment is present in this case. As $`\theta _w`$ is turned on the $`U_Y(1)`$ field is excited and the spherical symmetry is reduced to an axial symmetry . A very good approximation to the exact solution is obtained using the Ansatz by Klinkhamer and Laterveer , which requires four scalar functions of $`r`$ only,
$`g^{}a_idx^i=(1f_0(\xi ))F_3,`$
$`gW_i^a\sigma ^adx^i=(1f(\xi ))(F_1\sigma ^1+F_2\sigma ^2)+(1f_3(\xi ))F_3\sigma ^3,`$
$`๐ฝ={\displaystyle \frac{v}{\sqrt{2}}}\left(\begin{array}{c}0\\ h(\xi )\end{array}\right),`$ (10)
where $`g`$ and $`g^{}`$ are the $`SU(2)_L`$ and $`U(1)_Y`$ gauge couplings, $`v`$ is the higgs VEV such that $`M_W=gv/2`$, $`M_h=\sqrt{2\lambda }v`$, $`\xi =gvr`$, $`\sigma ^a`$ ($`a=1,2,3`$) are the Pauli matrices, and the $`F_a`$โs are 1-forms defined in Ref. . The boundary conditions for the four scalar functions are
$`f(\xi ),f_3(\xi ),h(\xi )0f_0(\xi )1`$ $`\mathrm{f}or\xi 0`$
$`f(\xi ),f_3(\xi ),h(\xi ),f_0(\xi )1`$ $`\mathrm{f}or\xi \mathrm{}.`$ (11)
It is known that for $`\theta _w0`$ the sphaleron has some interesting electromagnetic properties. In fact, differently from the pure $`SU(2)`$ case, in the physical case a nonvanishing hypercharge current $`J_i`$ comes-in. At the first order in $`\theta _w`$, $`J_i`$ takes the form
$$J_i^{(1)}=\frac{1}{2}g^{}v^2\frac{h^2(\xi )[1f(\xi )]}{r^2}ฯต_{3ij}x_j,$$
(12)
where $`h`$ and $`f`$ are the solutions in the $`\theta _w0`$ limit, giving for the dipole moment
$$\mu ^{(1)}=\frac{2\pi }{3}\frac{g^{}}{g^3v}_0^{\mathrm{}}๐\xi \xi ^2h^2(\xi )[1f(\xi )].$$
(13)
The reader should note that the dipole moment is a true electromagnetic one because in the broken phase only the electromagnetic component of the hypercharge field survives at long distances.
Following Ref. we will now consider what happens to the sphaleron when an external hypercharge field, $`A_i^Y`$, is turned on. Not surprisingly, the energy functional is modified as
$$E=E_0E_{\mathrm{d}ip},$$
(14)
with
$$E_0=d^3x\left[\frac{1}{4}F_{ij}^aF_{ij}^a+\frac{1}{4}f_{ij}f_{ij}+(D_i๐ฝ)^{}(D_i๐ฝ)+V(๐ฝ)\right]$$
(15)
and
$$E_{\mathrm{d}ip}=d^3xJ_iA_i^Y=\frac{1}{2}d^3xf_{ij}f_{ij}^c$$
(16)
with $`f_{ij}_iA_j^Y_jA_i^Y`$. We will consider here a constant external hypermagnetic field $`B_Y^c`$ directed along the $`x_3`$ axis. In the $`\theta _w0`$ limit the sphaleron has no hypercharge contribution and then $`E_{\mathrm{d}ip}^{(0)}=0`$. At $`O(\theta _w)`$, using (12) and (13) we get a simple dipole interaction
$$E_{\mathrm{d}ip}^{(1)}=\mu ^{(1)}B_Y^c.$$
(17)
In order to assess the range of validity of the approximation (17) one needs to go beyond the leading order in $`\theta _w`$ and look for a nonlinear $`B_Y^c`$โdependence of $`E`$. This requires to solve the full set of equations of motion for the gauge fields and the Higgs in the presence of the external magnetic field. Fortunately, a uniform $`B_Y^c`$ does not spoil the axial symmetry of the problem. Furthermore, the equation of motion are left unchanged ($`_if_{ij}^c=0`$) with respect to the free field case. The only modification induced by $`B_Y^c`$ resides in the boundary conditions since โ as $`\xi \mathrm{}`$ โ we now have
$$f(\xi ),h(\xi )1,f_3(\xi ),f_0(\xi )1B_Y^c\mathrm{sin}2\theta _w\frac{\xi ^2}{8gv^2}$$
(18)
wheras the boundary condition for $`\xi 0`$ are left unchanged.
The solution of the sphaleron equation of motions with the boundary conditions in the above were determined numerically by the authors of Ref.. They showed that in the considered $`B_Y^c`$โrange the corrections to the linear approximation
$$\mathrm{\Delta }E\mu ^{(1)}\mathrm{cos}\theta _WB_Y^c$$
are less than $`5\%`$. For larger values of $`B_Y^c`$ non-linear effects increase sharply. However, for such large magnetic fields the broken phase of the SM is believed to become unstable to the formation either of $`W`$-condensates or of a mixed phase . In such situations the sphaleron solution does not exist any more. Therefore, we will limit our considerations to safe values BYc
<
0.4T2
<
subscriptsuperscript๐ต๐๐0.4superscript๐2B^{c}_{Y}\mathrel{\vbox{\hbox{$<$}\nointerlineskip\hbox{$\sim$}}}0.4~{}T^{2}.
From the previous considerations it follows that the conclusion that the sphaleron freeze-out condition (5) is satisfied and the baryon asymmetry preserved was premature. Indeed, in an external magnetic field the relation between the VEV and the sphaleron energy is altered and Eq. (6) does not imply (5) any more. We can understand it by considering the linear approximation to $`E`$,
$$EE(B_Y^c=0)\mu ^{(1)}B_Y^c\mathrm{cos}\theta _W\frac{4\pi v}{g}\left(\epsilon _0\frac{\mathrm{sin}2\theta _w}{g}\frac{B_Y^c}{v^2}m^{(1)}\right)$$
(19)
where $`m^{(1)}`$ is the $`O(\theta _W)`$ dipole moment expressed in units of $`e/\alpha _WM_W(T)`$. From the Fig.1 we see that even if v(Tc)/Tc
>
1
>
๐ฃsubscript๐๐subscript๐๐1v(T_{c})/T_{c}\mathrel{\vbox{\hbox{$>$}\nointerlineskip\hbox{$\sim$}}}1 the washout condition E/Tc
>
35
>
๐ธsubscript๐๐35E/T_{c}\mathrel{\vbox{\hbox{$>$}\nointerlineskip\hbox{$\sim$}}}35 is far from being fulfilled.
It follows that the presence of a strong homogeneous hypermagnetic field does not seem to help the EWB.
## 4 Baryons from hypermagnetic helicity
A more interesting scenario may arise if hypermagnetic fields are inhomogeneous and carry a nontrivial topology. The topological properties of the hypermagnetic field are quantified by the, so called, hypermagnetic helicity, which coincide with the Chern-Simon number
$$N_{CS}=\frac{\alpha }{\pi }_Vd^3xB_YA_Y$$
(20)
where $`A_Y`$ is the hypercharge field. It is well known that the Chern-Simon number is related to the lepton and baryon number by the abelian anomaly. Recently it was noticed by Giovannini and Shaposhnikov that the magnetic helicity may have some non-trivial dynamics during the big-bang giving rise, through the anomaly equation, to a variation of the fermion and baryon contents of the Universe. Magnetic configurations with $`B_Y\stackrel{}{}\times B_Y`$ ( โmagnetic knotsโ ) may have been produced in the early Universe, for example, by the conformal invariance breaking coupling of a pseudoscalar field with the electromagnetic field which may arise, for example, in some superstring inspired models .
The interesting point that we would like to arise in the conclusions of this contribution is that a less exotic source of hypermagnetic helicity is provided by the Weiberg-Salam model itself. This source are electroweak strings. Electroweak strings are well known non-perturbative solutions of the Weinberg-Salam model (for a review see ). They generally carry a nonvanishing Chern-Simon number and, according to recent lattice simulations, they are copiously produced during the EWPT even if this transition is just a cross-over . Although electroweak string have been sometimes invoked for alternative mechanism of EWB, it was not always noticed in the literature that CP symmetry is naturally broken for twisted electroweak strings without calling for extension of the Higgs structure of the model. This is just because the twist give rise to non-orthogonal hyperelectric and hypermagnetic fields. We suggest that primordial magnetic fields, even if they are uniform, could provide a bias for the baryon number violation direction. Finally, electroweak string decay might provide the third Sacharov ingredient for EWB, namely an out-of-equilibrium condition.
In conclusions, we think that a more careful study of the possible role of electroweak strings for the EWB is worthwhile. |
warning/0002/hep-ph0002190.html | ar5iv | text | # Central Baryons in Dual Models and the Possibility of a Backward Peak in Diffraction
## 1 Introduction
In heavy ion collisions the stopping of incoming baryons is stronger than expected from a simple superposition of nucleon reactions. A number of mechanisms were introduced in a somewhat ad hoc way to repair this shortcoming and to enhance the slowing-down of baryons in heavy ion collisions and no real problem seems to remain. The aim of the present work is to better understand the basic mechanism.
In Topological models there is a definite mechanism which drastically increases the slowing-down of baryons in heavy ion reactions. Besides the cylinder of the conventional Pomeron scattering a membraned cylinder should appear. It could have a significant effect on the overall event structure in dense scattering processes. Just like the Pomeron it can be understood as a soft extrapolation of a known hard QCD exchange. If observed this second example of a connection between a soft and a hard object would be of considerable theoretical interest.
Unfortunately there is a large quantitative uncertainty. The relevant slope parameter is not sufficiently constrained by existing baryon exchange data to allow definite conclusions about the expected weight of this new contribution. However, there is a very specific effect in diffractive final states and possibly in electro-production. The existence of such a contribution can therefore be tested and its weight is accessible to direct experimental observation.
After a short excursion to hadron hadron scattering we will consider heavy ion reactions and discuss the proposed mechanism of baryon transport in this context. We will then turn to general consequences including the clear cut backward scattering peak in diffractive events.
## 2 Baryon transfer in particle scattering
### Available experimental data
To observe the slowing down of the incoming baryon charge to central or opposite rapidities, one needs to somehow identify the incoming contribution to the baryon spectra. It is usually assumed that the produced baryons and anti-baryons have identical distributions. A simple subtraction then provides the desired initial baryon distribution. This trick does not work for proton antiproton scattering, where the sea-baryon contribution cannot be determined from data without model assumptions. Unfortunately this precludes the use of post ISR data excepting HERA and diffractive systems with sufficiently large Pomeron-proton sub-energies.
Available are spectra from meson-baryon processes (compiled in ) and for a suitable combination of proton-proton and proton-antiproton processes (compiled in ). The incoming proton spectrum at fixed $`p_{}`$ (plotted in or ) are consistent with a slope of $`\alpha _{Transfer}\alpha _{Pomeron}=1`$ with large error. As the central data points are somewhat on the high side there is a hint of an eventual turnover to a flatter slope. The $`p_{}`$-integrated ratio of the incoming baryon and the sea-baryon distribution was given in to be:
$$A_{ISR}=\frac{\rho _{initialbaryoncharge}(y)}{\rho _{seabaryoncharge}(y)}=0.39\pm 0.05,\mathrm{\hspace{0.25em}0.33}\pm 0.05,\mathrm{\hspace{0.25em}0.23}\pm 0.05$$
for $`y=0.4`$, $`y=0`$ and $`y=+0.4`$<sup>1</sup><sup>1</sup>1 As the data have anyhow a large error the numbers are just taken from the figure 29. The subtraction assumed that the sea-antiproton distributions in proton-antiproton and proton-proton scattering are equal. . The central derivative of $`A_{ISR}`$ obtains no contribution from the (symmetric) sea-baryon distribution. Assuming the usual exponential distribution
$$A\mathrm{exp}[(\alpha _{Transfer}\alpha _{Pomeron})y]$$
(1)
the quantity
$$\frac{d/dyA_{ISR}}{A_{ISR}}|_{y=0}=\alpha _{Transfer}\alpha _{Pomeron}=0.49_{0.37}^{+0.42}$$
(2)
just yields the slope<sup>2</sup><sup>2</sup>2 A similar value was obtained in in a fit which required more assumptions. . While not in absolute contradiction with a slope around one the indicated value is again considerably less.
New preliminary data come from the H1 experiment at HERA. They observed the initial baryons asymmetry at laboratory rapidities
$$A_{H1}=\frac{\rho _{initialbaryoncharge}(y)}{\rho _{seabaryoncharge}(y)}=0.08\pm 0.01\pm 0.025$$
(3)
A simple extrapolation of the ISR value to the larger rapidity difference<sup>3</sup><sup>3</sup>3 It assumes that the sea stays constant. With an increasing sea baryon spectrum the conclusion about the flattening slope would be slightly stronger. would have "predicted":
$$A_{H1}=0.061_{0.046}^{+0.243}$$
(4)
Hence the H1 values lie within the expected range . The required rather flat value of the slope is
$$\alpha _{Transfer}\alpha _{Pomeron}=0.4\pm 0.2$$
The baryon stopping seen in the spectra is related by the Mueller-Kancheli relation to the annihilation cross section. With certain assumptions similar conclusions yielding a steep slope and a possible flattening at higher energies can be drawn from this data <sup>4</sup><sup>4</sup>4 Data on identified annihilation are available only at energies below $`20`$ GeV. The data fall of like a power of roughly $`a_{Transfer}1=1`$. In this range the difference in the proton-proton and proton-antiproton cross sections is saturated by annihilation. At high energies this difference in the cross sections turns to a flatter slope of roughly $`a_{Transfer}1=0.5`$. The data have smaller errors and the indication of a knee is stronger than in the inclusive ISR distribution. However the interpretation as annihilation process is not clear as mesonic trajectories ($`\omega +f_0`$ ) also contribute to $`p\overline{p}`$ scattering only. .
### Dual Topological picture
There are several Regge-pole contributions for the slowing down of baryons in hadron hadron scattering. The basic philosophy of the Dual Topological models in the classification of such exchanges involves โmaterializingโ or โsuppressedโ strings. โMaterializingโ means that the initial color fields are neutralized by a chain of hadronizing $`qq`$ pairs, โsuppressedโ means hadron-less neutralization by an exchange of a single quark. It is analogous to the Pomeron and the Reggeon exchange where in addition to a two chain Pomeron a one-chain Reggeon has to be considered. Phenomenologically contributions with various suppressed strings have to be considered as independent and additive. For each suppressed string an extra factor $`(\sqrt{1/M_{string}})`$ appears and restricts the suppressed contribution to low energies.
For a nuclear exchange one starts with a completely suppressed exchange, i. e. with the square of the quasi-elastic nucleon exchange amplitude
$$\alpha _{junction}^{III}1=2(\alpha _{Nucleon}1)=2$$
(5)
known from elastic backward scattering. Each of the three exchanged valence quarks can now be replaced by a โmaterializingโ string. Corresponding to three, two, one or zero strings there are four contributions with trajectories spaced by one half. At considered energies the first two of these โbaryoniumโ trajectories with two and three hadronizing strings
$$\alpha _{junction}^01=0.5,\alpha _{junction}^I1=1.0$$
(6)
will be relevant. They could be responsible for the initially steep ($`1.0`$) and then possible flattening ($`0.5`$) slope observed in the data discussed above.
The value of the final slope is rather uncertain. Even a value of $`\alpha _{junction}^I1=0`$ was proposed in the literature. The correspondence to the Odderon gives some support to such a value. Another uncertainty comes from the $`\omega `$-trajectory. The baryonium exchange has the same quantum numbers as the $`\omega `$-meson Reggeon. A simple estimate of an additional $`\omega `$-contribution leads to a too large non-annihilation contribution. One solution is to identify the initial trajectory as a mixture of both contributions was developed in literature. In this way it is possible to identify the turnover in the slope with onset of the sea baryon antibaryon production observed in the inclusive spectrum . The predicted strong correlation between baryon stopping and sea-baryon antibaryon density is not found in the data and we therefore assume here that the mixing if existing is very weak. An alternative solution to the $`\omega `$-problem will be given below.
### Implementation in Dual Parton model based Monte Carlo codes
The Dual Parton model was developed for high energies and it has in its present implementation no mesonic Regge-pole exchanges appear in the iterations which determine the cross sections.<sup>5</sup><sup>5</sup>5 To be precise, some Monte Carlo implementations (like DPMJET or PHOJET) include simple Regge exchanges to stay applicable at lower energies. Also if individual chains happen to come out too light to produce partons their parent quarks are annihilated to mimic a Regge-pol exchange as far as the final state is concerned. The neglect of the (flavor moving) Regge contribution is questionable in heavy ion scattering when multiple scattering processes are common and sub-energies of the involved constituents are often quite low. . Considering just the interplay of local baryonium exchanges within a global Pomeron exchange is straight forward. The factorization among strings allows to ignore the quark string which is common to both trajectories. The transition of the remaining diquark string (baryonium remnant) into an antiquark string (Pomeron remnant) can be implemented in a usual fragmentation scheme by a suitable choice of the splitting function for diquark-diquark, quark-diquark, diquark-quark and quark-quark transitions. It was implemented in most string models e.g. and it is part of the JETSET program (as diquarks or as pop-corn mechanism). Without relying on string factorization leading- and sea-baryon exchanges are also implemented in HIJING/$`B\overline{B}`$. .
## 3 Baryon enhancement <br>in dense heavy ion scattering
### Concepts for slowing-down initial baryons
There are a few completely conventional mechanisms of baryon transfer and central baryon production for multiple scattering processes in string models<sup>6</sup><sup>6</sup>6 Besides the purely kinematic โattenuationโ effect in multiple scattering events, the initial baryon can be slowed down if a secondary Pomeron exchange picks up a valence quark and leave the diquark with one (typically slower) sea quark. An obvious mechanism of central baryon production involves sea diquark-antidiquark pairs. Except for a limited number of valence partons, strings have to connect to sea partons which can also be diquark-antidiquark pairs. It is known from the analysis of the transverse momenta, i.e. of $`<p_{}>_{n_{\{charged\}}}`$ , that the partons of the string ends come from a harder (naturally more $`SU(3)`$ -symmetric) initial phase. Multiple additional sea quarks are therefore helpful in the understanding of the strangeness enhancement. . They were investigated numerically, they are helpful in some regions but not enough to explain the large stopping in heavy ion scattering.
To understand the data it seems necessary to include interplay of string if they get sufficiently dense in transverse space. It was proposed that there are new special strings. In contrast, we shall maintain here the general factorization hypothesis between initial scattering and the final hadronization within standard strings. We will just consider a more complex string structure.
The usual Pomeron exchange in the Dual Parton model leaves a quark and a diquark for the string ends. Diquarks are no special entities and multiple scattering processes have no reason not to split them in a suitable conventional two Pomeron interaction. It is natural to expect that diquark break-ups considerably slow down the baryons evolving. The probability for such an essentially un-absorbed process is:
$$[breakup]/[nobreakup][cutPomeronnumber]1$$
(7)
As required by the experimentally observed slowdown this is a drastic effect for heavy ion scattering while for hadron-hadron scattering multiple scattering are sufficiently rare (especially at energies studied in detail) to preserve the known hadron-hadron phenomenology. How such processes are affected was considered numerically in and no manifestly disturbing effects were found.
We emphasize here that the behavior of the baryon quantum number slowed down by such a break-up is not trivial. In topological models the baryon contains Y-shaped color electric fluxes. Two Pomerons intercepting two different branches will leave two โfreeโ valence quarks and a valence quark connected with the vortex line with the velocity of initial Baryon to form the end of the strings. Nothing is a priory known about the energy distribution of such quarks with vortex lines in the structure function. A simple identification with the attached quark distributions somewhat in the spirit of color evaporation models has no basis in a dual framework.
### Special baryon transfers in the Topological model
For a more detailed description of the slowing down we turn to the Dual Topological model introduced in section 2. A discussion of baryon transfers in such a framework was recently given by Kharzeev. We will here emphasize topological aspects.
In topological models a Pomeron exchange corresponds to a cylinder of in a certain way arranged gluon fields connecting the two scattering hadrons. If one considers an arbitrary plane intersecting this exchange (say at a fixed exchange-channel time) the intersection of the cylinder is topologically equivalent to a circle. More specifically in topological models amplitudes with clockwise respectively anticlockwise orientation have to be added or subtracted depending on the charge parity. The cylinders or the circles therefore come with two orientations. This distinction is usually not very important as it is always possible to attach hadrons in a matching way; except for C-parity conservation no special restrictions result.
Pomerons have a transverse extent and if they get close in transverse space they should interact. Hadronic interaction is sufficiently strong to be largely determined by geometry. It is therefore reasonable to expect that the coupling does not strongly depend on the orientation as long as there is no mechanism of suppression.
The two distinct configurations lead to different interactions. Two Pomerons with the same orientation can if they touch (starting locally at one point in the exchange-channel time) shorten their circumference and form a single circle. This then corresponds to the usual triple Pomeron coupling experimentally well known from diffractive processes.
For two Pomerons with opposite orientation the situation is more complicated. Like for soap bubbles the two surfaces which get in contact can merge and form a single membrane. The joining inverts the orientation of the membrane. On the intersecting plane one now obtains โ instead of the single circle โ three lines originating in a vortex point and ending in an anti-vortex point as shown in figure 1 .
Lacking a topological name for the three dimensional object the term membraned cylinder will be used in the following.
How do this membraned cylinder contribute to particle production? Similar to the triple Pomeron case there are three different ways to cut through a membraned cylinder. The cut which also intersects the membrane has vortex lines on both sides. They present a topological description of the baryon transfers considered in this paper. Cuts which intersect only two sheets contribute as an absorptive reduction of the two string contribution. Their negative contributions make the understanding of total cross sections not straightforward. It is possible that the membraned cylinder exchange has a vanishing or negative imaginary part. The $`\omega `$-exchange could indeed dominate the difference in the cross section while the contribution of the three string cut of the membraned cylinder could be more or less compensated by its negative two string cut. In the final states the annihilation process can be observed while the one string $`\omega `$-exchange contribution and the two string membraned cylinder reduction is hidden resp. contained in the usual two string contribution.
### The identification with the Odderon
Even though QCD cannot presently be used to calculate soft processes the typical absence of a abrupt changes in experimental distributions indicates that there is no discontinuous transition between soft and hard reactions both formulated on a partonic level. This provides the hope that hard processes can be used as a guide and that soft processes can be parameterized as an extrapolation of calculable hard processes.
The well known example is the connection between soft and hard Pomerons. To identify the hard partner of the soft Pomeron we first observe that the simplest representation of a Pomeron in PQCD involves the exchange of two gluons. As spin one particles exchanged gluons introduce no energy dependence and two gluons can form color singlets with the required positive charge parity. Following this concept it can be shown that a generalization of such an exchange gives in a rather well defined approximation the dominant contribution at very high energies. It involves a ladder of Reggeized gluons and is called the โhardโ or BFKL Pomeron. Topologically in a leading $`1/N`$ -expansion gluons can be represented by pairs of color lines drawn without crossing on a geometrical structure representing the considered contribution to the amplitude. In this expansion the leading structure of a BFKL Pomeron corresponds to a cylinder with the two basic gluons exchanged on opposites sides parallel to the axis. As they are in a color singlet state their matching color lines can be connected in front of the cylinder as shown in Fig.2a . Analogously the outer lines can be connected on the back of the cylinder.
Going back to the soft regime the basic assumption in topological models is that the $`1/N`$ -expansion stays valid <sup>7</sup><sup>7</sup>7 The $`1/N`$ expansion is in principle destroyed in the soft limit by huge combinatorial factors if the number of considered gluons increases. The hope is that preconfinement effects and a typical correlation of spatial coordinates and momenta create a situation where the (stray) long distance gluon exchanges which is responsible for the disturbing combinatorial factors cancel as color and anti-color lines are too closely neighbored when seen from a distance. and that the soft Pomeron therefore maintains its cylindrical structure. If cut, soft and hard Pomerons therefore lead to similar two string final states. As difference it remains that the slope of the soft Pomeron is just shifted downward roughly a third of a unit.
Can one find a similar connection for the membraned cylinder? The simplest representation spanning such a structure involves at least three gluons, one on each sheet again parallel to the axis. Any gluon connecting these exchanges has then to pass through a vortex line. In the $`1/N`$ expansion this means that the color lines have to cross like in Fig. 2b<sup>8</sup><sup>8</sup>8 In spite of the crossing of color lines it is a leading contribution in the $`1/N`$ expansion. The introduction of baryons in this expansion has to be done with care. For a given $`N`$ a membraned cylinder would actually require a structure with $`N2`$ membranes with $`N`$ gluons so that suppression ( $`1/N`$ ) associated with the crossed exchange would be compensated by the $`N1`$ connection choices. . A color singlet of three gluons can have the quantum numbers of a Pomeron or an Odderon. There is a simple topological property of the Odderon. A single gluon connection of type Fig.2a would project the color structure of the pair to that of a single gluon (or singlet) and the exchange would have to correspond to a Pomeron-like effective two gluon contribution. The Odderon will therefore have to involve only crossed connections shown in Fig.2b. Hence it has the topology of the membraned cylinder.
The identification with the Odderon fixes the C-parity of the membraned cylinder and a twisted membraned cylinder will have to contribute to $`pp`$ scattering at least asymptotically equally and with opposite sign as in $`p\overline{p}`$ scattering. For the twisted exchange only one type of cut exists, it will always involve a two step transition from a vortex-antivortex piece to a two string piece with both vortices on one side and to an antivortex-vortex piece. Depending on its sign it will absorb or add to the contribution in which the two string part is replaced by a usual Pomeron cut. In diffraction the situation is more complicated. As the imaginary (and real) part of the Pomeron-Pomeron-Odderon contribution has to vanish, the sum over all contributions will cancel. As some cuts have a negative sign no restriction on individual terms in the sum result.
In the same QCD approximation as the โhardโ Pomeron the properties of a โhardโ or BKP Odderon were calculated. The consensus is that the corrections to the initial gluon intercept are smaller than for the Pomeron and the intercept is rather firmly predicted to be close to $`0.96`$. There is a mismatch between this hard Odderon intercept and the (with large error) experimentally observed soft value again by about a third of a unit.
## 4 Experimental consequences <br>of membraned cylinder exchanges.
### Odderon in heavy ion scattering
In heavy ion scattering where the Pomerons are dense in transverse space they can join and form a Pomeron or a membraned cylinder. This helps with the problem of unphysical high string densities. The individual strings pairs are no longer independent but the general picture of particle production in separate universal strings survives. There should be a considerable probability of membraned cylinder exchanges growing proportional to the density:
$$\frac{[numberofmembranedcylinders]}{[Pomeronnumber]}\frac{[Pomeronnumber][Pomeronradius]^2}{[nucleusradius]^2}$$
(8)
The transition from a Pomeron pair to the centrally cut membraned-cylinder requires a baryon antibaryon pair production. Between a proton and a Pomeron the cut membraned-cylinder is a very efficient mechanism of baryon stopping.
As the slope is not well determined it is hard to obtain really reliable quantitative statements which can be tested with convincing results in heavy ion scattering. There is however a very specific qualitative prediction which should be testable.
### The backward peak in diffraction and possibly in electro production
We consider a diffractive system whose mass exceeds ISR energies. Usually the diffractively produced particles will originate in two strings of a cut Pomeron and the baryon charge will stay on the side of the initial proton. As usual there might be some migration to the center with a slope eventually corresponding to the difference of the Odderon and the Pomeron trajectory.
To accept the high soft Odderon slope suggested by the data on baryon transfers it is necessary to require a clear suppression from the coupling constants. It is quite natural to assume that the transition between a baryon exchange (cut membraned cylinder) and non baryon exchange (cut Pomeron) has no large overlap and a relatively small coupling. No such suppression is expected at a two Pomeron vertex. This has a direct consequence in rapidity space. At a certain distance it should be more favorable for an Odderon to span the total diffractive region and to utilize in this way the more favorable coupling to the two Pomerons. In consequence the initial baryon will end up exactly at the end of the string. In the usual presentation of rapidity plots this might be smeared out if different masses of diffractive system are included. This problem can be solved if one plots the rapidity distribution of initial baryons in relation to the inner end of the diffractive region, i.e. as function of
$$y_{\{Pomeron\}}=y_{\{CMS\}}\mathrm{ln}\frac{m\sqrt{s}}{M(diffr.)}$$
The expectation is that a small backward peak should then be visible. To substantiate this we show in Figure 3 the result of a calculation with the PHOJET Monte Carlo code with standard parameters. To select diffractive events a lower cutoff of $`x_F=0.95`$ was used.
PHOJET contains diquark exchanges and yields reasonable baryon spectra in the forward region. To obtain the postulated backward peak we just mix a suitable sample of such events. These special events are obtained by suitably inverting the rapidity distribution of usual events. For this inverted contribution the diquark exchanges were disabled to ensure that the backward baryon ends up on the last rank.
For a diffractive system with a mass of $`300`$ GeV the suppression from the slope is of the order of a factor of $`10`$. To account for the unknown ratio of coupling constants and for uncertainties connected with the somewhat simple implementation we add a factor $`6`$ and took $`30`$ million normal events and $`0.5`$ million inverted ones. It is important to stress that this is not an absolute prediction but an illustration of the reasonably expected effect. No effort was made to explicitely include the turnover in the baryon spectra discussed in section 2.
It is clear that suitable diffractive events should be available at the TEVATRON and that such data could be decisive. Similar measurements might also be possible at HERA. It is likely in the scattering of a virtual photon on a proton that the photon does not prefer a fixed topology and the coupling to an Odderon is also not disfavored. In this case the same backward peak might be observable in non diffractive $`ep`$ data.
## 5 Conclusion
The present paper wants to encourage the measurement of the initial baryon distribution in high mass diffractive systems. The known initial baryon distribution suggests that there is a contribution to the exchange which is most strongly suppressed by its coupling and not by its slope. It is argued that this should not be the case in diffractive events and in consequence there should be an observable tiny backward peak in the initial baryon distribution in the diffractive system. The prediction is important as it has manifest consequences for heavy ion processes, where it would be a strong mechanism for central baryon production and for the transport of initial baryons to the central and opposite region. It might also clarify the role of the Odderon.
## Acknowledgments
I thank Lech Szymanowski and L. Lukaszuk for discussion regarding the Odderon. I am indebted to J. Ranft and H. Anlauf for reading the manuscript and many useful comments. I also acknowledge his help in running the PHOJET code to obtain Figure 3. The work is based on work done with Patrick Aurenche. |
warning/0002/nlin0002039.html | ar5iv | text | # Front motion for phase transitions in systems with memory
## 1 Introduction
In this manuscript we consider the following partial integro-differential equation
$$ฯต^2\varphi _t=_0^ta(tt^{})[ฯต^2\mathrm{\Delta }\varphi +f(\varphi )+ฯตh](t^{})๐t^{},$$
(1)
where $`\varphi (x,t)`$ is a field, called an order parameter, defined in $`\mathrm{\Omega }\times [0,T]`$, where $`\mathrm{\Omega }R^2`$ and $`T>0`$, with Dirichlet and Neumann boundary conditions. The parameter $`ฯต`$ is assumed to be small, $`ฯต1`$, $`f(\varphi )`$ is a real odd function with a positive maximum equal to $`\varphi ^{}`$, a negative minimum equal to $`\varphi ^{}`$ and precisely three zeros in the closed interval $`[b,b]`$ located at $`0`$ and $`\pm b`$, where $`b`$ is a positive constant. For simplicity we will consider $`b=1`$. The operator $`\mathrm{\Delta }`$ is the Laplace operator. The kernel $`a`$ is assumed to be a piecewise continuous, differentiable at the origin and scalar-valued functions on $`(0,\mathrm{})`$ satisfying additional conditions to be described later.
Specific cases of equation (1) are $`a(t)=\delta (t)`$, the Allen-Cahn equation:
$$ฯต^2\varphi _t=ฯต^2\mathrm{\Delta }\varphi +f(\varphi )+ฯตh,$$
(2)
which is the simplest model of phase transition with a non-conserved order parameter, and the damped Klein-Gordon equation (used, e.g., in the theory of long Josephson junctions):
$$ฯต^2\varphi _{tt}+\gamma ฯต^2\varphi _t=ฯต^2\mathrm{\Delta }\varphi +f(\varphi )+ฯตh,$$
(3)
which is obtained for an exponentially decreasing kernel of type \[kn:rotnep1\]
$$a(t)=e^{\gamma t}.$$
(4)
Equation (1) can be thought of as a phenomenological equation based on energetic penalization driving the evolution of the system toward equilibrium states. More specifically, let us call $`F_ฯต=_{R^2}[\frac{1}{2}ฯต^2(\varphi )^2V(\varphi )]๐\overline{x}`$, with $`f(\varphi )+ฯตh=dV/d\varphi `$ the free energy of the system. The functional derivative of the free energy, $`\delta F_ฯต(\varphi )/\delta \varphi `$ is considered as a generalized force indicative of the tendency of the free energy to decay towards a minimum. The equation is obtained by assuming that the response of $`\varphi `$ to the tendency of the free energy to decay towards a minimum is given by \[kn:rotnep1, kn:rot1\],
$$ฯต^2\varphi _t=_0^ta(tt^{})\frac{\delta F_ฯต}{\delta \varphi }(t^{})๐t^{}.$$
A system of integro-differential equations closely related to (1), the phase field equations with memory, is \[kn:rotnep3, kn:rot1\]
$$\{\begin{array}{c}u_t+ฯต^2\varphi _t=_0^ta_1(tt^{})\mathrm{\Delta }u(t^{})๐t^{}\hfill \\ \\ ฯต^2\varphi _t=_0^ta_2(tt^{}[ฯต^2\mathrm{\Delta }\varphi +f(\varphi )+ฯตu](t^{})dt^{}\hfill \end{array}$$
(5)
in $`\mathrm{\Omega }\times [0,t]`$ where $`\mathrm{\Omega }R^2`$ and $`T>0`$ with Dirichlet and Neumann boundary conditions for the temperature field, $`u(x,t)`$ and the order parameter, $`\varphi (x,t)`$, respectively and for $`u`$ vanishing initially. The parameter $`ฯต`$ and $`f(\varphi )`$ are as for (1) and the kernels $`a_1`$ and $`a_2`$ are similar to $`a`$; i.e., piecewise continuous, differentiable at the origin, and scalar-valued functions on $`(0,\mathrm{})`$.
For $`a_1(t)=a_2(t)=\delta (t)`$, system (LABEL:eq:defas2) gives rise to the classical phase field equations with a non-conserved order parameter
$$\{\begin{array}{c}u_t+ฯต^2\varphi _t=\mathrm{\Delta }u,\hfill \\ \\ ฯต^2\varphi _t=ฯต^2\mathrm{\Delta }\varphi +f(\varphi )+ฯตu.\hfill \end{array}$$
(6)
System (LABEL:eq:defas2) also generalizes the hyperbolic phase field equations with a non-conserved order parameter:
$$\{\begin{array}{c}u_{tt}+ฯต^2\varphi _{tt}+\gamma _1u_t+ฯต^2\gamma _1\varphi _t=\alpha \mathrm{\Delta }u,\hfill \\ \\ ฯต^2\varphi _{tt}+ฯต^2\gamma _2\varphi _t=ฯต^2\mathrm{\Delta }\varphi +f(\varphi )+ฯตu,\hfill \end{array}$$
(7)
reducing to them when the kernels $`a_1`$ and $`a_2`$ are exponentially decreasing functions; i.e., $`a_i(t)=\alpha _ie^{\gamma _it}`$, $`i=1,2`$, with $`\alpha _i`$ and $`\gamma _i`$ are non-negative constants. System (LABEL:eq:defas1) is obtained by differentiating both equations in (LABEL:eq:defas2), rearranging terms and rescaling by means of the transformation $`t\alpha _2^{\frac{1}{2}}t`$, $`\gamma _2\alpha _2^{\frac{1}{2}}\gamma _2`$, $`\gamma _1\alpha _2^{\frac{1}{2}}\gamma _1`$ and $`\alpha :=\alpha _1/\alpha _2`$.
The phase field equations with memory (LABEL:eq:defas2) describe the process of phase transitions where memory effects are present. For a readable description of the classical phase field equations see \[kn:cag1\]. The first equation in (LABEL:eq:defas2) is based on the balance of heat equation for a non-Fourier process in which the expression for the heat flux is given by a convolution in time between the temperature gradient and the kernel $`a_1`$ \[kn:aizbar1, kn:nov1\]. The second equation is obtained in a similar way as equation (1) \[kn:rot1\].
In what follows we assume that there exists a solution $`\varphi (x,t)`$ of (1) or $`\{u(x,t),\varphi (x,t)\}`$ of (LABEL:eq:defas2) defined for all small $`ฯต`$, every $`x\mathrm{\Omega }`$ and every $`t[0,T]`$ which contains an internal layer. We also assume, for such solutions, that for all small $`ฯต0`$ and all $`t[0,T]`$, the domain $`\mathrm{\Omega }`$ can be divided into two open regions $`\mathrm{\Omega }_+(t;ฯต)`$ and $`\mathrm{\Omega }_{}(t,ฯต)`$ with a curve $`\mathrm{\Gamma }(t;ฯต)`$, separating between them. This interface defined by
$$\mathrm{\Gamma }(t;ฯต):=\{x\mathrm{\Omega }:\varphi (x,t;ฯต)=0\},$$
(8)
is assumed to be smooth, which implies that its curvature and its velocity are bounded independently of $`ฯต`$. The function $`\varphi `$ is assumed to vary continuously across the interface, far from the interface tending to $`1`$ when $`x\mathrm{\Omega }_+(t;ฯต)`$, -1 when $`x\mathrm{\Omega }_{}(t,ฯต)`$, with rapid spatial variation close to the interface. The problem is to derive a closed equation for the evolution of the interface asymptotically valid as $`ฯต0`$.
The latter problem has been studied formerly by means of formal asymptotic analysis for PDEs (2) and (3) \[kn:cagfif1, kn:fif1, kn:fifpen1\], as well as for systems (6) and (LABEL:eq:defas1) \[kn:rotnep1, kn:neu1\]. For (2) and (6), fronts evolve according to the mean curvature flow equation \[kn:rubste1, kn:cag3, kn:cag2\],
$$v=\kappa ,$$
(9)
where $`v`$ is the normal velocity of the interface and $`\kappa `$ its curvature. If $`y=S(x,t)`$ is the Cartesian description of the interface, equation (9) reads
$$S_t=\frac{S_{xx}}{1+S_x^2}.$$
For the undamped Klein-Gordon equation, (3) with $`\gamma =h=0`$, Neu \[kn:neu1\] proved that the evolution of the interface is governed by the Born-Infeld equation
$$(1S_t^2)S_{xx}+2S_xS_tS_{xt}(1+S_x^2)S_{tt}=0.$$
(10)
For (3) with $`\gamma >0`$ and for (LABEL:eq:defas1) Rotstein et. al. \[kn:rotnep1, kn:rot1\] showed that fronts move according to an extended version of the Born-Infeld equation given by
$$(1\alpha S_t^2)S_{xx}+2\alpha S_xS_tS_{xt}\alpha (1+S_x^2)S_{tt}$$
$$\gamma S_t(1+S_x^2\alpha S_t^2)\nu (1+S_x^2\alpha S_t^2)^{\frac{3}{2}}=0,$$
(11)
where the parameter $`\nu `$ (proportional to $`h`$ in (3)) is defined later.
In local (geometric) coordinates equation (11) reads
$$\frac{v_t}{1\alpha v^2}+\gamma v=\kappa +\nu (1\alpha v^2)^{\frac{1}{2}}.$$
(12)
where $`v`$, $`v_t`$ and $`\kappa `$ are the normal velocity, normal acceleration and curvature of the interface respectively. Equation (11) or (12) have been studied in \[kn:rot1\], and in \[kn:neu1\] for $`\gamma =\nu =0`$ .
The case of integrodifferential equations (1) and (LABEL:eq:defas2) is still less investigated. For (LABEL:eq:defas2) and for kernels which are Laplace transforms of suitable functions satisfying
$$_0^{\mathrm{}}a(t)๐t<\mathrm{},_0^{\mathrm{}}\overline{\alpha }(s)๐s<\mathrm{}\text{and}_0^{\mathrm{}}\overline{\alpha }(s)s๐s<\mathrm{}$$
where $`\overline{\alpha }(s)`$ is the inverse Laplace transform of $`a(t)`$, Rotstein et al. \[kn:rot1, kn:rotnep5\] showed that the equation governing the motion of the interface is also (11) where in this case
$$\gamma =\frac{a_2^{}(0)}{a_2(0)^2}.$$
In the present manuscript we show that in the asymptotic limit studied and for a certain class of kernels, which will be made clear below, equation (1) is equivalent to (3) from the point of view of fronts motion; i.e., the evolution of interfaces is described by the same equation, (11) with
$$\alpha =\left(_0^{\mathrm{}}\overline{\alpha }(\tau )๐\tau \right)^1,$$
where $`\overline{\alpha }(\tau )`$ is the inverse Laplace transform of $`a(t)`$, and
$$\gamma =\alpha ^2\left(_0^{\mathrm{}}\overline{\alpha }(\tau )\tau ๐\tau \right).$$
We also find the equivalence between (LABEL:eq:defas2) and (LABEL:eq:defas1) \[kn:rotnep5\] for a more general class of kernels, described below. In Section 2 we describe the class of kernels considered in this work, give some examples and explain our strategy in dealing with them. In Section 3 we derive the equation of motion for (1). We present the derivation treating the equation in Cartesian coordinates and representing the kernel as a Laplace tranform of a suitable function. For kernels represented as a Fourier transform of suitable functions the derivation is similar. In our derivation we assume that the interface has no oscillations in the sense that there are no points on the interface for which the velocity vanishes; i.e., the front is an advancing front. This assumption is not necessary for the case of exponentially decreasing kernels; i.e., for (3). In Section 4 we derive the equation of motion for (LABEL:eq:defas2) for kernels which are inverse Fourier transforms of suitable functions. Here we treat the problem in local coordinates.
## 2 Description of the class of kernels considered
In order to deal with the memory integral in deriving the equations of front motion for either (1) or (LABEL:eq:defas2) we represent the kernels $`a(t)`$ and $`a_i(t)`$, $`i=1,2`$, as Laplace transforms of suitable functions:
$$a(t)=_0^{\mathrm{}}\overline{\alpha }(\tau )e^{\tau t}๐\tau ,$$
(13)
or
$$a_i(t)=_0^{\mathrm{}}\overline{\alpha _i}(\tau )e^{\tau t}๐\tau .$$
(14)
for $`i=1,2`$. Substituting (13) into (1), rearranging terms and defining
$$\chi (t;\tau )=_0^te^{\tau (tt^{})}[ฯต^2\mathrm{\Delta }\varphi +f(\varphi )+ฯตh](t^{})๐t^{},$$
(15)
where it is understood that $`\chi `$ is also a function of the space variable $`x`$ and it is defined for all $`x\mathrm{\Omega }`$, we obtain
$$\{\begin{array}{c}ฯต^2\varphi _t=_0^{\mathrm{}}\overline{\alpha }(\tau )\chi (t;\tau )๐\tau ,\hfill \\ \\ \chi _t(t;\tau )+\tau \chi (t;\tau )=ฯต^2\mathrm{\Delta }\varphi +f(\varphi )+ฯตh,\hfill \end{array}$$
(16)
for all $`\tau [0,\mathrm{})`$. Note that the second equation in (LABEL:eq:memphasefield2comp) satisfies $`\chi (0,\tau )=0`$ for all $`x\mathrm{\Omega }`$ and for all $`\tau [0,\mathrm{})`$. We can see that, in spite of the fact that there is an integral involved in the first equation, that integral sums over the parameter $`\tau `$ not involving the time variable $`t`$ as in (1). For the phase field equations with memory, substituting (14) into (LABEL:eq:defas2), rearranging terms and calling
$$\chi (t;\tau )=_0^te^{\tau (tt^{})}[ฯต^2\mathrm{\Delta }\varphi +f(\varphi )+ฯตu](t^{})๐t^{},$$
(17)
and
$$v(t;\tau )=_0^te^{\tau (tt^{})}\mathrm{\Delta }u(t^{})๐t^{},$$
(18)
where, again, it is understood that $`\chi `$ and $`v`$ are also functions of the space variable $`x`$ and are defined for all $`x\mathrm{\Omega }`$, we obtain
$$\{\begin{array}{c}u_t+ฯต^2\varphi _t=_0^{\mathrm{}}\overline{\alpha _1}(\tau )v(t;\tau )๐\tau ,\hfill \\ \\ v_t(t;\tau )+\tau v(t;\tau )=\mathrm{\Delta }u,\hfill \\ \\ ฯต^2\varphi _t=_0^{\mathrm{}}\overline{\alpha }(\tau )\chi (t;\tau )๐\tau ,\hfill \\ \\ \chi _t(t;\tau )+\tau \chi (t;\tau )=ฯต^2\mathrm{\Delta }\varphi +f(\varphi )+ฯตu,\hfill \end{array}$$
(19)
for all $`\tau [0,\mathrm{})`$. Note that the second and fourth equations satisfy $`v(0,\tau )=\chi (0,\tau )=0`$ for all $`x\mathrm{\Omega }`$ and for all $`\tau [0,\mathrm{})`$.
For an exponentially decreasing kernel (4), $`\overline{\alpha }(\tau )=\delta (\tau \gamma )`$. For a kernel which is a linear combination of kernels of type (4), $`a(t)=_{k=1}^n\alpha _ke^{\gamma _kt}`$, for positive constants $`\alpha _k`$ and $`\gamma _k`$ ($`k=1,2,\mathrm{},n`$), we have $`\overline{\alpha }(\tau )=_{k=1}^n\alpha _k\delta (\tau \gamma _k)`$. In the latter case we can define for $`k=1,2,\mathrm{},n`$
$$\chi (t;\gamma _k)=_0^te^{\gamma _k(tt^{})}[ฯต^2\mathrm{\Delta }\varphi +f(\varphi )+ฯตh](t^{})๐t^{},$$
(20)
and substitute into (1) obtaining
$$\{\begin{array}{c}ฯต^2\varphi _t=_{k=1}^n\alpha _k\chi (t;\gamma _k),\hfill \\ \\ \chi _t(t;\gamma _k)+\gamma _k\chi (t;\gamma _k)=ฯต^2\mathrm{\Delta }\varphi +f(\varphi )+ฯตh,\hfill \end{array}$$
(21)
for $`k=1,2,\mathrm{},n`$. It is clear that (LABEL:eq:memphasefield2sum) is a particular (discrete) case of (LABEL:eq:memphasefield2comp) and the simplest generalization of (3). This discussion can be easily adapted to system (LABEL:eq:defas2).
The approach discussed up to now in this section allows us to consider kernels which are not necessarily exponentially decreasing. Examples are, for positive $`\alpha `$, $`\beta `$ and $`\gamma `$: $`a(t)=1/(\gamma t+\alpha )`$, for which $`\overline{\alpha }(\tau )=1/\gamma e^{\frac{\alpha }{\gamma }t}`$ or, more generally $`a(t)=1/(\gamma t+\alpha )^n`$ (with $`n`$ a positive integer), for which $`\overline{\alpha }(\tau )=1/(\gamma ^n(n1)!)\tau ^{n1}e^{\frac{\alpha }{\gamma }t}`$; $`a(t)=\beta /((\gamma t+\alpha )^2+\beta ^2)`$ and $`a(t)=(\gamma t+\alpha )/((\gamma t+\alpha )^2+\beta ^2)`$ for which $`\overline{\alpha }(\tau )=1/\gamma ^2e^{\frac{\alpha }{\gamma }t}\mathrm{sin}(\beta \tau )`$ and $`\overline{\alpha }(\tau )=1/\gamma ^2e^{\frac{\alpha }{\gamma }t}\mathrm{cos}(\beta \tau )`$ respectively. A class of kernels not included in (13) or (14) is $`a(t)=\alpha e^{\gamma t}\mathrm{cos}(\beta t)`$, or equivalently, $`a(t)=\alpha \frac{e^{z_1t}+e^{\overline{z}_1t}}{2}`$ where $`z_1=\gamma i\beta `$.
From a slightly different point of view, kernels of type $`a(t)=\alpha e^{\gamma t}\mathrm{cos}(\beta t)`$ can be represented as Fourier transforms of suitable functions.
$$a(t)=_{\mathrm{}}^{\mathrm{}}\overline{\alpha }(w)e^{iwt}๐w,$$
(22)
or
$$a_i(t)=_{\mathrm{}}^{\mathrm{}}\overline{\alpha _i}(w)e^{iwt}๐w.$$
(23)
for $`i=1,2`$ where
$$\overline{\alpha }(w)=\frac{1}{2\pi }_{\mathrm{}}^{\mathrm{}}a(t)e^{iwt}๐t,$$
and
$$\overline{\alpha _i}(w)=\frac{1}{2\pi }_{\mathrm{}}^{\mathrm{}}a_i(t)e^{iwt}dt.$$
for $`i=1,2`$. In this case the functions $`a(t)`$ and $`a_i(t)`$, $`i=1,2`$ are understood to be defined for all $`t`$ and being equal to zero for $`t<0`$. Substituting (22) into (1), rearranging terms and defining
$$\chi (t;w)=_0^te^{iw(tt^{})}[ฯต^2\mathrm{\Delta }\varphi +f(\varphi )+ฯตh](t^{})๐t^{},$$
(24)
where it is understood that $`\chi `$ is a function of the space variable $`x\mathrm{\Omega }`$, we obtain
$$\{\begin{array}{c}ฯต^2\varphi _t=_{\mathrm{}}^{\mathrm{}}\overline{\alpha }(w)\chi (t;w)๐w,\hfill \\ \\ \chi _t(t;w)+iw\chi (t;w)=ฯต^2\mathrm{\Delta }\varphi +f(\varphi )+ฯตh,\hfill \end{array}$$
(25)
for all $`w(\mathrm{},\mathrm{})`$. For the phase field equations with memory, substituting (23) into (LABEL:eq:defas2), rearranging terms and calling
$$\chi (t;w)=_0^te^{iw(tt^{})}[ฯต^2\mathrm{\Delta }\varphi +f(\varphi )+ฯตu](t^{})๐t^{},$$
(26)
and
$$v(t;w)=_0^te^{iw(tt^{})}\mathrm{\Delta }u(t^{})๐t^{},$$
(27)
where, again, it is understood that $`\chi `$ and $`v`$ are functions of the space variable $`x\mathrm{\Omega }`$. System (LABEL:eq:defas2) becomes, for all $`w`$
$$\{\begin{array}{c}u_t+ฯต^2\varphi _t=_{\mathrm{}}^{\mathrm{}}\overline{\alpha _1}(w)v(t;w)๐w,\hfill \\ \\ v_t(t;w)+iwv(t;w)=\mathrm{\Delta }u,\hfill \\ \\ ฯต^2\varphi _t=_{\mathrm{}}^{\mathrm{}}\overline{\alpha }(w)\chi (t;w)๐w,\hfill \\ \\ \chi _t(t;w)+iwchi(t;w)=ฯต^2\mathrm{\Delta }\varphi +f(\varphi )+ฯตu,\hfill \end{array}$$
(28)
for all $`w(\mathrm{},\mathrm{})`$. As we set at the beginning of the introduction, the kernels $`a(t)`$, are assumed to be piecewise continuous, differentiable at the origin, scalar-valued functions on $`(0,\mathrm{})`$. Moreover they are assumed to be independent of $`ฯต`$ and such that
$$_0^{\mathrm{}}a(t)๐t<\mathrm{},$$
and
$$_0^{\mathrm{}}\overline{\alpha }(s)๐s<\mathrm{}\text{and}_0^{\mathrm{}}\overline{\alpha }(s)s๐s<\mathrm{}.$$
when we use the Laplace transform, and
$$_{\mathrm{}}^{\mathrm{}}\overline{\alpha }(s)๐s<\mathrm{}\text{and}_{\mathrm{}}^{\mathrm{}}\overline{\alpha }(s)s๐s<\mathrm{}.$$
when we use the Fourier transform. In both cases this is equivalent to $`a(0)<\mathrm{}`$ and $`a^{}(0)<\mathrm{}`$.
## 3 The Allen-Cahn equation with memory - asymptotic analysis in Cartesian coordinates
For points outside the interface and for $`ฯต>0`$ we expand $`\varphi `$ as follows
$$\varphi =\varphi (x,t;ฯต)=\varphi ^0(x,t)+ฯต\varphi ^1(x,t)+ฯต^2\varphi ^2(x,t)+๐ช(ฯต^3),$$
and substitute into (1) obtaining the $`๐ช(1)`$ and $`๐ช(ฯต)`$ respectively
$$_0^ta(tt^{})f(\varphi ^0)๐t^{}=0,$$
whose solution is $`f(\varphi ^0)=0`$ (or $`\varphi ^0=\pm 1`$), and
$$_0^ta(tt^{})[f^{}(\varphi ^0)\varphi ^1+h]๐t^{}=0,$$
giving $`f^{}(\varphi ^0)\varphi ^1+h=0`$. Thus, for points which have not yet been reached by the moving front, we have
$$\chi ^0(t;\tau )=0,$$
(29)
and
$$\chi ^1(t;\tau )=0,$$
(30)
for all $`\tau [0,\mathrm{})`$.
For the asymptotic analysis using Cartesian coordinates the interface is represented by
$`y=S(x,t,ฯต)`$ for $`ฯต`$ sufficiently small and assume $`S_t0`$ for all $`x\mathrm{\Omega }`$ and all $`t0`$. We define a new variable
$$z:=\frac{yS(x,t,ฯต)}{ฯต}$$
which is assumed to be $`๐ช(1)`$ as $`ฯต0`$ near the interface. We call $`\mathrm{\Phi }`$ the asymptotic form of $`\varphi `$ as $`ฯต0`$ with $`z`$ fixed; i.e.,
$$\varphi =\mathrm{\Phi }(z,x,t;ฯต).$$
(31)
The field equations (LABEL:eq:memphasefield2comp) in $`(z,x,t)`$ coordinates, after differentiating the first equation with respect to $`t`$ and calling $`\chi =\mathrm{{\rm Y}}(x,z,t;\tau ,ฯต)`$, become (see Appendix A)
$$\{\begin{array}{c}ฯต^3\mathrm{\Phi }_{tt}2ฯต^2S_t\mathrm{\Phi }_{zt}+ฯตS_t^2\mathrm{\Phi }_{zz}ฯต^2S_{tt}\mathrm{\Phi }_z=_0^{\mathrm{}}\overline{\alpha }(\tau )[ฯต\mathrm{{\rm Y}}_tS_t\mathrm{{\rm Y}}_z]๐\tau ,\hfill \\ \\ ฯต\mathrm{{\rm Y}}_tS_t\mathrm{{\rm Y}}_z+ฯต\tau \mathrm{{\rm Y}}=ฯต^3\mathrm{\Phi }_{xx}2ฯต^2S_x\mathrm{\Phi }_{zx}+ฯต(1+S_x^2)\mathrm{\Phi }_{zz}ฯต^2S_{xx}\mathrm{\Phi }_z+ฯตf(\mathrm{\Phi })+ฯตh,\hfill \end{array}$$
(32)
for $`\tau [0,\mathrm{})`$.
The asymptotic expansion of $`\mathrm{\Phi }`$ is assumed to have the form
$$\mathrm{\Phi }\mathrm{\Phi }^0+ฯต\mathrm{\Phi }^1+ฯต^2\mathrm{\Phi }^2+๐ช(ฯต^3),\text{as}ฯต0.$$
(33)
Substituting into (LABEL:eq:memphasefield2syscomp) and equating coefficients of the corresponding powers of $`ฯต`$, we obtain the following problems for $`๐ช(1)`$, $`๐ช(ฯต)`$ and $`๐ช(ฯต^2)`$ respectively:
$$\{\begin{array}{c}_0^{\mathrm{}}\overline{\alpha }(\tau )S_t\mathrm{{\rm Y}}_z^0๐\tau =0,\hfill \\ \\ S_t\mathrm{{\rm Y}}_z^0=0,\hfill \end{array}$$
(34)
$$\{\begin{array}{c}S_t^2\mathrm{\Phi }_{zz}^0=_0^{\mathrm{}}\overline{\alpha }(\tau )[\mathrm{{\rm Y}}_t^0S_t\mathrm{{\rm Y}}_z^1d\tau ],\hfill \\ \\ \mathrm{{\rm Y}}_t^0S_t\mathrm{{\rm Y}}_z^1+\tau \mathrm{{\rm Y}}^0=(1+S_x^2)\mathrm{\Phi }_{zz}^0+f(\mathrm{\Phi }^0),\hfill \end{array}$$
(35)
and
$$\{\begin{array}{c}2S_t\mathrm{\Phi }_{zt}^0+S_t^2\mathrm{\Phi }_{zz}^1S_{tt}\mathrm{\Phi }_z^0=_0^{\mathrm{}}\overline{\alpha }(\tau )[\mathrm{{\rm Y}}_t^1S_t\mathrm{{\rm Y}}_z^2]๐\tau ,\hfill \\ \\ \mathrm{{\rm Y}}_t^1S_t\mathrm{{\rm Y}}_z^2+\tau \mathrm{{\rm Y}}^1=2S_x\mathrm{\Phi }_{zx}^0+(1+S_x^2)\mathrm{\Phi }_{zz}^1S_{xx}\mathrm{\Phi }_z^0+f^{}(\mathrm{\Phi }^0)\mathrm{\Phi }^1+h,\hfill \end{array}$$
(36)
for all $`\tau [0,\mathrm{})`$. From (29) and (30) we have
$$\underset{z\mathrm{}}{lim}\mathrm{{\rm Y}}^0(x,z,t;\tau )=0,$$
(37)
and
$$\underset{z\mathrm{}}{lim}\mathrm{{\rm Y}}^1(x,z,t;\tau )=0.$$
(38)
From (LABEL:eq:probsyscart0) we have $`\mathrm{{\rm Y}}_z^0(x,z,t;\tau )=0`$ for all $`\tau [0,\mathrm{})`$. Integrating with respect to $`z`$ and applying condition (37) (assuming $`S_t0`$) yields
$$\mathrm{{\rm Y}}^00,$$
(39)
for all $`\tau [0,\mathrm{})`$. Substituting (39) into (LABEL:eq:probsyscart1) the second equation becomes
$$S_t\mathrm{{\rm Y}}_z^1=(1+S_x^2)\mathrm{\Phi }_{zz}^0+f(\mathrm{\Phi }^0).$$
(40)
We see that $`\mathrm{{\rm Y}}_z^1`$ does not depend on $`\tau `$. We multiply the second equation in (LABEL:eq:probsyscart1) by $`\alpha (\tau )`$ and integrate with respect to $`\tau `$. We obtain
$$S_t^2\mathrm{\Phi }_{zz}^0=[(1+S_x^2)\mathrm{\Phi }_{zz}^0+f(\mathrm{\Phi }^0)]_0^{\mathrm{}}\overline{\alpha }(\tau )๐\tau $$
since $`\mathrm{\Phi }`$ and $`S`$ do not depend on $`\tau `$. We call
$$\alpha =\left(_0^{\mathrm{}}\overline{\alpha }(\tau )๐\tau \right)^1,$$
(41)
obtaining
$$(1+S_x^2\alpha S_t^2)\mathrm{\Phi }_{zz}^0+f(\mathrm{\Phi }^0)=0.$$
(42)
From (40) we have that $`S_t\mathrm{{\rm Y}}_z^1=\alpha S_t^2\mathrm{\Phi }_{zz}^0`$. Integrating with respect to $`z`$ and applying condition (38) we obtain
$$\mathrm{{\rm Y}}^1=\alpha S_t\mathrm{\Phi }_z^0.$$
(43)
In order to solve (42) we define a new variable
$$\xi :=\frac{z}{(1+S_x^2\alpha S_t^2)^{\frac{1}{2}}}.$$
(44)
In terms of $`\xi `$, equation (42) reads
$$\mathrm{\Phi }_{\xi \xi }^0+f(\mathrm{\Phi }^0)=0,$$
(45)
whose solution is $`\mathrm{\Phi }^0=\mathrm{\Psi }(\xi )`$, the unique solution of $`\mathrm{\Psi }^{\prime \prime }+f(\mathrm{\Psi })=0,\mathrm{\Psi }(\pm \mathrm{})=\varphi ^\pm ,\mathrm{\Psi }(0)=0`$. Thus
$$\mathrm{\Phi }^0=\mathrm{\Phi }^0\left(\frac{z}{(1+S_x^2\alpha S_t^2)^{\frac{1}{2}}}\right).$$
(46)
Multiplying the second equation in (LABEL:eq:probsyscart2) by $`\alpha (\tau )`$, integrating with respect to $`\tau `$, replacing (43) into (LABEL:eq:probsyscart2), and rearranging terms we obtain
$$(1+S_x^2\alpha _2S_t^2)\mathrm{\Phi }_{zz}^1+f^{}(\mathrm{\Phi }^0)\mathrm{\Phi }^1=(S_{xx}\alpha S_{tt})\mathrm{\Phi }_z^02\alpha S_t\mathrm{\Phi }_{zt}^0+$$
$$+2S_x\mathrm{\Phi }_{zx}^0h\alpha ^2S_t\left(_0^{\mathrm{}}\overline{\alpha }(\tau )\tau ๐\tau \right)\mathrm{\Phi }_z^0.$$
(47)
Setting
$$\gamma =\alpha ^2\left(_0^{\mathrm{}}\overline{\alpha }(\tau )\tau ๐\tau \right).$$
(48)
In terms of $`\xi `$, $`x`$ and $`t`$ equation (47) reads (see Appendix A)
$$\mathrm{\Phi }_{\xi \xi }^1+f^{}(\mathrm{\Phi }^0)\mathrm{\Phi }^1=\frac{S_{xx}\alpha S_{tt}\gamma S_t}{(1+S_x^2\alpha _2S_t^2)^{\frac{1}{2}}}\mathrm{\Phi }_\xi ^0\frac{2S_x(S_xS_{xx}\alpha S_tS_{xt})}{(1+S_x^2\alpha _2S_t^2)^{\frac{3}{2}}}(\xi \mathrm{\Phi }_{\xi \xi }^0+\mathrm{\Phi }_\xi ^0)+$$
$$+\frac{2\alpha S_t(S_xS_{xt}\alpha S_tS_{tt})}{(1+S_x^2\alpha _2S_t^2)^{\frac{3}{2}}}(\xi \mathrm{\Phi }_{\xi \xi }^0+\mathrm{\Phi }_\xi ^0)h.$$
(49)
It is straightforward to check that $`\mathrm{\Psi }^{}(\xi )`$ satisfies the homogeneous equation
$$\mathrm{\Phi }_{\xi \xi }^1+f^{}(\mathrm{\Phi }^0)\mathrm{\Phi }^1=0.$$
(50)
This implies that the operator $`\mathrm{\Lambda }`$ defined by
$$\mathrm{\Lambda }:=\frac{^2}{\xi ^2}+f^{}(\mathrm{\Psi }^{}(\xi ))$$
(51)
has a simple eigenvalue at the origin with $`\mathrm{\Psi }^{}`$ as the corresponding eigenfunction. Now the solvability condition for equation (49) gives
$$\frac{S_{xx}\alpha S_{tt}\gamma S_t}{(1+S_x^2\alpha _2S_t^2)^{\frac{1}{2}}}_{\mathrm{}}^{\mathrm{}}(\mathrm{\Psi }^{})^2๐\xi h_{\mathrm{}}^{\mathrm{}}\mathrm{\Psi }^{}๐\xi +$$
$$\frac{2S_x(S_xS_{xx}\alpha S_tS_{xt})2\alpha S_t(S_xS_{xt}\alpha S_tS_{tt})}{(1+S_x^2\alpha _2S_t^2)^{\frac{3}{2}}}_{\mathrm{}}^{\mathrm{}}(\xi \mathrm{\Psi }^{\prime \prime }+\mathrm{\Psi }^{})\mathrm{\Psi }^{}๐\xi =0.$$
(52)
A simple calculation shows that
$$_{\mathrm{}}^{\mathrm{}}\xi \mathrm{\Psi }^{}\mathrm{\Psi }^{\prime \prime }๐\xi =\frac{1}{2}_{\mathrm{}}^{\mathrm{}}(\mathrm{\Psi }^{})^2๐\xi .$$
(53)
Defining
$$\nu :=h\frac{\mathrm{\Psi }(+\mathrm{})\mathrm{\Psi }(\mathrm{})}{_{\mathrm{}}^{\mathrm{}}(\mathrm{\Psi }^{})^2๐\xi },$$
(54)
Substituting (53) and (54) into (52), multiplying equation (52) by $`(1+S_x^2\alpha _2S_t^2)^{\frac{3}{2}}`$ and rearranging terms we obtain (11). Note that for $`f(\varphi )=(\varphi \varphi ^3)/2`$ (Ginzburg-Landau theory), $`\mathrm{\Psi }(\xi )=\mathrm{tanh}\frac{\xi }{2}`$ and $`\nu :=3h`$, whereas for $`f(\varphi )=\mathrm{sin}\varphi `$, $`\mathrm{\Psi }(\xi )=4\mathrm{tan}^1e^\xi \pi `$ and $`\nu :=\frac{\pi }{4}h`$.
## 4 The phase field equations with memory - asymptotic analysis in local coordinates
### 4.1 Assumptions and definitions
For the asymptotic analysis, using local coordinates, we set $`x=(x_1,x_2)`$ and call $`d(x)`$ the distance from $`x`$ to $`\mathrm{\Gamma }`$; i.e., $`d(x)=dist(x,\mathrm{\Gamma })`$. We next define a local orthogonal coordinate system $`(r,s)`$ in a neighborhood of $`\mathrm{\Gamma }`$ in the following way
$$r(x,t;ฯต)=\{\begin{array}{ccc}d(x)\hfill & if\hfill & \varphi (x)>0\hfill \\ & & \\ d(x)\hfill & if\hfill & \varphi (x)<0,\hfill \end{array}$$
(55)
and $`s(x,t;ฯต)`$, a smooth function of $`t`$, such that on $`\mathrm{\Gamma }(t;ฯต)`$ it measures arclength from some point which moves normally to $`\mathrm{\Gamma }`$ as $`t`$ varies. The assumed initial smoothness of $`\mathrm{\Gamma }(t;ฯต)`$ implies that $`r`$ is a smooth function, at least, in a sufficiently small neighborhood of $`\mathrm{\Gamma }`$.
The outer expansions of $`\varphi `$ and $`u`$ are assumed to have the form
$$\varphi =\varphi (x,t;ฯต)=\varphi ^0(x,t)+ฯต\varphi ^1(x,t)+ฯต^2\varphi ^2(x,t)+๐ช(ฯต^3)$$
(56)
and
$$u=u(x,t;ฯต)=u^0(x,t)+ฯตu^1(x,t)+ฯต^2u^2(x,t)+๐ช(ฯต^3).$$
(57)
In order to determine the inner expansions we first define the inner variable
$$z(x,t;ฯต):=\frac{r(x,t;ฯต)}{ฯต}$$
(58)
and then assume the inner expansions to be given by
$$\varphi =\mathrm{\Phi }(z,s,t;ฯต)=\mathrm{\Phi }^0(z,s,t)+ฯต\mathrm{\Phi }^1(z,s,t)+๐ช(ฯต^2)$$
(59)
and
$$u=U(z,s,t;ฯต)=U^0(z,s,t)+ฯตU^1(z,s,t)+๐ช(ฯต^2).$$
(60)
The very definition of $`\mathrm{\Gamma }`$ requires $`\mathrm{\Phi }(0,s,t;ฯต)=0`$. In what follows we will use the following notation to refer to any variable $`g`$ evaluated by approaching $`\mathrm{\Gamma }`$ from either side ($`r>0`$ or $`r<0`$):
$$g_{\mathrm{\Gamma }^\pm }=\underset{r0^\pm }{lim}g(r,s,t;ฯต),$$
(61)
$$g_r_{\mathrm{\Gamma }^\pm }=\underset{r0^\pm }{lim}g_r(r,s,t;ฯต).$$
(62)
The following relations between the inner and outer variables obtained in \[kn:fif1\] are assumed to hold as $`\rho \pm \mathrm{}`$.
$$G^0(\rho ,s,t)=g^0(0^\pm ,s,t),$$
(63)
$$G^1(\rho ,s,t)=g^1(0^\pm ,s,t)+\rho g_r^0(0^\pm ,s,t).$$
(64)
### 4.2 Derivation of the equations of motion
Let us look at the system (LABEL:eq:defas2) assuming that $`u=0`$ initially and along the boundary.
#### 4.2.1 Outer problems
Substituting (56) and (57) into (LABEL:eq:defas2) and equating coefficients of the corresponding powers of $`ฯต`$ we obtain the $`๐ช(1)`$ and $`๐ช(ฯต)`$ outer problems respectively for points where the interface has not yet arrived:
$$\{\begin{array}{c}u_t^0=a_1\mathrm{\Delta }u^0,\hfill \\ \\ a_2f(\varphi ^0)=0,\hfill \end{array}$$
(65)
and
$$\{\begin{array}{c}u_t^1=a_1\mathrm{\Delta }u^1,\hfill \\ \\ a_2[f^{}(\varphi ^0)\varphi ^1+u^0]=0.\hfill \end{array}$$
(66)
The solution of (65), given the assumed initially and boundary conditions for $`u`$, is $`u^00`$, $`\varphi ^0=\pm 1`$. The solution of (66) is $`u^10`$, $`\varphi ^10`$.
#### 4.2.2 Inner problems
From the solution of the outer problems (65) and (66) and for points where the moving front has not yet arrived, we have
$$\chi ^0(t;w)=0,\text{and}v^0(t;w)=0,$$
(67)
and
$$\chi ^1(t;w)=0.\text{and}v^1(t;w)=0,$$
(68)
for all $`x\mathrm{\Omega }`$ and $`w(\mathrm{},\mathrm{})`$.
System (28) expressed in the $`(z,s,t)`$ coordinates (see appendix B), after differentiating the first and third equations with respect to $`t`$ and calling $`v=V(z,s,t;w,ฯต)`$ and $`\chi =\mathrm{{\rm Y}}(z,s,t;w,ฯต)`$, becomes
$$\{\begin{array}{c}r_t^2U_{zz}+ฯต[2r_tU_{zt}+r_{tt}U_z]+ฯต^2[U_{tt}+2s_tU_{st}+s_t^2U_{ss}+s_{tt}U_s+2r_ts_tU_{zs}+r_t^2\mathrm{\Phi }_{zz}]=\hfill \\ \\ =_{\mathrm{}}^{\mathrm{}}\overline{\alpha _1}(w)[ฯต^2V_t+ฯตr_tV_z+ฯต^2s_tV_s]๐w+๐ช(ฯต^3),\hfill \\ \\ ฯตr_tV_z+ฯต^2[V_t+s_tV_s+iwV]=U_{zz}+ฯต\mathrm{\Delta }rU_z+ฯต^2[U_{ss}|s|+U_s\mathrm{\Delta }s],\hfill \\ \\ ฯตr_t^2\mathrm{\Phi }_{zz}+ฯต^2[2r_t\mathrm{\Phi }_{zt}+r_{tt}\mathrm{\Phi }_z+2r_ts_t\mathrm{\Phi }_{zs}]=_{\mathrm{}}^{\mathrm{}}\overline{\alpha _2}(w)[ฯต\mathrm{{\rm Y}}_t+r_t\mathrm{{\rm Y}}_z+ฯตs_t\mathrm{{\rm Y}}_s]๐w+๐ช(ฯต^3),\hfill \\ \\ r_t\mathrm{{\rm Y}}_z+ฯต[\mathrm{{\rm Y}}_t+s_t\mathrm{{\rm Y}}_s+iw\mathrm{{\rm Y}}]=ฯต[\mathrm{\Phi }_{zz}+f(\mathrm{\Phi })]+ฯต^2[\mathrm{\Delta }r\mathrm{\Phi }_z+U]+๐ช(ฯต^3),\hfill \end{array}$$
(69)
for $`w(\mathrm{},\mathrm{})`$. From (67) and (68) we have
$$\underset{z\mathrm{}}{lim}\mathrm{{\rm Y}}^0(x,z,t;w)=0,\text{and}\underset{z\mathrm{}}{lim}V^0(x,z,t;w)=0,$$
(70)
and
$$\underset{z\mathrm{}}{lim}\mathrm{{\rm Y}}^1(x,z,t;w)=0,\text{and}\underset{z\mathrm{}}{lim}V^1(x,z,t;w)=0.$$
(71)
Substituting (59) and (60) into (69) and equating coefficients of the corresponding powers of $`ฯต`$ we obtain the $`๐ช(1)`$, $`๐ช(ฯต)`$ and $`๐ช(ฯต^2)`$ problems respectively.
$`๐ช(1)`$:
$$\{\begin{array}{c}r_t^2U_{zz}^0=0,\hfill \\ \\ U_{zz}^0=0,\hfill \\ \\ _{\mathrm{}}^{\mathrm{}}\overline{\alpha _2}(w)r_t\mathrm{{\rm Y}}_z^0๐w=0,\hfill \\ \\ r_t\mathrm{{\rm Y}}_z^0=0,\hfill \end{array}$$
(72)
$`๐ช(ฯต)`$:
$$\{\begin{array}{c}2r_tU_{zt}^0+r_t^2U_{zz}^1r_{tt}U_z^0=_{\mathrm{}}^{\mathrm{}}\overline{\alpha _1}(w)r_tV_z^0๐w=0,\hfill \\ \\ r_tV_z^0=U_{zz}^1+\mathrm{\Delta }rU_z^0,\hfill \\ \\ r_t^2\mathrm{\Phi }_{zz}^0=_{\mathrm{}}^{\mathrm{}}\overline{\alpha _2}(w)[\mathrm{{\rm Y}}_t^0+r_t\mathrm{{\rm Y}}_z^1+s_t\mathrm{{\rm Y}}_s^0]๐w,\hfill \\ \\ \mathrm{{\rm Y}}_t^0+r_t\mathrm{{\rm Y}}_z^1+s_t\mathrm{{\rm Y}}_s^0+iw\mathrm{{\rm Y}}^0=\mathrm{\Phi }_{zz}^0+f(\mathrm{\Phi }^0),\hfill \end{array}$$
(73)
$`๐ช(ฯต^2)`$:
$$\{\begin{array}{c}r_t^2U_{zz}^2+2r_tU_{zt}^1+r_{tt}U_z^1+U_{tt}^0+2s_tU_{st}^0+s_t^2U_{ss}^0+s_{tt}U_s^0+2r_ts_tU_{zs}^0\hfill \\ \\ +r_t^2\mathrm{\Phi }_{zz}^0=_{\mathrm{}}^{\mathrm{}}\overline{\alpha _1}(w)[V_t^0+r_tV_z^1+s_tV_s^0]๐w,\hfill \\ \\ r_tV_z^1+V_t^0+s_tV_s^0+iwV^0=U_{zz}^2+\mathrm{\Delta }rU_z^1+U_{ss}^0|s|+U_s^0\mathrm{\Delta }s,\hfill \\ \\ r_t^2\mathrm{\Phi }_{zz}^1+2r_t\mathrm{\Phi }_{zt}^0+r_{tt}\mathrm{\Phi }_z^0+2r_ts_t\mathrm{\Phi }_{zs}^0=_{\mathrm{}}^{\mathrm{}}\overline{\alpha _2}(w)[\mathrm{{\rm Y}}_t^1+r_t\mathrm{{\rm Y}}_z^2+s_t\mathrm{{\rm Y}}_s^1]๐w,\hfill \\ \\ \mathrm{{\rm Y}}_t^1+r_t\mathrm{{\rm Y}}_z^2+s_t\mathrm{{\rm Y}}_s^1+iw\mathrm{{\rm Y}}^1=\mathrm{\Phi }_{zz}^1+f^{}(\mathrm{\Phi }^0)\mathrm{\Phi }^1+\mathrm{\Delta }r\mathrm{\Phi }_z^0+U^0,\hfill \end{array}$$
(74)
for $`w(\mathrm{},\mathrm{})`$. We call
$$\alpha _i=\left(_{\mathrm{}}^{\mathrm{}}\overline{\alpha }_i(w)๐w\right)^1,$$
(75)
for $`i=1,2`$. A bounded solution of the first two equations in (72) satisfying the matching conditions (63) is $`U^10`$. From the third and fourth equations in (72) we have $`\mathrm{{\rm Y}}_z^0(x,z,t;w)=0`$ for $`w(\mathrm{},\mathrm{})`$. Integrating with respect to $`z`$ (assuming that $`r_t0`$) and applying condition (70) yields
$$\mathrm{{\rm Y}}^00,$$
(76)
for $`w(\mathrm{},\mathrm{})`$. Substituting $`U^10`$ in the first two equations in (73), multiplying the second and fourth equations in (73) by $`\overline{\alpha _1}`$ and $`\overline{\alpha _2}`$ respectively, and integrating with respect to $`w`$ yields
$$\{\begin{array}{c}(\alpha _1r_t^2)U_{zz}^1=0,\hfill \\ \\ (1\alpha _2r_t^2)\mathrm{\Phi }_{zz}^0+f(\varphi ^0)=0,\hfill \end{array}$$
(77)
Assuming that $`r_t^2\alpha _1`$, the bounded solution of the first equation in (77) satisfying the matching conditions (63) is $`U^10`$. To solve the second equation in (77), we assume that $`r_t^2\alpha _2`$ and we define the new variable
$$\xi :=\frac{z}{(1\alpha _2r_t^2)^{\frac{1}{2}}}.$$
(78)
In terms of $`(\xi ,s,t)`$, equation the second equation in (77) reads
$$\mathrm{\Phi }_{\xi \xi }^0+f(\mathrm{\Phi }^0)=0,$$
(79)
whose solution is $`\mathrm{\Phi }^0=\mathrm{\Psi }(\xi )`$, the unique solution of $`\mathrm{\Psi }^{\prime \prime }+f(\mathrm{\Psi })=0,\mathrm{\Psi }(\pm \mathrm{})=\pm 1,\psi (0)=0`$. Thus
$$\mathrm{\Phi }^0=\mathrm{\Phi }^0\left(\frac{z}{(1\alpha _2r_t^2)^{\frac{1}{2}}}\right),$$
(80)
which satisfies (63). From the fourth equation in (73) we have $`r_t\mathrm{{\rm Y}}_z^1=\alpha _2r_t^2\mathrm{\Phi }_{zz}^0`$. Integrating with respect to $`z`$ and applying condition (71) we get
$$\mathrm{{\rm Y}}^1=\alpha _2r_t^2\mathrm{\Phi }_z^0.$$
(81)
Substituting (80) into the third and fourth equations in (74), multiplying the fourth equation by $`\overline{\alpha _2}`$ and integrating with respect to $`w`$ yields
$$(1\alpha _2r_t^2)\mathrm{\Phi }_{zz}^1+f^{}(\mathrm{\Phi }^0)\mathrm{\Phi }^1=$$
$$=(\alpha _2r_{tt}+\gamma _2\alpha _2r_t\mathrm{\Delta }r)\mathrm{\Phi }_z^0++2r_t\mathrm{\Phi }_{zt}^0,$$
(82)
where
$$\gamma _2=i\alpha _2^2\left(_{\mathrm{}}^{\mathrm{}}\overline{\alpha _2}(w)w๐w\right).$$
(83)
Equation (82) expressed in the ($`\xi ,s,t`$) coordinate system reads (see appendix B)
$$\mathrm{\Phi }_{\xi \xi }^1+f^{}(\mathrm{\Phi }^0)\mathrm{\Phi }^1=$$
$$\frac{2r_t^2r_{tt}}{(1\alpha _2r_t^2)^{\frac{3}{2}}}(\xi \mathrm{\Phi }_{\xi \xi }^0+\mathrm{\Phi }_\xi ^0)+\frac{r_{tt}+\gamma _2r_t\mathrm{\Delta }r}{(1\alpha _2r_t^2)^{\frac{1}{2}}}\mathrm{\Phi }_\xi ^0.$$
(84)
It is straightforward to check that $`\mathrm{\Psi }^{}(\xi )`$ satisfies the homogeneous equation
$`\mathrm{\Phi }_{\xi \xi }^1+f^{}(\mathrm{\Phi }^0)\mathrm{\Phi }^1=0`$. That means that the operator $`\mathrm{\Lambda }:=\frac{^2}{\xi ^2}+f^{}(\mathrm{\Psi }^{}(\xi ))`$ has a simple eigenvalue at the origin with $`\mathrm{\Psi }^{}`$ as the corresponding eigenfuction. The solvability condition for the equation (84) now gives
$$\frac{2r_t^2r_{tt}}{(1\alpha _2r_t^2)^{\frac{3}{2}}}_{\mathrm{}}^{\mathrm{}}(\xi \mathrm{\Psi }^{\prime \prime }+\mathrm{\Psi }^{})\mathrm{\Psi }^{}๐\xi +$$
$$\frac{r_{tt}+\gamma _2r_t\mathrm{\Delta }r}{(1\alpha _2r_t^2)^{\frac{1}{2}}}_{\mathrm{}}^{\mathrm{}}(\mathrm{\Psi }^{})^2๐\xi =0.$$
(85)
A simple calculation shows that $`_{\mathrm{}}^{\mathrm{}}\xi \mathrm{\Psi }^{}\mathrm{\Psi }^{\prime \prime }๐\xi =\frac{1}{2}_{\mathrm{}}^{\mathrm{}}(\mathrm{\Psi }^{})^2๐\xi `$. Hence multiplying equation (85) by $`(1\alpha _2r_t^2)^{\frac{3}{2}}`$ and rearranging terms one obtains
$$\frac{r_{tt}}{1\alpha _2r_t^2}+\gamma _2r_t=\kappa .$$
(86)
Taking into consideration that on the interface $`\mathrm{\Delta }r=\kappa `$, the curvature of the interface, and that $`r_t=v`$, its normal velocity \[kn:fifpen1\], equation (86) becomes (12).
## 5 Conclusions
In this paper we showed that to the leading order and for a large class of kernels $`a`$, $`a_1`$ and $`a_2`$ under suitable assumptions on them, the law governing the evolution of interfaces for the integro-differential equation (1) is the same as for the differential equation (3). It is easy to see that $`\gamma _2`$ is given by
$$\gamma =\frac{a^{}(0)}{(a(0))^2}.$$
For equation (LABEL:eq:defas2) the result is similar with $`\gamma `$ and $`a`$ changed by $`\gamma _2`$ and $`a_2`$. Our derivation is only valid for advancing fronts; i.e., non-oscillating fronts in the sense that there exist points on the interface with vanishing velocity. This restriction prevent us from analyzing system where those oscillations may be relevant. This assumption is not necessary when the kernels are exponentially decreasing \[kn:rotnep1, kn:rotnep2\].
To solve the $`๐ช(ฯต)`$ problem we have assumed that $`|r_t|\sqrt{\alpha _1}`$ ($`\alpha _1`$ being similar to $`\alpha _2`$). On the other hand, assuming that $`1\alpha _2v^2>0`$ initially and that there is not change of sign in this direction during the evolution, a fact which is true for the circular case, we can easily see from equation (12) that $`|r_t|<\frac{1}{\sqrt{\alpha _2}}`$. Therefore, if $`\alpha _1\alpha _21`$ the former assumption does not add further restrictions on the interfacial motion.
Acknowledgement: The authors are indebted to A. Novick-Cohen for the formulation of the problem and to Ch. Charach for valuable discussions.
## Appendix A Cartesian coordinates
From Cartesian to ($`z,x,t`$) coordinates:
To go from Cartesian to ($`z,x,t`$) coordinates we transform derivatives as follows
$$\varphi _t=\mathrm{\Phi }_t\frac{1}{ฯต}S_t\mathrm{\Phi }_z,$$
$$\varphi _{tt}=\mathrm{\Phi }_{tt}\frac{1}{ฯต}2S_t\mathrm{\Phi }_{zt}+\frac{1}{ฯต^2}S_t^2\mathrm{\Phi }_{zz}\frac{1}{ฯต}S_{tt}\mathrm{\Phi }_z,$$
$$\varphi _{xx}=\mathrm{\Phi }_{xx}\frac{1}{ฯต}2S_x\mathrm{\Phi }_{zx}+\frac{1}{ฯต^2}S_x^2\mathrm{\Phi }_{zz}\frac{1}{ฯต}S_{xx}\mathrm{\Phi }_z,$$
$$\varphi _{yy}=\frac{1}{ฯต^2}\mathrm{\Phi }_{zz},$$
$$\varphi _{xxt}=\mathrm{\Phi }_{xxt}\frac{1}{ฯต}S_t\mathrm{\Phi }_{xxz}\frac{1}{ฯต}2S_{xt}\mathrm{\Phi }_{zx}\frac{1}{ฯต}2S_x\mathrm{\Phi }_{xzt}+\frac{1}{ฯต^2}2S_xS_t\mathrm{\Phi }_{xzz}+$$
$$+\frac{1}{ฯต^2}2S_xS_{xt}\mathrm{\Phi }_{zz}+\frac{1}{ฯต^2}S_x^2\mathrm{\Phi }_{zzt}\frac{1}{ฯต^3}S_x^2S_t\mathrm{\Phi }_{zzz}\frac{1}{ฯต}S_{xxx}\mathrm{\Phi }_z\frac{1}{ฯต}S_{xx}\mathrm{\Phi }_{zt}+\frac{1}{ฯต^2}S_{xx}S_t\mathrm{\Phi }_{zz},$$
$$\varphi _{yyt}=\frac{1}{ฯต^2}\mathrm{\Phi }_{zzt}\frac{1}{ฯต^3}S_t\mathrm{\Phi }_{zzz}.$$
From ($`z,x,t`$) to ($`\xi ,x,t`$) coordinates:
To go from ($`z,x,t`$) to ($`\xi ,x,t`$) coordinates, derivatives are transformed as follows
$$\xi _z=(1+S_x^2\alpha _2S_t^2)^{\frac{1}{2}},$$
$$\xi _x=z(1+S_x^2\alpha _2S_t^2)^{\frac{3}{2}}(S_xS_{xx}\alpha S_tS_{xt}),$$
$$\xi _t=z(1+S_x^2\alpha _2S_t^2)^{\frac{3}{2}}(S_xS_{xt}\alpha S_tS_{tt}),$$
$$\xi _{zx}=(1+S_x^2\alpha _2S_t^2)^{\frac{3}{2}}(S_xS_{xx}\alpha S_tS_{xt}),$$
$$\xi _{zt}=(1+S_x^2\alpha _2S_t^2)^{\frac{3}{2}}(S_xS_{xt}\alpha S_tS_{tt}).$$
$$\mathrm{\Phi }_z^0=\mathrm{\Phi }_\xi ^0\xi _z=\frac{1}{(1+S_x^2\alpha _2S_t^2)^{\frac{1}{2}}}\mathrm{\Phi }_\xi ^0$$
(87)
$$\mathrm{\Phi }_{zt}^0=\mathrm{\Phi }_{\xi \xi }^0\xi _z\xi _t+\mathrm{\Phi }_\xi ^0\xi _{zt}=\frac{S_xS_{xt}\alpha S_tS_{tt}}{(1+S_x^2\alpha _2S_t^2)^{\frac{3}{2}}}(\xi \mathrm{\Phi }_{\xi \xi }^0+\mathrm{\Phi }_\xi ^0).$$
(88)
$$\mathrm{\Phi }_{zx}^0=\mathrm{\Phi }_{\xi \xi }^0\xi _z\xi _x+\mathrm{\Phi }_\xi ^0\xi _{zx}=\frac{S_xS_{xx}\alpha S_tS_{xt}}{(1+S_x^2\alpha _2S_t^2)^{\frac{3}{2}}}(\xi \mathrm{\Phi }_{\xi \xi }^0+\mathrm{\Phi }_\xi ^0).$$
(89)
## Appendix B Local coordinates
### B.1 From Cartesian to ($`r,s,t`$) and $`(z,s,t)`$ coordinates
:
To go from Cartesian to $`(r,s,t)`$ coordinates we transform derivatives as follows ( using the fact that $`\left|r\right|1`$ )
$$\varphi _t=\mathrm{\Phi }_t+\mathrm{\Phi }_rr_t+\mathrm{\Phi }_ss_t,$$
$$\varphi _{tt}=\mathrm{\Phi }_{tt}+2\mathrm{\Phi }_{rt}r_t+2\mathrm{\Phi }_{st}s_t+\mathrm{\Phi }_{rr}r_t^2+\mathrm{\Phi }_{ss}s_t^2+\mathrm{\Phi }_rr_{tt}+\mathrm{\Phi }_ss_{tt}+2\mathrm{\Phi }_{rs}r_ts_t,$$
and
$$\mathrm{\Delta }\varphi =\mathrm{\Phi }_{rr}+\mathrm{\Phi }_r\mathrm{\Delta }r+\mathrm{\Phi }_{ss}\left|s\right|^2+\mathrm{\Phi }_s\mathrm{\Delta }s,$$
where the operators $``$ and $`\mathrm{\Delta }`$ refers only to the spatial variable $`x`$. These derivatives expressed in terms of $`(z,s,t)`$ are
$$\varphi _t=\mathrm{\Phi }_t+\frac{1}{ฯต}\mathrm{\Phi }_zr_t+\mathrm{\Phi }_ss_t,$$
$$\varphi _{tt}=\mathrm{\Phi }_{tt}+2\frac{1}{ฯต}\mathrm{\Phi }_{zt}r_t+2\mathrm{\Phi }_{st}s_t+\frac{1}{ฯต^2}\mathrm{\Phi }_{zz}r_t^2+\mathrm{\Phi }_{ss}s_t^2+\frac{1}{ฯต}\mathrm{\Phi }_zr_{tt}+\mathrm{\Phi }_ss_{tt}+2\frac{1}{ฯต}\mathrm{\Phi }_{zs}r_ts_t,$$
and
$$\mathrm{\Delta }\varphi =\frac{1}{ฯต^2}\mathrm{\Phi }_{zz}+\frac{1}{ฯต}\mathrm{\Phi }_z\mathrm{\Delta }r+\mathrm{\Phi }_{ss}\left|s\right|^2+\mathrm{\Phi }_s\mathrm{\Delta }s,$$
### B.2 From ($`z,s,t`$) to ($`\xi ,s,t`$) coordinates
To go from ($`z,s,t`$) to ($`\xi ,s,t`$) coordinates, derivatives are transformed as follows
$$\xi _z=(1\alpha r_t^2)^{\frac{1}{2}},$$
$$\xi _t=z(1\alpha r_t^2)^{\frac{3}{2}}\alpha r_tr_{tt}$$
$$\xi _{zt}=(1\alpha r_t^2)^{\frac{3}{2}}\alpha r_tr_{tt}.$$
$$\mathrm{\Phi }_z^0=\mathrm{\Phi }_\xi ^0\xi _z=\frac{1}{(1\alpha r_t^2)^{\frac{1}{2}}}\mathrm{\Phi }_\xi ^0$$
(90)
$$\mathrm{\Phi }_{zt}^0=\mathrm{\Phi }_{\xi \xi }^0\xi _z\xi _t+\mathrm{\Phi }_\xi ^0\xi _{zt}=\frac{\alpha r_tr_{tt}}{(1\alpha r_t^2)^{\frac{3}{2}}}(\xi \mathrm{\Phi }_{\xi \xi }^0+\mathrm{\Phi }_\xi ^0).$$
(91) |
warning/0002/cond-mat0002335.html | ar5iv | text | # Competing neural networks: Finding a strategy for the game of matching pennies
## I Introduction
Since the connection between disordered spin systems and symmetric binary neural networks was drawn intensive theoretical, numerical and experimental research has been devoted to this field within physics, and in the boundary of physics with biology and information theory, among others . From the viewpoint of the study of dynamical systems, neural networks are a special kind of distributed active systems , which in their most impressive realization โthe brainโ are able to display extremely sophisticated collective behavior. Actual models have of course much more modest scopes but, in spite of their simplicity, they have been able to imitate some basic features of cognitive processes. These models have also been extended to perform specific tasks, such as for instance process control and forecasting .
A basic capability of a wide class of neural-network models is that of learning, i.e. the possibility of modifying the internal architecture of the network to adapt its dynamics to an expected response. This process can take a variety of forms, to be chosen according to the aims of the model. Pattern storing and recognition โthe so-called associative memoryโ is perhaps the best known . Another well known instance is learning by generalization. In this case, the network is exposed to some input information and the output is compared with the expected response. Errors are usually backpropagated to modify the network dynamics through a change in its architecture. The network thus learns from experience. It is expected that after a certain learning transient the network is able to produce the correct output even from inputs not included in the learning sample. This kind of learning can be carried on under supervision, or the system can be designed to learn in an unsupervised manner, by means of a selforganization mechanism .
In this paper, we explore a neural-network model of the learning that takes place during a competitive game. Competitive games have recently attracted a great deal of attention among physicists as simple models of adaptive evolution and selforganization in biological, social, and economical systems . Neural networks have been designed and trained to play some highly complex games such as chess and backgammon . The complexity of these games, however, does not allow a systematic analysis of the learning process or a statistical evaluation of the performance accomplished. On the other hand, too simple games โsuch as those that admit a pure optimal strategy โ should be readily solved by a suitably designed neural network. In fact, finding a pure strategy can be associated with a maximization problem.
Here, we focus the attention at an intermediate level, choosing a competitive zero-sum game with very simple rules but lacking a pure optimal strategy, v.g. the game of matching pennies. Two neural networks are left to repeatedly play the game against each other. The successive game results are used on-line to feed the learning mechanism of the two players. As in the case of human players, each network tries to guess the strategy of its opponent and, thus, competition becomes a kind of mutual supervision. The optimal strategy for the game of matching pennies is a purely stochastic one. Thus, the challenge for the networks, whose dynamics is fully deterministic, consists in approximating as close as possible a random evolution. Our analysis of the time series generated during the game shows that even small networks with simple architectures do quite well โprobably better than any human being (not using a randomizing device) .
In the next section we describe in detail the game of matching pennies, and specify the architecture and learning dynamics of the competing neural networks. Section III is devoted to the study of the model as a time-discrete dynamical system โa mappingโ with emphasis in its phase-space evolution. In Sect. IV, we analyze statistical properties of the dynamics during the game, evaluating the performance of the networks within an information-theory approach. Finally, we discuss our results and consider some possible extensions.
## II The game and the players
In the game of matching pennies, player I chooses among two possible instances, say โheadsโ or โtails.โ Player II, not knowing player Iโs choice, also chooses either โheadsโ or โtails.โ Then, the two choices are disclosed โfor example, each player showing a pennyโ and, if they are the same, player I pays one cent to player II. If, on the contrary, the choices have been different, II pays one cent to I. The procedure is then repeated a large number of rounds, which has for instance been defined by a previous agreement between the players. In a less symmetric but very well-known realization of the same game, player II must guess in which hand has player I hidden a coin or any other small object. The pay-off rules are the same as for the game of matching pennies. Since at each round player Iโs loss (or gain) equals player IIโs gain (or loss), this is a zero-sum game. In game theory, a two-player zero-sum game is said to be a โstrictly competitiveโ game .
As the game proceeds, we expect the two players trying to outguess each other, keeping their own strategies secret. Due to the high symmetry of the game of matching pennies, however, there is no optimal pure strategy for either player. Of course, it would be a most poor strategy for any player to choose the same instance at every time step. But, moreover, any deterministic way of deciding which instance should be chosen at a given time step could be disclosed by the opponent in the long run. On the other hand, trying to guess the opponentโs strategy could lead an unsolvable, infinitely involved problem. As illustrated in , we may picture player I as thinking: โPeople usually choose heads; hence II will expect me to choose heads and choose heads himself, and so I should choose tails. But perhaps II is reasoning along the same line: heโll expect me to choose tails, and so Iโd better choose heads. But perhaps that is IIโs reasoning, soโฆโ In this way, it becomes impossible to determine a strategy in which either player could be confident. It follows that it is necessary for both players to introduce a mixed stochastic strategy where, at each time step, each player chooses an instance at random, with a certain probability distribution. The symmetry of the present game indicates clearly that the best strategy for both players is to choose heads or tails with equal probability. In the long run, this insures a zero average gain, whereas any other strategy implies a net gain for the opponent.
Our aim here is to study, as a dynamical system, a pair of competing neural networks playing the game of matching pennies. In particular, we are interested at analyzing whether the dynamics implies learning of an efficient strategy โon-line,โ i.e. as the game proceeds. Since the network dynamics and the learning algorithm considered in the following are deterministic, it cannot be expected that the networks will find the optimal (stochastic) strategy. However, it could be possible that the networks were able to approximate it by means of a complex deterministic dynamics over a sufficiently long period. The basic idea in the learning process is that the playing strategy of each network should emerge from trying to guess the opponentโs strategy. This is in fact the mechanism expected to drive the game between human players: though a general analysis of the game shows that the best way of playing is at random, each player tries to outguess the other assuming a deterministic strategy, at least, in the short term. The way of playing derives therefore from a (somewhat paradoxical) cooperative mechanism during the contest, where each player โsupervisesโ the learning of the other.
As for the architecture of each neural network, we take the simplest model, namely, the perceptron, introduced in and reviewed in standard books on neural networks (see, for example, ). It consists of a collection of $`N`$ inputs $`s_i(t)`$ and of $`N`$ synaptic weights $`w_i(t)`$ ($`i=1,\mathrm{},N`$) which define, at each time step, a single output $`\sigma (t)`$ as
$$\sigma (t)=S\left[\underset{i}{}w_i(t)s_i(t)\right].$$
(1)
Here $`S`$ is a step-shaped function, that we choose to be $`S(x)=\text{sign}(x)`$. Thus, $`\sigma =\pm 1`$. We associate each of this two possible values of the output with the instance chosen by the network at a given time step, say, $`\sigma (t)=+1`$ for heads and $`\sigma (t)=1`$ for tails.
We consider now two of these perceptrons (see Fig. 1), both with $`N`$ inputs. At each time step, the output of one of the perceptrons should be determined by the outputs of the other at the precedent steps. Indeed, this is the information available to each player on the strategy of the opponent. We associate therefore the inputs $`s_i^1`$ of perceptron I with the previous outputs $`\sigma ^2`$ of perceptron II and vice versa, as
$$s_i^{1,2}(t)=\sigma ^{2,1}(ti),$$
(2)
$`i=1,\mathrm{},N`$. Time steps are of unitary length.
Learning is a consequence of the comparison of the outputs of the two perceptrons at each time. If the outputs are identical perceptron II wins, and the synaptic weights $`w_i^1`$ of perceptron I are modified to produce a better prediction of the opponentโs output at the next round. Meanwhile, the synaptic weights $`w_i^2`$ of perceptron II can be left invariant, as they have led this perceptron to win the round. If, on the other hand, the outputs have been different, $`w_i^2`$ are modified and $`w_i^1`$ are maintained. A suitable algorithm for implementing this mechanism is the standard perceptron learning rule , which in our case implies
$$w_i^1(t+1)=w_i^1(t)\eta \mathrm{\Theta }[\sigma ^1(t)\sigma ^2(t)]s_i^1(t)\sigma ^2(t)$$
(3)
and
$$w_i^2(t+1)=w_i^2(t)+\eta \mathrm{\Theta }[\sigma ^1(t)\sigma ^2(t)]\sigma ^1(t)s_i^2(t),$$
(4)
with $`i=1,\mathrm{},N`$ and $`\eta =(1+N)^1`$. The Heaviside function $`\mathrm{\Theta }`$ โwhere $`\mathrm{\Theta }(x)=1`$ for $`x0`$ and $`\mathrm{\Theta }(x)=0`$ for $`x<0`$โ acts here as a mask, by selecting the perceptron whose synaptic weights are to be modified.
Suppose that the successive outputs of perceptron I are replaced by a periodic series of $`\pm 1`$ . From the viewpoint of perceptron II, this is interpreted as the opponentโs choice of a trivial strategy. As a matter of fact, the perceptron convergence theorem insures that if $`N`$ is large enough, i.e. if perceptron IIโs memory is sufficiently long-ranged, the learning procedure stops and, from then on, perceptron II wins all rounds. When the period of the output series of perceptron I is lower than $`N`$, in fact, it can be straightforwardly shown that there is at least one set of synaptic weights $`w_i^2`$ that make perceptron II able to win at every round. The number of steps needed to compute these synaptic weights is of order $`N^3`$ , and can be tested numerically in our system. It is therefore not expected that when two large perceptrons are left to play freely one of them will adopt a short-period strategy.
## III The system as a mapping: phase-space dynamics
Equations (1) to (4) define the dynamics of our system. They can be resumed in a $`4N`$-dimensional recursive mapping for the perceptron inputs and the synaptic weights only. The recursion equations are
$$\begin{array}{cc}\hfill s_1^{1,2}(t+1)=& S[_iw_i^{2,1}(t)s_i^{2,1}(t)],\hfill \\ \hfill s_i^{1,2}(t+1)=& s_{i1}^{1,2}(t)(i=2,\mathrm{},N),\hfill \\ \hfill w_i^1(t+1)=& w_i^1(t)\eta \mathrm{\Theta }[s_1^1(t+1)s_1^2(t+1)]\hfill \\ & \times s_1^1(t+1)s_i^1(t)(i=1,\mathrm{},N),\hfill \\ \hfill w_i^2(t+1)=& w_i^2(t)+\eta \mathrm{\Theta }[s_1^1(t+1)s_1^2(t+1)]\hfill \\ & \times s_1^2(t+1)s_i^2(t)(i=1,\mathrm{},N).\hfill \end{array}$$
(5)
The phase space corresponding to this mapping is discrete. In fact, the inputs $`s_i^{1,2}`$ can adopt the two values $`\pm 1`$ only. Moreover, $`w_i^{1,2}`$ can have real values but they vary on a discrete set, since according to Eqs. (3) and (4) the variation of the synaptic weights has always the same modulus, $`\left|\mathrm{\Delta }w^{1,2}\right|=\eta `$. Once the initial synaptic weights have been fixed, the discrete set of their possible future values is completely determined.
During the evolution, the synaptic weights can in principle run over an infinite set. However, though the synaptic weights are not expected to converge to fixed values but to continuously evolve as the game proceeds, it is reasonable to conjecture that they will not perform arbitrarily long excursions in phase space. To prove this conjecture, let us consider in detail the evolution of the synaptic weights, given by the two last equations in (5) or, equivalently, by Eqs. (3) and (4). These two equations can be written, respectively, as
$$\{\begin{array}{cc}w_i^1(t+1)=w_i^1(t)\eta \sigma ^1(t)s_i^1(t)\hfill & \text{if }\sigma ^1(t)=\sigma ^2(t)\hfill \\ w_i^1(t+1)=w_i^1(t)\hfill & \text{if }\sigma ^1(t)=\sigma ^2(t),\hfill \end{array}$$
(6)
and
$$\{\begin{array}{cc}w_i^2(t+1)=w_i^2(t)\eta \sigma ^2(t)s_i^2(t)\hfill & \text{if }\sigma ^1(t)=\sigma ^2(t)\hfill \\ w_i^2(t+1)=w_i^2(t)\hfill & \text{if }\sigma ^1(t)=\sigma ^2(t).\hfill \end{array}$$
(7)
We now select one of the perceptrons and restrict the dynamics of its synaptic weights to the time steps where they are effectively modified, by simply ignoring the steps where no changes occur. The evolution equations can be written in vectorial form as
$$๐ฐ(t+1)=๐ฐ(t)\eta S[๐ฐ(t)๐ฌ(t)]๐ฌ(t),$$
(8)
where the components of $`๐ฐ`$ and $`๐ฌ`$ are the synaptic weights and the inputs of the selected perceptron, respectively. We recall that $`S(x)`$ is the sign function. The scalar product $`๐ฐ๐ฌ`$ is defined in the usual way, cf. Eq. (1). Note that (8) holds for both perceptrons.
Let us now consider for a moment that, in Eq. (8), the vector $`๐ฌ`$ is independent of time. Under this assumption it is possible to reduce the system (8) to two equations for the quantities $`p(t)=๐ฌ๐ฐ(t)`$ and $`q(t)=|๐ฐ(t)|^2`$, namely,
$$\begin{array}{cc}p(t+1)\hfill & =p(t)(1\eta )S[p(t)]\hfill \\ q(t+1)\hfill & =q(t)2\eta |p(t)|+\eta (1\eta ).\hfill \end{array}$$
(9)
It can be easily seen from the first equation that $`p(t)`$ converges, after a certain transient, to a period-2 cycle. The two values of $`p`$ on this cycle, $`p_1`$ and $`p_2`$, satisfy the relation $`p_2=p_11+\eta `$. They depend on the initial conditions, but are always restricted to the intervals $`0<p_1<1\eta `$ and $`\eta 1<p_2<0`$. Accordingly, $`q(t)`$ oscillates between two values, $`q_1`$ and $`q_2`$, defined by the initial conditions and related by $`q_2=q_12\eta p_1+\eta (1\eta )`$. After the transient, the modulus of the vector $`๐ฐ`$ is therefore restricted to vary within the interval $`[W,W]`$ with $`W=\mathrm{max}\{\sqrt{q_1},\sqrt{q_2}\}`$.
In summary, for fixed $`๐ฌ`$ the evolution given by Eq. (8) drives the synaptic weights towards a bounded domain whose size depends on the initial condition but which is always finite. We stress that this is valid for any choice of $`๐ฌ`$. Coming now back to the case of variable inputs, we note that the number of possible values for $`๐ฌ(t)`$ is also finite, and equals $`2^N`$. Equation (8) can therefore be thought of as the application, at each time step, of one of the $`2^N`$ transformations just studied. Since each of them contracts the space of synaptic weights towards a bounded region, after the transient $`๐ฐ(t)`$ will always evolve within the union of all those regions. Disregarding transient effects, the space of synaptic weights is then finite. Hence, the accessible phase space of mapping (5) is finite and discrete.
As an illustration of the evolution of synaptic weights, we show in Fig. 2 the time dependence of $`|๐ฐ|`$ for both perceptrons. The initial weights were uniformly chosen at random in $`(0.2,0.2)`$, and $`N=10`$. The horizontal lines in the plot stand for the theoretical values of the temporal average of $`|๐ฐ|`$ for $`N=10`$ (full line) and $`N\mathrm{}`$ (dashed line). These can be calculated by taking the square of Eq. (8), namely,
$$q(t+1)=q(t)2\eta \sqrt{q(t)}|\widehat{๐ฐ}(t)๐ฌ(t)|+\eta (1\eta )$$
(10)
with $`\widehat{๐ฐ}=๐ฐ/|๐ฐ|`$. This recursion equation is analogous to the second of Eqs. (9). It can be seen that, for sufficiently large $`N`$, the average of $`|\widehat{๐ฐ}(t)๐ฌ(t)|`$ over time โor, equivalently, over random realizations of the vectors $`\widehat{๐ฐ}`$ and $`๐ฌ`$โ becomes independent of $`N`$ and approaches the limit $`|\widehat{๐ฐ}๐ฌ|=\sqrt{2/\pi }0.798`$. From Eq. (10), this implies that for large $`N`$:
$$|๐ฐ|=\sqrt{q}=\sqrt{\frac{\pi }{8}}0.627,$$
(11)
cf. . A better approximation for finite $`N`$ is $`|๐ฐ|=\sqrt{\pi N(N1)/8(N+1)^2}`$. For $`N=10`$ this gives $`|๐ฐ|0.540`$, which is the value plotted in Fig. 2. The average value of $`|๐ฐ|`$ provides an estimate for the size of the domain of phase space where the synaptic weights evolve after transients have elapsed. Note that the fact that $`|๐ฐ|`$ approaches a constant for large $`N`$ implies that, in average, the synaptic weights are $`w_i^{1,2}1/\sqrt{N}`$.
The main byproduct of the fact that for our system phase space is finite and discrete is that, after the transient has elapsed, the orbits will be periodic. It becomes therefore relevant to determine the length of the periods. In fact, if it resulted that orbits get typically trapped in short cycles, the problem would at once get uninteresting. We have measured the periods numerically, carrying out extensive series of $`100`$ to $`1000`$ realizations $`6.4\times 10^5`$ steps long, with $`N`$ ranging from $`2`$ to $`10`$. Initial conditions were chosen at random, with the synaptic weights uniformly distributed in $`(0.2,0.2)`$. The system has always been found to reach a periodic orbit for $`N<7`$. For a fixed value of $`N`$, periods show typically a broad distribution. The average period has been found to increase exponentially with $`N`$, as shown in Fig. 3. For $`N7`$, not all the realizations displayed periodicity, indicating the occurrence of periods longer than our numerical realizations. Some test realizations for $`N=10`$ suggest that periods could grow beyond $`10^8`$ steps.
The system thus seems to have two well-differentiated time scales. On the one hand, there should be a time scale associated with learning, of order $`N^3`$. As stated above, in the case of a single perceptron being trained to predict a periodic series this is in fact the number of steps needed to compute all the synaptic weights. For the competing perceptrons, the length of the initial transient during which the system explores phase space to find the bounded region where it will evolve later, should be of the same order. On the other hand, we have a much longer โrecursionโ time scale, of order $`A^N`$ ($`A5.05`$, Fig. 3), associated with the periods of orbits inside that region. Though the two-perceptron dynamics is dissipative, it resembles in this aspect that of Hamiltonian systems with many degrees of freedom. Indeed, according to Poincarรฉโs theorem , Hamiltonian systems are recurrent and, at sufficiently long times, they visit an arbitrarily small neighborhood of their initial state. However, in a statistical description of their evolution, it is possible to identify much shorter time scales, related to the relaxation of fast variables .
At the level of recursion time scales, the dynamics of the two-perceptron system is in a sense trivial. Orbits are in fact periodic at long times, and the results of successive game rounds will be repeated ad infinitum. When, during a whole period, one of the perceptrons is able to gain even the smallest advantage over the other, this small difference will continuously accumulate producing, in the long run, an arbitrarily large bias in the result of the game. As in the case of large Hamiltonian systems, however, recursion times are far beyond the reach of our (numerical) experience as the size of the perceptrons increases. Therefore, most of the realizations of the two-perceptron game analyzed bellow will always be restricted to the transient period, previous to the appearance of periodicity. In this stage, the relevant time scale is the learning time, of order $`N^3`$. Within such times we expect the system to reach a kind of stationary playing regime where, if the learning algorithm is efficient, the outputs of the two perceptrons should imitate a random series of $`\pm 1`$. In the next section we study the statistical properties of these output series.
## IV Statistical analysis of the game dynamics
Random properties in time series can be characterized in a variety of ways. In our case, where the relevant series are arrays of $`\pm 1`$, a suitable measure of time correlations is an information-like quantity . As shown below, this quantity can be used to characterize the correlation between different series and, consequently, the correlation of a series with itself. It has the advantage of being additive, and is therefore appropriate when comparing numerical results. We thus begin by defining the mutual information of two time series.
Consider two dichotomic stochastic processes $`S_1`$ and $`S_2`$ that, at each time step, can adopt the values $`\pm 1`$ with certain probability distributions. Let $`P(S_1,S_2)`$ be the joint probability for the processes, and $`P_1(S_1)=_{S_2}P(S_1,S_2)`$ and $`P_2(S_2)=_{S_1}P(S_1,S_2)`$ their individual (marginal ) probabilities. A measure of the correlation between the two processes is given by the mutual information , defined as
$$I=\underset{S_1=\pm 1}{}\underset{S_2=\pm 1}{}P(S_1,S_2)\mathrm{log}_2\left[\frac{P(S_1,S_2)}{P_1(S_1)P_2(S_2)}\right].$$
(12)
It can be shown that $`I0`$. For two uncorrelated processes, where $`P(S_1,S_2)=P_1(S_1)P_2(S_2)`$, the mutual information reaches its minimum, $`I=0`$. The maximal value of the mutual information is obtained for two identical stochastic processes, $`S_1=S_2`$, where $`I=P_1(+1)\mathrm{log}_2P_1(+1)P_1(1)\mathrm{log}_2P_1(1)`$. In particular, if $`P_1(+1)=P_1(1)=1/2`$, we get $`I=1`$.
The definition of mutual information, Eq. (12), suggests immediately a way of introducing a measure of autocorrelation for a single dichotomic stochastic process $`S`$ at different times. In fact, associating $`S_1(t)`$ and $`S_2(t)`$ with $`S(t)`$ and $`S(t+\tau )`$, respectively, we can introduce the (two-time) autoinformation as
$$\begin{array}{cc}\hfill I(t,\tau )=_{S(t)}_{S(t+\tau )}& P[S(t),S(t+\tau )]\hfill \\ & \\ & \times \mathrm{log}_2\left\{\frac{P[S(t),S(t+\tau )]}{P[S(t)]P[S(t+\tau )]}\right\}.\hfill \end{array}$$
(13)
If $`S`$ is a stationary stochastic process the autoinformation depends on the time interval $`\tau `$ only, $`II(\tau )`$. If the successive values of $`S`$ are uncorrelated we have $`I=0`$, whereas for $`\tau =0`$ we get the maximal value $`I(t,0)=P(+1)\mathrm{log}_2P(+1)P(1)\mathrm{log}_2P(1)`$.
In practice, for a finite realization of the stochastic processes, the probabilities involved in Eqs. (12) and (13) are approximated by the corresponding frequencies, which can be computed by simple counting of the relevant occurrences. This approximation implies that in the case of uncorrelated processes the information can differ from zero, due to fluctuations in the finite sample under consideration. It can be shown that for a $`T`$-step realization of uncorrelated stochastic processes where the individual probabilities of the two possible values $`\pm 1`$ are equal, $`P(+1)=P(1)=1/2`$, the probability distribution for the information to have a value $`I`$ is
$$p_T(I)=\sqrt{\frac{T\mathrm{ln}2}{\pi I}}\mathrm{exp}(TI\mathrm{ln}2),$$
(14)
for small $`I`$. The resulting mean value of the information is
$$I=_0^{\mathrm{}}Ip_T(I)๐I=\frac{1}{2T\mathrm{ln}2},$$
(15)
which decreases as $`T^1`$ as the series size grows. For large $`T`$, $`p_T(I)\delta (I)`$, as expected. Thus, the distribution of values for the information computed from finite samples of size $`T`$ and its average are to be respectively compared with $`p_T(I)`$ and $`I`$ in order to detect the presence of correlations.
We now consider two playing perceptrons with $`N=10`$, and apply the definition of autoinformation (13) to any of the two series of outputs, $`S(t)\sigma ^{1,2}(t)`$. The outputs are recorded after the first $`10^4`$ steps have elapsed, in order to avoid nonstationary transient effects during the first stage of learning (of order $`N^3`$ ). The recorded series are $`T=10^4`$ steps long, and the results presented below correspond to averages over $`5\times 10^4`$ realizations.
Figure 4 shows the measured average autoinformation as a function of $`\tau `$. The horizontal line corresponds to the average autoinformation (15) expected for an uncorrelated series with the present value of $`T`$, i.e. $`I7.2\times 10^5`$. We first note that, except for $`\tau =2`$ and $`4`$, the autoinformation of the output signal is always less that twice the value of $`I`$ for a random series. This implies that each perceptron exhibits a quite good performance in generating a random sequence. There are however certain regular patterns that suggest the presence of small but nontrivial correlations. Indeed, the average autoinformation oscillates strongly for small $`\tau `$, reaching high levels for even values of $`\tau `$ and dropping abruptly for odd values of $`\tau `$. On average, these oscillations decrease as $`\tau `$ grows, but they reappear near $`\tau =20`$ and $`30`$. Realizations for other values of $`N`$ indicate that the oscillation amplitude decreases as $`N`$ grows, and that the โburstsโ at which oscillations reappear occur when $`\tau `$ approaches integer multiples of $`N`$. The amplitude of these bursts decreases for larger multiples.
A more detailed description of the appearance of correlations in the output signals of the perceptrons is provided by the distribution of autoinformation values. Figure 5 displays the normalized frequencies of autoinformation values resulting from our sets of $`5\times 10^4`$ realizations of $`10^4`$-step series for various values of $`\tau `$. The curve corresponds to $`p_T(I)`$ for an uncorrelated series, Eq. (14). For $`\tau =1`$ practically no correlations are detected by the autoinformation. We note only a slight overpopulation for large $`I`$. On the other hand, for $`\tau =2`$, which corresponds to the largest deviation in the average autoinformation (see Fig. 4), the distribution is qualitatively different. It exhibits a maximum at a rather large value of the autoinformation ($`I2\times 10^3`$) and, except for small values of $`I`$, it is systematically much larger than the distribution expected for a random series. At $`\tau =10`$, the distribution has a profile similar to that observed for $`\tau =1`$, but the overpopulation at the tail is noticeably larger. This overpopulation grows further during the bursts where oscillations reappear. The plot for $`\tau =20`$ shows the distribution at the first of these bursts. In contrast, for the intermediate values at which the average autoinformation plotted in Fig. 4 reaches the information of a random series, the corresponding distribution cannot be distinguished from $`p_T(I)`$.
We have found that the oscillations of the average autoinformation shown in Fig. 4 are essentially a byproduct of the internal dynamics of each perceptron. In fact, if instead of using the opponentโs output, a perceptron is fed with a random series of $`\pm 1`$, the autoinformation of its own output oscillates as well. A detailed analysis of the output series reveals that, for even $`\tau `$, the product $`\sigma (t)\sigma (t+\tau )`$ is more frequently negative than positive. For instance, for $`\tau =2`$, the respective frequencies are about $`0.52`$ and $`0.48`$. We remark in passing that this small relative difference โof the order of a few percentโ produces an increment larger than one order of magnitude in the autoinformation, which evidences the sensibility of this quantity as a measure of correlations. For larger values of $`\tau `$, the difference is even smaller. On the other hand, for odd $`\tau `$ no differences are detected.
In order to trace the origin of the correlations observed for even $`\tau `$, a careful analysis of the learning algorithm has to be carried out. We consider first the case of $`\tau =2`$. After two time steps, the vector of synaptic weights can be written as
$$\begin{array}{cc}\hfill ๐ฐ(t+2)=& ๐ฐ(t)\eta \theta (t)\sigma (t)๐ฌ(t)\hfill \\ & \eta \theta (t+1)\sigma (t+1)๐ฌ(t+1),\hfill \end{array}$$
(16)
where $`\theta (t)=1`$ if the weights have been modified at time $`t`$, and $`\theta (t)=0`$ otherwise \[cf. Eq. (8)\]. When the perceptron if fed with a random signal, $`\theta (t)`$ can be seen as a stochastic process with equal probabilities for its two values. The product of the outputs two steps apart is
$$\begin{array}{ccc}\sigma (t)\sigma (t+2)\hfill & =\hfill & \sigma (t)S[๐ฐ(t+2)๐ฌ(t+2)]\hfill \\ & =\hfill & S[\sigma (t)๐ฐ(t)๐ฌ(t+2)\eta \theta (t)๐ฌ(t)๐ฌ(t+2)\hfill \\ & & \eta \theta (t+1)\sigma (t)\sigma (t+1)๐ฌ(t+1)๐ฌ(t+2)]\hfill \end{array}$$
(17)
Numerical measurements of the right-hand side (r.h.s.) of this equation show that the first two terms in the argument of the sign function have zero mean and do not produce a net contribution to the sign of $`\sigma (t)\sigma (t+2)`$. The only contribution to the correlation is originated in the third term. To verify this fact analytically, we first note that
$$\begin{array}{c}๐ฐ(t+1)=๐ฐ(t)[1+๐ช(1/\sqrt{N})],\hfill \\ ๐ฌ(t+1)๐ฌ(t+2)=๐ฌ(t)๐ฌ(t+1)[1+๐ช(1/\sqrt{N})].\hfill \end{array}$$
(18)
The first of these identities results from the fact that, as shown in the previous section, $`w_i1/\sqrt{N}`$ whereas, according to Eq. (5), its variation in one time step is given by $`\eta 1/N`$. The second identity can be readily proven from the evolution of $`s_i(t)`$, also given in Eq. (5). Consequently, neglecting terms of order $`1/\sqrt{N}`$, the sign of the product $`\sigma (t)\sigma (t+1)๐ฌ(t+1)๐ฌ(t+2)`$ can be approximated as follows:
$$\begin{array}{c}S[\sigma (t)\sigma (t+1)๐ฌ(t+1)๐ฌ(t+2)]\hfill \\ S\{[๐ฐ(t)๐ฌ(t)][๐ฐ(t)๐ฌ(t+1)][๐ฌ(t)๐ฌ(t+1)]\}.\hfill \end{array}$$
(19)
Note that the argument of the sign function in the r.h.s. of this equation is likely to be positive, since it is given by the product of the projections of two vectors, $`๐ฌ(t)`$ and $`๐ฌ(t+1)`$, along the direction of $`๐ฐ(t)`$ times their mutual scalar product. More explicitly,
$$\begin{array}{c}[๐ฐ(t)๐ฌ(t)][๐ฐ(t)๐ฌ(t+1)][๐ฌ(t)๐ฌ(t+1)]\hfill \\ =s_w^2(t)s_w^2(t+1)+s_w(t)s_w(t+1)๐ฌ^{}(t)๐ฌ^{}(t+1),\hfill \end{array}$$
(20)
with $`s_w=๐ฐ๐ฌ`$ and $`๐ฌ^{}=๐ฌs_w\widehat{๐ฐ}`$. The first term in the r.h.s. of this equation is always positive, whereas the second term is not expected to have a definite sign on average. Note moreover that the first term is of the order of unity, whereas the second term is of order $`\sqrt{N}`$. This implies that the relative importance of the positive contribution decreases as $`N`$ grows. Coming now back to Eq. (17) through Eq(19) it is clear that, when the synaptic weights are modified at time $`t+1`$ (i.e. $`\theta (t+1)=1`$), there is in average a negative contribution to $`\sigma (t)\sigma (t+2)`$, in agreement with numerical results. According to the above analysis, this correlation should become less important as $`N`$ grows. In fact, the autoinformation peak at $`\tau =2`$ is observed to decrease in the simulations.
For arbitrary $`\tau `$, the analysis can be repeated mutatis mutandis. We have
$$\begin{array}{cc}\hfill \sigma (t)\sigma (t+\tau )=S[& \sigma (t)๐ฐ(t)๐ฌ(t+\tau )\hfill \\ & \eta _{t^{}=0}^{\tau 1}\theta (t+t^{})\sigma (t)\sigma (t+t^{})\hfill \\ & \times ๐ฌ(t+t^{})๐ฌ(t+\tau )].\hfill \end{array}$$
(21)
Taking now into account that
$$\begin{array}{c}๐ฐ(t+t^{})=๐ฐ(t)[1+๐ช(\sqrt{t^{}/N})],\hfill \\ ๐ฌ(t+t^{})๐ฌ(t+\tau )=๐ฌ(t)๐ฌ(t+\tau t^{})[1+๐ช(\sqrt{t^{}/N})],\hfill \end{array}$$
(22)
the sign of the product $`\sigma (t)\sigma (t+t^{})๐ฌ(t+t^{})๐ฌ(t+\tau )`$ in the sum of Eq. (21) can be approximately written as
$$\begin{array}{c}S[\sigma (t)\sigma (t+t^{})๐ฌ(t+t^{})๐ฌ(t+\tau )]\hfill \\ S\{[๐ฐ(t)๐ฌ(t)][๐ฐ(t)๐ฌ(t+t^{})][๐ฌ(t)๐ฌ(t+\tau t^{})]\}.\hfill \end{array}$$
(23)
The argument of the sign function in the r.h.s. of this equation has a positive contribution of the same type as in Eq. (19) when $`t+t^{}=t+\tau t^{}`$, i.e. for $`\tau =2t^{}`$. Therefore, in the realizations where $`\theta (t+\tau /2)=1`$ a negative contribution to $`\sigma (t)\sigma (t+\tau )`$ appears. This of course requires $`\tau `$ to be even. Since other contributions have no definite sign, peaks in the average autoinformation are expected for even values of $`\tau `$, as observed. Note moreover that the order $`\sqrt{t^{}/N}`$ of the terms neglected in Eq. (23) increases with $`t^{}`$, i.e. with $`\tau `$. This explains why the height of the peaks decreases as $`\tau `$ grows.
Along the same line of analysis, it is possible to explain the bursts where the autoinformation peaks reappear. Now, however, it is necessary to take into account both perceptrons. In fact, the output of a single perceptron fed with a random signal does not exhibit such bursts. They are rather a consequence of the interaction between the two perceptrons during the game. The analysis, whose details we omit here, shows that bursts are originated by a kind of bouncing effect in the transmission of information between the opponents. This bouncing effects is attenuated as $`\tau `$ grows, and decreases for larger perceptrons, as observed in the numerical simulations.
In summary, the statistical analysis of perceptron outputs at time scales larger than the learning stage but much shorter than the recursion times, reveals that the perceptrons are quite efficient players of the game of matching pennies. Even with a relatively small number of inputs, i.e. with a relatively short-ranged memory, their dynamics is able to generate quasi-random mixed strategies. We recall that this behavior originates spontaneously from the deterministic learning algorithm with which each player is endowed to outguess its opponent. Remaining correlations, which could in principle be exploited by a โsmarterโ opponent to obtain a net gain during the game, are overall small and can in fact be reduced systematically by increasing the memory range.
## V Discussion
We have here considered an example of a fully deterministic learning system and explored its ability to behave stochastically. Concretely, we have coupled two deterministic perceptrons in such a way that they imitate two players of the game of matching pennies, trying to outguess each other. Since the optimal strategy for this game is a purely stochastic sequence of outputs, the learning process should lead the network dynamics to approach a random signal.
In the first place, we have observed that a perceptron producing a periodic signal can always be defeated by a sufficiently โsmartโ opponent, i.e. by a perceptron with a sufficiently large number of neurons. This kind of โdummyโ player provides in fact a linearly separable set of examples for the learning of its opponent . The learning task is thus to find a plane in the input space that separates the input states into two groups, namely those whose expected outputs are either $`+1`$ or $`1`$. On the other hand, when the two competing perceptrons are allowed to learn the situation is pretty much different. Since both networks are looking for the best performance, they both change their strategies on line and, thus, they may well provide not only a nonlinearly separable set of examples, but also an inconsistent one. That is to say, at two different times any perceptron can give two different outputs from the same input state. This is the reason why the learning process does, in fact, not converge, and why the system is expected to spontaneously develop stochastic-like dynamics.
Our main conclusion is that, despite the fact that the overall dynamics is in the long run periodic, the perceptrons do learn to behave quasi-stochastically over moderately long time intervals. An information-theoretical statistical analysis of the output signals shows slight time correlations, to be ascribed to the deterministic coupling between the learning mechanism and the outputs themselves, which act as the inputs of the respective opponents. The effect of these correlations is observed to decrease gradually as the number of neurons in each perceptron grows. Two seemingly paradoxical aspects of this learning process deserve to be pointed out, because of their suggestive similarity with learning in humans (or other animals) entrained in a systematic activity such as a repetitive competition game. In the first place, the mutual search for regularity in the opponentโs behavior leads the whole system to develop highly irregular evolution over long times, which can hardly be distinguished from purely random dynamics. In the second place, we stress that competition can here be interpreted as a form of mutually supervised learning and, thus, results in a kind of collaboration between the opponents.
Some natural extensions of the present model are worth considering for future work. An important question to be addressed regards the case where the entangled perceptrons are not equal in size, i.e. they have different numbers of neurons. In such a situation, in fact, the above quoted correspondence of competition and collaboration could fail to hold. Preliminary results along this line (not presented in this paper) suggest however that the advantage of a larger perceptron is relatively small. Only very small networks ($`N2`$) are systematically defeated by larger opponents, as they typically fall in short-period cyclic orbits.
The perceptron-like structure of our networks is probably the simplest instance among a large class of possible architectures. Fully connected networks and multilayer structures have been shown to exhibit very high performance in learning tasks . It would therefore be interesting to study how these more complex networks respond to mutually supervised learning. Finally, from the viewpoint of game theory, it would be relevant to analyze the dynamics of competing networks engaged in other games, especially, when ordinary optimization procedures do not lead to the optimal playing strategy. We mention, in particular, the iterated prisonerโs dilemma , which is attracting a great deal of attention as a paradigm of competition-collaboration interplay, and multiplayer minority games, recently studied by means of ensembles of globally coupled perceptrons . Competing neural networks could contribute to a better understanding of the complex learning mechanisms involved in such kind of social interactions. |
warning/0002/hep-th0002205.html | ar5iv | text | # 1 Introduction
## 1 Introduction
Over the few past years, there has been an increasing interest in studying the moduli space of vacua of the Coulomb and Higgs branches of supersymmeric gauge theories with eight supercharges in various dimensions. This interest is mostly due to the fact that extended supersymmetry severely restricts the quantum corrections to the moduli space metric and allows to make many exact computations . A large class of these gauge theories can be realized by brane configurations in type II strings on Calabi Yau manifolds by using either the Hanany -Witten method or the geometric engineering approach introduced and developed by Klemm ,Lerche , Mayr, Vafa and Warner and their collaborators; see ; see also .
Recently, a special interest has been given to the analysis of the hypermultiplet gauge invariant moduli space near the Higgs branch singularity. This analysis has been shown to be relevant for the study of many aspects in supersymmetric gauge theories with eight supercharges; in particular in the understanding of the assymptotic regions of the infrared IR low energy limits of the $`N=(4,4)`$ supersymmetric gauge theories in two dimensions especially within the so called throats of the Coulomb and Higgs branches where the theories are typically described by two $`2d`$ $`N=4`$ conformal Liouville theories with different central charges. In this context, it was shown in ; see also , that the IR limits of the Coulomb and Higgs branches have isomorphic throat regions associated with different small $`2d`$ $`N=4`$ subalgebras of the $`N=4`$ superconformal symmetry in two dimensions see also \[24-28\].
Vector and hypermultiplet moduli spaces have been also much studied in strings compactification to four dimensions where the low energy supergravity has scalar fields in both vector multiplets and hypermultiplets. The basic example of such compactification is given by type IIA string on $`R^{1,3}`$ times Calabi Yau threefolds which is believed to be dual to heterotic on $`R^{1,3}\times K3\times T^2`$, where the type IIA dilaton is in a hypermultiplet and the heterotic one is in a vector multiplet . Using local mirror symmetry and toric geometry methods, the absence of the type IIA dilaton in the vector multiplet has been exploited in to derive exact results in the Coulomb branch of type IIA on local Calabi Yau threefolds. In the same spirit, the hypermultiplet moduli space is independant of the heterotic string coupling and hence can be determined exactly from heterotic conformal field theory near the $`K3`$ ADE singularity. This issue has been analysed recently in and it was suggested that, in absence of small instanons, the hyperKahler moduli space for the heterotic string near the $`K3`$ ADE singularity is just the moduli space of vacua of a pure supersymmetric gauge theory in three dimensions with eight supercharges and ADE gauge group . Matter adjunction has been considered in . In the presence of small instantons, the hyperKahler moduli space has singularities which generally are interpreted in terms of singular conformal field theories or non perturbative massless particles. Note that contrary to the singularity of the coulomb branch generated by the one loop quantum corrections, the singularitiy of the Higgs branch of supersymmetric gauge theories with eight supercharges is not generated by quantum mechanics. It has been first suspected when tempting to understand the breakdown of string perturbation theory in type IIA on $`A_r`$ ALE surface, where the non perturbative phenomenon cannot be avoided by making the string coupling constant smaller . In lower dimensions, the Higgs branch singularity was also motivated by using duality between Higgs and Coulomb of $`N=(4,4)`$ supersymmetric gauge theories in two dimensions . More convincing and rigourous arguments for the existence of the Higgs branch throat are obtained from the study of the low energy physics of the D1/D5 system on $`X`$, which is equivalent to type IIB String theory on $`AdS_3\times S^3\times X`$ where $`X`$ is either $`T^4`$ or $`K3`$ .
In four dimensions, hypermultiplets moduli are moreover involved in the study of stringy instanton moduli space which is given by hyperkahler deformations of the classical instanton moduli space with non zero B field and small instanton singularites eliminated. Stringy instantons with non zero B field, including hyperkahler deformations resolving small instanton singularities, have been suggested in to be equally described as instantons on a non commutative space \[40,41 \]. Furthermore analysis involving hypermultiplets moduli is also encountered in the study of CFTโs obtained from compactifications of superstrings on Calabi Yau fourfolds and too particularly in the compactification of M-Theory on a Calabi Yau four-folds near the so called hyperkahler singularity .
In most of all of these studies, one of the basic eqs describing the moduli space of the hypermultiplets vacua reads, in the sigma model approach, as:
$$\underset{i=1}{\overset{n}{}}q_a^i[\phi _i^\alpha \overline{\phi }_{i\beta }+\phi _{i\beta }\overline{\phi }_i^\alpha ]=\stackrel{}{\xi _a}\stackrel{}{\sigma }_\beta ^\alpha ;a=1,\mathrm{},r.$$
(1)
Other basic eqs describing the throat region of the Higgs branch are given in section 5; see eqs(65,66). In eqs(1), the $`\phi _i^\alpha `$โs form a set $`\left\{\phi _i^\alpha ;1in\right\}`$ of $`n`$ component fields doublets $`\phi _i^\alpha `$ belonging to hypermultiplets and transforming in the $`(n,2)`$ representations of $`G\times SU\left(2\right)_R`$ group, where the group G will be specified later on. The $`\stackrel{}{\xi }_a`$โs are a collection of $`r`$ FI coupling 3-vectors, each of it transforms as a triplet under the usual $`SU\left(2\right)_R`$ symmetry rotating the eight supercharges. The $`q_a^i`$ parameters are the charges of $`\phi _i^\alpha `$โs under the $`U\left(1\right)^r`$ gauge group of the underlying supersymmetric gauge theory. For later use it is interesting to note that eqs(1) have a formal analogy with the following sigma model vaccum eqs of $`2d`$ $`N=2`$ supersymmetric $`U\left(1\right)^r`$ gauge theory involved in the analysis of the Coulomb branch of IIA superstrings on Calabi Yau threefolds with ADE singularities
$$\underset{i}{}q_a^i\left|X_i\right|^2=R_a;a=1,\mathrm{},r.$$
(2)
In eqs (2), the $`X_i`$โs are complex scalar fields , the $`R_a`$โs are FI couplings and the $`q_a^i`$โs are the $`U\left(1\right)^r`$ charges of the $`X_i`$โs which, for reference , read in the case of $`SU\left(n\right)`$ singularity as :
$$q_a^i=2\delta _a^i+\delta _a^{i1}+\delta _a^{i+1},$$
(3)
with the remarkable equality
$$\underset{i}{}q_a^i=0.$$
(4)
It is interesting to note here that in the $`2d`$ gauge theories with four supercharges, the above constraint eqs (4) is the condition under which the gauge theory flows in the infrared to $`2d`$ $`N=2`$ superconformal field theory . It is also the condition to have Kahler Calabi Yau backgrounds involved in superstring compactifications. Concerning eqs(1), we will see in section 4 that, under some assymptions, one can works out two remarkable classes of gauge invariant solutions of eqs (1) preserving the eight supercharges. The first class leads to the obtention of new singularities extending the usual $`N=2`$ ADE ones which are recovered as a special solutions by partial breaking of $`2d`$ $`N=4`$ supersymmetry down to $`2d`$ $`N=2`$. As preliminary results we find the following singular hypersurfaces:
$$\begin{array}{ccc}A_{n1}:U^{+\frac{n\left(n+1\right)}{2}}V^{+\frac{n\left(n+1\right)}{2}}=\left[Z^{+\left(n+1\right)}\right]^n\hfill & & \\ D_n:\left(x^{++}\right)^n+x^{++}\left(y^{+\left(n1\right)}\right)^2+\left(z^{+n}\right)^2=0\hfill & & \\ E_6:\left(x^{+6}\right)^2+\left(y^{+4}\right)^3+\left(z^{+3}\right)^4=0\hfill & & \\ E_7:\left(x^{+9}\right)^2+\left(y^{+6}\right)^3+y^{+6}\left(z^{+4}\right)^3=0\hfill & & \\ E_8:\left(x^{+15}\right)^2+\left(y^{+10}\right)^3+\left(z^{+6}\right)^5=0.\hfill & & \end{array}$$
(5)
In these eqs the charges carried by the various gauge invariants $`U,V,Z,x,y`$ and $`z`$ are Cartan charges of the $`U\left(1\right)_R`$ abelian subsymmetry of $`SU\left(2\right)_R`$ group. Eqs (5) extend the usual ADE complex surfaces of the $`N=2`$ backgrounds. For more details, see eqs (35-36). In the second class of gauge invariant solutions, we will show that eqs (4) are no longer constraint eqs; they are replaced by the following remarkable identity:
$$\underset{i}{}q_a^i+\underset{i}{}\left(q_a^i\right)=0,$$
(6)
which is usually fulfilled whatever the values of the $`q_a^i`$โs are. In the low energy limit, the above eqs lead then to $`2d`$ $`N=4`$ conformal models going beyond the $`N=2`$ ADE conformal ones.
The resemblance between eqs (1) and (2) is only formal, but turns out however to be very useful for the analysis of eqs(1) as well as their solving. Eqs(1) and (2) carry different meanings amongst which we quote the four followings:
(a) Eqs (2) describe the Coulomb branch leading to the well known gauge invariant Kahler moduli space while eqs (1) deal with the hypermultiplet branch and give a gauge invariant hyperkahler moduli space.
(b) For each $`U\left(1\right)`$ factor of the $`U\left(1\right)^r`$ gauge group, eqs (1) involve a triplet of FI parameters whereas eq (2) has only one. This feature goes with the previous one as it is related with the number of complex stucutres of Kahler and hyperkahler manifolds.
(c) Eqs (1) have a manifest $`SU\left(2\right)_R`$ symmetry which is absent in eqs (2). The latters have a $`U\left(1\right)`$ R symmetry. Eqs (1) are more restrictive than eqs(2) since they are the vaccum eqs of $`2d`$ $`N=4`$ supersymmetric linear $`\sigma `$ models. More precisely $`N=4`$ models in two dimensions involve three times the number of the D-flatness eqs of the $`N=2`$ models. This feature is easily seen on the space of the FI couplings which, for $`N=4`$, belong to $`\left(R^3\right)^rR_+^r\times \left(S^2\right)^r`$ while, for $`N=2`$, belong to $`R^r`$. The extra eqs associated with the moduli space $`\left(S^2\right)^r`$ are necessary conditions to have $`N=4`$ supersymmetry; without these constraints $`N=4`$ supersymmetry is partially broken down to $`N=2`$.
(d) It is now quite well established that eqs (1) and (2) hide moreover a comparable behaviour between the Coulomb and Higgs branches even near the singularity where eqs (1) and eqs (2) ceasse to be valid. We have already mentioned the duality between the two branches and their algebraic descriptions in terms of subalgebras of $`2d`$ $`N=4`$ conformal invariance. Later on, we shall give other arguments, geometrical and field theoretical, showing that in absence of FI couplings, $`\theta `$ terms and RR fields , both Coulomb and Higgs branches are described by singular CFTโs which seems to have something to do with $`2d`$ $`N=4`$ superconformal ADE Toda theories.
The formal similarity between eqs (1) and (2) together with the abovementioned features show that one may obtain new solutions of the gauge invariant moduli space of vacua of eqs (1) by using $`SU\left(2\right)_R`$ harmonic analysis and generalisations of methods of $`2d`$ $`N=2`$ supersymmetric linear sigma models.
The aim of this paper is to study these solutions and give interpretations in terms of blown up of singularities given by intersections of cotangent complex $`n`$ dimensional weighted projective spaces. Actually this study extends the results obtained for Coulomb branch of supersymetric gauge theories with four supercharges. Concerning the infrared dynamics of two dimensional $`N=(4,4)`$ gauge theories, we give also comments on the $`N=4`$ conformal Liouville description of the region in the viccinity of the singularity of the metric of the $`2d`$ $`N=4`$ Higgs branch generally interpreted as a semi infinite throat where the string coupling constant $`g_s=e^\varphi `$ blows up as the Liouville field $`\varphi `$ goes to infinity . Moreover, in an attempt towards an interpretation of the degenerate $`A_r`$ singularity carried by eqs (1), we give a field theoretical argument in favor to the hypothesis according to which the metric of the moduli space near the Higgs singularity maight be described by a $`N=4`$ conformal $`SU\left(r+1\right)`$ Toda theory in two dimensions. Of course this is just an observation which deserves in its own right a detailed study.
The presentation of this paper is as follows: In section 2, we review brefly the standard way used in handling eqs (1) where only a $`A_1`$ singularity has been considered, using the standard $`2d`$ $`N=2`$ supersymmetric analysis in which half of the eight supersymmetries are manifest. In this way of doing the $`SU\left(2\right)_R`$ symmetry is broken down to $`U\left(1\right)_R`$, a feature which is exploited in by making an appropriate choice of the FI 3-vector coupling where only one parameter is non zero. The two others are put to zero. In section 3, we develop a new way of doing by keeping all the three FI parameters non zero and the eight supercharges manifest. In this approach the $`SU\left(2\right)_R`$ symmetry is apparent but explicitly broken by the non zero FI terms. Our method enables us to exhibit manifestly the role of the three Kahler structures of the gauge invariant hyperkahler moduli spaces and permits moreovver to go beyond the $`A_1`$ singularity analysis of \[31,42 \]. Our way in handling eqs(1) involves two steps based on a geometric realisation of the $`SU\left(2\right)_R`$ symmetry and on the separation of the charges of the gauge and $`R`$-symmetries. The first step of this programme is described in Sectin 3 while the second step is studied in Section 4. The gauge and $`SU\left(2\right)_R`$ charge separation of the hypermultiplets moduli involves a parameter $`\gamma `$ taking the values $`\gamma =0`$ or $`\gamma =1`$ which distinguish two classes of solutions of eqs(1) both preserving the eight supercharges. For $`\gamma =0`$, we obtain a generalisation of the ADE complex surfaces reproducing the standard ones by partial breaking of $`2d`$ $`N=4`$ supersymmetry down to $`2d`$ $`N=2`$. For $`\gamma =1`$, we find new models which flow in the infrared to $`2d`$ $`N=(4,4)`$ scale invaraint models. In section 5, we study the moduli space of vacua of models with $`\gamma =1`$ by distinguishing the two cases $`\underset{i}{}q_a^i=0`$ and $`\underset{i}{}q_a^i0`$. We show by explicit computation that the hyperKahler moduli space, associated with eqs (1), is given by weighted complex projective spaces. Moreover, we study the Liouville description of the small instanton conformal theory near the singularity by using the field theoretical approach of Aharony and Berkooz and make comments regarding the $`A_r`$ singularity of eqs (1). In section 6, we give our conclusion.
## 2 Hyperkahler moduli space
In this section we review brefly the example of the hyperkahler cotangent bundle of complex projective space: $`T^{}\left(CP^2\right)`$; considered in the study of M theory on Calabi Yau fourfold after what we give the 2n complex dimension hyperkahler space $`T^{}\left(CP^n\right),n1`$, describing the instantons moduli space of one instanton on $`R^4`$ with gauge group $`U\left(n\right)`$ . To begin, note first of all that a Calabi Yau fourfolds can develop singularities of many types; this includes the $`\frac{C^4}{Z_4}`$ orbifold, the ADE hypersurface singularities considered recently in in the context of derivations of $`2d`$ CFTโs from type IIA string compactifications on Calabi Yau fourfolds , and the so called hyperkahler singularity we are intersted in here . To describe the $`T^{}\left(CP^2\right)`$ bundle, we consider $`2d`$ $`N=4`$ supersymmetric $`U\left(1\right)`$ gauge theory with one isotriplet FI coupling $`\stackrel{}{\xi }=(\xi ^1,\xi ^2,\xi ^3)`$ and three hypermultiplets of charges $`q_a^i=q^i=1;i=1,2,3`$ and taking G as $`SU\left(3\right)`$. The zero energy states of this gauge model are obtained by solving
$$\underset{i=1}{\overset{3}{}}\left[\overline{\phi }_{i\alpha }\phi _i^\beta +\phi _{i\alpha }\overline{\phi }_i^\beta \right]=\stackrel{}{\xi }\stackrel{}{\sigma }_\alpha ^\beta ,$$
(7)
which by the way is just a special situation of eqs (1) where all gauge charges $`q_a^j`$ are equal to one. Eqs (6) is a system of three eqs which,up to replacing the Pauli matrices by their expressions and using the $`SU\left(2\right)_R`$ transformations $`\phi ^\alpha =\epsilon ^{\alpha \beta }\phi _\beta `$ with $`\epsilon _{12}=\epsilon ^{21}=1`$ and $`\overline{\left(\phi ^\alpha \right)}=\overline{\phi }_\alpha `$, split as follows:
$$\begin{array}{ccc}\underset{j=1}{\overset{3}{}}\left(\left|\phi _j^1\right|^2\left|\phi _j^2\right|^2\right)\hfill & =\xi ^3& \hfill \left(a\right)\\ \underset{j=1}{\overset{3}{}}\phi _j^1\overline{\phi }_{j2}\hfill & =\xi ^1+i\xi ^2& \hfill \left(b\right)\\ \underset{j=1}{\overset{3}{}}\phi _j^2\overline{\phi }_{j1}\hfill & =\xi ^1i\xi ^2& \hfill \left(c\right).\end{array}$$
(8)
The moduli space of zero energy states of the classical gauge theory is the space of the solutions of eqs (7-8) divided by the action of the $`U\left(1\right)`$ gauge group. The solutions of eqs (7) depend on the values of the FI couplings. For the case where $`\xi ^1=\xi ^2=\xi ^3=0`$, the moduli space has an $`SU\left(3\right)\times SU\left(2\right)_R`$ symmetry; it is a cone over a seven manifold described by the eqs:
$$\underset{i=1}{\overset{3}{}}(\phi _{\alpha i}\overline{\phi }_i^\beta \phi _i^\beta \overline{\phi }_{\alpha i})=\delta {}_{\alpha }{}^{\beta }.$$
(9)
For the case $`\stackrel{}{\xi }\stackrel{}{0}`$, the abovementioned $`SU\left(3\right)\times SU\left(2\right)_R`$ symmetry is explicitly broken down to $`SU\left(3\right)\times U\left(1\right)_R`$. In the remarkable case where $`\xi ^1=\xi ^2=0`$ and $`\xi ^3`$ positive definite, it is not difficult to see that eqs (7) describe the cotangent bundle of $`CP^2`$. Indeed making the change
$$\psi _i=\frac{\phi _i^1}{\left[\underset{j=1}{\overset{3}{}}\left|\overline{\phi }_{j2}\right|^2+\xi ^3\right]^{\frac{1}{2}}},$$
(10)
and putting back into eq (8.a), one discovers that $`\psi _i`$โs satisfy $`_i\left|\psi _i\right|^2=1`$. The $`\psi _i`$โs parametrize the $`CP^2`$ space. On the other hand with $`\xi ^1=\xi ^2=0`$ conditions, eqs (8.b-c) may be interpreted to mean that $`\overline{\phi }_{2j}`$ lies in the cotangent space to $`CP^2`$ at the point determined by $`\psi _i`$. Although we are usually allowed to make the choice $`\xi ^1=\xi ^2=0`$ by using an appropriate $`SU\left(2\right)_R`$ transformation, we shall consider in sections 4 and 5, the generic cases where $`\xi ^1,\xi ^2`$ and $`\xi ^3`$ are all of them non zero as they form altogether the three Kahler parameters of hyperkahler manifolds. For the time being let us note that the previous analysis may be extended to the cases of $`2d`$ $`N=4`$ supersymmetric $`U\left(1\right)`$ gauge linear sigma model involving $`n+1`$ hypermutiplets with charges $`q^i=1;i=1,\mathrm{},n+1`$ and transforming in the fundamental representation of $`SU\left(n+1\right)`$. The vaccum energy equations of this $`U\left(1\right)`$ gauge model read as :
$$\begin{array}{ccc}\underset{j=1}{\overset{n+1}{}}\left(\left|\phi _j^1\right|^2\left|\overline{\phi }_{2j}\right|^2\right)\hfill & =\xi ^3& \\ \underset{j=1}{\overset{n+1}{}}\left(\phi _j^1\overline{\phi }_{2j}\right)\hfill & =\xi ^1+i\xi ^2.& \end{array}$$
(11)
For $`\xi ^1=\xi ^2=0`$ and $`\xi ^3`$ positive definite, the classical moduli space of the classical gauge theory is given by the cotangent bundle of complex $`n`$ projective space: $`T^{}CP^n`$ . For $`\xi ^1=\xi ^2=\xi ^3=0`$, one has just the conifold singularity of $`n`$ dimensional complex manifolds. Note also that near this singularity the low energy limits of this gauge theory is described by 2d N=(4,4) superconformal field theory of central charge $`C=6\left(n+11\right)=6n`$. In section 5, we shall turn to this point and describe the nature of the conformal field theory one has in the nearby of the Higgs branch singularity.
## 3 More on the Eqs (1)
Eqs (1) is a system of $`3r`$ equations or more precisely $`r`$ isovector eqs, each of which shares some general features with the usual $`2d`$ $`N=2`$ supersymmetic D-flatness conditions eqs (2); but in addition to the gauge charges, it carries a $`SU\left(2\right)_R`$ charge. To solve eqs (1), we shall use a different method than that used in . This method reproduces the solutions of the abovementioned study as particular cases and offers moreover a possiblity to address the question of multi instantons in $`R^4`$. Our approach is done in two steps and is based on the two following: First we combine methods of $`2d`$ $`N=2`$ supersymmetric sigma models, as used in describing the Kahler Coulomb branch of type IIA string on $`K3`$, and $`SU\left(2\right)_R`$ harmonic analysis allowing us to interpret $`SU\left(2\right)_R`$ representations as special functions on $`S_R^3=SU\left(2\right)_R`$. The index R carried by $`S_R^3`$, $`S_R^2`$ and $`U\left(1\right)_R`$ refers to the $`SU\left(2\right)_R`$. This step allows us to put eqs(1) into a manageable form which exhibit many similarities with eqs (2). Second , we introduce a convenient change of variables based on separating the $`U\left(1\right)_G^r`$ gauge charges and the $`U\left(1\right)_R`$ ones. This change of variables allows us to benifit from the similarities with the $`2d`$ $`N=2`$ supersymmetric gauge invariant backgrounds in order to study and solve eqs (1). In this section, we describe the first step of this programme; the charge separation of $`U\left(1\right)_G`$ and $`U\left(1\right)_R`$ symmetry factors will be studied in the next section. Our main purpose in what follows is to establish first that, up to $`SU\left(2\right)_R`$ transformations; eqs (1) can be rewitten in the following remarkable form:
$$\underset{j}{}q_a^j\phi _j^+\overline{\phi }_j^+=i\xi _a^{++};a=1,\mathrm{},r;$$
(12)
which , abstraction done of the plus indices describing the $`U\left(1\right)_R`$ Cartan charges carried by the $`\phi _j`$โs and the FI couplings, is comparable to eqs (2). In eqs (12); the moduli $`\phi _j^+`$, $`\overline{\phi }_j^+`$ and $`\xi _a^{++}`$ are related to $`\phi _j^\alpha `$, $`\overline{\phi }_j^\alpha `$ and $`\xi _a^{\left(\alpha \beta \right)}`$; they will be specified later on. To establish eqs (12) let first note that one may use the isomorphisms $`SU\left(2\right)_R=S_R^3`$ and $`\frac{SU\left(2\right)_R}{U\left(1\right)_R}S_R^2`$ to describe $`SU\left(2\right)_R`$ representations (both reducible and irreducible ) as harmonic functions on the sphere $`S_R^2`$ with definite $`U\left(1\right)_R`$ charge. This way of doing is well known in the study of $`SU\left(2\right)_R`$ representation theory; the main idea behind this construction may be summarized as follows: First, consider the following $`2\times 2`$ matrix
$$U=\left(\begin{array}{cc}u_1^+& u_2^+\\ u_1^{}& u_2^{}\end{array}\right)$$
(13)
and solve the isospin $`\frac{1}{2}`$ $`SU\left(2\right)_R`$ representation constraints namely the unimodularity $`detU=1`$ and the unitarity $`U^+U=U^+U=I`$ conditions. Straightforward algebra leads to:
$$\begin{array}{ccc}u^{\pm \alpha }=ฯต^{\alpha \beta }u_\beta ^\pm ;\overline{u}^{+\alpha }=u_\alpha ^{};ฯต_{\alpha \beta }=ฯต_{\beta \alpha }\hfill & & \\ u^{+\alpha }u_\alpha ^{}=1,u^{+\alpha }u_\alpha ^+=u^\alpha u_\alpha ^{}=0.\hfill & & \end{array}$$
(14)
Recall in passing that the $`u_\alpha ^\pm `$ harmonic variables are bosonic $`SU\left(2\right)_R`$ doublets which parametrize the unit $`S_R^3`$ sphere; they may be solved in terms of the standard $`S_R^3`$ variables $`\psi `$, $`\theta `$ and $`\varphi `$ as
$$\begin{array}{ccc}u_1^+=\mathrm{cos}\frac{\theta }{2}\mathrm{exp}\frac{i}{2}\left(\psi +\varphi \right)\hfill & & \\ u_2^+=\mathrm{sin}\frac{\theta }{2}\mathrm{exp}\frac{i}{2}\left(\psi \varphi \right)\hfill & & \\ u_1^{}=\mathrm{sin}\frac{\theta }{2}\mathrm{exp}\frac{i}{2}\left(\psi +\varphi \right)\hfill & & \\ u_2^{}=\mathrm{cos}\frac{\theta }{2}\mathrm{exp}\frac{i}{2}\left(\psi \varphi \right).\hfill & & \end{array}$$
(15)
We shall not use this realization hereafter as we shall take $`u_\alpha ^\pm `$ as our basic variables. Moreover, using the $`u_\alpha ^\pm `$ variables, the $`SU\left(2\right)_R`$ algebra is realized as differential operators on the space of harmonic functions on $`S_R^3`$:
$$\begin{array}{ccc}D^{++}=u^{+\alpha }\frac{}{u^\alpha };D^{}=u^\alpha \frac{}{u^{+\alpha }}\hfill & & \\ 2D^{++}=[D^0,D^{++}];2D^{}=[D^0,D^{}]\hfill & & \\ D^0=[D^{++},D^{}]=u^{+\alpha }\frac{}{u^{+\alpha }}u^\alpha \frac{}{u^\alpha }\hfill & & \end{array}$$
(16)
To study the $`SU\left(2\right)_R`$ representations by using the harmonic variables, it is more convenient to consider harmonic functions $`F^q\left(u_\alpha ^\pm \right)`$ with definte $`U\left(1\right)_R`$ charge q ; that is functions $`F^q\left(u_\alpha ^\pm \right)`$ satisfying the eigenfunction eq
$$[D^0,F^q]=qF^q.$$
(17)
These functions $`F^q`$ have a global harmonic expansion of total charge $`q`$ and carry $`SU\left(2\right)_R`$ representations. For example, taking $`q=2`$ and choosing $`F^{++}`$ as:
$$F^{++}\left(u_\alpha ^\pm \right)=u_{(\alpha }^+u_{\beta )}^+F^{\left(\alpha \beta \right)},$$
(18)
one sees that $`F^{++}`$ is the highest state of the isovector representation of $`SU\left(2\right)_R`$. This is also seen from the following eqs defining the highest states of $`SU\left(2\right)_R`$ of $`U\left(1\right)_R`$ charge equal to $`q`$
$$[D^0,F^q]=qF^q$$
(19)
$$[D^{++},F^q]=0.$$
Thus the harmonic functions $`F^{++}`$ altogether with $`F^0`$ and $`F^{}`$, defined as
$$F^0=[D^{},F^{++}]=u_{(\alpha }^+u_{\beta )}^{}F^{\left(\alpha \beta \right)}$$
(20)
$$F^{}=[D^{},F^0]=u_{(\alpha }^{}u_{\beta )}^{}F^{\left(\alpha \beta \right)},$$
form the three states of the isotriplet representation of the algebra (16) . In connection with the isotriplet representation $`\{F^{++}`$,$`F^0`$,$`F^{}\}`$, there is an interesting feature that we want to give at this level and which we will use later on when studying the solutions of eqs (12). This feature concerns the fact that one can usually realize $`F^q;q=0,\pm 2`$ as bilinears of isospinors $`f^+`$ and $`\overline{f}^+`$ as follows
$$\begin{array}{ccc}F^{++}=if^+\overline{f}^+=iu_{(\alpha }^+u_{\beta )}^+f^{(\alpha }\overline{f}^{\beta )}\hfill & & \\ F^0=\frac{i}{2}\left(f^+\overline{f}^{}+f^{}\overline{f}^+\right)=\frac{i}{2}u_{(\alpha }^+u_{\beta )}^{}f^{(\alpha }\overline{f}^{\beta )}\hfill & & \\ F^{}=if^{}\overline{f}^{}=iu_{(\alpha }^{}u_{\beta )}^{}f^{(\alpha }\overline{f}^{\beta )}.\hfill & & \end{array}$$
(21)
The complex number $`i`$ in front of the the factor of the right hand side of the above eqs ensures the reality condition of the isotriplet representation. Moreover the realization of $`F^q;q=0,\pm 2`$ as bilinears of isospinors reflects too simply the fact that the $`SU\left(2\right)_R`$ isovectors may be built from the symmetric product of the isospin $`\frac{1}{2}`$ representation and its conjugate. After this digression on the $`SU\left(2\right)_R`$ harmonic analysis, we turn now to eqs (1) which we write as:
$$\begin{array}{ccc}\underset{j}{}q_a^j\phi _j^+\overline{\phi }_j^+\hfill & =i\xi _a^{++}& \hfill \left(a\right)\\ \underset{j}{}q_a^j\left(\phi _j^+\overline{\phi }_j^{}+\phi _j^{}\overline{\phi }_j^+\right)\hfill & =2i\xi _a^0& \hfill \left(b\right)\\ \underset{j}{}q_a^j\phi _j^{}\overline{\phi }_j^{}\hfill & =i\xi _a^{}& \hfill \left(c\right).\end{array}$$
(22)
These eqs are obtained from eqs (1) by multipllying their both sides by $`u_{(\alpha }^+u_{\beta )}^+`$, $`u_{(\alpha }^+u_{\beta )}^{}`$ and $`u_{(\alpha }^{}u_{\beta )}^{}`$ respectively. Eqs (22) are also the D-flatness eqs one gets if one is using the $`2d`$ $`N=(4,4)`$ harmonic superspace formulation of $`2d`$ $`N=4`$ gauge theories . Thus like for eqs (1), eqs (22) form altogether a system of $`r`$ isovector eqs of the $`SU\left(2\right)_R`$ algebra (16); but with the remarkable difference that now it is enough to focuss attention on the highest weight states eqs (22.a). Knowing the solutions $`\phi _j^+`$ and $`\overline{\phi }_j^+`$ of eqs (20-a), one can also get the solutions of $`\phi _j^{}`$ and $`\overline{\phi }_j^{}`$ by acting on $`\phi ^+`$ and $`\overline{\phi }^+`$ by $`D^{}`$; namely:
$$\begin{array}{ccc}\phi _j^{}=[D^{},\phi _j^+]\hfill & & \\ \overline{\phi }_j^{}=[D^{},\overline{\phi }_j^+].\hfill & & \end{array}$$
(23)
In the end of this section, we would like to make two comments. The first comment is that one can use the isospinor bilinear realization of isotriplets eqs (21) to represent the Kahler parameters $`\xi _a^{++}`$ as follows:
$$\begin{array}{ccc}\xi _a^{++}=i\zeta _a^+\overline{\zeta }_a^+=iu_{(\alpha }^+u_{\beta )}^+\zeta _a^{(\alpha }\overline{\zeta }_a^{\beta )}\hfill & & \\ \zeta _a^\pm =u_\alpha ^\pm \zeta _a^\alpha ;\overline{\zeta }_a^\pm =u_\alpha ^\pm \overline{\zeta }_a^\alpha ,\hfill & & \end{array}$$
(24)
where $`\zeta _a^\alpha `$ and $`\overline{\zeta }_a^\alpha `$ may, roughly speaking, be viewed as the square roots of the FI couplings $`\xi _a^{\left(\alpha \beta \right)}`$. Similar relations involving $`\zeta _a^\pm `$ and $`\overline{\zeta }_a^\pm `$ for $`\xi _a^0`$ and $`\xi _a^{}`$ may be also written down. Putting back these relations in eqs (22.a), one gets
$$\underset{j}{}q_a^j\phi _j^+\overline{\phi }_j^+=\zeta _a^+\overline{\zeta }_a^+=u_{(\alpha }^+u_{\beta )}^+\zeta _a^{(\alpha }\overline{\zeta }_a^{\beta )}.$$
(25)
The second comment we want to do is that one can simplify further eqs (25) by making an extra change of variables which turns out to convenient when discussing the moduli space of gauge invariant vacua of eqs (1). This extra change consists to use the mapping $`R^3=R^+\times S^2`$ to write the FI isovectors $`\xi _a^{++}`$ as
$$\xi _a^{++}=R_a\eta _a^+\overline{\eta }_a^+=r_a^2\eta _a^+\overline{\eta }_a^+,$$
(26)
or equivalently by using the isospinors $`\zeta _a^\alpha `$ and $`\overline{\zeta }_a^\alpha `$ introduced previously:
$$\begin{array}{ccc}\zeta _a^\pm \hfill & =u_\alpha ^\pm \zeta _a^\alpha ;\overline{\zeta }_a^\pm =u_\alpha ^\pm \overline{\zeta }_a^\alpha & \hfill \left(a\right)\\ \zeta _a^\alpha \hfill & =r_a\eta _a^\alpha ;\overline{\zeta }_a^\alpha =r_a\overline{\eta }_a^\alpha & \hfill \left(b\right)\\ \zeta _a^\alpha \overline{\zeta }_{a\alpha }\hfill & =r_a^20& \hfill \left(c\right).\end{array}$$
(27)
Eqs (26) and (27) tell us that the $`R_a`$โs $`\left(R_a=r_a^20\right)`$ are the radial variables and the $`\eta _a^\alpha `$โs and $`\overline{\eta }_{a\alpha }`$โs, which satisfy
$$\eta _a^\alpha \overline{\eta }_{a\alpha }=1;\eta _a^\alpha \eta _{a\alpha }=\overline{\eta }_a^\alpha \overline{\eta }_{a\alpha }=0,$$
(28)
parametrize the two spheres $`S_a^2`$. The $`r_a^2`$ and $`\eta _a^\alpha `$ and $`\overline{\eta }_{a\alpha }`$ are in one to one correspondance with the $`r`$ FI isovectors. In other words eqs (28) describe a collection of $`r`$ unit two spheres which together with the $`r_a`$ conical variables of $`\left(R^3\right)^r`$ give the $`3r`$ parameters of the $`r`$ FI isovector couplings.
## 4 Separation of the charges of the $`U(1)_G`$ gauge and the $`U(1)_R`$ symmetries
Here we describe the separation of the gauge and $`U\left(1\right)_R`$ charges of the hypermultiplet scalar moduli $`\phi _j^+`$ and $`\overline{\phi }_j^+`$. As we mentioned earlier, this charge separation is the second step in our programme of finding the zero energy states of the classical $`2d`$ $`N=4`$ supersymmetric $`U\left(1\right)^r`$ gauge theory. Recall that the first step decribed in section 3 consists to interpret $`SU\left(2\right)_R`$ field representations as harmonic functions on $`S_R^2`$, fact which allowed us to put eqs (1) in a form similar to eqs(2) as shown on eqs (25). However the field variables of eqs( 25) still carry both $`U\left(1\right)_R`$ and $`U\left(1\right)^r`$ charges; these make their geometrical interpretations difficult and moreover do not allow to take advantage with their similarities with the $`N=2`$ supersymmetric D-flatness eqs (2) in looking for the solutions. Motivated by these two featurres, we have been lead to look for a way of rewriting eqs (25) so that the known results of $`2d`$ $`N=2`$ linear sigma models, including the geometric interpretation, can be exploited. We have found that this way may be achieved by introducing a parametrization of the hyperkahler moduli where the gauge charges and the $`U\left(1\right)_R`$ ones are separated but $`2d`$ $`N=4`$ supersymmetry still preserved. Note in passing that the idea of separating composite quantum numbers of fields is not a new idea; and corresponds just to a special feature of group representations theory. In the physics literature, the separation of the different charges of quantum fields has been used succesfully in various occasions in the past, in particular in coset models of $`2d`$ conformal field theory and in strongly correlated electrons models of low dimensional systems. One of the well known examples concerns complex spinors which may be separated into a spinon and a holon. For a review see . Thus our major aim in this section is to use this idea and try to solve these eqs by introducing the method of factorization of the two kinds of charges carried by the hyperKahler moduli. There are various ways one may follows in order to perform the separation of gauge and $`U\left(1\right)_R`$ charges carried by $`\phi _j^+`$โs and $`\overline{\phi }_j^+`$โs. A quite general way which allows to fulfill the four following requirements ((a), (b), (c) and (d)) is given by splitting $`\phi _j^+`$โs and $`\overline{\phi }_j^+`$โs as shown on eqs (29). These requirements are natural and may be stated as follows:
(a) The splitting should preserves supersymmetry, that is preserving the eight supercharges of the gauge theory.
(b)It should recover the results of \[31,42 \] summarized in section 2 but also extend the ADE models of $`2d`$ $`N=2`$ supersymmetric backgrounds , see also .
(c) It should give the standard ADE results up on breaking half of the eight supercharges.
(d) It should has a geometrical interpretation.
The factorization we propose is given by:
$$\begin{array}{ccc}\phi _j^+=X_j\eta _j^++\gamma Y_j\overline{\eta }_j^+\hfill & & \\ \overline{\phi }_j^+=\overline{X}_j\overline{\eta }_j^+\gamma \overline{Y}_j\overline{\eta }_j^+,\hfill & & \end{array}$$
(29)
where $`\eta _j^+`$ and $`\overline{\eta }_j^+`$ are as in eqs (27); that is:
$$\begin{array}{ccc}\eta _j^+=u_\alpha ^+\eta _j^\alpha ;\overline{\eta }_j^+=u_\alpha ^+\overline{\eta }_j^\alpha \hfill & & \\ \eta _j^\alpha \overline{\eta }_{\alpha j}=1;\eta _j^\alpha \eta _{\alpha j}=\overline{\eta }_j^\alpha \overline{\eta }_{\alpha j}=0\hfill & & \end{array}$$
(30)
and where $`X_j`$ and $`Y_j`$, $`j=1,\mathrm{},n`$ are complex fields carrying no $`U\left(1\right)_R`$ charges. The parameter $`\gamma `$ takes the values $`\gamma =0`$ or $`\gamma =1`$ and distinguish the two classes of solutions we will give hereafter. Note that similar decompositions to eqs (29) are also valid for $`\phi _j^{}`$ and $`\overline{\phi }_j^{}`$ and may be obtained from eqs (29) by acting on them by $`D^{}`$ as in eqs (23) . We will not use them in this discussion and then one can ignore them for the moment. Moreover the quantities $`X_j`$ , $`Y_j`$ and $`\eta _j^+`$ and $`\overline{\eta }_j^+`$ of the splitting (29) behave under $`U\left(1\right)^r`$ gauge and $`U\left(1\right)_R`$ transformations as follows:
$$\begin{array}{ccc}U\left(1\right)^r:X_jX_j^{}=\lambda ^{q_a^j}X_j\hfill & & \\ Y_jY_j^{}=\lambda ^{q_a^j}y_j\hfill & & \\ \eta _j^+\eta _j^+=\eta _j^+\hfill & & \\ \overline{\eta }_j^+\overline{\eta }_j^+=\overline{\eta }_j^+;\hfill & & \end{array}$$
(31)
and
$$\begin{array}{ccc}U\left(1\right)_R:X_jX_j^{}=X_j\hfill & & \\ Y_jY_j^{}=Y_j\hfill & & \\ \eta _j^+\eta _j^+=e^{i\theta }\eta _j^+\hfill & & \\ \overline{\eta }_j^+\overline{\eta }_j^+=e^{i\theta }\overline{\eta }_j^+.\hfill & & \end{array}$$
(32)
Actually eqs (31-32) define the factorization of the gauge charges and $`U\left(1\right)_R`$ ones. In what follows we shall use the splitting (29) to solve eqs (1) which, by help of the analysis of section 3, is also given by
$$\underset{j}{}q_a^j\phi _j^+\overline{\phi }_j^+=r_a^2\eta _a^+\overline{\eta }_a^+.$$
(33)
We shall give two classes of solutions of these eqs; the first class described a generalisation of the usual $`N=2`$ ADE ALE surfaces and the second class describes new models which flow in the infrared to $`2d`$ $`N=(4,4)`$ conformal field theories; the latters satisfy eqs (6).
### 4.1 Generalized ADE hypersurfaces
Here we would like to use the similarity betwen eqs (1) and eqs (2) in order to look for special solutions of eqs (1). These solutions are expected to describe generalisations of the standard eqs of ADE singularities associated with $`2d`$ $`N=4`$ linear sigma models. We find indeed that the moduli space of gauge invariant vacua of eqs (33) is appropriatly formulated in terms of harmonic variables of $`SU\left(2\right)_R`$ symmetry and the hypermultiplets vacua; see for instance eqs (35-36) for the case of a $`SU\left(n\right)`$ singularity. As a check consistency of of our results, we show that under some assumptions to be specified later on, the generalized ADE hypersurfaces we have obtained may be brought to the well known ADE models of $`2d`$ $`N=2`$ supersymmetric linear models. We show also, by explicit computation, that for the case of a $`SU\left(n\right)`$ singularity, the moduli space of gauge invariant vacua of eqs (33) is given by the usual ALE space with a $`SU\left(n\right)`$ singularity times the 2-sphere to power $`2n`$; i.e $`\left(S^2\right)^{2n}`$. Under the abovementioned assymption, this reduces to the usual ALE background times a 2-sphere. A similar result is also valid for the other singularities. To do so, let us first note that up on imposing the condition of ADE models
$$\begin{array}{ccc}_jq_a^j=0,a=1,\mathrm{},n1,\hfill & & \end{array}$$
(34)
one can imitate the analysis of $`2d`$ $`N=2`$ linear $`\sigma `$ models and build the gauge invariant moduli in terms of the $`\phi _j^+`$ fields. In the $`SU\left(n\right)`$ case for instance where $`q_a^j`$ is given by eq(3); there are three gauge invariant moduli; $`U^{+\frac{n\left(n+1\right)}{2}}`$ , $`V^{+\frac{n\left(n+1\right)}{2}}`$ and $`Z^{+\left(n+1\right)}`$ carrying $`\frac{n\left(n+1\right)}{2}`$ , $`\frac{n\left(n+1\right)}{2}`$ and $`\left(n+1\right)`$ $`U\left(1\right)_R`$ Cartan charges respectively. They are given by:
$$\begin{array}{ccc}U^{+\frac{n\left(n+1\right)}{2}}=\underset{j=0}{\overset{n}{}}\left(\phi _j^+\right)^{nj}\hfill & & \\ V^{+\frac{n\left(n+1\right)}{2}}=\underset{j=0}{\overset{n}{}}\left(\phi _j^+\right)^j\hfill & & \\ Z^{+\left(n+1\right)}=\underset{j=0}{\overset{n}{}}\left(\phi _j^+\right).\hfill & & \end{array}$$
(35)
They satisfy the following remarkable equation
$$U^{+\frac{n\left(n+1\right)}{2}}V^{+\frac{n\left(n+1\right)}{2}}=\left[Z^{+\left(n+1\right)}\right]^n$$
(36)
Eq (36) generalizes the usual equation of the ALE surface with $`SU\left(n\right)`$ singularity which, for later use, we recall it herebelow:
$$uv=z^n.$$
(37)
To better see the structure of eq (36), we use the splitting method of the charges of the $`\phi _j^+`$โs we have described earlier. Taking $`\gamma =0`$, the general splitting eqs (29) reduces to:
$$\phi _j^+=x_j\eta _j^+;\overline{\phi }_j^+=\overline{x}_j\overline{\eta }_j^+,$$
(38)
where $`X_j`$ and $`\eta _j^+`$ behave under gauge and $`U\left(1\right)_R`$ transformations as in eqs (31,32). Note by the way that like $`\phi _j^\pm `$, the realization $`X_j\eta _j^\pm `$ carries , for each value of j, four real degrees of freedom; two degrees come from $`X_j`$ and the two others from the parameters the 2- sphere described by $`\eta _j^\pm `$; eqs (30). Under $`2d`$ $`N=4`$ supersymmetric transformations which may be conveniently expressed as $`4d`$ $`N=2`$ supersymmetric transformations of fermionic parameters $`ฯต^\pm `$ and $`\overline{ฯต}^\pm `$, we have:
$$\delta \phi _j^+=ฯต^+\psi _j+\overline{ฯต}^+\overline{\chi }_j;$$
(39)
where $`\psi _j`$ and $`\overline{\chi }_j`$ are the Fermi partners of the $`\phi _j^\pm `$ scalars. $`(\phi _j^\pm ,\psi _j,\overline{\chi }_j)`$ constitute altogether the $`4d`$ $`N=2`$ free hypermultiplets. Using the splitting principle by factorizing $`ฯต^+`$ as $`ฯต\eta ^+`$ and $`\overline{ฯต}^+=\overline{ฯต}\overline{\eta }^+`$, and using eqs (38) and (39); we get
$$\eta _j^+\delta X_j+X_j\delta \eta _j^+=\overline{\eta }^+ฯต\psi _j+\overline{\eta }^+\overline{ฯต}\overline{\chi }_j,$$
(40)
or equivalenty
$$\eta _j^\alpha \delta X_j+X_j\delta \eta _j^\alpha =\overline{\eta }^\alpha ฯต\psi _j+\overline{\eta }^\alpha \overline{ฯต}\overline{\chi }_j.$$
(41)
Putting eqs (38) back into eqs (33), we get
$$\underset{j}{}q_a^j\left|X_j\right|^2\eta _j^+\overline{\eta }_j^+=r_a^2\eta _a^+\overline{\eta }_a^+$$
(42)
and
$$\begin{array}{ccc}U^{+\frac{n\left(n+1\right)}{2}}=uM^{+\frac{n\left(n+1\right)}{2}}\hfill & & \\ V^{+\frac{n\left(n+1\right)}{2}}=vN^{+\frac{n\left(n+1\right)}{2}}\hfill & & \\ Z^{+\left(n+1\right)}=zS^{+\left(n+1\right)};\hfill & & \end{array}$$
(43)
where $`u,v,z`$ and $`M^{+\frac{n\left(n+1\right)}{2}}`$, $`N^{+\frac{n\left(n+1\right)}{2}}`$ and $`S^{+\left(n+1\right)}`$ are gauge invariants given by
$$\begin{array}{ccc}u=\underset{j=0}{\overset{n}{}}X_j^j\hfill & ;N^{+\frac{n\left(n+1\right)}{2}}& \hfill =\underset{j=0}{\overset{n}{}}\left(\eta _j^+\right)^{nj}\\ v=\underset{j=0}{\overset{n}{}}X_j^{nj}\hfill & ;M^{+\frac{n\left(n+1\right)}{2}}& \hfill =\underset{j=0}{\overset{n}{}}\left(\eta _j^+\right)^j\\ z=\underset{j=0}{\overset{n}{}}X_j\hfill & ;S^{+\left(n+1\right)}& \hfill =\underset{j=0}{\overset{n}{}}\eta _j^+.\end{array}$$
(44)
Note that $`u,v`$ and $`z`$ verify the relation (37) and $`M^{+\frac{n\left(n+1\right)}{2}}`$ , $`N^{+\frac{n\left(n+1\right)}{2}}`$ and $`S^{+\left(n+1\right)}`$ satisfy eq (36). Eqs (42,43) may be brought to more familiar forms if we require moreover that the three following types of 2-spheres are identified:
(i) The $`\left(n+1\right)`$ 2- spheres parametrized by the $`\eta _j`$โs .
(ii) The $`\left(n1\right)`$ $`\eta _a`$ 2-spheres used in the parametrization of the FI couplings eqs (28).
(iii) The $`\eta ^+`$ 2-sphere involved in the factorization of the supersymmetric parameter $`ฯต^+`$ eq(40).
In other words, we require the following identity:
$$\eta _j^+=\eta _a^+=\eta ^+.$$
(45)
With this identification eqs, (42) reduce to the well known D- flatness conditions of the $`U\left(1\right)_R`$ gauge theory with four supercharges; namely:
$$\underset{j}{}q_a^j\left|X_j\right|^2=r_a^2.$$
(46)
Moreover eq (36) reduce to the usual ALE surface with $`SU\left(n\right)`$ singularity eq (37) since the $`M^{+\frac{n\left(n+1\right)}{2}}`$, $`N^{+\frac{n\left(n+1\right)}{2}}`$ and $`S^{+\left(n+1\right)}`$ gauge invariant become trivial as they are given by powers of $`\eta ^+`$ as shown here below.
$$\begin{array}{ccc}M^{+\frac{n\left(n+1\right)}{2}}=\left(\eta ^+\right)^{\frac{n\left(n+1\right)}{2}}=N^{+\frac{n\left(n+1\right)}{2}}\hfill & & \\ S^{+\left(n+1\right)}=\left(\eta ^+\right)^{n+1}.\hfill & & \end{array}$$
(47)
In the general case where the gauge charges $`q_a^j`$ of the $`X_j`$โs satisfy the consraints (4), eqs (46) is the vaccum energy of $`2d`$ $`N=2`$ supersymmetric linear $`sigma`$ models. Thus the classical moduli space M of the gauge invariant vacua of eqs (45,46) is then given by the 2-sphere parametrized by $`\eta ^+`$; eq (45), times the moduli space of the gauge invariant solutions of $`2d`$ $`N=2`$ supersymmetric vaccum energy states. In other words:
$$M=\frac{C^{n+1}}{C_{}^{}{}_{}{}^{n1}}\times S^2.$$
Note that the identification constraint eq (45) has a nice interpretation; it breaks explicitly half of the eight supersymmetries leaving then four supercharges preserved. These four supercharges are behind the reduction of eqs (42) down to eqs (46) leading to the standard ADE models. This feature is immediatly derived by combining eqs (41) and (45) as follows :
$$\eta ^\alpha \delta X_j+X_j\delta \eta ^\alpha =\overline{\eta }^\alpha ฯต\psi _j+\overline{\eta }^\alpha \overline{ฯต}\overline{\chi }_j.$$
(48)
Then multiplying both sides of this identity by $`\overline{\eta }_\alpha `$; one gets, after using eqs (30):
$$\delta X_j=ฯต\psi _j,$$
(49)
giving the usual supersymmetric transformations of the complex scalars of the $`2d`$ $`N=2`$ chiral multiplets. This completes the check of consistency of the generalised $`SU\left(n\right)`$ hypersurface singularity (36). Before going ahead let us summarize in few words what we have done until now. Starting from eqs (1) , we have shown that it is possible to put them into their equivalent form (33). The corresponding moduli space of gauge invaraint vacua is given by eq (36) which reduces to the standard ALE space with $`A_{n1}`$ singularity up on imposing the factorization eqs (38) and the conditions (45). The later breaks four supercharges among the original eight ones. The factorization (38) offers in turns a method for a geometric representation of hyperKahler backgrounds with eight supercharges. However we have not succeeded to solve directly eqs (33) nor (42) without breaking the eight supercharges. We will see later that it sitll possible to work out solutions with eight supercharges by using the general splitting (29) instead of the factorization (38) but still imposing eqs (45). Indeed to restore the eight supersymmetries by still using the constraint (45) we should take $`\gamma `$ non zero; say $`\gamma =1`$. Non zero $`\gamma `$ brings four extra supercharges which add to the old four existing ones carried by eq (38). This is easily seen from the above analysis and the splitting (38) where each part of the two terms of the right hand of eqs (39-41) carries four supersymmetries. We shall return to this feature with more details in the next subsection; for the time being we would like to make two comments regarding eqs (33).
(1) A naive analysis of eqs (33) suggests that the gauge invariant moduli space of vacua $`M`$ of eq (38) is given, for the generalized $`SU\left(n\right)`$ singularity eq (36), by the usual ALE space with $`SU\left(n\right)`$ singularity times $`2n`$ two-spheres. In other words:
$$M=\frac{C^{n+1}}{C_{}^{}{}_{}{}^{n1}}\times \left(S^2\right)^{2n},$$
where $`\left(n+1\right)`$ two spheres come from the $`\phi _j^+`$โs as shown in eqs (38) and $`\left(n1\right)`$ two spheres come from the FI couplings.
(2) As far eq (36) is concerned, one can also write down the generalized ADE models extending the usual $`N=2`$ ones . In addition to eq (36) which generalizes eq (37) , we have also
$$\left(x^{++}\right)^n+x^{++}\left(y^{+\left(n1\right)}\right)^2+\left(z^{+n}\right)^2=0,$$
describing the generalized $`D_n`$ singularity extending the standard ALE one namely:
$$x^n+xy^2+z^2=0.$$
More generally, we have the following results:
$$\left(x^{+6}\right)^2+\left(y^{+4}\right)^3+\left(z^{+3}\right)^4=0$$
$$\left(x^{+9}\right)^2+\left(y^{+6}\right)^3+y^{64}\left(z^{+4}\right)^3=0$$
$$\left(x^{+15}\right)^2+\left(y^{+10}\right)^3+\left(z^{+6}\right)^5=0.$$
These eqs extend respectively the following exceptional singulaerities
$$E_6:x^2+y^3+z^4=0$$
$$E_7:x^2+y^3+yz^3=0$$
$$E_8:x^2+y^3+z^5=0.$$
More informations about these extensions will be given in a future occasion.
### 4.2 Solutions with $`\gamma =1`$
Choosing $`\gamma =1`$ in the eqs (29) and putting back into eqs (33), we get a system of three eqs given by:
$$\begin{array}{ccc}\underset{j}{}q_a^j\left(\left|X_j\right|^2\left|Y_j\right|^2\right)\eta _j^+\overline{\eta }_j^+\hfill & =r_a^2\eta _a^+\overline{\eta }_a^+& \hfill \left(a\right)\\ \underset{j}{}q_a^j\left(X_j\overline{Y}_j\right)\eta _j^+\eta _j^+\hfill & =0& \hfill \left(b\right)\\ \underset{j}{}q_a^j\left(\overline{X}_jY_j\right)\overline{\eta }_j^+\overline{\eta }_j^+\hfill & =0.& \hfill \left(c\right)\end{array}$$
(50)
At this level no constraint has been imposed yet on the FI couplings contrary to the analysis of ref summarized in section 2. If moreover we require that all the two sphere $`\eta _j^+`$ and $`\eta _a^+`$ are identified as in eq (45); the above system reduces to
$$\begin{array}{ccc}\underset{j}{}q_a^j\left(\left|X_j\right|^2\left|Y_j\right|^2\right)\hfill & =r_a^2& \hfill \left(a\right)\\ _jq_a^j\left(X_j\overline{Y}_j\right)\hfill & =0& \hfill \left(b\right)\\ _jq_a^j\left(\overline{X}_jY_j\right)\hfill & =0.& \hfill \left(c\right)\end{array}$$
(51)
Eqs (51) have some remarkable features which have nice interpretations. Though the $`q_a^j`$ gauge charges of the hypermultiplets are not required to add to zero as in eq (4), eq (51.a) behave exactly as the D-flatness condition of $`2d`$ $`N=2`$ supersymmetric $`U\left(1\right)^r`$ gauge theory. The point is that eqs (51.a) involve twice the number of fields of eqs (2), but with opposite charges $`q_a^j`$ . Put differently; eq (51) involve two sets of fields $`X_j`$ and $`Y_j`$ of charge $`q_a^j`$ and ($`q_a^j`$) respectively. The sum of gauge charges of the $`X_j`$โs and $`Y_j`$โs add automatically to zero even though eq (4) is not fulfilled. Thus models with $`\gamma =1`$ flow in the IR to a $`2d`$ $`N=(4,4)`$ superconformal models extending the usual $`2d`$ $`N=(2,2)`$ ADE ones since the identities
$$\underset{j}{}q_a^j+\underset{j}{}\left(q_a^j\right)=0$$
(52)
go beyond the constraint eqs (4). Moreover eqs (51) may be fulfilled in different ways; either by taking all charges $`q_a^j`$ of the $`U\left(1\right)^r`$ gauge theory to be positive; say $`q_a^j=1;a=1,\mathrm{},r;j=1,\mathrm{},n`$, or part of the $`q_a^j`$โs are positive and the remaining ones are negative. In the case of a $`U\left(1\right)`$ gauge theory with $`\left(n+1\right)`$ hypermutiplets with gauge charges equal to one, eqs (51.a) describe a $`CP^n`$ manifold whereas eqs (51.b-c) which read as
$$\underset{j}{}X_j\overline{Y_j}=0,$$
(53)
together with their complex conjugate, show that the $`\overline{Y_j},`$โs are in the cotanget space of $`CP^n`$ at the point $`x_j=X_j/\left[\underset{i}{}\left|Y_i\right|^2+r_a^2\right]^{\frac{1}{2}}`$. Observe in passing that in case where some of the positive charges $`q_a^j`$ of $`U\left(1\right)^r`$ gauge theory are not equal to one, the corresponding moduli space is just the cotangent bundle of some weighted complex projective space, $`T^{}\left(WP^n\right)`$. Obesrve moreover that in the infrared limit this gauge theory flows to a $`2d`$ $`N=(4,4)`$ conformal field theory with central charge $`C=6n`$. In the next section we shall give some illustrating examples.
## 5 Moduli space of vacua of models with $`\gamma =1`$
In this section we want to study two types of vacua of the D-flatness conditions of $`2d`$ $`N=4`$ supersymmetric $`U\left(1\right)^r`$ gauge theory depending on the manner we deal with eqs(52). In other words starting from eqs(52),we develop hereafter two different, but equivalent, ways to solve them . These ways are associated with the value of the sum over the $`U\left(1\right)^r`$ charges of the hypermultiplet moduli that is ; $`\underset{i}{}q_a^i0`$ or $`\underset{i}{}q_a^i=0`$. To do so, we shall first study the case $`\underset{i}{}q_a^i0`$. We start by describing explicitly two examples after what we give the general sigma model result we have obtained and give also comments regarding the $`2d`$ $`N=4`$ Liouville description in the viccinity of these singularities. A similar analysis will be made for the other case $`\underset{i}{}q_a^i=0`$.
### 5.1 $`\underset{i}{}q_a^i0`$
A priori there are many ways to choose the $`q_a^i`$ charges such that $`\underset{i}{}q_a^i0`$; each of which corresponds to a definite model. A simple and instructif model is to consider a $`2d`$ $`N=4`$ supersymmetric abelian gauge theory with $`\left(r+1\right)`$ hypermultiplets whose scalar fields are denoted as $`\phi _j^+`$ and $`\overline{\phi }_j^+;j=0,1,\mathrm{},r`$. Using the splitting method described previously, we write the $`\phi _j^+`$โs and $`\overline{\phi }_j^+`$โs as
$$\begin{array}{ccc}\phi _j^+=X_j\eta ^++Y_j\overline{\eta }^+\hfill & & \\ \overline{\phi }_j^+=\overline{Y}_j\eta ^++\overline{X_j}\overline{\eta }^+,\hfill & & \end{array}$$
(54)
where their $`U\left(1\right)^r`$ charges are choosen as
$$q_i^j=\delta _{a1}^j+\delta _a^j.$$
(55)
Putting eqs (54) back into the D-flatness conditions (51) we get the following system of algebraic eqs
$$\begin{array}{ccc}\left|X_{a1}\right|^2+\left|X_a\right|^2\left(\left|Y_{a1}\right|^2+\left|Y_a\right|^2\right)=R_a\hfill & & \\ \underset{j}{}q_a^jX_j\overline{Y}_j=X_{a1}\overline{Y}_{a1}+X_a\overline{Y}_a=0\hfill & & \\ \underset{j}{}q_a^j\overline{X}_jY_j=\overline{X}_{a1}Y_{a1}+\overline{X}_aY_a=0.\hfill & & \end{array}$$
(56)
For later use , let us rewrite the two leading blocks of eqs of the above system, describing respectively models with $`U\left(1\right)`$ and $`U\left(1\right)^2`$ gauge groups associated with the values $`r=1`$ and $`r=2`$. For $`r=1`$, eqs (56) reduce to the three following eqs:
$$\begin{array}{ccc}\left|X_0\right|^2+\left|X_1\right|^2\left(\left|Y_0\right|^2+\left|Y_1\right|^2\right)=R_1\hfill & \left(a\right)& \\ X_0\overline{Y}_0+X_1\overline{Y}_1=0\hfill & \left(b\right)& \\ \overline{X}_0Y_0+\overline{X}_1Y_1=0.\hfill & \left(c\right)& \end{array}$$
(57)
Similarly we have, for the $`U\left(1\right)^2`$ gauge model, a system of six equations; three of them coincide with those given by eqs (57); the others are as follows:
$$\begin{array}{ccc}\left|X_1\right|^2+\left|X_2\right|^2\left(\left|Y_1\right|^2+\left|Y_2\right|^2\right)=R_2\hfill & \left(a\right)& \\ X_1\overline{Y}_1+X_2\overline{Y}_2=0\hfill & \left(b\right)& \\ \overline{X}_1Y_1+\overline{X}_2Y_2=0\hfill & \left(c\right)& \end{array}$$
(58)
To solve eqs (56), we shall adopt the following strategy. We shall first consider the solving of eqs (57), then we treat both eqs (57) and (58), after what we give the general solutions for eqs (56) and finally make some comments regarding the Liouville description of the singularities of the metric of the Higgs branch. For the $`2d`$ $`N=4`$ supersymmetric model with one $`U\left(1\right)`$ gauge factor, one should note first of all that the moduli space of gauge invariant vacua is a complex surface which becomes singular when $`R_1`$ vanishes. It is just the cotangent line bundle of the two- sphere $`S^2`$; $`T^{}\left(CP^1\right)`$. A naive way to see this feature is to set $`Y_0=Y_1=Y`$; a choice which reduces eq (57.a) to the following well known eq of $`N=2`$ linear sigma models with four supercharges
$$\left|X_0\right|^2+\left|X_1\right|^22\left|Y\right|^2=R.$$
(59)
This eq describes the blow up of the $`SU\left(2\right)`$ singularity of the ALE complex surface $`\frac{C^2}{Z_2}`$. An other way to deal with this singularity is to make the change of variables preserving the eight supercharges
$$\begin{array}{ccc}x_0=X_0\left[R_1+\left|Y_0\right|^2+\left|Y_1\right|^2\right]^{\frac{1}{2}}\hfill & & \\ x_1=X_1\left[R_1+\left|Y_0\right|^2+\left|Y_1\right|^2\right]^{\frac{1}{2}},\hfill & & \end{array}$$
(60)
leading to
$$\begin{array}{ccc}\left|x_0\right|^2+\left|y_1\right|^2=1\hfill & \left(a\right)& \\ x_0\overline{y}_0+x_1\overline{y}_1=0\hfill & \left(b\right)& \\ x_0y_0+\overline{x}_1y_1=0\hfill & \left(c\right).& \end{array}$$
(61)
Eqs (57.b-c), which by the way, are exchanged under complex conjugation, have a geometric meaning; they show that at each point $`(x_0,x_1)`$ of the base manifold $`B_1`$ there is an orthogonal fiber $`F_1`$ parametrized by $`(\overline{Y}_0,\overline{Y}_1)`$, defining altogether the cotangent line bundle $`T^{}CP^1`$. For non zero values of $`R_1`$ where the change (60) is well defined, eq (57.a) is a 2-sphere and then the bundle is smooth. For $`R_1`$ equals to zero, the change (60) falls down at the origin $`X_0=X_1=Y_0=Y_1=0`$ and the bundle becomes singular. Note that according to the ADHM construction, the moduli space of gauge invariant vacua of the $`U\left(1\right)`$ gauge model with two (or more) hypermultiplets is just the moduli space of small instantons on $`R^4`$. For $`R_1`$ positive definite, the small instanton singularity is blown up and in the limit $`R_1=0`$ the singularity is recovered. In two dimensions, it has been shown moreover that in this limit the fields $`X_0`$, $`X_1`$, $`Y_0`$ and $`Y_1`$ do not give a good description of the small instanton conformal theory near the singularity. The appropriate variables in this region turns out to be those of a $`2d`$ $`N=4`$ conformal Liouville field theory \[17,18 \]. To see this remarkable feature, it is interesting to use the field theoretical approach of Aharony and Berkooz \[18 \] regarding the study of the low energy lmits of $`2d`$ $`N=4`$ gauge theories whose lagrangian $`L=L_{gauge}+L_H`$ reads as:
$$\begin{array}{ccc}L=\frac{1}{4g_{YM}^2}d^2xtr\left(F_{\mu \nu }^2+\left(D_\mu V\right)+[V,V]^2+\overline{\psi }_V\gamma ^\mu D_\mu \psi _V+\overline{\psi }_V[V,\psi _V]+\stackrel{}{D}^2\right)+\hfill & & \\ d^2x\underset{hypermult}{}\left(\right|D_\mu \phi _H|^2+|V\phi _H|^2+\overline{\psi }_H\gamma ^\mu D_\mu \psi _H+\overline{\psi }_HV\psi _H+\overline{\psi }_VV\psi _H\hfill & & \\ +\overline{\phi }_HD\phi _H)\hfill & & \end{array}$$
(62)
In this formal eq $`D_\mu =\left(_\mu +A_\mu \right)`$ is the covariant derivative, $`(V_{A\overline{A}},A_\mu ,D^{\left(\alpha \beta \right)})`$ and $`(\psi _V^A,\psi _V^{\overline{A}})`$ are respectively the bosonic and fermionic fields of the vector multiplet carrying amongst others quantum charges of the $`SO\left(4\right)\times SU\left(2\right)_RSU\left(2\right)_r\times SU\left(2\right)_l\times SU\left(2\right)_R`$ . Note that $`SU\left(2\right)_r\times SU\left(2\right)_l\times SU\left(2\right)_R`$ is the R-symmetry of the gauge theory with eight supercharge in two dimensions. Note also that the indices $`A`$,$`\overline{A}`$ and $`\alpha `$ refer to the isospin $`\frac{1}{2}`$ representation of $`SU\left(2\right)_l,SU\left(2\right)_r`$ and $`SU\left(2\right)_R`$ respectively. Note Moreover that the $`\phi _H`$ scalars and their fermionic partners $`(\psi _H^L,\psi _H^R)`$ stand for the fields of the hypermutiplets. Following , the low energy limit of this gauge theory, which involves taking $`g_{YM}\mathrm{}`$, is described by two decoupled $`2d`$ $`N=(4,4)`$ superconformal field theories; one describing the Higgs branch whose central charge $`C_H=6\left(n_Hn_V\right)`$ where $`n_H`$ the number of hypermultiplets and $`n_V`$ is the number of vector multiplets. The other conformal field theory corresponds to the Coulomb Branch of central charge $`C_V=6`$. An argument supporting this paticular feature comes from the analysis of the R-symmetries of the $`N=(4,4)`$ superconformal algebra which includes left and right moving $`su\left(2\right)`$ Kac Moody subalgebras. The R-symmetry of the Higgs branch is exactly $`SU\left(2\right)_l\times SU\left(2\right)_rSO\left(4\right)`$ encountered earlier while the R-symmetry of the Coulomb branch is given by a non visible group $`SO\left(4\right)SU\left(2\right)_l\times SU\left(2\right)_r`$ containing $`SU\left(2\right)_R`$ as a diagonal subgroup. Since the Coulomb and Higgs superconformal theories have different R-symmetries; they cannot be identified. Moreover taking the naive limit $`g_{YM}\mathrm{}`$ in eq (62), one sees $`L_{gauge}`$ is removed and the lagrangian of the low energy gauge theory is reduced to $`L_H`$; the lagrangian of the Higgs branch namely:
$$\begin{array}{ccc}L_H=d^2x\underset{hypermult}{}\left[\right|D_\mu \phi _H|^2+|V\phi _H|^2+\overline{\psi }_H\gamma ^\mu D_\mu \psi _H+\overline{\psi }_HV\psi _H+\overline{\psi }_VV\psi _H\hfill & & \\ +\overline{\phi }_HD\phi _H]\hfill & & \end{array}$$
(63)
where now the vector multiplet fields are auxiliary fields which may be eliminated through their eqs of motion. However following see also , it more useful to regard the vector multiplet fields as the basic objects instead of the matter fields and integrating over the hypermultiplet fields in order to describe the behavior near the singularity in the moduli space. In this lagrangian approach, one obtains an induced effective action of the vector mutiplet fields which describe the region near the singularity of the Higgs branch $`(\phi 0`$ or $`V\mathrm{})`$. In the case of supersymmetric $`U\left(1\right)`$ gauge theory with $`N_f`$ hypermultiplets and one vector multtiplet , supersymmetry and $`SO\left(4\right)`$ symmetry constraint the metric of the four gauge scalar fields $`\left(V_i\right)=(V_1,V_2,V_3,V_4)`$ in the vector multiplet to be of the form:
$$ds^2=N_f\frac{1}{\left[\left(V_1\right)^2+V_2^2+V_3^2+V_4^2\right]}\left[\left(dV_{1}^{}{}_{}{}^{2}\right)+\left(dV_{2}^{}{}_{}{}^{2}\right)+\left(dV_{3}^{}{}_{}{}^{2}\right)+\left(dV_{4}^{}{}_{}{}^{2}\right)\right].$$
(64)
or equivalently by changing to radial cordinates $`\underset{m=1}{\overset{4}{}}\left(dV_m\right)^2=dv^2+v^2\underset{i=1}{\overset{3}{}}\left(d\mathrm{\Omega }_i\right)^2`$, and defining a new variable $`\varphi =\sqrt{\frac{N_f}{2}}\mathrm{log}\left(\frac{v}{M}\right)`$ for some mass scale $`M`$:
$$ds^2=d\varphi ^2+\frac{N_f}{2}\underset{i=1}{\overset{3}{}}\left(d\mathrm{\Omega }_i\right)^2,$$
(65)
together with the 3-form torsion $`H`$ given by $`\left(N_f\right)`$ times the volume form of the 3-sphere namely:
$$H=N_fd\mathrm{\Omega }_1d\mathrm{\Omega }_2d\mathrm{\Omega }_3=N_f๐๐$$
(66)
The effective theory in the region of large $`V`$ is described by a Liouville field $`\varphi `$, its fermionic partner $`\psi _\varphi `$ and a supersymmetric level $`N_f`$ $`SU\left(2\right)`$ WZW model generated by the usual currents $`J^\pm `$ and $`J^3`$ which may be rewritten as the sum of a bosonic level $`\left(N_f2\right)`$ $`SU\left(2\right)`$ WZW model plus three free fermions. $`\psi _{SU\left(2\right)}^\pm `$ and $`\psi _{SU\left(2\right)}^3`$ Altogether these fields give a realization of the $`N=4`$ conformal field theory of the central charge $`C=6\left(N_f1\right)`$ as shown on the following central charge counting
$$6\left(N_f1\right)=2+\frac{3\left(N_f2\right)}{N_f}+\left(1+3Q^2\right);$$
(67)
where $`Q=\left(N_f1\right)\sqrt{\frac{2}{N_f}}`$. In the end of this digression on the physics in the throat of the Higgs branch, note that the Liouville field $`\varphi `$ is intimately related with the four scalars $`\left\{V_m\right\}`$ of the vector multiplet and then with the abelian $`U\left(1\right)`$ gauge factor as shown on the following eq.
$$\mathrm{exp}\left(\frac{\sqrt{2}}{N_f}\varphi \right)\sqrt{V_1^2+V_2^2+V_3^2+V_4^2}.$$
(68)
Therefore there is one to one correspondance between the liouville field $`\varphi `$ the vector multiplet of the $`2d`$ $`N=4`$ $`U\left(1\right)`$ gauge theory In other words the field $`\varphi `$ is one to one correspondance with the $`U\left(1\right)`$ factor of the gauge theory. In this regards, one ask the following question. What happens if, instead of $`2d`$ $`N=4`$ supersymmetric $`U\left(1\right)`$ gauge theory , we consider a $`U\left(1\right)^r`$ gauge theory involving $`r`$ abelian $`U\left(1\right)`$ factors and then $`4r`$ scalars $`V_{m,a}`$; $`a=1,..,r;m=1,2,3,4`$? Before discussing the answer to this question, let us first consider the linear sigma model solutions for a typical $`U\left(1\right)^r`$ D-flatness eqs. This concerns for example of type eqs (57-58) which are associated with a $`U\left(1\right)\times U\left(1\right)`$ supersymmetric linear sigma model. Following the same steps we described above, one can solve these eqs in a similar way as for the $`U\left(1\right)`$ theory. The result is that eqs (57-58) describe a two dimensional complex surface given by two intersecting smooth $`T^{}CP^1`$โs of base
$$\begin{array}{ccc}\left|x_0\right|^2+\left|x_1\right|^2=1\hfill & & \\ \left|\overline{x}_1\right|^2+\left|\overline{x}_2\right|^2=1,\hfill & & \end{array}$$
(69)
where $`\left|\overline{x}_1\right|`$ and $`\left|\overline{x}_2\right|`$ are obtained from eqs (57-58) and analogous changes as in eq (60). Using the results of and the discussions made in the end of the previous example, one sees that here also the fields $`X_i`$ and $`Y_i`$ could not be the appropriate variables in the viccinity of the singularity. Since this singularity is a degenerate singularity of type $`A_2`$, we expect to have more than one Liouville mode in this region and then a more general $`2dN=4`$ conformal field theory with backgound charges. A naive way to see this feature is to use the radial coordinates change of the $`U\left(1\right)`$ gauge theory which allowed us to put eq(64) into its equivalent form eqs ( 65,66). Since in the $`U\left(1\right)\times U\left(1\right)`$ gauge theory we are discussing we have two kinds of scalar fields $`V_{1,m}`$ and $`V_{2,m}`$ corresponding to each $`U\left(1\right)`$ factor of the $`U\left(1\right)\times U\left(1\right)`$ group, one is tempted to extend the above radial charge to lead to a $`N=4`$ conformal $`su\left(3\right)`$ Toda theory . Indeed, starting from the radial parametrization
$$\underset{m=1}{\overset{4}{}}\left|dV_{\rho ,m}\right|^2=\left(dv_\rho \right)^2+v_{\rho }^{}{}_{}{}^{2}\underset{i=1}{\overset{3}{}}\left(d\mathrm{\Omega }_{\rho ,i}\right)^2;\rho =1,2$$
(70)
and introducing two scalar fields $`\varphi _\rho `$:
$$\varphi _\rho =a_\rho \mathrm{log}\frac{v_\rho }{M}$$
(71)
where the coefficients $`a_\rho `$ should be determined by $`2d`$ $`N=4`$ conformal invariance; one can write down an extension of eqs (65,66). Supersymmetry and $`SO\left(4\right)`$ invariance suggest the following extension :
$$\begin{array}{ccc}ds^2=\frac{1}{2}K_{\rho \sigma }\left[d\varphi _\rho d\varphi _\sigma +a_\rho a_\sigma \underset{i=1}{\overset{3}{}}d\mathrm{\Omega }_{\rho ,i}d\mathrm{\Omega }_{\sigma ,i}\right];\hfill & & \\ H=2a_\rho d\mathrm{\Omega }_\rho ,\hfill & & \end{array}$$
(72)
where $`K_{\rho \sigma }`$ is $`su\left(3\right)`$ the Cartan matrix . More generally, this analysis may be extended in a natural way to any $`2d`$ $`N=4`$ $`U\left(1\right)^r`$ gauge theory $`r1`$ in presence of $`N_{f,r}`$ hypermultiplets. To do so one should first note that for a $`2dN=4`$ supersymmetric $`U\left(1\right)^r`$ linear $`\sigma `$ model with $`\left(r+1\right)`$ hypermultiplets the moduli space is given by the intersection of $`r`$ $`T^{}CP^1`$โs. When all the FI coupling variables vanish simultaneously, the physics within the Higgs branch throat is expected to be described by a general $`2d`$ $`N=4`$ superconformal Toda theory. In this region the metric is expected to have a form like that given by eqs(65,66). Progress in this direction will be reported elsewhere .
### 5.2 $`\underset{i}{}q_a^i=0`$
This situation is the relevent one in the analysis of the moduli space of gauge invariant vacua of $`2dN=2`$ supersymmetric linear sigma models. It ensures that in the infrared, the gauge theory flows to a superconformal one and plays a crucial role in the study of superstrings compactifications on local Calabi Yau manifolds with ADE singularities. The $`q_a^j`$โs satisfying the relation $`\underset{i}{}q_a^i=0`$ are also one of the main ingredient in toric geometry especially in the toric construction of Calabi Yau manifold and their mirrors .
In the case of $`2d`$ $`N=4`$ supersymmetric linear sigma models we have been studing, the sum over the $`q_a^j`$ charges is automatically fulfilled as shown on eq (52) and then one maight conclude that it is not necessary to distinguish the two senarios described in paragraphs 5.1 and 5.2. Though this remark is partially true, there are however some remarkable subtilities we will comment in a moment. Moreover, distinguishing the two senarios is also relevant for studying $`N=4`$ supersymmetric backgrounds by imetating methods of $`2d`$ $`N=2`$ supersymmetric linear models as we have done in subsection 4.2 . In what follows we give two examples illustrating the above remarks . In the first example we consider $`2d`$ $`N=4`$ $`U\left(1\right)^2`$ linear sigma model with three hypermultiplets of $`q_a^j`$ charges chosen as:
$$q_a^j=\delta _{a1}^j\delta _a^j;\underset{j=0}{\overset{2}{}}q_a^j=0.$$
(73)
The D-flatness conditions, which may be deduced from eqs (56) and ( 73 ) read as
$$\begin{array}{ccc}\left(\left|X_0\right|^2\left|X_1\right|^2\right)\left(\left|Y_0\right|^2\left|Y_1\right|^2\right)=R_1\hfill & & \\ X_0\overline{Y}_0X_1\overline{Y}_1=\overline{X}_0Y_0\overline{X}_1Y_1=0,\hfill & & \end{array}$$
(74)
together with
$$\begin{array}{ccc}\left(\left|X_1\right|^2\left|X_2\right|^2\right)\left(\left|Y_1\right|^2\left|Y_2\right|^2\right)=R_2\hfill & & \\ X_1\overline{Y}_1X_2\overline{Y}_2=\overline{X}_1Y_1\overline{X}_2Y_2=0.\hfill & & \end{array}$$
(75)
A way to handle these eqs is to note that they are quite similar to eqs (57,58) up to permutating the roles of $`X_1`$, $`X_2`$ and $`\overline{Y}_1`$ and $`\overline{Y}_2`$ respectively. From this view point eqs (74 ) may be rewritten as
$$\begin{array}{ccc}\left(\left|X_0\right|^2+\left|\overline{Y}_1\right|^2\right)\left(\left|Y_0\right|^2+\left|X_1\right|^2\right)=R_1\hfill & & \\ X_0\overline{Y}_0+X_1\left(\overline{Y}_2\right)=\overline{X}_0Y_0+\overline{X}_1\left(Y_2\right)=0.\hfill & & \end{array}$$
(76)
For later use it is intersting to rename the feld variables of eqs (76)as $`X_0=Z_0`$,$`\left(\overline{Y}_1\right)=Z_1`$; $`Y_0=W_0`$ and $`X_1=W_1`$. Putting this change in the above eqs, one sees that the resulting relations are comparable to those given by eqs (58,61 ). Thus eqs (76) describe just a cotangent bundle of $`CP^1`$. The base $`B_1`$ and the fiber $`F_1`$ are respectively parametrized by the local coodinates $`(z_0,z_1)`$ and $`(w_0,w_1)`$ where the $`z_i`$โs and $`w_i`$โs are related to $`Z_i`$โs and $`W_i`$โs by analogous formulas to those given by eqs (60). In the case of the $`2d`$ $`N=4`$ supersymmetric $`U\left(1\right)\times U\left(1\right)`$ gauge theory, we have to solve the system of eqs (74) and (75) which we rewrite for convienience as
$$\begin{array}{ccc}\left(\left|Z_0\right|^2+\left|Z_1\right|^2\right)\left(\left|W_0\right|^2+\left|W_1\right|^2\right)=R_1\hfill & & \\ Z_0W_0+Z_1W_1=\overline{Z}_0\overline{W}_0+\overline{Z}_1\overline{W}_1=0,\hfill & & \end{array}$$
(77)
Eqs (77) coincide with eqs (76) we have considered above while eqs (75) read now as
$$\begin{array}{ccc}\left(\left|W_1\right|^2+\left|Z_2\right|^2\right)\left(\left|Z_1\right|^2+\left|W_2\right|^2\right)=R_2\hfill & & \\ Z_1W_1+Z_2W_2=\overline{Z}_1\overline{W}_1+\overline{Z}_2\overline{W}_2=0,\hfill & & \end{array}$$
(78)
where we have set $`\overline{Y}_2=Z_2`$ and $`X_2=W_2`$. For positive definite values of $`R_2`$, if we take $`Z_1=W_2=0`$, one sees that the complex coordinates $`(W_1,Z_2)`$ parametrize a $`CP^1`$ complex curve which is isomorphism to a real $`2`$-sphere of radius $`\sqrt{R_2}`$. In the limit when $`R_2`$ goes to zero, this two sphere collapse and one ends with a $`SU\left(2\right)`$ singularity. For generic values of $`Z_1`$ and $`W_2`$, eqs (78) describe a cotangent bundle: $`T^{}CP^1`$ exactly as for eqs (77), the radius of the base $`B_1`$ of $`T^{}CP^1`$ is proportional to $`\sqrt{R_1}`$. Eqs (77) and (78) describe then two intersecting cotangent $`CP^1`$โs whose bases $`B_1`$ and $`B_2`$ as well as fibers $`F_1`$ and $`F_2`$ are roughly speaking parametrized by $`(Z_0,Z_1)`$,$`(W_1,Z_2)`$, $`(W_0,Z_1)`$ and $`(W_2,Z_1)`$ respectively. The second example we want to give deals with the case of a $`2d`$ $`N=4`$ supersymmetric $`U\left(1\right)^r`$ gauge theory with $`r+2`$ hypermultiplets $`\phi _j^+`$ of charges $`q_a^j`$ given by
$$q_a^j=2\delta _a^j+\delta _a^{j1}+\delta _a^{j+1},$$
(79)
satisfying the identity
$$\underset{i=0}{\overset{r+1}{}}q_a^j=0.$$
(80)
Using the splitting (29) with $`\gamma =1`$, one sees that the $`X_j`$โs and $`Y_j`$โs transform under the $`C_{}^{}{}_{}{}^{r}`$ actions in the same manner as the $`\phi _j^+`$ namely:
$$\begin{array}{ccc}X_j\lambda ^{q_a^j}X_j\hfill & & \\ Y_j\lambda ^{q_a^j}Y_j;\hfill & & \end{array}$$
(81)
where $`\lambda `$ is non zero complex parameter. Putting eqs (79) into the D-flatness eqs (56 ), one gets the following system of $`3r`$ eqs:
$$\begin{array}{ccc}\left(\left|X_{a1}\right|^2+\left|X_{a+1}\right|^22\left|X_a\right|^2\right)\left(\left|Y_{a1}\right|^2+\left|Y_{a+1}\right|^22\left|Y_a\right|^2\right)=R_a\hfill & \left(a\right)& \\ X_{a1}\overline{Y}_{a1}+X_{a+1}\overline{Y}_{a+1}2X_a\overline{Y}_a=0\hfill & \left(b\right)& \\ \overline{X}_{a1}Y_{a1}+\overline{X}_{a+1}Y_{a+1}2\overline{X}_aY_a=0.\hfill & \left(c\right)& \end{array}$$
(82)
For the simple example of the $`U\left(1\right)`$ gauge theory, one can check by following the same procedure we described in the previous example that the system of eqs given herebelow
$$\begin{array}{ccc}\left(\left|X_0\right|^2+\left|X_2\right|^2+2\left|Y_1\right|^2\right)\left(\left|Y_0\right|^2+\left|Y_2\right|^2+2\left|X_1\right|^2\right)=R_1\hfill & & \\ X_0\overline{Y}_0+X_2\overline{Y}_22X_1\overline{Y}_1=0\hfill & & \\ \overline{X}_0Y_0+\overline{X}_2Y_22\overline{X}_1Y_1=0,\hfill & & \end{array}$$
(83)
describe the cotangent bundle of the complex two dimensions weigthed complex pojective space $`WP_{1,2,1}^2`$. For the general $`U\left(1\right)^r`$ gauge theory, if all the $`R_a`$โs are non zero, eqs (82) describe the intersection of $`r`$ $`WP_{1,2,1}^2`$ cotangent bundles.
In the end of this discussion we would like to make a comment regading the second example. This concerns the link between our present analysis and the usual ADE $`N=2`$ syupersymetric models. Starting from the splitting (29) of the hypermultilpet moduli, one may view the $`X_j`$โs and $`Y_j`$โs as vacua of the moduli space of two orthogonal copies of $`2d`$ $`N=2`$ supersymmetric $`A_r`$ models. This feature is easily seen on the $`C_{}^{}{}_{}{}^{r}`$ action on these fields as shown on eqs (81) and may be rendered more manifest by analysing the D-flatness eqs (82). Ignoring for the moment eqs (82.b-c) and setting $`R_a=A_aB_a`$ with $`A_aB_a`$, one may put eqs (82.a) into the following remarkable form describing two copies of $`A_r`$ models
$$\begin{array}{ccc}\left(\left|X_{a1}\right|^2+\left|X_{a+1}\right|^22\left|X_a\right|^2\right)=A_a\hfill & & \\ \left(\left|Y_{a1}\right|^2+\left|Y_{a+1}\right|^22\left|Y_a\right|^2\right)=B_a\hfill & & \end{array}$$
(84)
For $`A_a=0`$; $`B_a`$ positive definite or $`A_a`$ positive definite; $`B_a=0`$, one of the $`A_r`$ models is singular while $`A_a=B_a=0`$ both of them are singular. For $`A_a`$ positive definite; $`B_a`$ positive definite, eqs( 84) describe the blown up of the two $`A_r`$ singularities. Note that each of the $`A_r`$ models has $`N=2`$ supersymmetry whereas the origonal eq (82) from which they come have $`N=4`$ supersummetry. This means that eqs( 82.b-c) are the lacking piece that rotates the two orthogonal $`N=2`$ supersymmetries. Eqs (82.b-c) reflect the neceessary conditions to get $`N=4`$ supersymmetry in two dimensions starting from two orthogonal blocks of $`N=2`$ supersymmetric models. From this naive parametization of the two intersecting $`T^{}CP^1`$โs, one sees that the base $`B_1`$ intersects the fiber $`F_2`$ along the $`Z_1`$ direction and the fiber $`F_1`$ intersects the base $`B_2`$ along the $`\omega _1`$ direction.
## 6 Conclusion
In this paper we have studied two main things. First, we have developed the analysis of the resolution of ADE singularities of hyperKahler manifolds involved in strings compactification. This concerns too particularly the moduli spaces of the Higgs branch of supersymmetric $`U\left(1\right)^r`$ gauge theories with eight supercharges. Second, we have initiated the analysis of singular CFTโs with higher order degeneracies by using the field theoretical approach of Aharony and Berkooz. Actually this study may be viewed as an extension of the recent works dealing with the leading $`A_1`$ singularity.
Concerning the first part, we have studied the solutions of the D-flatness eqs of supersymmetric $`U\left(1\right)^r`$ gauge theories with eight supercharges by using the linear sigma model approach. We have given, amongst others, a geometrical interpretation of the blown up singularities as a collection of intersecting cotangent complex dimensional weighted projective spaces depending on the number of hypermultiplets and gauge supermultiplets. This examination extends the standard linear sigma model analysis performed for the Kahler Coulomb branch of supersymmetric gauge theories with four supercharges. Our way of doing go beyond literature analysis where only half of the eight supersymmetries are manifest. Our method preserves manifesty all the eight supersymmetries and is realised in two steps based on a geometric realization of the $`SU\left(2\right)_R`$ symmetry on one hand and on a separation of the charges of the gauge and $`R`$-symmetries on the other hand. The factorisation of the gauge and $`SU\left(2\right)_R`$ charges of the hypermultiplets moduli involves a parameter $`\gamma `$ taking the values $`\gamma =0`$ or $`\gamma =1`$ which distinguish two classes of solutions of eqs(1) both preserving the eight supercharges. For $`\gamma =0`$, we have obtained a generalisation of the ADE complex surfaces reproducing the standard ones by partial breaking of $`2d`$ $`N=4`$ supersymmetry down to $`2d`$ $`N=2`$. For $`\gamma =1`$, we have found new models which flow in the infrared to $`2d`$ $`N=(4,4)`$ scale invaraint models. In this context several examples are given and classified according the manner one solves eqs(6). In the second part of this paper, we have studied the infrared dynamics of two dimensional $`N=(4,4)`$ gauge theories using field theoretical methods. We have made comments regarding the $`N=4`$ conformal Liouville description of the region in the viccinity of the singularity of the metric of the $`2d`$ $`N=4`$ Higgs branch. In this region, the string coupling constant $`g_s=e^\varphi `$ blows up as the Liouville field $`\varphi `$ goes to infinity . In an attempt towards an interpretation of the degenerate $`A_r`$ singularity carried by eqs (1), we have given field theoretical arguments suggesting that the metric of the moduli space near the Higgs singularity maight be described by a $`N=4`$ conformal $`SU\left(r+1\right)`$ Toda theory in two dimensions. This observation needs however a detailed study. In this regards a project of checking this observation for the case of the $`sp\left(2\right)`$ gauge group is understudy .
Aknowledgements
One of us(EHS) would like to thank Profs J.Shinar and J.Vary for kind hospitality at IITAP, Iowa State University; where a part of this work is done. AB would like to thank the organizers of the Spring Workshop on Superstrings and related Matters (March 1999), The Abdus Salam International Centere for Theoretical Physics Trieste ,Italy, for hospitality.
This research work has been supported by the program PARS number 372-98 CNR.
References
1. \] N. Seiberg and E. Witten , Nucl. Phys (1994)B426-19,hep-th/9407087.
2. \] N. Seiberg and E. Witten , Nucl. Phys (1994)B431-498,hep-th/9408099.
3. \] A. ELHassouni, T.Lhallabi, E.Oudrhiri and E.H.Saidi; Class.Quant.Grav 5 (1988)287-298.
4. \] A. Hanany and E. Witten ; Nucl. Phys B492 152-190 (1997),hep-th/9611230.
5. \] E. Witten ; Nucl. Phys B500 3-42 (1997),hep-th/9703166.
6. \] A. Klemm, W. Lerche , P Mayr , C Vafa and N. Warner; Nucl Phys B477(19997)
7. \] S. Katz,A.Klemm and C. Vafa; Nucl. Phys (1997) B497 173-195.
8. \] S. Katz, P. Mayr and C. Vafa, Mirror symmetry and exact solution of 4d N=2 gauge theories I Adv,Theor. Math.Phys(1998)1.53.
9. \] N.C. Leung and C. Vafa ; Adv .Theo . Math. Phys 2(1998) 91, hep-th/9711013
10. \] C. Vafa, Exact Results of N=2 Compactifications of Heterotic Strings, Summer School in High Energy Physics and Cosmology (ICTP, Trieste,1996).
11. \] C. Vafa, Geometric Physics ; Talk presented at the ICMโ 98.
12. \] P. Mayr , Geometric Construction of N=2 gauge theories, Spring Workshop on Superstrings and Related Matters, (ICTP, Trieste, March 1999).
13. \] A.Belhaj, A. EL Fallah and E.H. Saidi, On the affine$`D_4`$ mirror geometry; Class. Quantum. Grav. 16 (1999)3297-3306.
14. \] A.Belhaj, A. EL Fallah and E.H. Saidi On non simply laced mirror geometries in type II strings; Class. Quantum. Grav. 17 (1999)1-18.
15. A. Shapere and C. Vafa, BPS Structrure of ARgyress -Douglas Superconformal theories ;hep-th/9910182.
16. \] T. R. Taylor and C. Vafa, RR flux on Calabi Yau and Partial Supersymmetry Breeking ;hep-th/9912152.
17. \]N. Seiberg and E. Witten, D1-D5 and singular CFT ; JHEP 9904 (1999) 017, hep-th/9903224
18. \] Aharoni, Berkooz IR Dynamics of d=2 , N=(4,4) Gauge Theories and DLCQ of Little String theories ; JHEP 9910 (1999) 0030, hep-th/9909101.
19. \] A. Given , D. Kutasov and N. Seiberg , More Comments on String on $`ADS_3`$ ; hep-th/9806194.
20. \] D. Kutasov and N. Seiberg, Comments on String on $`ADS_3`$ ; Adv .Theo . Math. Phys 2(1998) 733, hep-th/9806194.
21. \] E. Diakonosko and N. Seiberg , The Coulomb branch of (4,4) supersymmetric Field theories in two dimensions , hep-th/9707158 ,JHEP 9707(1997)1.
22. \] J. Brodie Two Dimensions Mirror Symmetry from M-theory ; Nucl. Phys B517(1998)36-52, hep-th/9709228.
23. \] S. Sehti, The Matrix Formulation of type II Five branes ; Nucl. Phys B52(1998) 158,hep-th/9710005.
24. \] K. Schoutens, O(N) Extended Superconformal Field Theory In Superspace ; Nucl. Phys B295(1998)634.
25. \] A. Servin, W. Troost and A. Van Prooeyen, Superconformal Algebras In Two dimensions with N=4 ; Nucl. Phys B304(1988)447.
26. \] E. A.Ivanov , S. O. Krivonos and V. M. Leviant , A New Class Of Superconformal Sigma models With the Wess -Zumino action ; Nucl. Phys B304(1988)601.
27. \] C. Kounass , M. Porrati and B. Rostand, On N=4 extend super-Liouville Theory ; Physs. Lett B225 (1991)611.
28. \] E. H. Saidi and M. Zakkari, Class Quant Grav 7 (1991) 123; E. H. Saidi and M. Zakkari, Inter Jour Mod A6 (1991) 3175.
29. \] S. Kachru and C. Vafa, Exact Results for N=2 compactifications of Heterotic Strings ; Nucl Phys B 450(1995) 69, hep-th/9505105.
30. \] S. Ferarra, J.A. Harvey, A . Strominger, and C Vafa, Second -Quantized mirror Symmetry ; Physs. Lett B361 (1995)59-65. hep-th/9505162.
31. \] E .Witten, Heteroric String Conformal Field Theory ad ADE Singularities , hep-th/9909299.
32. \]N. Seiberg and E .Witten Gauge Dynamics and Compactification to Three Dimensions In The mathematique Beauty of Physics , hep-th/9609122 .
33. \]K. Hori, H. Ooguri and C. Vafa, Non -Abelian Conifold Trasitions and N=4 in Three dimensions ; Nucl. Phys B504 117 (1997), hep-th/9705220.
34. \] Moshe Rozali, Hypermultiplet Moduli Space and Three Dimensional Gauge Theories ; hep-th/9910238
35. \] P. Mayr, Conformal Field Theory on $`K3`$ and Three dimensional gauge theories ; hep-th/9910268.
36. \]E. Witten, String Theory Dynamics in Various Dimensions ; Nucl Phys B 433(1995) 85, hep-th/9503124.
37. \] P.S. Aspinwall, Enhanced Gauge Symmetries and $`K3`$ surfaces ; Phys . Lett B443 (1995) 329, hep-th/9507012.
38. \] J. Maldacena , J. Michelson and A. Strominger, $`ADS_3`$ Black Holes and Stringy Exclusion Principale ; JHEP 12 (1998) 05 ,hep-th/9908142.
39. \]N. Seiberg and E. Witten, String Duality and Non Commutative geometry ; JHEP 9909(1999) 032, hep-th/9908142.
40. \] N. Nekarasov and A. Schwartz, Instanton on Non Commutative $`R^4`$ and (2,0) Superconformal Six dimensional theory , hep-th/9802068.
41. \] Shiraz Minwalla, Mark. Van. Raamsdonk and Nathan Seiberg, Non commutative perturbative Dynamics , hep-th/9912072.
42. \] S. Gukov , C. Vafa and E. Witten , CFTโs From Calabi Yau Four-folds ; hep-th/9906070
43. \] E. Silverstein and E. Witten Global U(1) R-symmetry and Conformal Invariance of(0,2) Models , hep-th/9403054
44. \] P.Aspinwall, B. Green and D.R. Plesser; Nucl Phys B416 414.
45. \] P.Aspinwall and B.R. Greene; Nucl Phys B437 (1995) 205-230.
46. \] E. Witten, Small Instantons In String Theory ; Nucl . Phys B460, 541 (1996),hep-th/9511030.
47. \] A. Galperin, E. Ivanov, S. Kalitzin, V. Ogievetsky and E. Sokatchev; Class . Quant . Grav 1(1984) 469.
48. \] E. Sahraoui and E. H. Saidi, On the completly integrable four dimensions N=2 hypermultiplet self - couplings ; Class . Quant . Grav 16(1999) 1605.
49. \] I. Benkadour , M. Bennai, E. Y. Diaf and E. H. Saidi , On Matrix model Compactification on non commutative $`F_0`$ geometry ; to appear in Class . Quant . Grav 17(2000)
50. \] T. Lhalabi and E.H.Saidi; Inter Jour Mod Phys3 (1988) 187-201
51. \] E. Sokatchev and K. Stelle; Class . Quant . Grav 4(1987)501.
52. \] k. Schoutens ; Phys . Rev. lett 79 (1997) 2608 , Cond-mat / 9706166.
53. \] M. Bershadsky , T. M. Chiang , B. Grenne , A. Johansen and C.I. Lazariou F-theory and linear Sigma Models ; Nucl . Phys . B527(1998) 531-570, hep-th/9712023.
54. \]M. Berkooz and H. Verlinde; Matrix Therory, ADS/CFT and Higgs branch, hep-th /9907100
55. \] A. Belhaj and E.H. Saidi, work in progress.
56. \] Batyrev,(1994)J.Alg. Geo. 3 493 |
warning/0002/hep-th0002049.html | ar5iv | text | # Finite temperature Casimir effect for a dilute ball satisfying ฯตโข๐=1 Revised Version
## 1 Introduction
The Casimir effect has in recent years attracted considerable interest. For reviews on this topic see, for example, Plunien *et al.*, Mostepanenko and Trunov, and Milton . The problem has been rather difficult to solve, once curved boundaries are present, because of divergences in the formalism and insufficient contact with experiments. In this paper we will consider a compact ball of radius a, permittivity $`ฯต_1`$ and permeability $`\mu _1`$, situated in an infinite medium whose corresponding material constants are $`ฯต_2`$ and $`\mu _2`$.
We shall take
$$ฯต_1\mu _1=ฯต_2\mu _2=1$$
(1)
so that the photon velocity is everywhere equal to $`c`$.
The relativistic condition (1) was, in the context of the Casimir effect, to our knowledge first introduced by Brevik and Kolbenstvedt . The reason why this particular condition was imposed was that when calculating the Casimir surface force density on the sphere, in this paper called $`f`$, one avoids thereby the subtraction of a contact term: the contact term turns out to be zero just as in the case of a singular spherical conducting sphere in a vacuum . Later, this kind of medium has been considered from various viewpoints \[7 - 19\]. A common feature of most of these papers is the use of macroscopic electrodynamics of continuous media.
The main purpose of the present paper is to consider the Casimir effect for a compact ball satisfying the condition (1), at finite temperatures. Finite temperature Casimir theory has been considered previously by several authors, e.g., Lifshitz , Fierz , Sauer , Mehra , Brown and Maclay , Milton et al. , Revzen et al. , and Klich . Also the very recent extensive treatise of Feinberg et al. ought to be mentioned, since it describes in detail the characteristic properties of the high temperature limit, and contains also an extensive reference list.
We start in the next section by considering, as a first step, the $`T=0`$ theory, and demonstrate the equivalence between the Green function method and the mode summation method. The equivalence has of course to hold if the theory is to make sense physically, but the equivalence is not so easy to see by a mere inspection of the mathematics. A detailed calculation is desirable. In section 3 we generalize the Casimir energy expression obtained in to the fourth order in the diluteness parameter $`\xi `$; cf. Eq. (15) and the definition in Eq. (7). In section 4 we consider the Casimir free energy $`F(T)`$ at finite $`T`$, emphasizing the difference between $`F(T)`$ and the Casimir energy $`E(T)`$, and give an order of magnitude estimate of $`F(T)`$ in the case of a dilute medium endowed with a โsquareโ dispersion relation implying a sharp high frequence cutoff at imaginary frequency $`\widehat{\omega }=\omega _0`$. The presence of a dispersion-induced strong, negative, cutoff-dependent contribution to the Casimir energy (or free energy, for finite $`T`$), seems first to have been pointed out by Candelas , for the case of a perfectly conducting shell.
At high temperatures we find $`F(T)`$ to be negative, decreasing linearly with $`T`$. This agrees with results found previously in and . The Casimir energy $`E(T)`$ itself tends to zero at high temperatures.
## 2 Equivalence between the Green function approach and the mode sum approach when $`T=0`$. Nondispersive medium
In this section we consider first the Casimir surface force density $`f`$ at zero temperature, for a nondispersive medium. There are in principle two different methods for calculating $`f`$ (more precisely, the part of $`f`$ that is independent of the high frequency cutoff):
The first method is to make use of the Green functions. There are two scalar Green functions, corresponding to two different electromagnetic modes. Again, this method can be carried out in one of two different ways:
1. First, one can evaluate $`f`$ as the difference between the radial Maxwell stress tensor components on the outside and inside of the surface $`r=a`$. We have to subtract off the volume terms in the scalar Green functions, so that only the boundary dependent terms are left. This regularization ensures that the surface force goes to zero when $`a\mathrm{}`$, which is as we expect on physical grounds. Also, this regularization permits us to calculate the Casimir energy $`E`$ once $`f`$ is known, by the formula $`E/a=4\pi a^2f`$. (At $`T=0`$, $`E`$ is the same as the free energy $`F`$.)
2. Alternatively, one can calculate the energy $`E`$ by integrating the Maxwell energy density over the volume, retaining, as above, only the boundary dependent terms in the Green functions. The outcome of these two approaches are of course in agreement.
The Green function method was first employed by Milton, for the case of an ordinary dielectric ball , and was later applied to the case of an $`ฯต\mu =1`$ ball in Ref. .
The result of the calculation can be given as follows. After a complex frequency rotation $`\omega i\widehat{\omega }`$ we can express all physical quantities as integrals over frequencies $`\widehat{\omega }`$ along the imaginary frequency axis.<sup>3</sup><sup>3</sup>3 This frequency rotation is a bit more involved than it might seem at first sight, since the Hankel functions imply that there are singularities in the lower half of the complex frequency plane. Physically, these singularities occur because we assume an infinite exterior region (thus, no large exterior sphere making all eigenfrequencies real). The legitimacy of the complex frequency rotation in the presence of these complexities follows from a more careful derivation; cf. the recent discussion in Ref. . The cutoff dependence is omitted (it disappears in the regularization procedure). We define Riccati-Bessel functions $`s_l`$ and $`e_l`$ by ($`x=\widehat{k}a)`$
$$s_l(x)=\sqrt{\frac{\pi x}{2}}I_\nu (x),e_l(x)=\sqrt{\frac{2x}{\pi }}K_\nu (x),$$
(2)
corresponding to the Wronskian $`W\{s_l,e_l\}=1`$. Here $`\nu =l+\frac{1}{2}`$; $`I_\nu `$ and $`K_\nu `$ are modified Bessel functions. We introduce two dispersion functions $`\mathrm{\Delta }_l^{\mathrm{TE}}`$and $`\mathrm{\Delta }_l^{\mathrm{TM}}`$corresponding to the TE and TM modes:
$`\mathrm{\Delta }_l^{\mathrm{TE}}(x)=\sqrt{ฯต_1\mu _2}s_l^{}(x)e_l(x)\sqrt{ฯต_2\mu _1}s_l(x)e_l^{}(x),`$ (3)
$`\mathrm{\Delta }_l^{\mathrm{TM}}(x)=\sqrt{ฯต_2\mu _1}s_l^{}(x)e_l(x)\sqrt{ฯต_1\mu _2}s_l(x)e_l^{}(x).`$ (4)
The Casimir energy as derived in on the basis of the Green function can be written as
$$E=\frac{1}{2a}\frac{(\mu _1\mu _2)^2}{\mu _1\mu _2}_{\mathrm{}}^{\mathrm{}}\frac{dy}{2\pi }e^{i\delta y}\underset{l=1}{\overset{\mathrm{}}{}}\frac{2l+1}{4}\frac{x}{\mathrm{\Delta }_l^{\mathrm{TE}}(x)\mathrm{\Delta }_l^{\mathrm{TM}}(x)}\left[\lambda _l^2(x)\right]^{}.$$
(5)
Here $`y=\widehat{\omega }a`$, $`x=|y|`$; the time difference $`\tau =tt^{}`$ has been Euclidean rotated, $`\tau i\widehat{\tau }`$, and $`\delta =\widehat{\tau }/a`$ is the time-splitting parameter. Moreover, we have defined the function $`\lambda _l(x)`$ as $`\lambda _l(x)=[s_l(x)e_l(x)]^{}`$ \[In Ref we used the symbols $`D_l(x)`$ and $`\stackrel{~}{D}_l(x)`$; the relations between these quantities and those defined by (3) and (4) are $`D_l(x)=\sqrt{\mu _1\mu _2}\mathrm{\Delta }_l^{\mathrm{TE}}(x)`$ and $`\stackrel{~}{D}_l(x)=\sqrt{\mu _1\mu _2}\mathrm{\Delta }_l^{\mathrm{TM}}(x)`$. \]
Consider next the second method, which consists in summing over the modes by means of contour integration. This method was recently used in Ref. , and is in turn based on the technique developed by Lambiase and Nesterenko and Nesterenko and Pirozhenko . One gets (when we supply the same parameter $`\delta `$ as above)
$$E=\frac{1}{2a}\underset{l=1}{\overset{\mathrm{}}{}}(2l+1)_{\mathrm{}}^{\mathrm{}}\frac{dy}{2\pi }e^{i\delta y}x\frac{d}{dx}\mathrm{ln}\left[\mathrm{\Delta }_l^{\mathrm{TE}}(x)\mathrm{\Delta }_l^{\mathrm{TM}}(x)\right].$$
(6)
Introducing the parameter $`\xi `$ by
$$\xi =\frac{ฯต_1ฯต_2}{ฯต_1+ฯต_2}=\frac{\mu _2\mu _1}{\mu _2+\mu _1},$$
(7)
the expression (6) can alternatively be written as
$$E=\frac{1}{2a}\underset{l=1}{\overset{\mathrm{}}{}}(2l+1)_{\mathrm{}}^{\mathrm{}}\frac{dy}{2\pi }e^{i\delta y}x\frac{d}{dx}\mathrm{ln}\left[1\xi ^2\lambda _l^2(x)\right].$$
(8)
The two methods thus lead to two expressions for the Casimir energy, given by Eqs. (5) and (6). The expressions, of course, have to be equal, but to see the equality is not quite trivial by mere inspection. The equality is tantamount to the equality
$$\left[\mathrm{\Delta }_l^{\mathrm{TE}}(x)\mathrm{\Delta }_l^{\mathrm{TM}}(x)\right]^{}=\frac{(\mu _1\mu _2)^2}{4\mu _1\mu _2}\left[\lambda _l^2(x)\right]^{}.$$
(9)
Now it turns out that this equality is not so difficult to prove after all if we insert the expressions (3) and (4) on the left hand side and take into account the above mentioned Wronskian between $`s_l`$ and $`e_l`$. Actually, the expression (6) was also made use of in the Green-function calculation of Milton and Ng .
## 3 A remark on the dilute ball when T=0
Before leaving the zero temperature theory let us make a brief remark on the nondispersive dilute ball, meaning that
$$|\xi |1.$$
(10)
The most practical expression to make use of for the Casimir energy, is that of Eq.(8). We expand the logarithm to fourth order in $`\xi `$:
$$\mathrm{ln}(1\xi ^2\lambda _l^2)=\xi ^2\lambda _l^2(x)\frac{1}{2}\xi ^4\lambda _l^4+๐ช(\xi ^6),$$
(11)
and make the following steps:
1. put $`\delta =0`$;
2. perform a partial integration in Eq.(8);
3. interchange summation and integration signs.
Then, omitting the $`๐ช(\xi ^6)`$ terms in (11),
$$E=\frac{\xi ^2}{\pi a}\underset{l=1}{\overset{\mathrm{}}{}}\nu _0^{\mathrm{}}๐x\lambda _l^2(x)\frac{\xi ^4}{2\pi a}\underset{l=1}{\overset{\mathrm{}}{}}\nu _0^{\mathrm{}}๐x\lambda _l^4(x),$$
(12)
which can be processed further using the uniform asymptotic expansion for $`s_l`$ and $`e_l`$ up to $`๐ช(1/\nu ^4)`$, in the same way as in Refs. . There is one single divergent term in Eq. (12), which can be regularized by means of the Riemann zeta function, $`\zeta (s)`$. The only formula needed in practice is
$$\underset{l=1}{\overset{\mathrm{}}{}}\nu ^s=\left(2^s1\right)\zeta (s)2^s.$$
(13)
Moreover, we use the summations $`_1^{\mathrm{}}\nu ^2=\frac{1}{2}\pi ^24`$, $`_1^{\mathrm{}}\nu ^4=\frac{1}{6}\pi ^416`$, as well as the integral
$$_0^{\mathrm{}}(1+z^2)^{p/2}๐z=\frac{1}{2}\mathrm{B}(\frac{1}{2},\frac{p1}{2})$$
(14)
($`\mathrm{B}`$ is the beta function), to get
$$\begin{array}{cc}\hfill E=& \frac{3\xi ^2}{64a}\left[1+\frac{9}{128}\left(\frac{1}{2}\pi ^24\right)\frac{423}{16384}\left(\frac{1}{6}\pi ^416\right)\right]\hfill \\ & \frac{63\xi ^4}{16384a}\left[\left(\frac{1}{2}\pi ^24\right)\frac{351}{1792}\left(\frac{1}{6}\pi ^416\right)\right].\hfill \end{array}$$
(15)
This expression generalizes that of equation (3.10) in up to the fourth order in $`\xi `$. Numerically, the expressions between square parentheses in (15) are respectively 1.05966 and 0.88880. The fourth order contribution to the Casimir energy is thus weak and negative, corresponding to an attractive force component.
In connection with this calculation, we wish to emphasise two points. First, we may avoid the time-splitting parameter $`\delta `$ from the beginning, using zeta function regularization (here the Riemann function) instead. Secondly, the convergence parameter $`s`$ that was for mathematical reasons introduced in , can simply be omitted. At least from a pragmatic point of view, the nondispersive medium becomes quite simple.
## 4 Dispersive medium: Finite temperatures
We now take into account dispersive properties of the two media, meaning that $`ฯต=ฯต(\omega )`$, $`\mu =\mu (\omega )`$. We shall still assume that $`ฯต(\omega )\mu (\omega )=1`$, so that the velocity of photons in either medium is equal to the velocity of light. After frequency rotation, along the imaginary frequency axis, $`ฯต=ฯต(i\widehat{\omega })`$, $`\mu =\mu (i\widehat{\omega })`$.
Some care ought to be taken when considering the dispersive formalism, since frequency derivatives of the material constants are no longer ignorable straightaway. Let us first consider the $`T=0`$ total zero-point energy, which we shall call $`E_{I+II}`$, of the inner (I) and the outer (II) regions. This energy is simply given by Eq. (6), which we repeat here with a slightly generalized notation:
$$E=\frac{1}{2a}\underset{l=1}{\overset{\mathrm{}}{}}(2l+1)_{\mathrm{}}^{\mathrm{}}\frac{dy}{2\pi }e^{i\delta y}x\frac{d}{dx}\mathrm{ln}\left[\mathrm{\Delta }_l^{\mathrm{TE}}(i\widehat{\omega },a)\mathrm{\Delta }_l^{\mathrm{TM}}(i\widehat{\omega },a)\right].$$
(16)
The reason for this equality is that the two dispersion functions $`\mathrm{\Delta }_l^{TE}`$ and $`\mathrm{\Delta }_l^{TM}`$, as given by the dispersive generalization of Eqs. (3) and (4), retain their physical meaning also in the presence of dispersion. Thus the TE and TM modes are still determined by $`\mathrm{\Delta }_l^{TE}=0`$ and $`\mathrm{\Delta }_l^{TM}=0`$ respectively. The expression (16) needs to be regularized: we subtract off the zero-point energy $`E_{uniform}`$ corresponding to the limit $`a\mathrm{}`$; cf. . As $`s_l(x)=\frac{1}{2}e^x`$ and $`e_l(x)=e^x`$ for large $`x`$, we get from (3) and (4)
$$\mathrm{\Delta }_l^{TE}(i\widehat{\omega },a\mathrm{})=\mathrm{\Delta }_l^{TM}(i\widehat{\omega },a\mathrm{})=\frac{1}{2}(\sqrt{ฯต_1\mu _2}+\sqrt{ฯต_2\mu _1}).$$
(17)
The Casimir energy $`E=E_{I+II}E_{uniform}`$ becomes accordingly, at $`T=0`$,
$$E=\frac{1}{2a}\underset{l=1}{\overset{\mathrm{}}{}}(2l+1)_{\mathrm{}}^{\mathrm{}}\frac{dy}{2\pi }e^{i\delta y}x\frac{d}{dx}\mathrm{ln}\left[\frac{4\mathrm{\Delta }_l^{\mathrm{TE}}(i\widehat{\omega },a)\mathrm{\Delta }_l^{\mathrm{TM}}(i\widehat{\omega },a)}{(\sqrt{ฯต_1\mu _2}+\sqrt{ฯต_2\mu _1})^2}\right].$$
(18)
Observing that $`1+\lambda _l=2s_l^{}e_l`$, $`1\lambda _l=2s_le_l^{}`$ we find (when reverting to the notation $`\mathrm{\Delta }_l^{TE,TM}(x)`$ instead of $`\mathrm{\Delta }_l^{TE,TM}(i\widehat{\omega },a)`$) that
$$\mathrm{\Delta }_l^{TE}(x)\mathrm{\Delta }_l^{TM}(x)=\frac{1}{4}(\sqrt{ฯต_1\mu _2}+\sqrt{ฯต_2\mu _1})^2[1\xi ^2\lambda _l^2(x)],$$
(19)
which finally means that, at $`T=0`$,
$$E=\frac{1}{2a}\underset{l=1}{\overset{\mathrm{}}{}}(2l+1)_{\mathrm{}}^{\mathrm{}}\frac{dy}{2\pi }x\frac{d}{dx}\mathrm{ln}\left[1\xi ^2(x)\lambda _l^2(x)\right].$$
(20)
It is rather remarkable that this simple formula continues to hold even in the presence of arbitrary frequency dispersion. Since dispersion implies that we introduce a physically based high frequency cutoff, we expect that there is no longer any need for a time splitting parameter. We have therefore put $`\delta =0`$ in the expression (20). We will henceforth restrict ourselves to the first ( i. e., the second order) term in the expansion (11).
In this section we adopt the simplest imaginable dispersion relation (a โsquareโ form): we let the inner medium (a ball) correspond to
$$\mu (i\widehat{\omega })=\{\begin{array}{cc}\mu =\text{const},\hfill & \widehat{\omega }\omega _0\hfill \\ 1,\hfill & \widehat{\omega }>\omega _0,\hfill \end{array}$$
(21)
whereas the outer medium is a vacuum. It is to be noted here that the general thermodynamic law saying that the susceptibility has to decrease monotonically along the positive imaginary frequency axis , is satisfied by this curve for $`ฯต(i\widehat{\omega })=1/\mu (i\widehat{\omega })`$ if it is given a small negative slope for $`\widehat{\omega }<\omega _0`$. We still assume dilute media, $`|\xi |1`$, as above.
Consider now the case of finite temperatures. Since the thermal Green function is periodic in imaginary time, we can replace the Fourier transform in the $`T=0`$ theory with Fourier series in imaginary time. The transition to finite temperature theory is accomplished by means of a discretization of the frequencies,
$$\widehat{\omega }\widehat{\omega }_n=2\pi n/\beta ,xx_n=\widehat{\omega }_na,$$
(22)
where $`n`$ is an integer and $`\beta =1/T`$ (we put $`k_\mathrm{B}=1`$). The rule for going from frequency integral to sum over Matsubara frequencies is
$$_0^{\mathrm{}}dxt\underset{n=0}{\overset{\mathrm{}}{^{}}},t=2\pi a/\beta .$$
(23)
Here $`t`$ is a nondimensional temperature, and the prime on the summation sign in (23) means that the $`n=0`$ term is counted with half weight.
Since $`x_n`$ can be written as $`nt`$, we can express the Casimir free energy at finite temperatures as
$$F(T)=\frac{\xi ^2t}{\pi a}\underset{l=1}{\overset{\mathrm{}}{}}\nu \underset{n=0}{\overset{\mathrm{}}{^{}}}\lambda _l^2(nt)$$
(24)
to $`๐ช(\xi ^2)`$. The corresponding zero temperature expression, still to be called $`E`$, is given by the first term to the right in Eq. (12).
The following delicate point ought to be noted. We know that the mean energy per mode of the electromagnetic field is $`\phi (\omega ,\beta )=\frac{1}{2}\omega \mathrm{coth}(\frac{1}{2}\beta \omega )`$. The function $`\phi (\omega ,\beta )`$ is to replace $`\frac{1}{2}\omega `$, when we construct the general contour integral for the Casimir energy $`E(T)`$ at finite temperatures. It is the function $`\mathrm{coth}(\frac{1}{2}\beta \omega )`$ which is responsible for the occurrence of the Matsubara frequencies along the imaginary frequency axis. Now, our formalism implies morover a transition from the energy $`E(T)`$ to the free energy $`F(T)`$. This is accomplished by means of a partial integration, making use of the relations
$$\frac{\phi (\omega ,\beta )}{\omega }=\frac{}{\beta }\left(\frac{\beta }{\omega }\phi (\omega ,\beta )\right),E(T)=\frac{}{\beta }\left(\beta F(T)\right).$$
For this reason we have changed the energy symbol from E to $`F`$ in Eq. (24). At $`T=0`$, obviously $`F(T=0)=E(T=0)E`$. We thank I. Klich (personal communication) for helpful comments regarding this point.
Let us in this context also recall how the surface force density $`f`$ is calculated. Generally, as mentioned in section 2, $`f`$ can be found as the difference between the radial Maxwell stress tensor components on the outside and inside of the surface $`r=a`$. The stress components are in turn constructed from the two-point functions like $`E_i(๐ซ)E_k(๐ซ^{})`$ when $`๐ซ`$ and $`๐ซ^{}`$ are close together, at a given temperature $`T`$. The appropriate energy function is found by integrating $`4\pi a^2f`$ over $`a`$, from initial position $`a=\mathrm{}`$, to the final position where the radius of the sphere terminates at the fixed value $`a`$. This process is to take place at constant temperature. The appropriate energy function is accordingly the free energy $`F(T)`$. Thus we see that, whereas in force considerations at finite temperatures one is lead naturally to the free energy, while starting from energy considerations one has instead to make use of thermodynamical considerations to calculate $`F(T)`$.
Mathematically, the case $`n=0`$ requires special attention. For this purpose we observe the approximate expression for $`s_l`$ and $`e_l`$ at low $`x`$:
$$s_l(x)=\frac{\sqrt{\pi }}{\mathrm{\Gamma }(\nu +1)}\left(\frac{x}{2}\right)^{\left(\nu +\frac{1}{2}\right)},e_l(x)=\frac{\mathrm{\Gamma }(\nu )}{\sqrt{\pi }}\left(\frac{x}{2}\right)^{\left(\nu +\frac{1}{2}\right)}.$$
(25)
From this it follows that $`s_l(x)e_l(x)=x/2\nu `$, $`\lambda _l(x)=1/2\nu `$, implying that the contribution to (24) from $`n=0`$ is
$$F(T,n=0)=\frac{\xi ^2t}{8\pi a}\underset{l=1}{\overset{\mathrm{}}{}}\frac{1}{\nu }.$$
(26)
We add the contribution from $`n1`$, restricting ourselves to the dominant term in the uniform asymptotic expansion of $`\lambda _l`$ :
$$\lambda _l(x)=\frac{(1+x^2/\nu ^2)^{3/2}}{2\nu }\left[1+๐ช\left(\frac{1}{\nu ^2}\right)\right].$$
(27)
Thereby
$$F(T)=\frac{\xi ^2t}{8\pi a}\underset{l=1}{\overset{l_0}{}}\frac{1}{\nu }\left[1+\underset{n=1}{\overset{\mathrm{}}{}}\frac{2\nu ^6}{(\nu ^2+n^2t^2)^3}\right].$$
(28)
We have here truncated the sum over $`l`$ at $`l=l_0`$ in order to avoid divergences. It is easy to give an order-of-magnitude estimate of $`l_0`$: Our dispersion relation (21) implies that photons having frequencies $`\omega `$ higher than $`\omega _0`$ do not โseeโ the sphere at all. A photon of limiting frequency $`\omega _0`$ just touching the surface of the sphere has an angular momentum equal to $`\omega _0a=x_0`$. This type of argument has repeatedly been used in Casimir calculations . One might put $`l_0=cx_0`$ where $`c`$ is a coefficient of order unity. However, in view if simplicity, and since we are able to give only an estimate of the order of magnitude anyhow, we put henceforth $`c=1`$. As upper limit in the summation we thus take
$$l_0=x_0.$$
(29)
Now taking into account the exact summation formula derived in Eq.(A3) in ( with $`p`$ an arbitrary constant):
$$\begin{array}{cc}\hfill \underset{n=1}{\overset{\mathrm{}}{}}\frac{1}{(n^2+p^2)^3}=& \frac{1}{2p^6}+\frac{3\pi }{16p^5}\mathrm{coth}\pi p\hfill \\ & +\frac{3\pi ^2}{16p^4}\frac{1}{\mathrm{sinh}^2\pi p}+\frac{\pi ^3}{8p^3}\frac{\mathrm{coth}\pi p}{\mathrm{sinh}^2\pi p},\hfill \end{array}$$
(30)
we can write Eq. (28) as
$$F(T)=\frac{3\xi ^2}{64a}\underset{l=1}{\overset{x_0}{}}\left[\mathrm{coth}\frac{\pi \nu }{t}+\frac{\pi \nu /t}{\mathrm{sinh}^2\pi \nu /t}+\frac{2}{3}\left(\frac{\pi \nu }{t}\right)^2\frac{\mathrm{coth}\pi \nu /t}{\mathrm{sinh}^2\pi \nu /t}\right];$$
(31)
the first term on the right hand side of (30) drops out.
This expression can be calculated, with $`x_0`$ as an input parameter, to show how $`F(T)`$ varies with $`t`$. It is of interest first to examine the case of low temperatures. Let us assume that
$$t1,$$
(32)
implying that $`\pi \nu /t1`$. Then the first term in (31) dominates, and can be replaced by unity. We obtain approximately
$$F(t0)=\frac{3\xi ^2}{64a}x_0.$$
(33)
This negative expression establishes a constant low-temperature plateau for the Casimir free energy.
It is of interest to compare this result with that obtained from the slightly more complicated though physically more realistic Lorentz (or Sellmeir) dispersion relation. The latter, when taken only along the imaginary frequency axis, can be written as
$$\mu (i\widehat{\omega })1=\frac{\chi _0}{1+\widehat{\omega }^2/\omega _0^2}=\frac{\chi _0}{1+x^2/x_0^2},$$
(34)
where $`\chi _0`$ is the zero frequency magnetic susceptibility. (The damping term is omitted.) Here $`\widehat{\omega }=\omega _0`$ represents a โsoftโ cutoff. This case was considered in , for the zero-temperature case. For a dilute medium, $`|\chi _0|1`$, the result of the calculation in can be written
$$F(t0)E=\frac{\xi ^2}{16a}\left(x_0\frac{3}{4}\right).$$
(35)
This formula assumes mathematically, strictly speaking, that $`x_01`$ in addition to $`t1`$. The first condition is however not very stringent; it turns out numerically that (35) holds to better than 1 percent accuracy even when $`x_0`$ is as low as about 4. We thus see that the dominant first term in Eq. (35) agrees qualitatively well with Eq. (33), the latter obtained on basis of the โsquareโ dispersion relation (21). The latter expression is $`3/4`$ of the Lorentz-based expression $`E(\xi ^2/16a)x_0`$. This is a satisfactory result, in view of the semi-quantitative agreement between the meaning of the parameters $`x_0`$ occuring in (16) and in (34).
The second term in Eq. (35), $`E(\mathrm{nondisp})=(3/64)\xi ^2`$, is the characteristic positive Casimir energy (corresponding to a repulsive force component) for a nondispersive dilute ball satisfying the condition $`ฯต\mu =1`$, at zero temperature. Cf., for instance, .
The other limiting case of interest is that of high temperatures. Let us assume that
$$tx_0,$$
(36)
implying that $`\pi \nu /t1`$ for all values of $`l`$ in Eq. (31). Making use of $`\mathrm{sinh}z=z`$ and $`\mathrm{cosh}z=1`$ for small $`z`$, we obtain as leading term
$$F(tx_0)=\frac{\xi ^2t}{8\pi a}\underset{l=1}{\overset{x_0}{}}\frac{1}{\nu }=\frac{\xi ^2}{4\beta }\underset{l=1}{\overset{\mathrm{}}{}}\frac{1}{\nu }.$$
(37)
This expression is independent of $`a`$. The high - temperature Casimir free energy, still negative, thus decreases proportionally with $`t`$. This is qualitatively in agreement with and .
Figures 1 - 3 show how the nondimensional reduced Casimir free energy $`F_{red}(t)`$, defined by
$$F_{red}(t)\frac{F(T)}{3\xi ^2x_0/(64a)},$$
(38)
varies with $`t`$, in the cases when $`x_0=\{10,100,1000\}`$. The low-temperature plateaus evidently lie at $`F_{red}(0)=1`$. These figures are given on a linear scale, thus exhibiting the linear behaviour predicted by Eq. (37) at high temperatures. It is of interest to check the accuracy of (37) in cases when the condition (36) is reasonably well, but not extremely well, satisfied. Let us choose $`x_0=10,t=100`$ as an example. From the expression
$$F_{red}(tx_0)=\frac{8t}{3\pi x_0}\underset{l=1}{\overset{x_0}{}}\frac{1}{\nu }$$
(39)
we obtain, taking into account that $`_1^{10}1/\nu =2.361749`$, that $`F_{red}(x_0=10,t=100)=20.047149.`$ The full series solution yields the number -20.047154. The approximate formula (39) is thus in this case extremely accurate. Choosing $`t=10`$, we get from Eq. (39) that $`F_{red}(x_0=10,t=10)=2.005`$, whereas the series (31) yields -2.091. Even in this case, where the condition (36) obviously is broken, the error in the formula (39) is thus only about 4 per cent. Generally speaking, at high temperatures when the diluteness of the medium becomes more pronounced, the Casimir free energy becomes small, as we would expect.
In connection with the temperature dependence of the free energy the following physical argument should be noticed. From the dispersion relations (21) or (34) we would expect that the influence from the dispersion effect becomes important when the temperature becomes so high that the most significant frequencies in the thermal radiation field are of the same order as $`\omega _0`$. And as most significant thermal frequency, it is natural to take the value $`\omega =\omega _m`$ corresponding to the maximum of the blackbody distribution. From Wienโs dispacement law we know that $`\omega _m=2.8/\beta `$. Thus, we expect that dispersion becomes significant when $`\omega \omega _02.8/\beta `$ which, in view of the definition $`t=2\pi a/\beta `$, means that $`t2x_0`$. Now, in Figs. 1-3 the linear scale chosen for the temperature implies that the intervals covered on the ordinate axes are small, and the mentioned effect becomes hidden (the โshoulderโ in each diagram appears to lie at a lower value of $`t`$). The effect becomes however clear if we replace the linear temperature scale by a logarithmic scale. Figure 4 shows, as an example, such a diagram for the case $`x_0=100`$. Then the ordinate interval becomes considerably larger, and we see that the prediction $`t2x_0`$ for the position of the โshoulder turns out to be reasonable. (In the diagram we have drawn the line $`t=2x_0`$ to emphasize the transitional region.) The same behaviour was previously found to hold in .
We note in passing that the high temperature entropy $`S`$ turns out to be constant: from Eq. (37) we get
$$S(tx_0)=\frac{F(T)}{T}=\frac{\xi ^2}{4}\underset{l=1}{\overset{x_0}{}}\frac{1}{\nu }.$$
(40)
## 5 An improved dispersion relation
Instead of the formal divergence encountered in the previous section and the necessity to truncate the sum over $`l`$ in Eq. (28) at an upper limit $`l=l_0`$, it is of interest to consider a mathematically more involved but physically more correct dispersion relation, that leads to a finite result for the free energy even after summing over all $`l`$ up to infinity.
Let us consider the following dispersion relation, being essentially of the Lorentz form but augmented by a term in the denominator that describes spatial dispersion:
$$\mu (i\widehat{\omega },l)1=\frac{\chi _0}{1+(x^2+l^2)/x_0^2},$$
(41)
$`\chi _0`$ being a constant. As in the previous section we assume that the ball of radius $`a`$ is dilute ($`|\chi _0|1`$), and that the exterior region is a vacuum. The spatial dispersion term $`l^2/x_0^2`$ ensures convergence when summing over all $`l`$. Since $`\xi =(1\mu )/(1+\mu )`$ according to (7), we obtain to first order in $`\chi _0`$, at finite temperatures,
$$\xi \xi _{ln}=\frac{\chi _0/2}{1+(n^2t^2+l^2)/x_0^2}.$$
(42)
Instead of the expression (28) we now get
$$F(T)=\frac{\xi _0^2x_0^4t}{8\pi a}\underset{l=1}{\overset{\mathrm{}}{}}\left[\frac{1}{\nu (x_0^2+l^2)^2}+\underset{n=1}{\overset{\mathrm{}}{}}\frac{2\nu ^5}{(x_0^2+l^2+n^2t^2)^2(\nu ^2+n^2t^2)^3}\right],$$
(43)
where we have introduced the symbol $`\xi _0=\chi _0/2`$. We resolve the last term into partial fractions, and introduce the symbol $`S(p;k)`$:
$$S(p;k)=\underset{n=1}{\overset{\mathrm{}}{}}(n^2+p^2)^k$$
(44)
Then, if $`p`$ and $`q`$ are arbitrary constants, we get
$`{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{(n^2+p^2)^2(n^2+q^2)^3}}`$ $`=`$ $`{\displaystyle \frac{3}{(p^2q^2)^4}}S(p;1){\displaystyle \frac{1}{(p^2q^2)^3}}S(p;2)`$ (45)
$`+`$ $`{\displaystyle \frac{3}{(p^2q^2)^4}}S(q;1){\displaystyle \frac{2}{(p^2q^2)^3}}S(q;2)`$
$`+`$ $`{\displaystyle \frac{1}{(p^2q^2)^2}}S(q;3).`$
From Appendix A in we write down the expressions for $`S(p;k)`$; $`k=1,2`$:
$$S(p;1)=\frac{1}{2p^2}+\frac{\pi }{2p}\mathrm{coth}\pi p,$$
(46)
$$S(p;2)=\frac{1}{2p^4}+\frac{\pi }{4p^3}\mathrm{coth}\pi p+\frac{\pi ^2}{4p^2}\frac{1}{\mathrm{sinh}^2\pi p}$$
(47)
($`S(p;3)`$ was given in Eq. (30)). We define, in analogy with Eq. (38), the nondimensional reduced free energy to be
$$F_{red}(t)\frac{F(T)}{3\xi _0^2x_0/(64a)},$$
(48)
and obtain after some manipulations the following expression:
$$\begin{array}{cc}\hfill F_{red}(t)& =\frac{16x_0^3t}{3\pi }\underset{l=1}{\overset{\mathrm{}}{}}\{\frac{1}{2\nu (x_0^2+l^2)^2}\hfill \\ & \frac{3\nu ^5}{t^2(x_0^2l\frac{1}{4})^4}\left[S(\frac{\sqrt{x_0^2+l^2}}{t};1)S(\frac{\nu }{t};1)\right]\hfill \\ & \frac{\nu ^5}{t^4(x_0^2l\frac{1}{4})^3}\left[S(\frac{\sqrt{x_0^2+l^2}}{t};2)+2S(\frac{\nu }{t};2)\right]\hfill \\ & +\frac{\nu ^5}{t^6(x_0^2l\frac{1}{4})^2}S(\frac{\nu }{t};3)\}.\hfill \end{array}$$
(49)
This expression is calculated numerically for two of the cases above, viz. $`x_0=\{10,100\}`$. The results are shown in Fig. 5.
It is seen that the curves have roughly the same form as in Figs. 1 - 2. There is a definite low - temperature plateau, and for high values of $`t`$ the curves tend towards a linear form. In any case, whatever we adopt a square dispersion relation or a spatial - dispersion relation, the sensitivity with respect to increasing values of $`t`$ is most pronounced when $`x_0`$ is smallest. The limiting values for $`t0`$ in Fig. 5 are in agreement with those of Figs. 1 and 2 to within about 70 per cent. This is the order of magnitude - agreement that we might expect, since the parameter $`\xi _0`$ made use of in Eq. (48) is similar to, but not identical with, the nondispersive quantity $`\xi `$ made use of in Eq. (38).
It is of interest to give analytic approximations for the case of low temperatures, $`t1`$. Then the quantities $`p=\sqrt{x_0^2+l^2}/t`$ and $`q=\nu /t`$ are large, and it becomes convenient to make use of the approximate formula
$$\underset{n=1}{\overset{\mathrm{}}{}}\frac{1}{(n^2+p^2)^k}=\frac{\sqrt{\pi }\mathrm{\Gamma }(k\frac{1}{2})}{2\mathrm{\Gamma }(k)p^{2k1}}\frac{1}{2p^{2k}}+\left(\frac{\pi }{p}\right)^k\frac{e^{2\pi p}}{\mathrm{\Gamma }(k)},$$
(50)
from which it follows, in our notation, that
$$S(p;1)=\frac{\pi }{2p}\frac{1}{2p^2}+\frac{\pi }{p}e^{2\pi p},$$
(51)
$$S(p;2)=\frac{\pi }{4p^3}\frac{1}{2p^4}+\left(\frac{\pi }{p}\right)^2e^{2\pi p},$$
(52)
$$S(p;3)=\frac{3\pi }{16p^5}\frac{1}{2p^6}+\frac{1}{2}\left(\frac{\pi }{p}\right)^3e^{2\pi p}.$$
(53)
Inserting Eqs. (51)-(53) into Eq. (49) we obtain a reasonable simple series expression for $`F_{red}(t)`$, at low temperatures. A further simplification can be obtained if we keep only the leading order terms in Eqs. (51)-(53). For $`p1`$ this means that we keep only the first term in each equation. We then get in this limit
$`F_{red}(t0)`$ $`=`$ $`x_0^3{\displaystyle \underset{l=1}{\overset{\mathrm{}}{}}}\{{\displaystyle \frac{8\nu ^5}{(x_0^2l\frac{1}{4})^4}}({\displaystyle \frac{1}{\sqrt{x_0^2+l^2}}}{\displaystyle \frac{1}{\nu }})`$ (54)
$`+`$ $`{\displaystyle \frac{8\nu ^5}{3(x_0^2l\frac{1}{4})^3}}({\displaystyle \frac{1}{2(x_0^2+l^2)^{3/2}}}+{\displaystyle \frac{1}{\nu ^3}}){\displaystyle \frac{1}{(x_0^2l\frac{1}{4})^2}}\}.`$
In this expression, the temperature $`t`$ is no longer present. Numerical trials show that Eq. (54) is very accurate for $`t1`$. Choosing $`x_0=10,t=1`$, we obtain from Eq. (54) $`F_{red}(x_0=10,t=1)=0.6526`$, whereas the full formula (49) yields the number -0.6530, i. e., an error of 0.06 per cent. Even with $`t=2`$, where Eq. (49) yields the number -0.678, the error in Eq. (54) is increased to only about 3 per cent.
## 6 Conclusion and final remarks
Let us summarize as follows:
(1) For a relativistic medium, i.e., a medium satisfying the condition $`ฯต\mu =1`$, the $`T=0`$ result for the Casimir energy is most conveniently written in the form (8). This expression holds for a nondispersive medium, $`\delta =\widehat{\tau }/a`$ being the time-splitting parameter. The permittivity is here arbitrary; the medium need not be dilute. The regularization is made by subtracting off the volume terms in the two scalar Green functions. This regularization is natural, among else things because it permits one to relate $`E`$ to the surface force density $`f`$ by the relation $`E/a=4\pi a^2f`$. The equivalence between the Green function approach and the mode sum approach is demonstrated explicitly.
(2) As an additional by-result at $`T=0`$, the expression (15) generalizes the second order result for $`E`$ found in up to the fourth order in the parameter $`\xi `$.
(3) Including the dispersive effect, the simplest dispersion relation one can imagine is the โsquareโ relation given in Eq. (21). The finite-temperature expression for the free energy on the basis of this dispersion relation is given in Eq. (31), with $`x_0=\omega _0a`$ being a high frequency cutoff which is in practice dependent on the detailed structure of the medium. With reasonable accuracy $`x_0`$ can represent the maximum $`l_0`$ of the angular momentum variable $`l`$. Typical results are shown in Figs. 1 - 4. In the limiting case $`t0`$ they are in accordance with Eq. (33), and for $`tx_0`$ they are in accordance with Eq. (39). In particular, the free energy is for high temperatures negative, and is a linear function of $`t`$. The entropy is correspondingly a (positive) constant. These results are in accordance with those of Refs. and . The transition region between low-temperature and high-temperature theory can be determined with reasonable accuracy from a simple physical argument.
These results are similar to those obtained on the basis of the Lorentz dispersion relation in .
(4) One can take spatial dispersion into account, as we have done in section 5. The basic dispersion relation is then Eq. (41). A physically appealing feature of this kind of approach is that the sum over angular momenta does not have to be truncated. Typical results for this case are shown in Fig. 5.
(5) The condition $`ฯต\mu =1`$ simplifies the theory, as noted already in Ref. . For an ordinary dielectric ball ($`\mu =1`$) the calculation becomes more complex and difficult to interpret .
(6) A final remark is called for, as regards the numerical computation of the series in the case of spatial dispersion. The evaluation of the expression (49) turned out to be more difficult than one might expect beforehand. Thus, a simple use of Matlab turned out to be insufficient. The problems seem to be associated with numerical noise, caused by the โcriticalโ terms for which $`lx_0^2`$. There are relatively large individual terms that almost, but not quite, compensate each other in the sum. We managed to do the calculation making use of double - precision MS-DOS Quick Basic. In practice, adequate precision was found with inclusion of less than $`x_0^2`$ terms.
Acknowledgment
We thank Michael Revzen and Israel Klich for valuable information about the finite temperature problem. |
warning/0002/hep-lat0002007.html | ar5iv | text | # 1 Introduction
## 1 Introduction
Quark masses are fundamental parameters of QCD that cannot be determined by theoretical considerations only and, due to the confinement of quarks inside hadrons, cannot be directly measured. Quark masses can, however, be introduced as short-distance effective couplings. As such, they are scale and scheme dependent quantities, the values of which depend on the adopted definition. Nonetheless, quark masses are very important for phenomenology since they enter many theoretical predictions of physical quantities such as CKM matrix elements, $`b`$-hadron inclusive semileptonic decays, total widths, etc. This is the reason why, in the last years, much effort has been devoted to accurately determine their values.
It is useful to classify quarks into two classes: light quarks, the masses of which are lower or of the order of $`\mathrm{\Lambda }_{QCD}`$ (the $`u`$, $`d`$ and $`s`$ quarks are light), and heavy quarks, with masses larger than $`\mathrm{\Lambda }_{QCD}`$ ( $`b`$ and $`t`$ quarks are heavy and, to some extent, the charm quark $`c`$ too). Light-quark masses are extracted from hadron spectroscopy using lattice QCD simulations and also QCD sum rules and $`\tau `$ decay data . Heavy-quark masses can be extracted from the properties of hadrons containing heavy quarks: the $`B`$-meson spectrum from the lattice HQET , the $`\mathrm{{\rm Y}}`$ (or $`J/\psi `$) spectrum with lattice NRQCD or QCD Sum Rules and mass effects in 3โjets $`b\overline{b}g`$ events .
In this paper, we present the first unquenched HQET lattice calculation of the $`b`$ quark mass. The idea (see also ) is to combine the HQET unquenched lattice computation of the $`B`$-meson binding energy with the recent next-to-next-to-leading ( NNLO ) perturbative calculation of the matching of the $`\overline{MS}`$ quark mass to its lattice HQET counterpart . We stress that both unquenched simulations and NNLO matching are necessary ingredients to improve the accuracy of previous results . The former is necessary to control potentially large vacuum-polarization contributions to the $`B`$-meson propagator. The latter is crucial to reduce renormalon ambiguities in the continuum-lattice matching . After a careful analisys of the systematics errors, our best result is
$$\overline{\textcolor[rgb]{1,0,0}{m}}_\textcolor[rgb]{1,0,0}{b}\textcolor[rgb]{1,0,0}{(}\overline{\textcolor[rgb]{1,0,0}{m}}_\textcolor[rgb]{1,0,0}{b}\textcolor[rgb]{1,0,0}{)}\textcolor[rgb]{1,0,0}{=}\textcolor[rgb]{1,0,0}{(}\textcolor[rgb]{1,0,0}{\mathrm{\hspace{0.17em}4.26}}\textcolor[rgb]{1,0,0}{\pm }\textcolor[rgb]{1,0,0}{\mathrm{\hspace{0.17em}0.03}}\textcolor[rgb]{1,0,0}{\pm }\textcolor[rgb]{1,0,0}{0.05}\textcolor[rgb]{1,0,0}{\pm }\textcolor[rgb]{1,0,0}{\mathrm{\hspace{0.17em}0.07}}\textcolor[rgb]{1,0,0}{)}\textcolor[rgb]{1,0,0}{\mathrm{GeV}}\textcolor[rgb]{1,0,0}{,}$$
where the first error is statistical; the second is the systematic error from the spread of values due to the use of different time intervals, fitting methods, smearing types and cube sizes for the interpolating operators, the dependence of the results on the mass of the sea quarks, the calibration of the lattice spacing and an evaluation of the $`1/m_b`$ corrections; the third is an estimate of the error due to the uncertainties in the values of $`\alpha _s`$ and to the effects of higher-order terms in eq. (13). A detailed discussion of the different errors can be found below.
The paper is organized as follows: in Sect. 2, we briefly describe our method and give the main formulae we used; in Sect. 3, we discuss the lattice computation of the binding energy; details of the simulation and numerical results are presented in Sect. 4, where we also discuss the procedure used for analyzing the unquenched lattice data; in Sect. 5, we carefully study the different sources of systematic errors in our results. Finally, in Sect. 6, we present our final numbers and compare them with other recent determinations.
## 2 The method.
The key idea to determine the $`b`$-quark mass consists in matching the propagator in QCD to its lattice HQET counterpart . As shown below, this matching allows us to relate the pole mass to the binding energy and to the physical mass of the $`B`$ meson. The renormalized $`\overline{MS}`$ $`b`$-quark mass at a given scale $`\mu `$ can then be obtained from the pole mass by using perturbation theory. In this section we briefly recall the formulae relevant to our study.
Lattice HQET is an effective theory of QCD. The relation between the inverse $`b`$-quark propagator in QCD, $`S^1`$, and its lattice HQET counterpart, $`S_L^1`$, can be written, to lowest order in $`1/m_b`$, as
$$\left(\frac{1+v\text{/}}{2}\right)S_P^1(p,m_b;\mu )\left(\frac{1+v\text{/}}{2}\right)S^1\left(\frac{1+v\text{/}}{2}\right)=C(\mu a,\alpha _s)\left(\frac{1+v\text{/}}{2}\right)S_L^1((vk)a),$$
(1)
where $`S_P^1`$ is the projected $`b`$-quark propagator, $`\mu `$ the renormalization point, $`a`$ the lattice spacing, $`p=m_bv+k`$ the momentum of the $`b`$-quark, $`v`$ its velocity and $`k`$ the residual momentum. $`C(\mu a,\alpha _s)`$ is the relevant Wilson coefficient. It contains all the mass dependence of the right hand side of eq. (1) since, by construction, the HQET propagator is independent of the $`b`$ quark mass. It should be noticed that in order for the HQET to be applicable, $`k`$ must satisfy the condition $`|k|m_b`$. In writing eq. (1) we have chosen as expansion parameter the quark mass appearing in the original propagator, namely $`m_b`$.
The procedure to find $`C`$ is well known: calculate the $`b`$-quark propagator in QCD and in the lattice HQET to a given order in $`\alpha _s`$, expand the former in inverse powers of $`m_b`$ to a given order (lowest order in our case), and finally compare both expressions at a fixed scale $`\mu `$ (with $`\mu \mathrm{\Lambda }_{QCD}`$) to extract $`C(\mu a,\alpha _s)`$. Renormalization group can then be used to evolve this function to any scale.
To illustrate and clarify the key points of this strategy, we briefly sketch the derivation of our master formula for the $`b`$-quark mass to $`O(\alpha _s)`$ and then we extend our equation to include higher orders in perturbation theory. The inverse quark propagator in QCD can be written in the form
$$iS^1(p,m_b;\mu )=p\text{/}m_b+\mathrm{\Sigma }_1(p^2,m_b)+(p\text{/}m_b)\mathrm{\Sigma }_2(p^2,m_b).$$
(2)
It is very easy to calculate the self-energy form factors $`\mathrm{\Sigma }_1`$ and $`\mathrm{\Sigma }_2`$ to one loop in some renormalization scheme and for a fixed gauge. By writing the $`b`$-quark momentum in the $`B`$ meson, $`p`$, as $`p=m_bv+k`$ and expanding in powers of $`1/m_b`$, one finds
$`S_P^1(p,m_b;\mu )`$ $`=`$ $`m_bm_b^{pole}`$ (3)
$`+`$ $`(vk)\left[1+\alpha _s(\mu )\left(\gamma _{Dk}\mathrm{ln}\left({\displaystyle \frac{\mu }{2(vk)}}\right)+c_2\right)+\mathrm{}\right],`$
where $`c_2`$ is a scheme dependent constant (the expression of which is irrelevant for our discussion) and $`\gamma _m`$ and $`\gamma _{Dk}`$ are the scheme-independent one-loop anomalous dimensions of the mass and operator $`Dk`$, respectively.
The pole mass, $`m_b^{pole}`$, is defined as the position of the pole of the propagator $`S^1`$, at a given order in perturbation theory,
$`S_P^1(p\text{/}=m_b^{pole},m_b;\mu )=\mathrm{\hspace{0.17em}0}.`$ (4)
To one loop, the explicit calculation of eq. (3) gives
$$m_b^{pole}=m_b\left[1+\alpha _s(\mu )\left(\gamma _m\mathrm{ln}\left(\frac{m_b}{\mu }\right)+c_1\right)\right].$$
(5)
In order to implement the matching $`S_L^1`$, the propagator of the lattice HQET, must be evaluated at the same order in perturbation theory
$$S_L^1((vk)a)=(vk)\left[1+\alpha _s(a)\left(\gamma _{Dk}\mathrm{ln}\left(\frac{1}{2(vk)a}\right)+d_2\right)\right]\alpha _s(a)\frac{X_0}{a}+\mathrm{}$$
(6)
where $`d_2`$ is a scheme dependent constant. From eq. (6) we learn that an additive, linearly divergent mass term is generated on the lattice: the so-called residual mass, $`\delta m`$. Inserting eq. (6) into eq. (3) and taking into account the expression of the pole mass in eq. (5), we obtain
$$S_P^1(p,m_b;\mu )=m_bm_b^{pole}+\alpha _s(\mu )\frac{X_0}{a}+C(\mu a,\alpha _s)S_L^1((vk)a),$$
(7)
where the Wilson coefficient has the form
$$C(\mu a,\alpha _s)=\mathrm{\hspace{0.17em}1}+\alpha _s(\mu )\left(\gamma _{Dk}\mathrm{ln}\left(\mu a\right)+c_2d_2\right)$$
(8)
and we have used the fact that the difference between $`\alpha _s(\mu )`$ and $`\alpha _s(a)`$ is $`๐ช(\alpha _s^2)`$. Comparing eq. (1) with eq. (7), the important relation between the HQET expansion mass parameter $`m_b`$ and the pole $`b`$-quark mass can be derived
$$m_b^{pole}=m_b+\alpha _s(\mu )\frac{X_0}{a}m_b+\delta m.$$
(9)
To lowest order in $`1/m_b`$, the HQET mass formula can now be used to eliminate the unknown expansion parameter, $`m_b`$, by expressing it in terms of the physical mass of a $`b`$-hadron, specifically the $`B`$-meson, and the non-perturbative binding energy, $``$, which is independent of $`m_b`$,
$$M_B=m_b++๐ช(1/m_b)$$
(10)
Using the equation above, we get
$$m_b^{pole}=M_B+\delta m+๐ช(1/m_b).$$
(11)
Finally, the pole mass is converted into the $`\overline{MS}`$ mass through the well-known one-loop perturbative relation
$`\overline{m}_b(\overline{m}_b)`$ $`=`$ $`m_b^{pole}\left[1{\displaystyle \frac{4}{3}}\left({\displaystyle \frac{\alpha _s(\overline{m}_b)}{\pi }}\right)\right]`$ (12)
$`=`$ $`\left[M_B+\alpha _s(\overline{m}_b){\displaystyle \frac{X_0}{a}}\right]\left[1{\displaystyle \frac{4}{3}}\left({\displaystyle \frac{\alpha _s(\overline{m}_b)}{\pi }}\right)\right]+๐ช(1/m_b).`$
We stress that $`\overline{m}_b(\overline{m}_b)`$ is obtained from the non-perturbative quantity $``$ combined with the perturbative calculation of lattice ($`\alpha _sX_0`$) and continuum ($`4/3\alpha _s(\overline{m}_b)/\pi `$) coefficients.
The generalization of eq. (12) to higher orders is straightforward. One gets
$`\overline{m}_b(\overline{m}_b)`$ $`=`$ $`m_b^{pole}\left[1+{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}\left({\displaystyle \frac{\alpha _s(\overline{m}_b)}{\pi }}\right)^{n+1}D_n\right]+๐ช(1/m_b)`$ (13)
$`=`$ $`\left[M_B+{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}(\alpha _s(\overline{m}_b))^{n+1}{\displaystyle \frac{X_n}{a}}\right]\left[1+{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}\left({\displaystyle \frac{\alpha _s(\overline{m}_b)}{\pi }}\right)^{n+1}D_n\right],`$
which is the master equation of our analysis. $`D_n`$ and $`X_n`$ are constants which depend on the number of flavours, $`n_f`$, the masses of the active quarks and the lattice action used for the light quarks (see below). The procedure used to calculate $`\overline{m}_b(\overline{m}_b)`$ is the following :
* compute the binding energy of the $`B`$ meson in lattice units, $`a`$, using the HQET for the heavy quark and a given discretized action (Wilson, Alpha, Staggered) to describe the light-quark dynamics (we discuss in detail this computation in Sects. 3 and 4).
* evaluate the value of the lattice spacing, $`a`$, from the light-hadron spectroscopy.
* take the experimental value of the $`B`$-meson mass, $`M_B`$, as input.
* insert the values of these quantities in eq. (13) and obtain $`\overline{m}_b(\overline{m}_b)`$ to a given order in perturbation theory and up to $`๐ช(1/m_b)`$ corrections.
A few important remarks are in order at this point:
1. The bare binding energy $``$ is not a physical quantity since it diverges linearly as $`a0`$ and needs to be subtracted . Since both the pole and the $`\overline{MS}`$ mass are finite quantities, instead, the divergence of $``$ is cancelled by the corresponding divergence of the residual mass $`\delta m`$, expressed in terms of the constants $`X_n`$.
2. In practice the cancellation is incomplete because we only know the values of few constants $`X_n`$. Therefore, at a given order of the perturbative expansion, we cannot take the lattice spacing too small.
3. The large-$`n_f`$ approximation shows that the perturbative series for $`\delta m`$, and hence for the pole mass, suffers from renormalon singularities . In other words, the coefficients $`X_n`$ are expected to grow as $`const.\times n!`$ as $`n\mathrm{}`$. These singularities give rise to ambiguities of $`๐ช(\mathrm{\Lambda }_{QCD})`$ in the sum of the perturbative series. A solution to this problem, which is the one adopted here, is to consider a short distance definition of the $`b`$-quark mass, such as the $`\overline{MS}`$, $`\overline{m}_b`$, because it is free of renormalon ambiguities (up to $`๐ช(\mathrm{\Lambda }_{QCD}^2/m_b)`$ at the order at which we are working).
4. In the expression of the $`\overline{MS}`$ mass, a delicate cancellation of renormalon singularities occurs: the renormalon of the series for $`\delta m`$ (with coeffcients $`X_n`$) is cancelled by the perturbative expansion of the coefficient relating the pole and the $`\overline{MS}`$ mass (with coeffcients $`D_n`$).
5. In order to achieve the cancellation of renormalon singularities in eq. (13), the same coupling constant has to be used in the expansion of $`\delta m`$ and in the relation between the pole mass and the $`\overline{MS}`$ mass . For this reason, although we work at a fixed order of perturbation theory, we believe that the most reasonable choice is to expand both the continuum and lattice series using the same coupling constant.
From the discussion above, it is clear that the precision of our results for the $`b`$-quark mass at given $`(,a)`$ is limited by the number of terms calculated in lattice ($`X_n`$) and in continuum ($`D_n`$) perturbation theory. The relation between the pole and the $`\overline{MS}`$ mass has been obtained to $`๐ช(\alpha _s^2)`$ by Gray *et al.* and, recently, to $`๐ช(\alpha _s^3)`$ by refs.
$`D_0`$ $`=`$ $`{\displaystyle \frac{4}{3}},`$
$`D_1`$ $`=`$ $`11.6656+\mathrm{\hspace{0.17em}1.0414}{\displaystyle \underset{i=1}{\overset{n_f}{}}}\left[1{\displaystyle \frac{\overline{m}_i}{\overline{m}_b}}\right],`$
$`D_2`$ $`=`$ $`157.116+\mathrm{\hspace{0.17em}23.8779}n_f0.6527n_f^2.`$ (14)
Note that the three-loop correction has been evaluated with massless quarks.
As for the residual mass, $`\delta m`$, it can be expressed in terms of the bare lattice coupling, $`\alpha _0`$, as
$$\delta m=\underset{n=0}{\overset{\mathrm{}}{}}(\alpha _0)^{n+1}\frac{\overline{X}_n}{a}$$
(15)
The constant $`\overline{X}_0`$ is simply $`\overline{X}_0=X_0`$ given by the three-dimensional integral
$$X_0=C_F\frac{1}{8\pi ^2}_\pi ^\pi d^3k\frac{1}{2_{i=1}^3\mathrm{sin}^2(k_i/2)}=\mathrm{\hspace{0.17em}2.1173}$$
(16)
where $`C_F=(N^21)/2N`$ and $`N`$ is the number of colors. Martinelli and Sachrajda have performed the calculation of $`\overline{X}_1`$ by extracting $`\delta m`$ from the exponential decrease of the expectation value of large Wilson loops with the perimeter . More recently, Burgio *et al.* have obtained a preliminary estimate of $`\overline{X}_2`$ from large Wilson loops computed with the Numerical Stochastic Perturbation Theory (NSPT) on a $`24^4`$ lattice in the quenched approximation . In summary, the results are
$`\overline{X}_0`$ $`=`$ $`2.1173`$
$`\overline{X}_1`$ $`=`$ $`11.152+n_f(\mathrm{\hspace{0.17em}0.282}+\mathrm{\hspace{0.17em}0.035}c_{SW}0.391c_{SW}^2)`$
$`\overline{X}_2`$ $`=`$ $`73.5(9.2)`$ (17)
where the value of the coefficient $`c_{SW}`$ depends on the lattice fermions used in the simulation: for Wilson fermions $`c_{SW}=0`$, for Clover-SW tree-level improved fermions $`c_{SW}=1`$ and for the non-perturbatively improved ones, $`c_{SW}`$ depens on $`\beta `$ (see for details). The numerical value of $`\overline{X}_2`$ has been obtained in the quenched approximation ($`n_f=0`$) and thus it does not include fermion-loop effects.
The next step is to express $`\delta m`$ in terms of the $`\overline{MS}`$ coupling $`\alpha _s`$, i.e. to calculate the coefficients $`X_n`$ from the $`\overline{X}_n`$ of eq. (17). The relation between $`\alpha _s`$ and $`\alpha _0`$ can be written, to $`๐ช(\alpha _0^3)`$, as:
$$\alpha _s\left(\mu \right)=\alpha _0+d_1(\mu a)\alpha _0^2+d_2(\mu a)\alpha _0^3+\mathrm{}$$
(18)
The pure gauge contributions to $`d_1`$ and $`d_2`$ have been calculated in ref. and respectively. The quark contribution to $`d_1`$ for Wilson fermions can be found in and for improved fermions with generic $`c_{SW}`$ in (see also ). Unfortunately, the quark contribution to the two-loop coefficient $`d_2`$ is still unknown. This calculation is necessary if $`X_2`$ is used to obtain the N<sup>3</sup>LO mass in the unquenched case. So far we have
$`d_1(x)`$ $`=`$ $`{\displaystyle \frac{\beta _0}{2\pi }}\mathrm{ln}(x){\displaystyle \frac{\pi }{2N}}+\mathrm{\hspace{0.17em}2.13573}N`$ (19)
$`+`$ $`n_f(0.08413+\mathrm{\hspace{0.17em}0.0634}c_{SW}\mathrm{\hspace{0.17em}0.3750}c_{SW}^2)`$
$`d_2(x)`$ $`=`$ $`d_1(x)^2{\displaystyle \frac{17N^2}{12\pi ^2}}\mathrm{ln}(x)+{\displaystyle \frac{3\pi ^2}{8N^2}}\mathrm{\hspace{0.17em}2.8626216}+\mathrm{\hspace{0.17em}1.249116}N^2`$
where $`\beta _0=\frac{1}{3}(11N\mathrm{\hspace{0.17em}2}n_f)`$ and $`d_2`$ is given for $`n_f=0`$. By inverting eq. (18) for $`\alpha _0`$ and inserting it in eq. (15), for $`N=3`$ we get the values of the constants $`X_{0,1,2}`$:
$`X_0`$ $`=`$ $`2.1173`$
$`X_1`$ $`=`$ $`1.30533+\mathrm{\hspace{0.17em}3.70677}\mathrm{ln}(\overline{m}_ba)`$
$`+`$ $`n_f(0.1038720.224653\mathrm{ln}(\overline{m}_ba)0.0992368c_{SW}+0.402988c_{SW}^2)`$
$`X_2`$ $`=`$ $`\overline{X}_2+\mathrm{\hspace{0.17em}6.48945}(3.57877+\mathrm{ln}(\overline{m}_ba))(3.29596+\mathrm{ln}(\overline{m}_ba))`$ (20)
Since the value of $`X_2`$ is preliminary, we cannot really use it to obtain our final result. We will only show that, even in the quenched case, the result of ref. is not precise enough to obtain a useful information.
## 3 Lattice computation of $``$.
The bare binding energy of the $`B`$ meson, $``$, is measured on the lattice by studying the large-time behaviour of the $`B`$ meson propagator . The $`b`$-quark is described by the discretized version of the HQET,
$$_{HQET}=\overline{b}(x)D_4b(x),$$
(21)
with the covariant derivative defined as
$$D_4b(x)=U_\mu (x)b(x+\widehat{\mu })b(x),$$
(22)
where $`\widehat{\mu }`$ indicates the $`\mu `$-direction and $`U_\mu (x)`$ is the link variable between the lattice sites $`x`$ and $`x+\widehat{\mu }`$. Light quarks, $`q`$, are simulated with some fermion action, in our case the Wilson action.
It is well known that correlation functions involving heavy quarks suffer from a large contamination from higher-mass excitations to the lightest-state contribution, to which we are interested in. In order to improve the isolation of the lightest state, we use cube and double-cube smeared axial-current operators of size $`L_s`$, as interpolating operators of the B meson :
$`A_\mu ^L(x)`$ $`=`$ $`\overline{b}(x)\gamma _\mu \gamma _5q(x),`$
$`A_\mu ^S(x)`$ $`=`$ $`{\displaystyle \underset{i}{\overset{L_s}{}}}\overline{b}(x_i)\gamma _\mu \gamma _5q(x),`$ (23)
$`A_\mu ^D(x)`$ $`=`$ $`{\displaystyle \underset{i,j}{\overset{L_s}{}}}\overline{b}(x_i)\gamma _\mu \gamma _5q(x_j).`$
From the above operators we construct the two-point correlation functions
$$C_{L_s}^{RR^{}}(t)=\underset{\stackrel{}{x}}{}0|A_4^R(\stackrel{}{x},t)A_4^R^{}(\stackrel{}{0},0)|0,$$
(24)
where $`R,R^{}`$ stands for $`L`$ (Local), $`S`$ (Single smeared) and $`D`$ (Double smeared) interpolating operators. The correlation functions are computed after rotating the links in the Coulomb gauge.
At large time distances, $`C_{L_s}^{RR^{}}(t)`$ behaves as:
$$C_{L_s}^{RR^{}}(t)Z_{L_s}^RZ_{L_s}^R^{}e^t,$$
(25)
where
$$Z_{L_s}^R=\frac{1}{\sqrt{2M_B}}0|A_4^R(\stackrel{}{0},0)|B.$$
(26)
By fitting the large time behaviour of $`C^{RR^{}}(t)`$ to eq. (25), the bare binding energy can be extracted. Due to contamination from excited states, the interpolating operators couple to the ground state in different ways so that the actual value of $``$ has some dependence on the cube size and the smearing type used in the analysis. To obtain our best estimate, we compare different methods and account the mixing as a systematic effect contributing to the final error.
We first use the Best Cube Method (BCM): we base our results on the best cube, defined as the one which yields the largest and flattest effective-mass plateau . In practice, we proceed as the following example illustrates. Consider the double smeared interpolating operators. For $`tt_{min}`$ and for each $`L_s`$, we fit $`C_{L_s}^{DD}(t)`$ and $`C_{L_S}^{LD}(t)`$ to a single-state propagator (25). $`t_{min}`$ is taken as the time at which we start observing a plateau for both the effective mass $`\mathrm{\Delta }E_{L_s}^{DD}`$ and the ratio $`\mathrm{R}_{L_s}^{DL/DD}`$ (we call BC DL/DD method the case where the smearing is in the sink) or, alternatively, for the effective mass $`\mathrm{\Delta }E_{L_s}^{LD}`$ and the ratio $`\mathrm{R}_{L_s}^{LD/DD}`$ (we call BC LD/DD method the case where the smearing is in the origin). Effective masses and ratios are defined by
$`\mathrm{\Delta }E_{L_s}^{RR^{}}(t)`$ $`=`$ $`\mathrm{ln}\left({\displaystyle \frac{C_{L_s}^{RR^{}}(t)}{C_{L_s}^{RR^{}}(t+1)}}\right),`$
$`\mathrm{R}_{L_s}^{RR^{}/NN^{}}(t)`$ $`=`$ $`{\displaystyle \frac{C_{L_s}^{RR^{}}(t)}{C_{L_s}^{NN^{}}(t)}}{\displaystyle \frac{Z_{L_s}^RZ_{L_s}^R^{}}{Z_{L_s}^NZ_{Ls}^N^{}}}.`$ (27)
For $`tt_{min}`$, the ground state is assumed to have been isolated. The next step is to combine the exponential fit for $`C_{L_s}^{DD}(t)`$ ($`C_{L_s}^{LD}(t))`$ and the average value of the ratio $`\mathrm{R}_{L_s}^{DL/DD}(t)`$ ($`\mathrm{R}_{L_s}^{LD/DD}(t)`$) in the plateau region to obtain both $``$ and $`Z^L`$, the matrix element for the local axial current. Similarly, the method is applied to single smeared operators.
A drawback of the BCM is that, in practice, we have only few different cubes at disposal (essentially only two cubes are really useful, $`L_s=79`$, as suggested by earlier lattice studies). Nevertheless, this method is able to give a reasonable isolation of the lightest state and an accurate determination of the binding energy (although the method is less efficient for the determination of the matrix elements of the local axial current).
In order to improve the accuracy of our analysis, we have also used the Multifit Method which consists in performing a global fit of the data for all cube sizes and smearing types by minimizing the total $`\chi ^2`$ . Consider, for example, the double-smeared operators. In this case, $`\chi _{total}^2`$ is defined by
$`\chi _{total}^2`$ $`=`$ $`{\displaystyle \underset{t=t_i}{\overset{t_f}{}}}\{{\displaystyle \underset{L_s=7,9}{}}\left({\displaystyle \frac{C_{L_s}^{LD}(t)_{DATA}C_{L_s}^{LD}(t)}{\sigma _{L_S}^{LD}(t)}}\right)`$ (28)
$`+`$ $`{\displaystyle \underset{L_s=7,9}{}}\left({\displaystyle \frac{C_{L_s}^{DD}(t)_{DATA}C_{L_s}^{DD}(t)}{\sigma _{L_S}^{DD}(t)}}\right)^2\}`$
where $`\sigma _{L_s}^{LD(DD)}`$ is the jacknife error of the data points. We also impose the consistency condition that the binding energy $``$ should be the same for all smearing types and cubes sizes. Moreover, in order to reduce the effect from excited states, it is convenient to fit the data by including the contribution of at least one excited state to $`C_{L_s}^{RR^{}}(t)`$.
## 4 Strategy of the unquenched analysis.
Before giving details on the simulation parameters, the calibration of the lattice spacing and the extraction of the binding energy $``$, we discuss the general strategy followed in the analysis of the unquenched data. This is a crucial issue due to the confusion which, we think, exists in the literature.
Our main point, which is justified below, is that the correct procedure consists in performing independent quenched-like calculations of all quantities for each fixed value of $`\textcolor[rgb]{1,0,0}{k}_{\textcolor[rgb]{1,0,0}{s}\textcolor[rgb]{1,0,0}{e}\textcolor[rgb]{1,0,0}{a}}`$, including the spectroscopy, the calibration of the lattice spacing $`\textcolor[rgb]{1,0,0}{a}`$ and the calculation of the relevant matrix elements . Only when all the quantities are expressed in physical units, the results can be extrapolated in the sea quark mass. Thus, one ends up with a different set of lattice parameters for each $`k_{sea}`$, such as the critical kappa, $`k_{cr}`$, the light quark masses $`k_u`$ and $`k_s`$, the lattice spacing and so on. Extrapolations in the valence quark mass, $`k_v`$, should also be performed at fixed $`k_{sea}`$, without ever mixing up the extrapolation in the valence and sea quark masses, which must remain distinct steps of the analysis. In all respects, the value of $`k_{sea}`$ is an external โfieldโ which controls the link dynamics.
The argument is the following. A change in the value of the sea quark mass(es) modifies the value of the effective coupling constant, because the latter receives contributions from virtual-quark loops. A change of the coupling constant may induce a rapid variation of the value of the lattice spacing which depends exponentially on $`\alpha _s`$. Therefore, strictly speaking, lattice results for different $`k_{sea}`$ correspond to different lattice dynamics and are not directly comparable. Only when the results have been converted to physical units, by using the lattice spacing extracted for each $`k_{sea}`$, comparisons and extrapolations are possible. We stress again that a combined (in valence and sea quark masses) chiral extrapolation of lattice quantities, as for example the quark masses, may produce incorrect results because in this way we are mixing results corresponding to different values of the lattice spacing and all the parameters of the extrapolation do depend on $`k_{sea}`$ through $`a`$.
Having explained our strategy, we turn to the numerical results for the binding energy. We have performed an unquenched computation of $``$ with two degenerate sea quarks at two values of their mass, $`k_{sea}=0.1575`$ and $`k_{sea}=0.1580`$. The heavy and light quark propagators have been computed using the set of unquenched link configurations generated by the T$`\chi `$L Collaboration. Details of the simulation can be found in ref. . Light quarks are simulated using the Wilson action whereas heavy quarks are described with the discretized HQET. The parameters of our run are given in Table 1. The calibration of the lattice spacing has been performed using the $`K^{}K`$ lattice-plane method of ref. .
With the BC method, we find that the flattest and largest plateaus for $`\mathrm{\Delta }E_{L_s}^{RR^{}}(t)`$ and $`\mathrm{R}_{L_s}^{RR^{}/NN^{}}(t)`$ correspond to the cube size $`L_s=7`$ and the correlation functions $`LD`$ (smearing in the origin) and $`DD`$. With the Multifit Method, we also obtain that a two-state global fit of the correlations $`LD`$ and $`DD`$ describes very well the data. The agreement for single smeared interpolating operators is, instead, much worse. Therefore we base our results on double smeared operators since this is the most efficient way of isolating the lightest state. Our best estimates of the values of the binding energy $``$ are
$`a_{B_d}`$ $`=`$ $`\{\mathrm{\hspace{0.17em}0.588}(11)(5),\mathrm{\hspace{0.17em}0.606}(15)(2)\}`$
$`a_{B_s}`$ $`=`$ $`\{\mathrm{\hspace{0.17em}0.620}(8)(5),\mathrm{\hspace{0.17em}0.632}(12)(2)\}`$ (29)
for the two values of $`k_{sea}=0.1575`$ and $`0.1580`$, respectively. The first error is statistical and the second systematic. The latter has been obtained from the spread of our results due to different time intervals of the fit, cube sizes and smearing types. Since a full account of the different methods and evaluation of the uncertainties can be found in ref. , we do not give further details here.
## 5 Sources of systematic error.
Using eqs. (13), (14), (20) and (29), we can readily obtain the value of the $`b`$ quark mass. As the value of $`X_2`$ is still preliminary and incomplete, we derive our results with the two-loop formula corresponding to the NNLO matching. In order to evaluate the systematic errors on these results, we carefully studied the different sources of uncertainties coming from the use of eq. (13) at this order: the value of $`\alpha _s`$, higher-order perturbative corrections, input meson mass, method of extracting $``$, value of $`k_{sea}`$ and $`1/m_b`$ corrections. In the following, unless stated otherwise, the central values in the tables correspond to $`k_{sea}=0.1580`$, $``$ has been obtained with the Multifit Method, the input mass is the $`B_s`$ meson mass and a linear chiral extrapolation in $`k_v`$ to $`k_s`$ has been performed. In the tables, the first error is statistical and the second the systematic one based on the spread of the results due to the uncertainty in the lattice spacing, the different forms of writing eq. (13) (see sect. 5.2) and the use of different time-intervals and fitting methods.
### 5.1 Dependence of $`\overline{m}_b(\overline{m}_b)`$ on $`\alpha _s`$.
In order to obtain $`\overline{m}_b(\overline{m}_b)`$ we have to choose the value of $`\alpha _s`$ to be used in the perturbative calculations. At all orders, to cancel the renormalon singularities in eq. (13), the same coupling constant has to be used for the lattice and continuum series. For this reason, although our calculation is truncated at $`๐ช(\alpha _s^2)`$, we prefer to take the same coupling constant for both the lattice and the continuum cases.
One possibility is to consider the physical value of $`\alpha _s`$ obtained by running the experimental coupling $`\alpha _s(M_Z)=0.118`$ down to $`\alpha _s(m_b)`$ with $`n_f=5`$.
A second possibility is to account that our simulation has been performed with $`n_f=2`$ and compute $`\alpha _s`$ at the NLO with a (still to be determined) $`\mathrm{\Lambda }_{QCD}^{n_f=2}`$. In the quenched case, the value of $`\mathrm{\Lambda }_{QCD}^{n_f=0}250`$ MeV has been measured in ref. Since the physical value of $`\mathrm{\Lambda }_{QCD}`$ is expected to be larger than the quenched one, in the second case we have used the NLO value of $`\alpha _s(m_b)`$ obtained by varying $`\mathrm{\Lambda }_{QCD}^{n_f=2}`$ in the range $`[250,\mathrm{\hspace{0.17em}350}]`$ MeV.
The results for different values of $`\alpha _s`$ are presented in Table 2. The dependence on $`\mathrm{\Lambda }_{QCD}^{n_f=2}`$ is very weak: the maximum spread of the values is less than $`20`$ MeV. The difference between the central values obtained with the $`n_f=5`$ and $`n_f=2`$ coupling is of about $`50`$ MeV. We have taken this as a very conservative estimate of the error due to the choice of the coupling constant.
Note that the value of $`\alpha _s(m_b)`$ for $`\mathrm{\Lambda }_{QCD}^{n_f=2}=300`$ MeV, $`\alpha _s(m_b)=0.182`$, corresponds with a very good approximation to the arithmetic (and geometric) average of the quenched ($`n_f=0`$) and physical ($`n_f=5`$) couplings ($`\alpha _s(m_b)=0.15`$ and $`\alpha _s(m_b)=0.22`$ respectively). For this reason our central value for $`\overline{m}_b(\overline{m}_b)`$ is that computed with $`\alpha _s(m_b)=0.182`$.
### 5.2 Dependence of $`\overline{m}_b(\overline{m}_b)`$ on higher orders.
Eq. (13), which is used to evaluate $`\overline{m}_b(\overline{m}_b)`$ to $`๐ช(\alpha _s^2)`$ consists in the product of two factors. We can, then, organize the formula including (not expanding) or excluding (expanding to $`๐ช(\alpha _s^2)`$) the $`๐ช(\alpha _s^3)`$ terms arising from the product. We take the difference between these (formally equivalent) procedures, as an estimate of unknown higher-order terms in perturbation theory (PT). In Table 3, the values of $`\overline{m}_b(\overline{m}_b)`$ obtained from the expanded and not expanded forms of eq. (13) are presented. For $`n_f=2`$ the dependence of our results on higher orders is $`30`$$`40`$ MeV, for $`n_f=5`$ the difference is $`60`$ MeV. From these spreads, we conclude that a fair estimate of the effect of expanding or not expanding eq. (13) is $`50`$ MeV. From this estimate we assume from higher-order terms an error of $`\pm 25`$ MeV.
As best estimate of $`\overline{m}_b(\overline{m}_b)`$, for each choice of the value of the coupling constant, we take the average of the results obtained with the not-expanded and expanded form of the master formula. In this way we have computed the central values given in tables 2, 4, 5, 6 and 7.
In order to get an independent estimate of the systematic uncertainty due to higher-order terms in the perturbative espansion, we also tried to compute $`\overline{m}_b(\overline{m}_b)`$ using for $`\overline{X}_2`$ the preliminary result of eq. (20). In the numerical calculations, obtained with the expanded form of eq. (13), we allowed $`\overline{X}_2`$ to vary in the $`1\sigma `$ interval $`[64.3,\mathrm{\hspace{0.17em}82.7}]`$, obtaining the range of masses given in Table 3 (only the central values are given). The values of $`\overline{m}_b(\overline{m}_b)`$ at order $`๐ช(\alpha _s^2)`$ are also given for comparison. There are huge numerical cancellations occurring in the calculation of $`X_2`$ from $`\overline{X}_2`$ in eq. (20). For this reason, the difference between the NNLO and the (approximate) N<sup>3</sup>LO results varies from about zero to $`180`$ MeV, depending on the value of $`\overline{X}_2`$. This quantity, even in the quenched case, is still affected by such a large uncertainty that it is impossible to use it for any realistic estimate. We urgently call for a more precise determination of $`\overline{X}_2`$ in both the unquenched and quenched cases.
### 5.3 Dependence of $`\overline{m}_b(\overline{m}_b)`$ on the input $`B`$ meson mass.
Consistent values of $`\overline{m}_b(\overline{m}_b)`$ should be obtained using as input either the $`B_d`$ or the $`B_s`$ meson masses. The corresponding values of the binding energy, $`_{B_d}`$ and $`_{B_s}`$, respectively, are used in the two cases. This checks the lattice value of the $`M_{B_d}`$-$`M_{B_s}`$ mass splitting and the smoothness of our chiral extrapolation. Actually, for the $`B_s`$ meson, the physical value of the strange quark mass, corresponding to $`k_s`$, is within the range of valence quark masses (see Table 1) and only a mild interpolation, rather than an extrapolation, is needed. For the $`B_d`$ meson, instead, we have extrapolated almost to the chiral limit. In order to compute the pole mass we have taken $`M_{B_d}=5.279`$ and $`M_{B_s}=5.375`$ GeV . In Table 4, we compare our results for the two cases. The results are nicely compatible and a small difference ($`30`$ MeV) is observed in the two cases.
### 5.4 Dependence of $`\overline{m}_b(\overline{m}_b)`$ on the method for extracting $``$.
In sect. 3, we discussed the two methods used to determine the binding energy $``$: the Best Cube method and the Multifit method. If the lightest state has been well isolated, both methods should give compatible results for the $`b`$ quark mass. In Table 5, we present the value of $`\overline{m}_b(\overline{m}_b)`$ for the BC method (obtained with $`LD/DD`$, smearing in the origin, and $`L_s=7`$, which gives the flattest and largest plateau) and for the Multifit method. Also in this case the results differ by $`30`$ MeV.
### 5.5 Dependence of $`\overline{m}_b(\overline{m}_b)`$ on $`k_{sea}`$.
As discussed before, in order to compare the values obtained for different $`k_{sea}`$ and attempt an extrapolation, we have first to convert the lattice quantities to physical units. In Table 6, the values of $`\overline{m}_b(\overline{m}_b)`$ for either values of $`k_{sea}`$ are given. The dependence of our results on $`k_{sea}`$, in the sea-quark mass region of our simulation, is small. Indeed, taking into account that they correspond to independent simulations, we are not able to observe any dependence on $`k_{sea}`$ within errors. Therefore, the only possible strategy is not to attempt an extrapolation in $`k_{sea}`$ and take the value at the lightest $`k_{sea}`$, i.e. $`k_{sea}=0.1580`$, as the best estimate of the physical value of the $`b`$ quark mass. The difference between the results at the two values of $`k_{sea}`$ is accounted as a systematic effect in the final error.
### 5.6 $`1/m_b`$ correcctions to $`\overline{m}_b(\overline{m}_b)`$.
Our whole analysis is performed to the lowest order of the expansion in $`1/m_b`$. This means that $`๐ช(1/m_b)`$ contributions to the relation between the QCD and the lattice HQET quantities in eqs. (1) and (10) have been neglected. We now make an estimate of the error introduced by these higher-order corrections.
The HQET pseudoscalar mass formula including $`1/m_b`$ corrections is given by
$$M_B=m_b+\frac{\lambda _1}{2m_b}\frac{3\lambda _2}{2m_b}+๐ช(1/m_b^2),$$
(30)
where the parameters $`\lambda _1`$ and $`\lambda _2`$ are matrix elements between $`B`$ states of the kinetic and chromomagnetic operators
$`\lambda _1`$ $``$ $`{\displaystyle \frac{B|\overline{b}\stackrel{}{D}^2b|B}{2M_B}},`$
$`\lambda _2`$ $``$ $`{\displaystyle \frac{B|\overline{b}(\stackrel{}{S}g\stackrel{}{B})b|B}{M_B}},`$ (31)
with $`\stackrel{}{S}`$ the spin operator of the $`B`$-meson and $`\stackrel{}{B}`$ the chromomagnetic field ($`B_i=\frac{1}{2}ฯต_{ijk}G_{jk}`$).
It is straightforward to estimate $`\lambda _2`$ because which is related to the vector-pseudoscalar mass splitting,
$$\lambda _2=\frac{1}{4}\left(M_B^{}^2M_B^2\right)\mathrm{\hspace{0.17em}0.12}GeV^2$$
(32)
The extraction of $`\lambda _1`$, is, instead, more difficult as demonstrated by the spread of values obtained with different approaches: the lepton-energy spectrum in inclusive semileptonic $`B`$ decays using Zero Recoil sum rules, QCD Sum rules, experiment data analysis and the HQET Virial Theorem (see and references therein). It has also been estimated on the lattice using the discretized HQET . Although biased by the lattice results, we prefer a small value for this parameter, in the absence of an accurate determination we let it to vary in the interval $`0.5`$$`0.0`$ GeV<sup>2</sup>. With this range, we find that the contribution of the $`1/m_b`$ corrections to the pole mass is at most $`30`$ MeV. Due to the theoretical uncertainties on $`\lambda _1`$, we do not attempt to correct the $`1/m_b`$ terms but include their effect as a sytematic error on the final result.
### 5.7 Continuum limit of $`\overline{m}_b(\overline{m}_b)`$.
To date, in quenched lattice simulations the binding energy $`a`$ has been computed at three values of $`\beta =6.0`$, $`6.2`$ and $`6.4`$ <sup>1</sup><sup>1</sup>1 Finite volume effects may be present in the results at $`\beta =6.4`$ since the lattice volume was rather small.. In Table 7 we give the values of the $`b`$-quark mass from these quenched simulations for different lattice spacings. The quenched results are computed using the values of the binding energy from the APE Collaboration and NNLO quenched master formula with the coupling constant $`\alpha _s(\overline{m}_b)=0.15`$. The quenched values are very close to our new result with $`n_f=2`$. Although one may argue that there is a (rather mild) tendency towards lower values as $`a`$ decreases, with the present uncertainties we cannot attempt any extrapolation in $`a`$ or realistic estimate of the discretization errors.
## 6 Final result for $`\overline{m}_b(\overline{m}_b)`$ and comparison with other determinations.
We consider as best estimate of $`\overline{m}_b(\overline{m}_b)`$ the value obtained with the pole mass extracted by using the mass of the $`B_s`$ meson and the binding energy $`_{B_s}`$ measured on the lattice through the Multifit method at $`k_{sea}=0.1580`$, by averaging the results of the expanded and not expanded form of eq. (13) and by taking the NLO coupling constant $`\alpha _s`$ computed at NLO with $`n_f=2`$ and $`\mathrm{\Lambda }_{QCD}^{n_f=2}=300`$ MeV. Using the estimate of the different errors discussed in the previous section we then obtain:
$$\overline{\textcolor[rgb]{1,0,0}{m}}_\textcolor[rgb]{1,0,0}{b}\textcolor[rgb]{1,0,0}{(}\overline{\textcolor[rgb]{1,0,0}{m}}_\textcolor[rgb]{1,0,0}{b}\textcolor[rgb]{1,0,0}{)}\textcolor[rgb]{1,0,0}{=}\textcolor[rgb]{1,0,0}{(}\textcolor[rgb]{1,0,0}{\mathrm{\hspace{0.17em}4.26}}\textcolor[rgb]{1,0,0}{\pm }\textcolor[rgb]{1,0,0}{\mathrm{\hspace{0.17em}0.03}}\textcolor[rgb]{1,0,0}{\pm }\textcolor[rgb]{1,0,0}{0.05}\textcolor[rgb]{1,0,0}{\pm }\textcolor[rgb]{1,0,0}{\mathrm{\hspace{0.17em}0.07}}\textcolor[rgb]{1,0,0}{)}\textcolor[rgb]{1,0,0}{\mathrm{GeV}}\textcolor[rgb]{1,0,0}{.}$$
(33)
The first error is statistical. The second is the systematic error obtained from the spread of values due to the use of different time intervals, fitting methods, smearing types and cube sizes for the interpolating operators, the dependence of the results on the $`k_{sea}`$, the calibration of the lattice spacing and the estimate of the $`1/m_b`$ corrections. Finally, the third is an estimate of the error due to the uncertainties in the values of $`\alpha _s`$ and to the effects of higher-order terms in eq. (13). We find that the latter is the most important source of error in the final result. For this reason a big effort must be done to compute the unknown N<sup>3</sup>LO contributions to the residual mass on the lattice and the NNLO matching coefficient between the lattice and continuum $`\alpha _s`$ in the unquenched case. On the numerical side, a non-perturbative calculation of $`\mathrm{\Lambda }_{QCD}^{n_f=2}`$ is also important.
Our new result (33) modifies and improve the previous one obtained from quenched lattice simulations with NLO matching only
$$\overline{m}_b(\overline{m}_b)=(\mathrm{\hspace{0.17em}4.15}\pm \mathrm{\hspace{0.17em}0.05}\pm 0.20)\mathrm{GeV}$$
(34)
where the first error is due to the lattice systematics and the second is an estimate of higher orders.
It is interesting to compare eq. (33) with recent values obtained with completely different approaches as mass effects in 3-jets $`b\overline{b}g`$ events, $`b\overline{b}`$ production cross-section and $`\mathrm{{\rm Y}}`$ spectroscopy. Our final result is in good agreement with most of NNLO estimates, as shown in fig.1. In the figure we also give our world average and error. This average has been obtained by using only the most recent NNLO determinations from $`\mathrm{{\rm Y}}`$ spectroscopy and lattice QCD, i.e. we did not use the results of refs. , either because they have been superseeded by more accurate calculations or because they are only computed at the NLO accuracy. The average is
$$\overline{\textcolor[rgb]{1,0,0}{m}}_\textcolor[rgb]{1,0,0}{b}\textcolor[rgb]{1,0,0}{(}\overline{\textcolor[rgb]{1,0,0}{m}}_\textcolor[rgb]{1,0,0}{b}\textcolor[rgb]{1,0,0}{)}\textcolor[rgb]{1,0,0}{=}\textcolor[rgb]{1,0,0}{\mathrm{\hspace{0.17em}4.23}}\textcolor[rgb]{1,0,0}{\pm }\textcolor[rgb]{1,0,0}{0.07}\textcolor[rgb]{1,0,0}{\mathrm{GeV}}$$
(35)
which corresponds to a relative error of less than $`2\%`$ comparable to the precision on the top quark mass. The masses of the quarks of the heaviest and last discovered generation are, and will probably remain, the most accurately determined quark masses.
## Acknowlegments
We are extremely grateful to all the members of the T$`\chi `$L Collaboration for providing us with the gauge configurations necessary to this study. We thank our collaborators D. Becirevic and V. Lubicz for illuminating discussions on the subject of this paper. V. G. has been supported by CICYT under Grant AEN-96-1718, by DGESIC under Grant PB97-1261 and by the Generalitat Valenciana under Grant GV98-01-80. L. G. has been supported in part under DOE Grant DE-FG02-91ER40676. G. M. and F. R. acknowledge the M.U.R.S.T. and the INFN for partial support. |
warning/0002/cond-mat0002129.html | ar5iv | text | # Cryptographical Properties of Ising Spin Systems
## Abstract
The relation between Ising spin systems and public-key cryptography is investigated using methods of statistical physics. The insight gained from the analysis is used for devising a matrix-based cryptosystem whereby the ciphertext comprises products of the original message bits; these are selected by employing two predetermined randomly-constructed sparse matrices. The ciphertext is decrypted using methods of belief-propagation. The analyzed properties of the suggested cryptosystem show robustness against various attacks and competitive performance to modern cyptographical methods.
Public-key cryptography plays an important role in many aspects of modern information transmission, for instance, in the areas of electronic commerce and internet-based communication. It enables the service provider to distribute a public key which may be used to encrypt messages in a manner that can only be decrypted by the service provider. The on-going search for safer and more efficient cryptosystems produced many useful methods over the years such as RSA (by Rivest, Shamir and Adleman), elliptic curves, and the McEliece cryptosystem to name but a few.
In this Letter, we employ methods of statistical physics to study a specific cryptosystem, somewhat similar to the one presented by McEliece. These methods enable one to study the typical performance of the suggested cryptosystem, to assess its robustness against attacks and to select optimal parameters.
The main motivation for the suggested cryptosystem comes from previous studies of Gallager-type error-correcting codes and their physical properties. The analysis exposes a significantly different behaviour for the two-matrix based codes (such as the MN code) and single-matrix codes, which may be exploited for constructing an efficient cryptosystem.
In the suggested cryptosystem, a plaintext represented by an $`N`$ dimensional Boolean vector $`๐(0,1)^N`$ is encrypted to the $`M`$ dimensional Boolean ciphertext $`๐ฑ`$ using a predetermined Boolean matrix $`G`$, of dimensionality $`M\times N`$, and a corrupting $`M`$ dimensional vector $`๐ป`$, whose elements are 1 with probability $`p`$ and 0 otherwise, in the following manner
$$๐ฑ=G๐+๐ป,$$
(1)
where all operations are (mod 2). The matrix $`G`$ and the probability $`p`$ constitute the public key; the corrupting vector $`๐ป`$ is chosen at the transmitting end. The matrix $`G`$, which is at the heart of the encryption/decryption process is constructed by choosing two randomly-selected sparse matrices $`A`$ and $`B`$ of dimensionality $`M\times N`$ and $`M\times M`$ respectively, defining
$`G=B^1A\text{(mod 2)}.`$
The matrices $`A`$ and $`B`$ are generally characterised by $`K`$ and $`L`$ non-zero unit elements per row and $`C`$ and $`L`$ per column respectively; all other elements are set to zero. The finite, usually small, numbers $`K`$, $`C`$ and $`L`$ define a particular cryptosystem; both matrices are known only to the authorised receiver. Suitable choices of probability $`p`$ will depend on the maximal achievable rate for the particular cryptosystem as discussed below.
The authorised user may decrypt the received ciphertext $`๐ฑ`$ by taking the (mod 2) product $`B๐ฑ=A๐+B๐ป`$. Solving the equation
$$A๐บ+B๐=A๐+B๐ป\text{(mod 2)},$$
(2)
is generally computationally hard. However, decryption can be carried out for particular choices of $`K`$ and $`L`$ via the iterative methods of Belief Propagation (BP), where pseudo-posterior probabilities for the decrypted message bits, $`P(S_i=1|๐ฑ)1iN`$ (and similarly for $`๐`$), are calculated by solving iteratively a set of coupled equations for the conditional probabilities of the ciphertext bits given the plaintext and vice versa. For details of the method used and the explicit equations see .
The unauthorised receiver, on the other hand, faces the task of decrypting the ciphertext $`๐ฑ`$ knowing only $`G`$ and $`p`$. The straightforward attempt to try all possible $`๐ป`$ constructions is clearly doomed, provided that $`p`$ is not vanishingly small, giving rise to only a few corrupted bits; decomposing $`G`$ to the matrices $`A`$ and $`B`$ is known to be a computationally hard problem, even if the values of $`K,C`$ and $`L`$ are known. Another approach to study the problem is to exploit the similarity between the task at hand and the error-correcting model suggested by Sourlas, which we will discuss below.
The treatment so far was completely general. We will now make use of insight gained from our analysis of Gallager-type and Sourlas error-correcting codes to suggest a specific cyptosystem construction and to assess its performance and capabilities. The method used in both analyses is based on mapping the problem onto an Ising spin system Hamiltonian, in the manner discovered by Sourlas, which enables one to analyse typical properties of such systems.
To facilitate the mapping we employ binary representations $`(\pm 1)`$ of the dynamical variables $`๐บ`$ and $`๐`$, the vectors $`๐ฑ`$, $`๐ป`$ and $`๐`$, and the matrices $`A`$, $`B`$ and $`G`$, rather than the Boolean $`(0,1)`$ ones.
The binary ciphertext $`๐ฑ`$ is generated by taking products of the relevant binary plaintext message bits $`J_{i_1,i_2\mathrm{}}=\xi _{i_1}\xi _{i_2}\mathrm{}\zeta _{i_1,i_2\mathrm{}}`$, where the indices $`i_1,i_2\mathrm{}`$ correspond to the non-zero elements of $`B^1A`$, and $`\zeta _{i_1,i_2\mathrm{}}`$ is the corresponding element of the corrupting vector (the indices $`i_1,i_2\mathrm{}`$ corresponds to the specific choice made for each ciphertext bit). As we use statistical mechanics techniques, we consider both plaintext ($`N`$) and ciphertext ($`M`$) dimensionalities to be infinite, keeping the ratio between them $`N/M`$ finite. Using the thermodynamic limit is quite natural here as most transmitted ciphertexts are long and finite size corrections are likely to be small.
An authorised user may use the matrix $`B`$ to obtain Eq.(2). To explore the systemโs capabilities one examines the Gibbs distribution, based on the Hamiltonian
$``$ $`=`$ $`{\displaystyle \underset{<i_1,..,i_K;j_1,..,j_L>}{}}๐_{<i_1,..,i_K;j_1,..,j_L>}\delta [1;๐ฅ_{<i_1,..,i_K;j_1,..,j_L>}`$ (3)
$``$ $`S_{i_1}\mathrm{}S_{i_K}\tau _{j_1}\mathrm{}\tau _{j_L}]{\displaystyle \frac{F_s}{\beta }}{\displaystyle }_{i=1}^NS_i{\displaystyle \frac{F_\tau }{\beta }}{\displaystyle }_{j=1}^M\tau _j.`$ (4)
The tensor product $`๐_{<i_1,..,i_K;j_1,..,j_L>}๐ฅ_{<i_1,..,i_K;j_1,..,j_L>}`$, where $`๐ฅ_{<i_1,..,j_L>}=\xi _{i_1}\xi _{i_2}..\xi _{i_K}\zeta _{j_1}\zeta _{j_2}..\zeta _{j_L}`$, is the binary equivalent of $`A๐+B๐ป`$, treating both signal ($`๐บ`$ and index $`i`$) and the corrupting noise vector ($`๐`$ and index $`j`$) simultaneously. Elements of the sparse connectivity tensor $`๐_{<i_1,..,j_L>}`$ take the value 1 if the corresponding indices of both signal and noise are chosen (i.e., if all corresponding elements of the matrices $`A`$ and $`B`$ are 1) and 0 otherwise; it has $`C`$ unit elements per $`i`$-index and $`L`$ per $`j`$-index, representing the systemโs degree of connectivity. The $`\delta `$ function provides $`1`$ if the selected sitesโ product $`S_{i_1}..S_{i_K}\tau _{j_1}..\tau _{j_L}`$ is in disagreement with the corresponding element $`๐ฅ_{<i_1..j_L>}`$, recording an error, and $`0`$ otherwise. Notice that this term is not frustrated, and can therefore vanish at sufficiently low temperatures ($`T=1/\beta 0`$), imposing the restriction of Eq.(2), while the last two terms, scaled with $`\beta `$, survive. The additive fields $`F_s`$ and $`F_\tau `$ are introduced to represent our prior knowledge on the signal and noise distributions, respectively.
The random selection of elements in $`๐`$ introduces disorder to the system which is treated via methods of statistical physics. More specifically, we calculate the partition function $`๐ต(๐,๐ฑ)=\text{Tr}_{\{๐บ,๐\}}\mathrm{exp}[\beta ]`$, which is then averaged over the disorder and the statistical properties of the plaintext and noise, using the replica method, to obtain the related free energy $`=\mathrm{ln}๐ต_{\xi ,\zeta ,๐}`$. The overlap between the plaintext and the dynamical vector $`m=\frac{1}{N}_{i=1}^N\xi _iS_i`$ will serve as a measure for the decryption success.
Studying this free energy for the case of $`K=L=2`$ and in the context of error-correcting codes, indicates the existence of paramagnetic and ferromagnetic solutions depicted in the inset of Fig.1. For corruption probabilities $`p>p_s`$ one obtains either a dominant paramagnetic solution or a mixture of ferromagnetic ($`m=\pm 1`$) and paramagnetic ($`m=0`$) solutions as shown in the inset; thin and thick lines correspond to higher and lower free energies respectively, dashed lines represent unstable solutions. Lines between the $`m=\pm 1`$ and $`m=0`$ axes correspond to sub-optimal ferromagnetic solutions. Reliable decryption may only be obtained for $`p<p_s`$, which corresponds to a spinodal point, where a unique ferromagnetic solution emerges at $`m=1`$ (plus a mirror solution at $`m=1`$).
The most striking result is the division of the complete space of $`๐บ`$ and $`๐`$ values to two basins of attraction for the ferromagnetic solutions, for $`p<p_s`$, implying convergence from any initialisation of the BP equations. Critical corruption rate values for $`M/N=2`$ were obtained from the analysis and via BP solutions as shown in Fig.1, in comparison to the rate obtainable from Shannonโs channel capacity (solid line). The priors assumed for both the plaintext (unbiased in this case, $`F_s=0`$) and the corrupting vector ($`F_\tau =(1/2)\mathrm{ln}[(1p)/p]`$) correspond to Nishimoriโs condition , which is equivalent to having the correct prior within the Bayesian framework
The initial conditions for the BP-based decryption were chosen almost at random, with a very slight bias of $`๐ช(10^{12})`$ in the initial magnetisation, corresponding to typical statistical fluctuation for a system size of $`10^{24}`$. Cryptosystems with other $`K`$ and $`L`$ values, e.g., $`K,L3`$, seem to offer optimal performance in terms of the corruption rate they accommodate theoretically, but suffer from a decreasingly small basin of attraction as $`K`$ and $`L`$ increase. The co-existence of stable ferromagnetic and paramagnetic solutions implies that the system will converge to the undesired paramagnetic solution from most initial conditions which are typically of close-to-zero magnetisation. It may still be possible to use successfully specific matrices with higher $`K`$ and $`L`$ values (such as in); however, these cannot be justified theoretically and there is no clear adventage in using them.
To conclude, for the authorised user, the $`K=L=2`$ cryptosystem offers a guaranteed convergence to the plaintext solution, in the thermodynamic limit $`N\mathrm{}`$, as long as the corruption process has a probability below $`p_s`$. The main consequence of finite plaintexts would be a decrease in the allowed corruption rate with little impact on the decoding success.
The task facing the unauthorised user, i.e., finding the plaintext from Eq. (1) was investigated via similar methods by considering the Hamiltonian
$`={\displaystyle \underset{i_1,..i_K^{}}{}}๐ข_{i_1,..i_K^{}}J_{i_1,..i_K^{}}S_{i_1}..S_{i_K^{}}{\displaystyle \frac{F_s}{\beta }}{\displaystyle \underset{k=1}{\overset{N}{}}}S_k,`$
using Nishimoriโs temperature $`\beta =(1/2)\mathrm{ln}[(1p)/p]`$. The number of plaintext bits in each product is denoted $`K^{}`$, $`๐บ`$ is the $`N`$ dimensional binary vector of dynamical variables and $`๐ข`$ is a dense tensor with $`C^{}`$ unit elements per index (setting the rest of the elements to zero) and is the binary equivalent of the Boolean matrix $`G`$. The latter, together with the statistical properties of the corrupting vector $`๐ป`$ constitutes the public key used to determine the ciphertext $`๐ฑ`$. The last term on the right is required in the case of sparse or biased messages and will require assigning a certain value to the additive field $`F_s`$.
The matrix $`G`$ generated in the case of $`K=L=2`$ is dense and has a certain distribution of unit elements per row. The fraction of rows with a low (finite, not of $`๐ช(N)`$) number of unit elements vanishes as $`N`$ increases, allowing one to approximate this scenario by the diluted Random Energy Model studied in where $`K^{},C^{}\mathrm{}`$ while keeping the ratio $`C^{}/K^{}`$ finite.
To investigate the typical properties of this (frustrated) model, we calculate again the partition function and the free energy by averaging over the randomness in choosing the plaintext, the corrupting vector and the choice of the random matrix $`G`$ (being generated by a product of two sparse random matrices). To assess the likelihood of obtaining spin-glass/ferromagnetic solutions, we calculated the free energy landscape (per plaintext bit - $`f`$) as a function of overlap $`m`$. This can be carried out straightforwardly using the analysis of , and provides the energy landscape shown in Fig.2. This has the structure of a golf-course with a relatively flat area around the one-step replica symmetry breaking (frozen) spin-glass solution and a very deep but extremely narrow area, of $`๐ช(1/N)`$, around the ferromagnetic solution. To validate the use of the random energy model we also added numerical data ($`+`$, with error-bars), obtained by BP, which are consistent with the theoretical results.
This free-energy landscape may be related directly to the marginal posterior $`P(S_i=1|๐ฑ)1iN`$ and is therefore indicative of the difficulties in obtaining ferromagnetic solutions when the starting point for the search is not infinitesimally close to the original plaintext (which is clearly highly unlikely). It is plausible that any local search method, starting at some distance from the ferromagnetic solution, will fail to produce the original plaintext. Similarly, any probabilistic method, such as simulated annealing, will require an exponentially long time for converging to the $`m=1`$ solution. Numerical studies of similar energy landscapes show that the time required increases exponentially with the system size.
Most attacks on this cryptosystems, by an unauthorised user, will face the same difficulty: without explicit knowledge of the current plaintext and/or the decomposition of $`G`$ to the matrices $`A`$ and $`B`$ it will require an exponentially long time to decipher a specific ciphertext. Partial or complete knowledge of the ciphertext/plaintext as well as partial knowledge of the matrix $`B`$ (while $`๐ช(N)`$ of the elements remain unknown) will not be helpful for decomposing $`G`$ which will still require an exponentially long time to carry out.
We will consider here two attacks on specific plaintexts with partial knowledge of the corrupting vector $`๐ป`$ or of the matrix $`B`$. In the first case, knowing $`p_aM`$ of the $`pM`$ corrupting bits may allow one to subtract the approximated vector $`\widehat{๐ป}`$ from the ciphertext and take the product of $`G^1`$ and the resulting ciphertext. This attack is similar to the task facing an unauthorised user with a reduced corruption rate of $`(pp_a)`$. For any non-vanishing difference between $`p_a`$ and $`p`$, deciphering the transmitted message remains a difficult task.
A second attack is that whereby the matrix $`B`$ is known to some degree; for instance, the location of a fraction of the unit elements, say $`1\rho `$ is known. From Eq.(2) one can identify the absent information as having a higher effective corruption level of $`p+g(\rho )`$, where $`g()`$ is some non-trivial function that depends on the actual scenario. To secure the transmission one may work sufficiently close to the critical corruption level $`p_s`$ such that the additional effective noise $`\rho `$ will bring the system beyond the critical corruption rate and into the paramagnetic/spin-glass regime. However, the drawbacks of working very close to $`p_s`$ are twofold: Firstly, average decryption times using BP methods ($`\tau `$) will diverge proportionally to $`1/(p_sp)`$ as demonstrated in the inset of Fig.2. Secondly, finite-size effects are expected to be larger close to $`p_s`$, which practically means that the system may not converge to the attractive optimal solutions in some cases.
We will end this presentation with a short discussion on the advantages and drawbacks of the suggested method in comparison with existing techniques. Firstly, we would like to point out the differences between this method and the McEliece cryptosystem. The latter is based on Goppa codes and is limited to relative low corruption levels. These may allow for decrypting the ciphertext using (many) estimates of the corruption vector. Our suggestion allows for a significant corruption level, thus increasing the security of the cryptosystem. In addition, the suggested construction, $`K=L=2`$, is not discussed in the information theory literature (e.g. in ) which typically prefers higher $`K,L`$ value systems. Secondly, in comparison to RSA where decryption takes $`๐ช(N^3)`$ operations, our method only requires $`๐ช(N)`$ operations, multiplied by the number of BP iterations (which is typically smaller than 100 for most system sizes examined except very close $`p_s`$). Encryption costs are of $`๐ช(N^2)`$ (as in RSA) while the inversion of the matrix $`B`$ is carried out only once and requires $`O(N^3)`$ operations.
The two obvious drawbacks of our method are: 1) The transmission of the public key, which is a dense matrix of dimensionality $`M\times N`$. However, as public key transmission is carried out only once for each user we do not expect it to be of great significance. 2) The ciphertext to plaintext bit ratio is greater than one to allow for corruption, in contrast to RSA where it equals 1. Choosing the $`N/M`$ ratio is in the hands of the user and is directly related to the security level required; we therefore do not expect it to be problematic as the increased transmission time is compensated by a very fast decryption.
We examine the typical performance of a new cryptosystem, based on insight gained from our previous studies, by mapping it onto an Ising spin system; this complements the information theory approach which focuses on rigorous worst-case bounds. We show that authorised decryption is fast and simple while unauthorised decryption requires a prohibitively long time. Important aspects that are yet to be investigated include finite size effects and methods for alleviating the drawbacks of the new method.
Acknowledgement Support by JSPS-RFTF (YK), The Royal Society and EPSRC-GR/L52093 (DS) is acknowledged. We would like to thank Manfred Opper and Hidetoshi Nishimori for reading the manuscript. |
warning/0002/physics0002027.html | ar5iv | text | # Femtosecond soliton amplification in nonlinear dispersive traps and soliton dispersion management
## I INTRODUCTION
In 1973 Hasegawa and Tappert showed theoretically that an optical pulse in a dielectric fibers forms an envelope solitons, and in 1980 Mollenauer, Stolen and Gordon demonstrated the effect experimentally. This discovery is significant in its application to optical communications. Today the optical soliton is regarded as an important alternative for the next generation of high speed telecommunication systems.
The theory of NSE solitons was developed for the first time in 1971 by Zakharov and Shabad . The concept of the soliton involves a large number of interesting problems in applied mathematics since it is an exact analytical solution of a nonlinear partial differential equations. The theory of optical solitons described by the nonlinear Schrรถdinger equation has produced perfect agreement between theory and experiment .
In this paper we present mathematical description of solitary waves propagation in a nonlinear dispersive medium with varying parameters.
The soliton spectral tunneling effect was theoretically predicted in . This is characterized in the spectral domain by the passage of a femtosecond soliton through a potential barrier-like spectral inhomogeneity of the group velocity dispersion (GVD), including the forbidden band of a positive GVD. It is interesting to draw an analogy with quantum mechanics where the solitons are considered to exhibit particle-like behavior. The soliton spectral tunneling effect also can be considered as an example of the dynamic dispersion soliton management technique. In the first part of the paper we will concentrate on the problem of femtosecond solitons amplification. We will show that spectral inhomogeneity of GVD allows one to capture a soliton in a sort of spectral trap and to accumulate an additional energy during the process of the soliton amplification. In the second part we will consider the problem of the short soliton pulse propagation in the nonlinear fiber with static non-uniform inhomogeneity of GVD. The methodology developed does provide a systematic way to generate infinite โoceanโ of the chirped soliton solutions of NSE model with varying coefficients.
## II FEMTOSECOND SOLITON AMPLIFICATION
It is well known that due to the Raman self-scattering effect (called soliton self-frequency shift ) the central femtosecond soliton frequency shifts to the red spectral region and so-called colored solitons are generated. This effect decreases significantly the efficiency of resonant amplification of femtosecond solitons. The mathematical model we consider based on the modified NSE including the effects of molecular vibrations and soliton amplification processes (see details in ):
$$i\frac{\psi }{z}=\frac{1}{2}\frac{^2\psi }{\tau ^2}+i\sigma \frac{^3\mathrm{\Psi }}{\tau ^3}+(1\beta )|\psi |^2\psi +\beta Q\psi +\frac{G}{2}P$$
(1)
$$\mu ^2\frac{^2Q}{t^2}+2\mu \delta \frac{Q}{t}+Q=|\psi |^2,\text{and},\gamma _a\frac{P}{\tau }+P(1+i\gamma _a\mathrm{\Delta }\mathrm{\Omega })=i\psi ,$$
(2)
As numerical experiments showed the GVD inhomogeneity as a potential well allows one to capture a soliton in a sort of spectral trap. Figure 1 shows the nonlinear dynamics of the soliton spectral trapping effect in the spectral domain. As soliton approaches the well, it does not slow down but speeds up, and then, after it has got into the well, the soliton is trapped. There exists a long time of soliton trapping in internal region of the well .This effect opens a controlled possibility to increase the energy of a soliton. As follows from our computer simulations the capture of moving in the frequency space femtosecond colored soliton by a dispersive trap formed in an amplifying optical fiber makes it possible to accumulate an additional energy in the soliton dispersive trap and to reduce significantly the soliton pulse duration.
## III DISPERSION MANAGEMENT: CHIRPED SOLITONS
Let us consider the propagation of a nonlinear pulse in the anomalous (or normal) group velocity dispersion fiber of length $`Z_1.`$ The complex amplitude $`q`$ of the light wave in a fiber with variable parameters $`D_2(Z)`$ , $`N_2(Z)`$ and $`\mathrm{\Gamma }(Z)`$ is described by the nonlinear Schrodinger equation
$$i\frac{q}{Z}+\frac{1}{2}D_2(Z)\frac{^2q}{T^2}+N_2(Z)q^2q=i\mathrm{\Gamma }(Z)q$$
(3)
Theorem 1. Consider the NSE (3) with varying dispersion, nonlinearity and gain. Suppose that Wronskian W\[N<sub>2</sub>,D<sub>2</sub>\] of the functions N<sub>2</sub>(Z) and D<sub>2</sub>(Z) is nonvanishing, thus two functions N<sub>2</sub>(Z) and D<sub>2</sub>(Z) are linearly independent. There are then infinite number of solutions of Eq. (3) in the form of Eq.4
$$q(Z,T)=\sqrt{\frac{D_2(Z)}{N_2(Z)}}\text{ }P(Z)\text{ }Q\left[P(Z)T\right]\text{ }\mathrm{exp}\left[i\frac{P(Z)}{2}\text{ }T^2+i\underset{0}{\overset{Z}{}}K(Z^{^{}})๐Z^{^{}}\right]$$
(4)
where function $`Q`$ describes fundamental functional form of bright $`Q=\text{sech}(P(Z)T)`$ or dark $`Q=\text{th}(P(Z)T)`$ NSE solitons and the real functions P(Z), D$`{}_{2}{}^{}(Z)`$, N$`{}_{2}{}^{}(Z)`$ and $`\mathrm{\Gamma }(Z)`$ are determined by the following nonlinear system of equations :
$$\frac{1}{P^2(Z)}\frac{P(Z)}{Z}+D_2(Z)=0\text{ ; }\frac{1}{2}D_2(Z)P(Z)+\frac{W[N_2(Z),D_2(Z)]}{2D_2(Z)N_2(Z)}=\mathrm{\Gamma }(Z)$$
(5)
Theorem 2. Consider the NSE (3) with varying dispersion, nonlinearity and gain. Suppose that Wronskian W\[N<sub>2</sub>,D<sub>2</sub>\] of the functions N<sub>2</sub>(Z) and D<sub>2</sub>(Z) is vanishing, thus two functions N<sub>2</sub>(Z) and D<sub>2</sub>(Z) are linearly dependent. There are then infinite number of solutions of Eq. (3) of the following form Eq. 6
$$q(Z,T)=C\text{ }P(Z)\text{ }Q\left[P(Z)T\right]\text{ }\mathrm{exp}\left[i\frac{P(Z)}{2}\text{ }T^2+i\underset{0}{\overset{Z}{}}K(Z^{^{}})๐Z^{^{}}\right]$$
(6)
where function $`Q`$ describes the fundamental form of bright (or dark) NSE soliton and the real functions P(Z), D$`{}_{2}{}^{}(Z)`$, N$`{}_{2}{}^{}(Z)`$ and $`\mathrm{\Gamma }(Z)`$ are determined by the following nonlinear system of equations :
$$D_2(Z)=\frac{1}{P^2(Z)}\frac{P(Z)}{Z}\text{ ; }\mathrm{\Gamma }(Z)=\frac{1}{2}\frac{1}{P}\frac{P(Z)}{Z}\text{ ; }N_2(Z)=D_2(Z)/C^2$$
(7)
The function P(Z) is required only to be once-differentiable, but otherwise arbitrary function, there is no restrictions.
To prove Theorems 1 and 2 we first construct a stationary localized solution of Eq. (3) by introducing Kumar- Hasegawaโs quasi-soliton concept through
$$q(Z,T)=\sqrt{\frac{D_2(Z)}{N_2(Z)}}\text{ }P(Z)\text{ }Q\left[P(Z)T\right]\text{ }\mathrm{exp}\left[i\frac{P(Z)}{2}\text{ }T^2+i\underset{0}{\overset{Z}{}}K(Z^{^{}})๐Z^{^{}}\right]$$
(8)
where $`D_2`$(Z), $`N_2`$(Z), $`P(Z)`$ and $`K(Z)`$ are the real functions of $`Z.`$ Substituting expression (8) into Eq. (3) and separating real and imaginary parts we obtain the system of two equations
$$\frac{1}{2}\text{sign}(D_2)\frac{^2Q}{S^2}+Q^3+\left(E\frac{S^2}{2}\mathrm{\Omega }^2(Z)\right)\text{ }Q=0$$
(9)
$$\frac{P}{Z}Q+P\frac{Q}{S}\frac{S}{Z}+\frac{1}{2}\frac{1}{D_2(Z)}\frac{D_2}{Z}PQ\frac{1}{2}\frac{1}{N_2(Z)}\frac{N_2}{Z}PQ+\frac{1}{2}D_2P^2Q+D_2P^2T\frac{Q}{S}\frac{S}{T}=\mathrm{\Gamma }PQ$$
(10)
Where
$$S(Z,T)=P(Z)T\text{ ; }\frac{S}{Z}=T\frac{P}{Z}\text{ ; }\frac{S}{T}=P(Z)$$
(11)
In Eq. (9) the parameters $`E`$ and $`\mathrm{\Omega }`$ are โthe energyโ and โfrequencyโ of ordinary quantum mechanical harmonic ocsillator
$$\mathrm{\Omega }^2(Z)=\frac{D_2^1(Z)}{P^2(Z)}\left(\frac{1}{P^2(Z)}\frac{P}{Z}+D_2(Z)\right)\text{ ; }E(Z)=K(Z)/P^2(Z)/D_2(Z)$$
(12)
Eq. (9) represents the nonlinear Schrodinger equation for the harmonic ocsillator. As must be in the case of Hamiltonian system Eq. (9) may be written in the form
$$\frac{\delta H}{\delta Q^{}}=0$$
(13)
$$H=\text{ }\left[\frac{1}{2}\text{sign}(D_2)\left|\frac{Q}{X}\right|^2+\frac{1}{2}\alpha \left|Q\right|^4+\left(E\frac{X^2}{2}\mathrm{\Omega }^2(Z)\right)\text{ }\left|Q\right|^2\right]๐X$$
(14)
The derivative in (13) is functional derivative. For the first time this equation was solved numerically by Kumar and Hasegawa in and gave rise a new concept of quasi-solitons . Now we make the important assumption about the solution of Eq. (9).
Let us consider the complete nonlinear regime when Eq. (9) represents the ideal NLS equation, i.e. we will allow $`\mathrm{\Omega }(Z)0`$ , then from (12) follows that
$$\frac{1}{P^2(Z)}\frac{P(Z)}{Z}+D_2(Z)=0$$
(15)
We now look for a solution of Eq. (10) which satisfies the condition (15). Substituting the expression (15) and relations (11) into Eq. (10) we obtain
$$\frac{1}{2}D_2(Z)P(Z)+\frac{1}{2}\frac{1}{D_2(Z)}\frac{D_2(Z)}{Z}\frac{1}{2}\frac{1}{N_2(Z)}\frac{N_2(Z)}{Z}=\mathrm{\Gamma }(Z)$$
(16)
Using notation
$`W\{N_2,D_2\}=N_2{\displaystyle \frac{D_2(Z)}{Z}}D_2{\displaystyle \frac{N_2(Z)}{Z}}`$
one can obtain the soliton solution of Eq. 3 in the form of the chirped solitons Eqs. 4-5 and Eqs. 6-7.. Consequently, we have found the infinite โoceanโ of solutions. The methodology developed does provide a systematic way of new and new chirped soliton solutions generation.
## IV DIFFERENT REGIMES OF SOLITON MANAGEMENT
Lemma 1: Soliton GVD management. Consider the NSE (3) with constant nonlinear coefficient N$`{}_{2}{}^{}=const`$ and with varying along Z-coordinate GVD parameter. Suppose that dispersion management function is known arbitrary analytical function :D<sub>2</sub>(Z)=$`\mathrm{\Phi }(Z)`$ . The function $`\mathrm{\Phi }(Z)`$ is required only to be once-differentiable and once integrable, but otherwise arbitrary function, there is no restrictions. There are then infinite number of solutions of Eq. (3) of the form of the chirped dispersion managed dark and bright solitons Eq. 4, where the main functions P and $`\mathrm{\Gamma }`$ are given by
$$D_2(Z)=\mathrm{\Phi }(Z)\text{ };\text{ }P(Z)=\frac{1}{\left[C\mathrm{\Phi }(Z)๐Z\right]}$$
(17)
$$\mathrm{\Gamma }(Z)=\frac{1}{2}\frac{\mathrm{\Phi }(Z)}{\left[C\mathrm{\Phi }(Z)๐Z\right]}+\frac{1}{2}\frac{1}{\mathrm{\Phi }(Z)}\frac{\mathrm{\Phi }(Z)}{Z}$$
(18)
Lemma 2: Soliton intensity management. Consider the NSE (3) with constant nonlinear coefficient N$`{}_{2}{}^{}=const`$ and with varying along Z-coordinate the dispersion and gain. Suppose that intensity of the soliton pulse is determined by the known management function: D<sub>2</sub>(Z)P<sup>2</sup>(Z)=$`\mathrm{\Theta }(Z),`$where the function $`\mathrm{\Theta }(Z)`$ is required only to be once-differentiable and once integrable, but otherwise arbitrary function, there is no restrictions. There are then infinite number of solutions of Eq. (3) of the form of the chirped dispersion managed dark and bright solitons Eq. 4 with parameters given by
$`D_2(Z)P^2(Z)=\mathrm{\Theta }(Z)\text{ ; }P(Z)={\displaystyle \mathrm{\Theta }(Z)๐Z}+C\text{ ; }D_2(Z)={\displaystyle \frac{\mathrm{\Theta }(Z)}{\left[C\mathrm{\Theta }(Z)๐Z\right]^2}}`$
$$\mathrm{\Gamma }(Z)=\frac{1}{2}\frac{\mathrm{\Theta }(Z)}{\left[C\mathrm{\Theta }(Z)๐Z\right]}+\frac{1}{2}\frac{1}{\mathrm{\Theta }(Z)}\frac{\mathrm{\Theta }(Z)}{Z}$$
(19)
Lemma 3: Soliton pulse duration management: optimal soliton compression. Consider the NSE (3) with constant nonlinear coefficient N$`{}_{2}{}^{}=const`$ and with varying along Z-coordinate the dispersion and gain coefficients. Suppose that pulse duration of a soliton is determined by the known analytical function: P(Z)=$`\mathrm{{\rm Y}}(Z)`$, where the function $`\mathrm{{\rm Y}}(Z)`$ is required only to be two-differentiable , but otherwise arbitrary function, there is no restrictions. There are then infinite number of solutions of Eq. (3) of the form of the chirped dispersion managed dark and bright solitons Eq. 4 with the main parameters given by
$$D_2(Z)=\frac{1}{\mathrm{{\rm Y}}^2(Z)}\frac{\mathrm{{\rm Y}}(Z)}{Z}\text{ ; }\mathrm{\Gamma }(Z)=\frac{1}{2}\left(\frac{\mathrm{{\rm Y}}(Z)}{Z}\right)^1\frac{}{Z}\left(\frac{1}{\mathrm{{\rm Y}}(Z)}\frac{\mathrm{{\rm Y}}(Z)}{Z}\right)$$
(20)
Lemma 4: Soliton amplification management: optimal soliton compression. Consider the NSE (3) with constant nonlinear coefficient N$`{}_{2}{}^{}=const`$ and with varying along Z-coordinate the dispersion and gain coefficients. Suppose that the gain coefficient is determined by the known control function: $`\mathrm{\Gamma }`$(Z)=$`\mathrm{\Lambda }(Z),`$ where the function $`\mathrm{\Lambda }(Z)`$ is required only to be once integrable , but otherwise arbitrary function, there is no restrictions. There are then infinite number of solutions of Eq. (3) of the form of the chirped dispersion managed dark and bright solitons of the Eq. 4 where
$$\left|P(Z)\right|=\mathrm{exp}\left[\mathrm{exp}\left(2\mathrm{\Lambda }(Z^{^{\prime \prime }})๐Z^{^{\prime \prime }}\right)๐Z^{^{}}\right]$$
(21)
$$\left|D_2(Z)\right|=\frac{\mathrm{exp}\left(2\mathrm{\Lambda }(Z)๐Z\right)}{\mathrm{exp}\left[\mathrm{exp}\left(2\mathrm{\Lambda }(Z^{^{\prime \prime }})๐Z^{^{\prime \prime }}\right)๐Z^{^{}}\right]}$$
(22)
Lemma 5: Combined dispersion and nonlinear soliton management. Consider the NSE (3) with varying nonlinear coefficient N$`{}_{2}{}^{}(Z)`$ and with varying along Z-coordinate the dispersion and gain coefficients too. Suppose that Wronskian W\[N<sub>2</sub>,D<sub>2</sub>\] is vanishing, or that the functions N<sub>2</sub>(Z) and D<sub>2</sub>(Z) are linearly dependent. Suppose also that the function D<sub>2</sub>(Z) is determined by the initial control function D<sub>2</sub>(Z)=$`\mathrm{\Xi }(Z),`$where the function $`\mathrm{\Xi }(Z)`$ is required only to be once integrable, but otherwise arbitrary function, there is no restrictions. There are then infinite number of solutions of Eq. (3) of the form of the chirped dispersion managed dark and bright solitons of the Eq. 6 where
$$P(Z)=\frac{1}{\left[C\mathrm{\Xi }(Z)๐Z\right]}\text{ ; }N_2(Z)=D_2(Z)/C^2$$
(23)
$$\mathrm{\Gamma }(Z)=\frac{1}{2}\frac{\mathrm{\Xi }(Z)}{\left[C\mathrm{\Xi }(Z)๐Z\right]}$$
(24)
The analytical solutions for the different regimes of the main soliton parameters management (intensity, pulse duration, amplification or absorption ) in the case of W\[N<sub>2</sub>,D<sub>2</sub>\]=0 can be obtained by using theorem 2.
Let us consider some examples. The case of $`\mathrm{\Gamma }(Z)0`$ and N<sub>2</sub>(Z)=N$`{}_{2}{}^{}(0)`$ corresponds to the problem of ideal GVD soliton management. The soliton solution in this case is:
$`q(Z,T)`$ $`=`$ $`\eta N_2^{1/2}(0)\mathrm{exp}(\text{ }{\displaystyle \frac{C}{2}}Z\text{ })\text{ sech }\left[\eta T\mathrm{exp}(CZ)\right]`$ (26)
$`\mathrm{exp}\left[iT^2{\displaystyle \frac{C}{2}}\mathrm{exp}(CZ)i{\displaystyle \frac{1}{2}}\eta ^2Z\mathrm{exp}(CZ)\right]`$
$`q(Z,T)`$ $`=`$ $`\eta N_2^{1/2}(0)\mathrm{exp}(\text{ }{\displaystyle \frac{C}{2}}Z\text{ })\text{ th }[\eta T\mathrm{exp}(CZ]`$ (28)
$`\mathrm{exp}\left[iT^2{\displaystyle \frac{C}{2}}\mathrm{exp}(CZ)i\eta ^2Z\mathrm{exp}(CZ)\right]`$
Here $`T`$ and $`Z`$ are ordinary variables and $`C`$ is arbitrary constant. If we use the expressions D<sub>2</sub>(Z)=constant and N<sub>2</sub>=N$`{}_{2}{}^{}(0)`$ then we obtain the following solutions of Eq. (3) in the form of hyperbolically growing ideal bright and dark solitons (for the first time reported in
$$q(Z,T)=\frac{\chi \text{ }N_2^{1/2}(0)}{(12\mathrm{\Gamma }(0)Z)}\text{ sech}\left[\frac{\chi T}{(12\mathrm{\Gamma }(0)Z)}\right]\mathrm{exp}\left[i\frac{T^2\mathrm{\Gamma }(0)}{(12\mathrm{\Gamma }(0)Z)}i\frac{\chi ^2Z}{2(12\mathrm{\Gamma }(0)Z)}\right]$$
(29)
$$q(Z,T)=\frac{\chi \text{ }N_2^{1/2}(0)}{(12\mathrm{\Gamma }(0)Z)}\text{ th}\left[\frac{\chi T}{(12\mathrm{\Gamma }(0)Z)}\right]\mathrm{exp}\left[i\frac{T^2\mathrm{\Gamma }(0)}{(12\mathrm{\Gamma }(0)Z)}i\frac{\chi ^2Z}{(12\mathrm{\Gamma }(0)Z)}\right]$$
(30)
In the case of $`\mathrm{\Gamma }(Z)G_0`$ and N<sub>2</sub>=$`N_2(0)`$ the solution of Eq. 3 is given by:
$$\text{ }Q(P(Z)T)=\eta N_2^{1/2}(0)\text{ sech }\left[\eta P(Z)T\right]$$
(31)
$$\text{ }Q(P(Z)T)=\eta N_2^{1/2}(0)\text{ th }\left[\eta P(Z)T\right]$$
(32)
$$P(Z)=P(0)\mathrm{exp}(\frac{1}{2G_0}(\mathrm{exp}(2G_0Z)1))$$
(33)
$$D_2(Z)=D_2(0)\mathrm{exp}(2G_0Z\frac{1}{2G_0}(\mathrm{exp}(2G_0Z)1))$$
(34)
When GVD is a hyperbolically decreasing function of Z
$$D_2(Z)=\frac{1}{1+\beta Z}$$
(35)
then from Lemma 1 follows the explicit soliton solution in the form of Eq. 4
$$P(Z)=\frac{1}{1\frac{1}{\beta }\mathrm{ln}(1+\beta Z)}$$
(36)
$$\mathrm{\Gamma }(Z)=\frac{1}{2(1+\beta Z)}\left[\frac{1\mathrm{ln}(1+\beta Z)}{1\frac{1}{\beta }\mathrm{ln}(1+\beta Z)}\right]$$
(37)
Let us consider the soliton intensity management problem. Chirped soliton pulse of Eq. 3 with the constant intensity can be obtained by using Lemma 2
$$P(Z)=CZ1;\text{ }D_2(Z)=C/(1+CZ)^2;\text{ }\mathrm{\Gamma }(Z)=C/2/(1+CZ)$$
(38)
Let us consider some periodical chirped soliton solutions of Eq. 3. Suppose that the soliton intensity varies periodically as
$$D_2(Z)P^2(Z)=\mathrm{\Theta }(Z)=1+\delta \mathrm{sin}^{2n}Z$$
(39)
Then soliton solution in the case of n=2 is determined by Eq. 4 with parameters:
$$D_2(Z)=\mathrm{\Theta }(Z)/P^2(Z);\text{ }P(Z)=C\left[Z+\delta \left(\frac{3Z}{8}\frac{\mathrm{sin}2Z}{4}+\frac{\mathrm{sin}4Z}{32}\right)\right]$$
(40)
$$\mathrm{\Gamma }(Z)=\frac{1}{2}\frac{\left(1+\delta \mathrm{sin}^4Z\right)}{C\left[Z+\delta \left(\frac{3Z}{8}\frac{\mathrm{sin}2Z}{4}+\frac{\mathrm{sin}4Z}{32}\right)\right]}+\frac{1}{2}\frac{2\text{ }\mathrm{sin}2Z\text{ }\mathrm{sin}^2Z}{\left(1+\delta \mathrm{sin}^4Z\right)}$$
(41)
Let us consider some periodical solutions of Eq. 3 in the case of the linearly dependent parameters of the media. The simplest solution of Eq. 3 in the form of Eq. 6 is:
$`P(Z)`$ $`=`$ $`\mathrm{{\rm Y}}(Z)=\left(1+\delta \mathrm{sin}^2Z\right);\text{ }N_2(Z)=D_2(Z)={\displaystyle \frac{\delta \mathrm{sin}2Z}{\left(1+\delta \mathrm{sin}^2Z\right)^2}};`$ (42)
$`\mathrm{\Gamma }(Z)`$ $`=`$ $`{\displaystyle \frac{\delta }{2}}{\displaystyle \frac{\mathrm{sin}2Z}{\left(1+\delta \mathrm{sin}^2Z\right)}}`$ (43)
The next periodical soliton solution is given by
$$D_2(Z)=N_2(Z)=\mathrm{cos}Z;\text{ }P(Z)=\frac{1}{\left(C\mathrm{sin}Z\right)};\text{ }\mathrm{\Gamma }(Z)=\frac{\mathrm{cos}Z}{2\left(C\mathrm{sin}Z\right)}$$
(44)
The main soliton features of the solutions given by theorem 1 and theorem 2 were investigated by using direct computer simulations. We have investigated the interaction dynamics of particle-like solutions obtained, their soliton-like character was calculated with the accuracy as high as 10<sup>-9</sup>. We also have investigated the influence of high-order effects on the dynamics of dispersion and amplification management. As follows from numerical investigations elastic character of chirped solitons interacting does not depend on a number of interacting solitons and their phases. Figure 2 shows the computer simulation dynamics of three hyperbolically growing solitons Eq. 29. NSE solution with periodic dispersion coefficient is shown in Figure 3. Here the dispersion management function is
$$D_2(Z)=1+\delta \mathrm{sin}^2(Z)$$
(45)
and the soliton solution is given by Eqs. 17-18. In Figure 3 parameters C=200 and $`\delta =0.9.`$ Figure 4 represents the two dispersion managed solitons interaction in the case of equal phases and in Figure 5 the interaction dynamics of two solitons is shown in the case of opposite phases. Figure 6 shows the intensity managed solitons dynamics of the form presented by Eq. 38. Figures 7-9 show the nonlinear propagation and interaction of the dispersion and nonlinear managed solitons of Eqs. 42-43. The main parameters in computer simulations were C=200; $`\delta =\pm 0.9`$. Figure 10 illustrates the dynamics of the fission of the bound states of two hyperbolically growing solitons Eqs. 29-30 produced by self-induced Raman scattering effect given by Eqs.2-3. This remarkable fact also emphasize the full soliton features of solutions discussed. They not only interact elastically but they can form the bound states and these bound states split under perturbations. The possibility to find the plethora of soliton solutions in the case of strong dispersion management is reported in the recent paper of Zakharov and Manakov . |
warning/0002/hep-ex0002018.html | ar5iv | text | # SOLAR NEUTRINOS:WHAT NEXT?
## 1 Introduction
The reader who is familiar with solar neutrino research may wish to skip directly to the last section entitled: What Next?
Solar neutrinos have been detected experimentally with fluxes and energies that are qualitatively consistent with solar models that are constructed assuming that the sun shines by nuclear fusion reactions. The first experimental result, obtained by Ray Davis and his collaborators in 1968, has now been confirmed by four other beautiful experiments, Kamiokande, SAGE, GALLEX, and SuperKamiokande. The observation of solar neutrinos with approximately the predicted energies and fluxes establishes empirically the theory that main sequence stars derive their energy from nuclear fusion reactions in their interiors and has inaugurated what we all hope will be a flourishing field of observational neutrino astronomy.
Although the calculated neutrino fluxes depend upon high powers of the central temperature of the solar model, the experiments and the solar model theory are so precise that persistent quantitative discrepancies have existed between the model predictions and the solar model calculations for over thirty years
Important experiments are underway that will provide diagnostic information about the physical properties of neutrinos that are created in the center of the sun and detected on earth in really long baseline experiments. At this workshop, we will hear discussions of the SuperKamiokande, SNO, BOREXINO, HELLAZ, HERON, ICARUS, LENS, and KamLAND experiments.
I will discuss predictions of the combined standard model in the main part of this review. By โcombinedโ standard model, I mean the predictions of the standard solar model and the predictions of the standard electroweak model. We need a solar model to tell us how many neutrinos of what energy are produced in the sun and we need an electroweak theory to tell us how the number and flavor content of the neutrinos are changed as they make their way from the center of the sun to detectors on earth. For all practical purposes, the standard electroweak model states that nothing happens to solar neutrinos after they are created in the deep interior of the sun. Using standard electroweak theory and fluxes from the standard solar model, one can calculate the rates of neutrino interactions in different terrestrial detectors with a variety of energy sensitivities. The combined standard model also predicts that the energy spectrum from a given neutrino source should be the same for neutrinos produced in terrestrial laboratories and in the sun and that there should not be measurable time-dependences (other than the seasonal dependence caused by the earthโs orbit around the sun). The spectral and temporal departures from standard model expectations are expected to be small in all currently operating experiments and have not yet yielded definitive results. Therefore, I will concentrate here on inferences that can be drawn by comparing the total rates observed in solar neutrino experiments with the combined standard model predictions.
I will begin by reviewing in Section 2 the quantitative predictions of the combined standard solar model and then describe in Section 3 the three solar neutrino problems that are established by the chlorine, Kamiokande, SAGE, GALLEX, and SuperKamiokande experiments. In Section 4, I detail the uncertainties in the standard model predictions and then show in Section 5 that helioseismological measurements indicate that the standard solar model predictions are accurate for our purposes. In Section 5, I discuss the implications for solar neutrino research of the precise agreement between helioseismological measurements and the predictions of standard solar models. Next, ignoring all knowledge of the sun, I cite analyses in Section 6 that show that one cannot fit the existing experimental data with neutrino fluxes that are arbitrary parameters, unless one invokes new physics to change the shape or flavor content of the neutrino energy spectrum. I summarize in Section 7 the characteristics of the best-fitting neutrino oscillation descriptions of the experimental data. Finally, I will discuss and summarize the results in Section 8.
If you want to obtain numerical data or subroutines that are discussed in this talk, or to see relevant background information, you can copy them from my Web site: http://www.sns.ias.edu/$``$jnb.
Before we begin the detailed discussion, I want to make just a brief historical diversion. Nearly all of the current interest in solar neutrinos centers around the opportunity to use the sun as a neutrino source in a very long baseline oscillation experiment. In preparing for a talk in honor of Fred Reines a few months ago, I ran across a long forgotten 1972 letter from Bruno Pontecorvo, the originator of the hypothesis that oscillations may be observed in solar neutrino experiments. For your interest, I enclose a reproduction of this letter in Fig. 1.
For the benefit of neutrino pioneers of today, it is perhaps worth remarking that Ray Davis and I never considered the possibility that solar neutrinos could be used to learn more about neutrinos when, in the early 1960โs, we were first analyzing the potentialities of a practical chlorine experiment. We sold the experiment as a fundamental test of the hypothesis that the sun shines by nuclear fusion reactions in its interior. Only after the first results of the chlorine experiment showed in 1968 a conflict with the solar model calculations and Gribov and Pontecorvo published their epochal 1969 paper on vacuum oscillations of solar neutrinos did we begin to consider the possibility that solar neutrinos might tell us something new about particle physics. Maybe there are previously unimagined physics treasures to be discovered in future neutrino experiments.
## 2 Standard Model Predictions
Table 1 gives the neutrino fluxes and their uncertainties for our best standard solar model, hereafter BP98. Figure 2 shows the predicted neutrino fluxes from the dominant $`p`$-$`p`$ fusion chain.
The BP98 solar model includes diffusion of heavy elements and helium, makes use of the nuclear reaction rates recommended by the expert workshop held at the Institute of Nuclear Theory, recent (1996) Livermore OPAL radiative opacities, the OPAL equation of state, and electron and ion screening as determined by the recent density matrix calculation. The neutrino absorption cross sections that are used in constructing Table 1 are the most accurate values available and include, where appropriate, the thermal energy of fusing solar ions and improved nuclear and atomic data. The validity of the absorption cross sections has recently been confirmed experimentally using intense radioactive sources of $`{}_{}{}^{51}\mathrm{Cr}`$. The ratio, $`R`$, of the capture rate measured (in GALLEX and SAGE) to the calculated $`{}_{}{}^{51}\mathrm{Cr}`$ capture rate is $`R=0.95\pm 0.07(\mathrm{exp})+{}_{0.03}{}^{+0.04}(\mathrm{theory})`$ and was discussed extensively at Neutrino 98 by Gavrin and by Kirsten. The neutrino-electron scattering cross sections, used in interpreting the Kamiokande and SuperKamiokande experiments, now include electroweak radiative corrections.
Figure 3 shows for the chlorine experiment all the predicted rates and the estimated uncertainties ($`1\sigma `$) published by my colleagues and myself since the first measurement by Ray Davis and his colleagues in 1968. This figure should give you some feeling for the robustness of the solar model calculations. Many hundreds and probably thousands of researchers have, over three decades, made great improvements in the input data for the solar models, including nuclear cross sections, neutrino cross sections, measured element abundances on the surface of the sun, the solar luminosity, the stellar radiative opacity, and the stellar equation of state. Nevertheless, the most accurate predictions of today are essentially the same as they were in 1968 (although now they can be made with much greater confidence). For the gallium experiments, the neutrino fluxes predicted by standard solar models, corrected for diffusion, have been in the range $`120`$ SNU to $`141`$ SNU since 1968. A SNU is a convenient unit with which to describe the measured rates of solar neutrino experiments: $`10^{36}`$ interactions per target atom per second.
There are three reasons that the theoretical calculations of neutrino fluxes are robust: 1) the availability of precision measurements and precision calculations of input data; 2) the connection between neutrino fluxes and the measured solar luminosity; and 3) the measurement of the helioseismological frequencies of the solar pressure-mode ($`p`$-mode) eigenfrequencies. I have discussed these reasons in detail in another talk.
Figure 4 displays the calculated $`{}_{}{}^{7}\mathrm{Be}`$ and $`{}_{}{}^{8}\mathrm{B}`$ neutrino fluxes for all $`19`$ standard solar models which have been published in the last $`10`$ years in refereed science journals. The fluxes are normalized by dividing each published value by the flux from the BP98 solar model; the abscissa is the normalized $`{}_{}{}^{8}\mathrm{B}`$ flux and the ordinate is the normalized $`{}_{}{}^{7}\mathrm{Be}`$ neutrino flux. The rectangular box shows the estimated $`3\sigma `$ uncertainties in the predictions of the BP98 solar model.
All of the solar model results from different groups fall within the estimated 3$`\sigma `$ uncertainties in the BP98 analysis (with the exception of the Dar-Shaviv model whose results have not been reproduced by other groups). This agreement demonstrates the robustness of the predictions since the calculations use different computer codes (which achieve varying degrees of precision) and involve a variety of choices for the nuclear parameters, the equation of state, the stellar radiative opacity, the initial heavy element abundances, and the physical processes that are included.
The largest contributions to the dispersion in values in Figure 4 are due to the choice of the normalization for $`S_{17}`$ (the production cross-section factor for $`{}_{}{}^{8}\mathrm{B}`$ neutrinos) and the inclusion, or non-inclusion, of element diffusion in the stellar evolution codes. The effect in the plane of Fig. 4 of the normalization of $`S_{17}`$ is shown by the difference between the point for BP98 (1.0,1.0), which was computed using the most recent recommended normalization, and the point at (1.18,1.0) which corresponds to the BP98 result with the earlier (CalTech) normalization.
Helioseismological-observations have shown that element diffusion is occurring and must be included in solar models, so that the most recent models shown in Fig. 4 now all include helium and heavy element diffusion. By comparing a large number of earlier models, it was shown that all published standard solar models give the same results for solar neutrino fluxes to an accuracy of better than 10% if the same input parameters and physical processes are included.
Bahcall, Krastev, and Smirnov have compared the observed rates with the calculated, standard model values, combining quadratically the theoretical solar model and experimental uncertainties, as well as the uncertainties in the neutrino cross sections. Since the GALLEX and SAGE experiments measure the same quantity, we treat the weighted average rate in gallium as one experimental number. We adopt the SuperKamiokande measurement as the most precise direct determination of the higher-energy $`{}_{}{}^{8}\mathrm{B}`$ neutrino flux.
Using the predicted fluxes from the BP98 model, the $`\chi ^2`$ for the fit to the three experimental rates (chlorine, gallium, and SuperKamiokande, see Fig. 5) is
$$\chi _{\mathrm{SSM}}^2\text{(3 experimental rates)}=61.$$
(1)
The result given in Eq. (1), which is approximately equivalent to a $`20\sigma `$ discrepancy, is a quantitative expression of the fact that the standard model predictions do not fit the observed solar neutrino measurements.
## 3 Three Solar Neutrino Problems
I will now compare the predictions of the combined standard model with the results of the operating solar neutrino experiments.
We will see that this comparison leads to three different discrepancies between the calculations and the observations, which I will refer to as the three solar neutrino problems.
Figure 5 shows the measured and the calculated event rates in the five ongoing solar neutrino experiments. This figure reveals three discrepancies between the experimental results and the expectations based upon the combined standard model. As we shall see, only the first of these discrepancies depends in an important way upon the predictions of the standard solar model.
### 3.1 Calculated versus Observed Absolute Rate
The first solar neutrino experiment to be performed was the chlorine radiochemical experiment, which detects electron-type neutrinos that are more energetic than $`0.81`$ MeV. After more than a quarter of a century of operation of this experiment, the measured event rate is $`2.56\pm 0.23`$ SNU, which is a factor of three less than is predicted by the most detailed theoretical calculations, $`7.7_{1.0}^{+1.2}`$ SNU. Most of the predicted rate in the chlorine experiment is from the rare, high-energy <sup>8</sup>B neutrinos, although the <sup>7</sup>Be neutrinos are also expected to contribute significantly. According to standard model calculations, the $`pep`$ neutrinos and the CNO neutrinos (for simplicity not discussed here) are expected to contribute less than 1 SNU to the total event rate.
This discrepancy between the calculations and the observations for the chlorine experiment was, for more than two decades, the only solar neutrino problem. I shall refer to the chlorine disagreement as the โfirstโ solar neutrino problem.
### 3.2 Incompatibility of Chlorine and Water Experiments
The second solar neutrino problem results from a comparison of the measured event rates in the chlorine experiment and in the Japanese pure-water experiments, Kamiokande and SuperKamiokande. The water experiments detect higher-energy neutrinos, most easily above $`7`$ MeV, by by observing the Cerenkov radiation from neutrino-electron scattering: $`\nu +e\nu ^{}+e^{}.`$ According to the standard solar model, <sup>8</sup>B beta decay, and possibly the $`hep`$ reaction, are the only important source of these higher-energy neutrinos.
The Kamiokande and SuperKamiokande experiments show that the observed neutrinos come from the sun. The electrons that are scattered by the incoming neutrinos recoil predominantly in the direction of the sun-earth vector; the relativistic electrons are observed by the Cerenkov radiation they produce in the water detector. In addition, the water Cerenkov experiments measure the energies of individual scattered electrons and therefore provide information about the energy spectrum of the incident solar neutrinos.
The total event rate in the water experiments, about $`0.5`$ the standard model value (see Fig. 5), is determined by the same high-energy <sup>8</sup>B neutrinos that are expected, on the basis of the combined standard model, to dominate the event rate in the chlorine experiment. I have shown elsewhere that solar physics changes the shape of the <sup>8</sup>B neutrino spectrum by less than 1 part in $`10^5`$ . Therefore, we can calculate the rate in the chlorine experiment (threshold $`0.8`$ MeV) that is produced by the <sup>8</sup>B neutrinos observed in the Kamiokande and SuperKamiokande experiments at an order of magnitude higher energy threshold.
If no new physics changes the shape of the <sup>8</sup>B neutrino energy spectrum, the chlorine rate from <sup>8</sup>B alone is $`2.8\pm 0.1`$ SNU for the SuperKamiokande normalization ($`3.2\pm 0.4`$ SNU for the Kamiokande normalization), which exceeds the total observed chlorine rate of $`2.56\pm 0.23`$ SNU.
Comparing the rates of the SuperKamiokande and the chlorine experiments, one findsโassuming that the shape of the energy spectrum of <sup>8</sup>B $`\nu _e`$โs is not changed by new physicsโthat the net contribution to the chlorine experiment from the $`pep`$, <sup>7</sup>Be, and CNO neutrino sources is negative: $`0.2\pm 0.3`$ SNU. The contributions from the $`pep`$, <sup>7</sup>Be, and CNO neutrinos would appear to be completely missing; the standard model prediction for the combined contribution of $`pep`$, <sup>7</sup>Be, and CNO neutrinos is a relatively large $`1.8`$ SNU (see Table 1). On the other hand, we know that the <sup>7</sup>Be neutrinos must be created in the sun since they are produced by electron capture on the same isotope (<sup>7</sup>Be) which gives rise to the <sup>8</sup>B neutrinos by proton capture.
Hans Bethe and I pointed out that this apparent incompatibility of the chlorine and water-Cerenkov experiments constitutes a โsecondโ solar neutrino problem that is almost independent of the absolute rates predicted by solar models. The inference that is usually made from this comparison is that the energy spectrum of $`{}_{}{}^{8}\mathrm{B}`$ neutrinos is changed from the standard shape by physics not included in the simplest version of the standard electroweak model.
### 3.3 Gallium Experiments: No Room for <sup>7</sup>Be Neutrinos
The results of the gallium experiments, GALLEX and SAGE, constitute the third solar neutrino problem. The average observed rate in these two experiments is $`73\pm 5`$ SNU, which is accounted for in the standard model by the theoretical rate of $`72.4`$ SNU that is calculated to come from the basic $`p`$-$`p`$ and $`pep`$ neutrinos (with only a 1% uncertainty in the standard solar model $`p`$-$`p`$ flux). The <sup>8</sup>B neutrinos, which are observed above $`6.5`$ MeV in the Kamiokande experiment, must also contribute to the gallium event rate. Using the standard shape for the spectrum of $`{}_{}{}^{8}\mathrm{B}`$ neutrinos and normalizing to the rate observed in Kamiokande, $`{}_{}{}^{8}\mathrm{B}`$ contributes another $`6`$ SNU. (The contribution predicted by the standard model is $`12`$ SNU, see Table 1.) Given the measured rates in the gallium experiments, there is no room for the additional $`34\pm 3`$ SNU that is expected from <sup>7</sup>Be neutrinos on the basis of standard solar models (see Table 1).
The seeming exclusion of everything but $`p`$-$`p`$ neutrinos in the gallium experiments is the โthirdโ solar neutrino problem. This problem is essentially independent of the previously-discussed solar neutrino problems, since it depends strongly upon the $`p`$-$`p`$ neutrinos that are not observed in the other experiments and whose theoretical flux can be calculated accurately.
The missing <sup>7</sup>Be neutrinos cannot be explained away by a change in solar physics. The <sup>8</sup>B neutrinos that are observed in the Kamiokande experiment are produced in competition with the missing <sup>7</sup>Be neutrinos; the competition is between electron capture on <sup>7</sup>Be versus proton capture on <sup>7</sup>Be. Solar model explanations that reduce the predicted $`{}_{}{}^{7}\mathrm{Be}`$ flux generically reduce much more (too much) the predictions for the observed $`{}_{}{}^{8}\mathrm{B}`$ flux.
The flux of <sup>7</sup>Be neutrinos, $`\varphi ({}_{}{}^{7}\mathrm{Be})`$, is independent of measurement uncertainties in the cross section for the nuclear reaction $`{}_{}{}^{7}\mathrm{Be}(p,\gamma )^8`$B; the cross section for this proton-capture reaction is the most uncertain quantity that enters in an important way in the solar model calculations. The flux of <sup>7</sup>Be neutrinos depends upon the proton-capture reaction only through the ratio
$$\varphi ({}_{}{}^{7}\mathrm{Be})\frac{R(e)}{R(e)+R(p)},$$
(2)
where $`R(e)`$ is the rate of electron capture by <sup>7</sup>Be nuclei and $`R(p)`$ is the rate of proton capture by <sup>7</sup>Be. With standard parameters, solar models yield $`R(p)10^3R(e)`$. Therefore, one would have to increase the value of the $`{}_{}{}^{7}\mathrm{Be}(p,\gamma )^8`$B cross section by more than two orders of magnitude over the current best-estimate (which has an estimated experimental uncertainty of $``$ 10%) in order to affect significantly the calculated <sup>7</sup>Be solar neutrino flux. The required change in the nuclear physics cross section would also increase the predicted neutrino event rate by more than 100 in the Kamiokande experiment, making that prediction completely inconsistent with what is observed.
I conclude that either: 1) at least three of the five operating solar neutrino experiments (the two gallium experiments plus either chlorine or the two water Cerenkov experiments, Kamiokande and SuperKamiokande) have yielded misleading results, or 2) physics beyond the standard electroweak model is required to change the energy spectrum of $`\nu _e`$ after the neutrinos are produced in the center of the sun.
## 4 Uncertainties in the Flux Calculations
I will now discuss uncertainties in the solar model flux calculations.
Table 2 summarizes the uncertainties in the most important solar neutrino fluxes and in the Cl and Ga event rates due to different nuclear fusion reactions (the first four entries), the heavy element to hydrogen mass ratio (Z/X), the radiative opacity, the solar luminosity, the assumed solar age, and the helium and heavy element diffusion coefficients. The $`{}_{}{}^{14}\mathrm{N}+p`$ reaction causes a 0.2% uncertainty in the predicted pp flux and a 0.1 SNU uncertainty in the Cl (Ga) event rates.
The predicted event rates for the chlorine and gallium experiments use recent improved calculations of neutrino absorption cross sections. The uncertainty in the prediction for the gallium rate is dominated by uncertainties in the neutrino absorption cross sections, $`+6.7`$ SNU ($`7`$% of the predicted rate) and $`3.8`$ SNU ($`3`$% of the predicted rate). The uncertainties in the chlorine absorption cross sections cause an error, $`\pm 0.2`$ SNU ($`3`$% of the predicted rate), that is relatively small compared to other uncertainties in predicting the rate for this experiment. For non-standard neutrino energy spectra that result from new neutrino physics, the uncertainties in the predictions for currently favored solutions (which reduce the contributions from the least well-determined <sup>8</sup>B neutrinos) will in general be less than the values quoted here for standard spectra and must be calculated using the appropriate cross section uncertainty for each neutrino energy.
The nuclear fusion uncertainties in Table 2 were taken from Adelberger et al., the neutrino cross section uncertainties from Bahcall (1997) and Bahcall et al. (1996), the heavy element uncertainty was taken from helioseismological measurements, the luminosity and age uncertainties were adopted from BP95, the $`1\sigma `$ fractional uncertainty in the diffusion rate was taken to be $`15`$%, which is supported by helioseismological evidence, and the opacity uncertainty was determined by comparing the results of fluxes computed using the older Los Alamos opacities with fluxes computed using the modern Livermore opacities. To include the effects of asymmetric errors, the now publicly-available code for calculating rates and uncertainties (see discussion in previous section) was run with different input uncertainties and the results averaged. The software contains a description of how each of the uncertainties listed in Table 2 were determined and used.
The low energy cross section of the $`{}_{}{}^{7}\mathrm{Be}+p`$ reaction is the most important quantity that must be determined more accurately in order to decrease the error in the predicted event rates in solar neutrino experiments. The <sup>8</sup>B neutrino flux that is measured by the Kamiokande, Super-Kamiokande, and SNO experiments is, in all standard solar model calculations, directly proportional to the $`{}_{}{}^{7}\mathrm{Be}+p`$ cross section. If the $`1\sigma `$ uncertainty in this cross section can be reduced by a factor of two to 5%, then it will no longer be the limiting uncertainty in predicting the crucial <sup>8</sup>B neutrino flux (cf. Table 2).
## 5 How Large an Uncertainty Does Helioseismology Suggest?
Could the solar model calculations be wrong by enough to explain the discrepancies between predictions and measurements for solar neutrino experiments? Helioseismology, which confirms predictions of the standard solar model to high precision, suggests that the answer is probably โNo.โ
Figure 6 shows the fractional differences between the most accurate available sound speeds measured by helioseismology and sound speeds calculated with our best solar model (with no free parameters). The horizontal line corresponds to the hypothetical case in which the model predictions exactly match the observed values. The rms fractional difference between the calculated and the measured sound speeds is $`1.1\times 10^3`$ for the entire region over which the sound speeds are measured, $`0.05R_{}<R<0.95R_{}`$. In the solar core, $`0.05R_{}<R<0.25R_{}`$ (in which about $`95`$% of the solar energy and neutrino flux is produced in a standard model), the rms fractional difference between measured and calculated sound speeds is $`0.7\times 10^3`$.
Helioseismological measurements also determine two other parameters that help characterize the outer part of the sun (far from the inner region in which neutrinos are produced): the depth of the solar convective zone (CZ), the region in the outer part of the sun that is fully convective, and the present-day surface abundance by mass of helium ($`Y_{\mathrm{surf}}`$). The measured values, $`R_{\mathrm{CZ}}=(0.713\pm 0.001)R_{}`$ and $`Y_{\mathrm{surf}}=0.249\pm 0.003`$ are in satisfactory agreement with the values predicted by the solar model BP98, namely, $`R_{\mathrm{CZ}}=0.714R_{}`$, and $`Y_{\mathrm{surf}}=0.243`$. However, we shall see below that precision measurements of the sound speed near the transition between the radiative interior (in which energy is transported by radiation) and the outer convective zone (in which energy is transported by convection) reveal small discrepancies between the model predictions and the observations in this region.
If solar physics were responsible for the solar neutrino problems, how large would one expect the discrepancies to be between solar model predictions and helioseismological observations? The characteristic size of the discrepancies can be estimated using the results of the neutrino experiments and scaling laws for neutrino fluxes and sound speeds.
All recently published solar models predict essentially the same fluxes from the fundamental pp and pep reactions (amounting to $`72.4`$ SNU in gallium experiments, cf. Table 1), which are closely related to the solar luminosity. Comparing the measured gallium rates and the standard predicted rate for the gallium experiments, the <sup>7</sup>Be flux must be reduced by a factor $`N`$ if the disagreement is not to exceed $`n`$ standard deviations, where $`N`$ and $`n`$ satisfy $`72.4+(34.4)/N=72.2+n\sigma `$. For a $`1\sigma `$ ($`3\sigma `$) disagreement, $`N=6.1(2.05)`$. Sound speeds scale like the square root of the local temperature divided by the mean molecular weight and the <sup>7</sup>Be neutrino flux scales approximately as the $`10`$th power of the temperature. Assuming that the temperature changes are dominant, agreement to within $`1\sigma `$ would require fractional changes of order $`0.09`$ in sound speeds ($`3\sigma `$ could be reached with $`0.04`$ changes), if all model changes were in the temperature<sup>1</sup><sup>1</sup>1I have used in this calculation the GALLEX and SAGE measured rates reported by Kirsten and Gavrin at Neutrino 98. The experimental rates used in BP98 were not as precise and therefore resulted in slightly less stringent constraints than those imposed here. In BP98, we found that agreement to within $`1\sigma `$ with the then available experimental numbers would require fractional changes of order $`0.08`$ in sound speeds ($`3\sigma `$ could be reached with $`0.03`$ changes.). This argument is conservative because it ignores the <sup>8</sup>B and CNO neutrinos which contribute to the observed counting rate (cf. Table 1) and which, if included, would require an even larger reduction of the <sup>7</sup>Be flux.
I have chosen the vertical scale in Fig. 6 to be appropriate for fractional differences between measured and predicted sound speeds that are of order $`0.04`$ to $`0.09`$ and that might therefore affect solar neutrino calculations. Fig. 6 shows that the characteristic agreement between solar model predictions and helioseismological measurements is more than a factor of $`40`$ better than would be expected if there were a solar model explanation of the solar neutrino problems.
## 6 Fits Without Solar Models
Suppose (following the precepts of Hata et al., Parke, and Heeger and Robertson ) we now ignore everything we have learned about solar models over the last $`35`$ years and allow the important $`pp`$, $`{}_{}{}^{7}\mathrm{Be}`$, and $`{}_{}{}^{8}\mathrm{B}`$ fluxes to take on any non-negative values. What is the best fit that one can obtain to the solar neutrino measurements assuming only that the luminosity of the sun is supplied by nuclear fusion reactions among light elements (the so-called โluminosity constraintโ)?
The answer is that the fits are bad, even if we completely ignore what we know about the sun. I quote the results from Bahcall, Krastev and Smirnov (1998).
If the CNO neutrino fluxes are set equal to zero, there are no acceptable solutions at the $`99`$% C. L. ($`3\sigma `$ result). The best-fit is worse if the CNO fluxes are not set equal to zero. All so-called โsolutionsโ of the solar neutrino problems in which the astrophysical model is changed arbitrarily (ignoring helioseismology and other constraints) are inconsistent with the observations at much more than a $`3\sigma `$ level of significance. No fiddling of the physical conditions in the model can yield the minimum value, quoted above, that was found by varying the fluxes independently and arbitrarily.
Figure 4 shows, in the lower left-hand corner, the best-fit solution and the $`1\sigma `$$`3\sigma `$ contours. The $`1\sigma `$ and $`3\sigma `$ limits were obtained by requiring that $`\chi ^2=\chi _{\mathrm{min}}^2+\delta \chi ^2`$, where for $`1\sigma `$ $`\delta \chi ^2=1`$ and for $`3\sigma `$ $`\delta \chi ^2=9`$. All of the standard model solutions lie far from the best-fit solution and even lie far from the $`3\sigma `$ contour.
Since standard model descriptions do not fit the solar neutrino data, we will now consider models in which neutrino oscillations change the shape of the neutrino energy spectra.
## 7 Neutrino Oscillations
The experimental results from all five of the operating solar neutrino experiments (chlorine, Kamiokande, SAGE, GALLEX, and SuperKamiokande) can be fit well by descriptions involving neutrino oscillations, either vacuum oscillations (as originally suggested by Gribov and Pontecorvo )or resonant matter oscillations (as originally discussed by Mikheyev, Smirnov, and Wolfenstein (MSW) ).
Table 3 summarizes the four best-fit solutions that are found in the two-neutrino approximation. Only the SMA and vacuum oscillation solutions fit well the recoil electron energy spectrum measured in the SuperKamiokande experimentโif the standard value for the $`hep`$ production reaction cross section ($`{}_{}{}^{3}\mathrm{He}+p{}_{}{}^{4}\mathrm{He}+e^++\nu _e`$) is used. However, for over a decade I have not given an estimated uncertainty for this cross section. The transition matrix element is essentially forbidden and the actual quoted value for the production cross section depends upon a delicate cancellation between two comparably sized terms that arise from very different and hard to evaluate nuclear physics. I do not see anyway at present to determine from experiment or from first principles theoretical calculations a relevant, robust upper limit to the $`hep`$ production cross section (and therefore the $`hep`$ solar neutrino flux).
The possible role of $`hep`$ neutrinos in solar neutrino experiments is discussed extensively in Bahcall and Krastev (1998) The most important unsolved problem in theoretical nuclear physics related to solar neutrinos is the range of values allowed by fundamental physics for the $`hep`$ production cross section.
## 8 Discussion
When the chlorine solar neutrino experiment was first proposed, the only stated motivation was โโฆto see into the interior of a star and thus verify directly the hypothesis of nuclear energy generation in stars.โ This goal has now been achieved,
The focus has shifted to using solar neutrino experiments as a tool for learning more about the fundamental characteristics of neutrinos as particles. Experimental effort is now concentrated on answering the question: What are the probabilities for transforming a solar $`\nu _e`$ of a definite energy into the other possible neutrino states? Once this question is answered, we can calculate what happens to $`\nu _e`$โs that are created in the interior of the sun. Armed with this information from weak interaction physics, we can return again to the original motivation of using neutrinos to make detailed, quantitative tests of nuclear fusion rates in the solar interior. Measurements of the flavor content of the dominant low energy neutrino sources, $`p`$-$`p`$ and <sup>7</sup>Be neutrinos, will be crucial in this endeavor and will require another generation of superb solar neutrino experiments (see the comments in Section 9).
Three decades of refining the input data and the solar model calculations has led to a predicted standard model event rate for the chlorine experiment, $`7.7`$ SNU, which is very close to $`7.5`$ SNU, the best-estimate value obtained in 1968. The situation regarding solar neutrinos is, however, completely different now, thirty years later. Four experiments have confirmed the original chlorine detection of solar neutrinos. Helioseismological measurements are in excellent agreement with the standard solar model predictions and very strongly disfavor (by a factor of $`40`$ or more) hypothetical deviations from the standard model that are require to fits the neutrino data (cf. Fig. 6). Just in the last two years, improvements in the helioseismological measurements have resulted in a five-fold improvement in the agreement between the calculated standard solar model sound speeds and the measured solar velocities (cf. Figure 2 of the Neutrino 96 talk with Figure 6 of this talk).
## 9 What next?
More than $`98`$% of the calculated standard model solar neutrino flux lies below $`1`$ MeV. The rare <sup>8</sup>B neutrino flux is the only solar neutrino source for which measurements of the energy have been made, but <sup>8</sup>B neutrinos constitute a fraction of less than $`10^4`$ of the total solar neutrino flux.
The next goal of solar neutrino astronomy is to measure neutrino fluxes below $`1`$ MeV. We should begin today preparing for experiments that will measure the <sup>7</sup>Be neutrinos (energy of $`0.86`$ MeV) and the fundamental $`p`$-$`p`$ neutrinos ($`<0.43`$ MeV). Indeed, we have heard at this workshop some marvelously exciting descriptions of how such low energy experiments could be carried out. The BOREXINO observatory, which can detect $`\nu e`$ scattering, is the only approved solar neutrino experiment which can measure energies less than $`1`$ MeV.
The $`p`$-$`p`$ neutrinos are overwhelmingly the most abundant source of solar neutrinos, carrying about $`91`$% of the total flux according to the standard solar model. The <sup>7</sup>Be neutrinos constitute about $`7`$% of the total standard model flux.
If we want to test and to understand neutrino oscillations with high precision using solar neutrino sources, then we have to measure the neutrino-electron scattering rate with <sup>7</sup>Be neutrinos, as will be done with the BOREXINO experiment, and also the CC (neutrino-absorption) rate with <sup>7</sup>Be neutrinos (no approved experiment). With a neutrino line as provided by <sup>7</sup>Be electron-capture in the sun, unique and unambiguous tests of neutrino oscillation models can be carried out if one knows both the charged-current and the neutral current reaction rates .
I believe we have calculated the flux of $`p`$-$`p`$ neutrinos produced in the sun to an accuracy of $`\pm 1`$%. Unfortunately, we do not yet have a direct measurement of this flux. The gallium experiments only tell us the rate of capture of all neutrinos with energies above $`0.23`$ MeV.
The most urgent need for solar neutrino research is to develop a practical experiment to measure directly the $`p`$-$`p`$ neutrino flux and the energy spectrum of electrons produced by target interactions with $`p`$-$`p`$ neutrinos. Such an experiment can be used to test the precise and fundamental standard solar model prediction of the $`p`$-$`p`$ neutrino flux. Moreover, the currently favored neutrino oscillation solutions all predict a strong influence of oscillations on the low-energy flux of $`\nu _e`$.
Figure 7 shows the calculated neutrino survival probability as a function of energy for three global best-fit MSW oscillation solutions. You can see directly from this figure why we have to have accurate measurements for the $`p`$-$`p`$ and <sup>7</sup>Be neutrinos: the currently favored solutions exhibit their most characteristic and strongly energy dependent features below $`1`$ MeV. In all of these solutions, the survival probability shows a dramatic increase with energy below $`1`$ MeV, whereas in the region above $`5`$ MeV (accessible to SuperKamiokande and to SNO) the energy dependence of the survival probability is at best modest.
The $`p`$-$`p`$ neutrinos are the gold ring of solar neutrino astronomy. Their measurement will constitute a simultaneous and critical test of stellar evolution theory and of neutrino oscillation solutions.
The most exciting result of this workshop for me has been the possibility discussed here of a synergistic experiment involving a huge (megaton?) nucleon decay detector in which an inner region is reserved for solar neutrino experiments(see, for example, the talks by Jung, Nakahata, and Ypsilantis at this workshop). The most straightforward solar neutrino experiments that could be carried out with this detector would be precision measurements of the temporal dependences of the relatively high-energy <sup>8</sup>B neutrinos. One could measure with such a detector the zenith-angle dependence of the solar neutrino-event rate (the generalization of the day-night difference) and the seasonal dependence (generalization of the winter-summer difference). The design could relax somewhat the precise requirements for energy calibration and for energy resolution used for the SuperKamiokande and SNO experiments and concentrate instead on limiting the systematic uncertainties in the detector that could contribute to the error budget in the day-night or seasonal dependences. After three years of very careful measurements, the SuperKamiokande experiment has, as we have heard at this conference, about a $`2\sigma `$ result for the day-night difference. They do not yet have the statistics to report a meaningful measurement of the full zenith-angle dependence or the seasonal dependence. The predicted temporal effects are small, generally of order a percent, with the currently favored neutrino oscillation solutions.
Nature has provided us with many different baselines and with many different matter column densities with which to do Very Long Baseline (VLB) studies of neutrino oscillations. The earth-sun distance varies continuously during the year between $`1.496(1.0\pm 0.017)10^{13}`$ cm and the column density through the earth to a terrestrial detector varies from $`0\mathrm{gm}\mathrm{cm}^2`$ during the day to more than $`10^9\mathrm{gm}\mathrm{cm}^2`$ at night.
A solar neutrino detector ten or more times the volume of the current SuperKamiokande experiment, as discussed in concept at this workshop, could measure precisely the results of many different VLB neutrino oscillation experiments. This would be a fantastic โSmoking Gunโ detector. I have a hard time sitting down when imagining such an exciting possibility.
## Acknowledgments
I acknowledge support from NSF grant #PHY95-13835.
## References |
warning/0002/astro-ph0002013.html | ar5iv | text | # 1 Preamble
## 1 Preamble
The subject of large-scale structure is in a period of very rapid development. For many years, this term would have meant only one thing: the distribution of galaxies. However, we are increasingly able to probe the primordial fluctuations through the CMB, so that the problem of galaxy formation and clustering is now only one aspect of the general picture of structure formation. The rationale for studying the large-scale distribution of galaxies is therefore altering. Ten years ago, we were happy to produce samples based on a rather sparse random sampling of the galaxy distribution, with the main aim of tying down statistics such as the large-scale power spectrum of number-density fluctuations. A major goal of the subject remains the measurement of the fluctuation spectrum for wavelengths $`\stackrel{>}{}100`$ Mpc, and the demonstration that this agrees in shape with what can be inferred from the CMB. Nevertheless, we are now increasingly interested in studying the pattern of galaxies with the highest possible fidelity โ demanding deep, fully-sampled surveys of the local universe. Such studies will tell us much about the processes by which galaxies formed and evolved within the distribution of dark matter. The aim of these lectures is therefore to look both backwards and forwards: reviewing the foundations of the subject and looking forward to the future issues.
## 2 The CDM family album
### 2.1 The linear spectrum
The basic picture of inflationary models (but also of cosmology before inflation) is of a primordial power-law spectrum, written dimensionlessly as the logarithmic contribution to the fractional density variance, $`\sigma ^2`$:
$`\mathrm{\Delta }^2(k)={\displaystyle \frac{d\sigma ^2}{d\mathrm{ln}k}}k^{3+n},`$ (1)
where $`n`$ stands for $`n_\mathrm{S}`$ hereafter. This undergoes linear growth
$`\delta _k(a)=\delta _k(a_0)\left[{\displaystyle \frac{D(a)}{D(a_0)}}\right]T_k,`$ (2)
where the linear growth law is
$`D(a)=ag[\mathrm{\Omega }(a)]`$ (3)
in the matter era, and the growth suppression for low $`\mathrm{\Omega }`$ is
$`g(\mathrm{\Omega })`$ $`\mathrm{\Omega }^{0.65}(\mathrm{open})`$ $`\mathrm{\Omega }^{0.23}(\mathrm{flat})`$ (4)
The transfer function $`T_k`$ depends on the dark-matter content as shown in figure 1.
Note the baryonic oscillations in figure 1; these can be significant even in CDM-dominated models when working with high-precision data. Eisenstein & Hu (1998) are to be congratulated for their impressive persistence in finding an accurate fitting formula that describes these wiggles. This is invaluable for carrying out a search of a large parameter space.
The state of the linear-theory spectrum after these modifications is illustrated in figure 2. The primordial power-law spectrum is reduced at large $`k`$, by an amount that depends on both the quantity of dark matter and its nature. Generally the bend in the spectrum occurs near $`1/k`$ of order the horizon size at matter-radiation equality, $`(\mathrm{\Omega }h^2)^1`$. For a pure CDM universe, with scale-invariant initial fluctuations ($`n=1`$), the observed spectrum depends only on two parameters. One is the shape $`\mathrm{\Gamma }=\mathrm{\Omega }h`$, and the other is a normalization. On the shape front, a government health warning is needed, as follows. It has been quite common to take $`\mathrm{\Gamma }`$-based fits to observations as indicating a measurement of $`\mathrm{\Omega }h`$, but there are three reasons why this may give incorrect answers:
(1) The dark matter may not be CDM. An admixture of HDM will damp the spectrum more, mimicking a lower CDM density.
(2) Even in a CDM-dominated universe, baryons can have a significant effect, making $`\mathrm{\Gamma }`$ lower than $`\mathrm{\Omega }h`$. An approximate formula for this is given in figure 2 (Peacock & Dodds 1994; Sugiyama 1995).
(3) The strongest (and most-ignored) effect is tilt: if $`n1`$, then even in a pure CDM universe a $`\mathrm{\Gamma }`$-model fit to the spectrum will give a badly incorrect estimate of the density (the change in $`\mathrm{\Omega }h`$ is roughly $`0.3(n1)`$; Peacock & Dodds 1994).
### 2.2 Normalization
The other parameter is the normalization. This can be set at a number of points. The COBE normalization comes from large angle CMB anisotropies, and is sensitive to the power spectrum at $`k10^3h\mathrm{Mpc}^1`$. The alternative is to set the normalization near the quasilinear scale, using the abundance of rich clusters. Many authors have tried this calculation, and there is good agreement on the answer:
$`\sigma _8(0.50.6)\mathrm{\Omega }_m^{0.6}.`$ (5)
(White, Efstathiou & Frenk 1993; Eke et al. 1996; Viana & Liddle 1996). In many ways, this is the most sensible normalization to use for LSS studies, since it does not rely on an extrapolation from larger scales.
Within the CDM model, it is always possible to satisfy both these normalization constraints, by appropriate choice of $`\mathrm{\Gamma }`$ and $`n`$. This is illustrated in figure 3. Note that vacuum energy affects the answer; for reasonable values of $`h`$ and reasonable baryon content, flat models require $`\mathrm{\Omega }_m0.3`$, whereas open models require $`\mathrm{\Omega }_m0.5`$.
### 2.3 The nonlinear spectrum
On smaller scales ($`k\stackrel{>}{}0.1`$), nonlinear effects become important. These are relatively well understood so far as they affect the power spectrum of the mass (e.g. Hamilton et al. 1991; Jain, Mo & White 1995; Peacock & Dodds 1996). Based on a fitting formula for the similarity solution governing the evolution of scale-free initial conditions, it is possible to predict the evolved spectrum in CDM universes to a few per cent precision (e.g. Jenkins et al. 1998).
These methods can cope with most smoothly-varying power spectra, but they break down for models with a large baryon content. Figure 1 shows that rather large oscillatory features would be expected if the universe was baryon dominated. The lack of observational evidence for such features is one reason for believing that the universe might be dominated by collisionless nonbaryonic matter (consistent with primordial nucleosynthesis if $`\mathrm{\Omega }_m\stackrel{>}{}0.1`$).
Nevertheless, baryonic fluctuations in the spectrum can become significant for high-precision measurements. Figure 4 shows that order 10% modulation of the power may be expected in realistic baryonic models (Eisenstein & Hu 1998; Goldberg & Strauss 1998). Most of these features are however removed by nonlinear evolution. The highest-$`k`$ feature to survive is usually the second peak, which almost always lies near $`k=0.05\mathrm{Mpc}^1`$ (no $`h`$, for a change). This feature is relatively narrow, and can serve as a clear proof of the past existence of baryonic oscillations in forming the mass distribution (Meiksin, White & Peacock 1999). However, figure 4 emphasizes that the easiest way of detecting the presence of baryons is likely to be through the CMB spectrum. The oscillations have a much larger โvisibilityโ there, because the small-scale CMB anisotropies come directly from the coupled radiation-baryon fluid, rather than the small-scale dark matter perturbations.
## 3 Statistics
Statistical measures of the cosmological density field relate to properties of the dimensionless density perturbation field
$`\delta (๐ฑ){\displaystyle \frac{\rho (๐ฑ)\rho }{\rho }},`$ (6)
although $`\delta `$ need not be assumed to be small.
### 3.1 Correlation functions
The simplest measure is the autocorrelation function of the density perturbation
$`\xi _\mathrm{A}(๐ซ)\delta (๐ฑ)\delta (๐ฑ+๐ซ),`$ (7)
This is a straightforward statistical measure that can also be computed for the dark-matter distribution in $`N`$-body simulations. Formally, the averaging operator here is an ensemble average, but one generally appeals to the ergodic nature of the density field to replace this with a volume average.
However, galaxies are a point process, so what astronomers can measure in practice is the two-point correlation function, which gives the excess probability for finding a neighbour a distance $`r`$ from a given galaxy. By regarding this as the probability of finding a pair with one object in each of the volume elements $`dV_1`$ and $`dV_2`$,
$`dP=\rho _0^2[1+\xi _2(r)]dV_1dV_2.`$ (8)
Is it true that $`\xi _\mathrm{A}(r)=\xi _2(r)\mathrm{?}`$ Life would certainly be simple if so, and much work on large-scale structure has implicitly assumed the Poisson clustering hypothesis, in which galaxies are assumed to be sampled at random from some continuous underlying density field. Many of the puzzles in the field can however be traced to the fact that this hypothesis is probably false, as discussed below.
A related quantity is the cross-correlation function. Here, one considers two different classes of object (a and b, say), and the cross-correlation function $`\xi _{ab}`$ is defined as the (symmetric) probability of finding a pair in which $`dV_1`$ is occupied by an object from the first catalogue and $`dV_2`$ by one from the second. Both cross- and auto-correlation functions are readily extended to higher orders and considerations of $`n`$-tuples of points in a given geometry.
### 3.2 Fourier space
For the Fourier counterpart of this analysis, we assume that the field is periodic within some box of side $`L`$, and expand as a Fourier series:
$`\delta (๐ฑ)={\displaystyle \delta _๐คe^{i๐ค๐ฑ}}.`$ (9)
For a real field, $`\delta _๐ค(๐ค)=\delta _๐ค^{}(๐ค)`$. Using this definition in the correlation function, most cross terms integrate to zero through the periodic boundary conditions, giving
$`\xi (๐ซ)={\displaystyle \frac{V}{(2\pi )^3}}{\displaystyle |\delta _๐ค|^2e^{i๐ค๐ซ}d^3k}.`$ (10)
In short, the correlation function is the Fourier transform of the power spectrum.
We shall usually express the power spectrum in dimensionless form, as the variance per $`\mathrm{ln}k`$ ($`\mathrm{\Delta }^2(k)=d\delta ^2/d\mathrm{ln}kk^3P[k]`$):
$`\mathrm{\Delta }^2(k){\displaystyle \frac{V}{(2\pi )^3}}\mathrm{\hspace{0.17em}4}\pi k^3P(k)={\displaystyle \frac{2}{\pi }}k^3{\displaystyle _0^{\mathrm{}}}\xi (r){\displaystyle \frac{\mathrm{sin}kr}{kr}}r^2๐r.`$ (11)
This gives a more easily visualizable meaning to the power spectrum than does the quantity $`VP(k)`$, which has dimensions of volume: $`\mathrm{\Delta }^2(k)=1`$ means that there are order-unity density fluctuations from modes in the logarithmic bin around wavenumber $`k`$. $`\mathrm{\Delta }^2(k)`$ is therefore the natural choice for a Fourier-space counterpart to the dimensionless quantity $`\xi (r)`$.
In the days before inflation, the primordial power spectrum was chosen by hand, and the minimal assumption was a featureless power law:
$`|\delta _k|^2P(k)k^n`$ (12)
The index $`n`$ governs the balance between large- and small-scale power. Similarly, a power-law spectrum implies a power-law correlation function. If $`\xi (r)=(r/r_0)^\gamma `$, with $`\gamma =n+3`$, the corresponding 3D power spectrum is
$`\mathrm{\Delta }^2(k)={\displaystyle \frac{2}{\pi }}(kr_0)^\gamma \mathrm{\Gamma }(2\gamma )\mathrm{sin}{\displaystyle \frac{(2\gamma )\pi }{2}}\beta (kr_0)^\gamma `$ (13)
($`=0.903(kr_0)^{1.8}`$ if $`\gamma =1.8`$). This expression is only valid for $`n<0`$ ($`\gamma <3`$); for larger values of $`n`$, $`\xi `$ must become negative at large $`r`$ (because $`P(0)`$ must vanish, implying $`_0^{\mathrm{}}\xi (r)r^2๐r=0`$). A cutoff in the spectrum at large $`k`$ is needed to obtain physically sensible results.
The most interesting value of $`n`$ is the scale-invariant spectrum, $`n=1`$, i.e. $`\mathrm{\Delta }^2k^4`$. To see how the name arises, consider a perturbation $`\delta \mathrm{\Phi }`$ in the gravitational potential:
$`^2\delta \mathrm{\Phi }=4\pi G\rho _0\delta \delta \mathrm{\Phi }_k=4\pi G\rho _0\delta _k/k^2.`$ (14)
The two powers of $`k`$ pulled down by $`^2`$ mean that, if $`\mathrm{\Delta }^2k^4`$ for the power spectrum of density fluctuations, then $`\mathrm{\Delta }_\mathrm{\Phi }^2`$ is a constant. Since potential perturbations govern the flatness of spacetime, this says that the scale-invariant spectrum corresponds to a metric that is a fractal: spacetime has the same degree of โwrinklinessโ on each resolution scale. The total curvature fluctuations diverge, but only logarithmically at either extreme of wavelength.
### 3.3 Error estimates
A key question for these statistical measures is how accurate they are โ i.e. how much does the result for a given finite sample depart from the ideal statistic averaged over an infinite universe? Terminology here can be confusing, in that a distinction is sometimes made between sampling variance and cosmic variance. The former is to be understood as arising from probing a given volume only with a finite number of galaxies (e.g. just the bright ones), so that $`\sqrt{N}`$ statistics limit our knowledge of the mass distribution within that region. The second term concerns whether we have reached a fair sample of the universe, and depends on whether there is significant power in density perturbation modes with wavelengths larger than the sample depth. Clearly, these two aspects are closely related.
The quantitative analysis of these errors is most simply performed in Fourier space, and was given by Feldman, Kaiser & Peacock (1994). The results can be understood most simply by comparison with an idealized complete and uniform survey of a volume $`L^3`$, with periodicity scale $`L`$. For an infinite survey, the arbitrariness of the spatial origin means that different modes are uncorrelated:
$`\delta _k(๐ค_i)\delta _k^{}(๐ค_j)=P(k)\delta _{ij}.`$ (15)
Each mode has an exponential distribution in power (because the complex coefficients $`\delta _k`$ are 2D Gaussian-distributed variables on the Argand plane), for which the mean and rms are identical. The fractional uncertainty in the mean power measured over some $`k`$-space volume is then just determined by the number of uncorrelated modes averaged over:
$`{\displaystyle \frac{\delta \overline{P}}{\overline{P}}}={\displaystyle \frac{1}{N_{\mathrm{modes}}^{1/2}}};N_{\mathrm{modes}}=\left({\displaystyle \frac{L}{2\pi }}\right)^3{\displaystyle d^3k}.`$ (16)
The only subtlety is that, because the density field is real, modes at $`k`$ and $`k`$ are perfectly correlated. Thus, if the $`k`$-space volume is a shell, the effective number of uncorrelated modes is only half the above expression.
Analogous results apply for an arbitrary survey selection function. In the continuum limit, the Kroneker delta in the expression for mode correlation would be replaced a term proportional to a delta-function, $`\delta [๐ค_i๐ค_j]`$). Now, multiplying the infinite ideal survey by a survey window, $`\rho (๐ซ)`$, is equivalent to convolution in the Fourier domain, with the result that the power per mode is correlated over $`k`$-space separations of order $`1/D`$, where $`D`$ is the survey depth.
Given this expression for the fractional power, it is clear that the precision of the estimate can be manipulated by appropriate weighting of the data: giving increased weight to the most distant galaxies increases the effective survey volume, boosting the number of modes. This sounds too good to be true, and of course it is: the above expression for the fractional power error applies to the sum of true clustering power and shot noise. The latter arises because we transform a point process. Given a set of $`N`$ galaxies, we would estimate Fourier coefficients via $`\delta _k=(1/N)_i\mathrm{exp}(i๐คx_i)`$. From this, the expectation power is
$`|\delta _k|^2=P(k)+1/N.`$ (17)
The existence of an additive discreteness correction is no problem, but the fluctuations on the shot noise hide the signal of interest. Introducing weights boosts the shot noise, so there is an optimum choice of weight that minimizes the uncertainty in the power after shot-noise subtraction. Feldman, Kaiser & Peacock (1994) showed that this weight is
$`w=(1+\overline{n}P)^1,`$ (18)
where $`\overline{n}`$ is the expected galaxy number density as a function of position in the survey.
Since the correlation of modes arises from the survey selection function, it is clear that weighting the data changes the degree of correlation in $`k`$ space. Increasing the weight in low-density regions increases the effective survey volume, and so shrinks the $`k`$-space coherence scale. However, the coherence scale continues to shrink as distant regions of the survey are given greater weight, whereas the noise goes through a minimum. There is thus a trade-off between the competing desirable criteria of high $`k`$-space resolution and low noise. Tegmark (1996) shows how weights may be chosen to implement any given prejudice concerning the relative importance of these two criteria. See also Hamilton (1997b,c) for similar arguments.
### 3.4 Karhunen-Loรจve and all that
Given these difficulties with correlated results, it is attractive to seek a method where the data can be decomposed into a set of statistics that are completely uncorrelated with each other. Such a method is provided by the Karhunen-Loรจve formalism. Vogeley & Szalay (1996) argued as follows. Define a column vector of data d d d ๐\hbox to0.0pt{\hbox to5.20486pt{\hfil\text@underline{\phantom{\hbox{$d$}}}\hglue 1.0pt\hfil}\hss}d; this can be quite abstract in nature, and could be e.g. the numbers of galaxies in a set of cells, or a set of Fourier components of the transformed galaxy number counts. Similarly, for CMB studies, d d d ๐\hbox to0.0pt{\hbox to5.20486pt{\hfil\text@underline{\phantom{\hbox{$d$}}}\hglue 1.0pt\hfil}\hss}d could be $`\delta T/T`$ in a set of pixels, or spherical-harmonic coefficients $`a_\mathrm{}m`$. We assume that the mean can be identified and subtracted off, so that d d=0delimited-โจโฉ d ๐0\langle\hbox to0.0pt{\hbox to5.20486pt{\hfil\text@underline{\phantom{\hbox{$d$}}}\hglue 1.0pt\hfil}\hss}d\rangle=0 in ensemble average. The statistical properties of the data are then described by the covariance matrix
$`C_{ij}d_id_j^{}`$ (19)
(normally the data will be real, but it is convenient to keep things general and include the complex conjugate).
Suppose we seek to expand the datavector in terms of a set of new orthonormal vectors:
d d=iai ฯ ฯi; ฯ ฯi ฯ ฯj=ฮดij.formulae-sequence d ๐subscript๐subscript๐๐ ฯ subscript๐๐ ฯ subscriptsuperscript๐๐ ฯ subscript๐๐subscript๐ฟ๐๐\displaystyle{\hbox to0.0pt{\hbox to5.20486pt{\hfil\text@underline{\phantom{\hbox{$d$}}}\hglue 1.0pt\hfil}\hss}d}=\sum_{i}a_{i}{\hbox to0.0pt{\hbox to6.51392pt{\hfil\text@underline{\phantom{\hbox{$\psi$}}}\hglue 1.0pt\hfil}\hss}\psi}_{i};\quad\quad{\hbox to0.0pt{\hbox to6.51392pt{\hfil\text@underline{\phantom{\hbox{$\psi$}}}\hglue 1.0pt\hfil}\hss}\psi}^{*}_{i}\cdot{\hbox to0.0pt{\hbox to6.51392pt{\hfil\text@underline{\phantom{\hbox{$\psi$}}}\hglue 1.0pt\hfil}\hss}\psi}_{j}=\delta_{ij}. (20)
The expansion coefficients are extracted in the usual way: aj= d d ฯ ฯjsubscript๐๐ d ๐ ฯ superscriptsubscript๐๐\smash{a_{j}=\hbox to0.0pt{\hbox to5.20486pt{\hfil\text@underline{\phantom{\hbox{$d$}}}\hglue 1.0pt\hfil}\hss}d\cdot\hbox to0.0pt{\hbox to6.51392pt{\hfil\text@underline{\phantom{\hbox{$\psi$}}}\hglue 1.0pt\hfil}\hss}\psi_{j}^{*}}. Now require that these coefficients be statistically uncorrelated, $`a_ia_j^{}=\lambda _i\delta _{ij}`$ (no sum on $`i`$). This gives
ฯ ฯi d d d d ฯ ฯj=ฮปiฮดij, ฯ superscriptsubscript๐๐delimited-โจโฉ d ๐ d superscript๐ ฯ subscript๐๐subscript๐๐subscript๐ฟ๐๐\displaystyle\hbox to0.0pt{\hbox to6.51392pt{\hfil\text@underline{\phantom{\hbox{$\psi$}}}\hglue 1.0pt\hfil}\hss}\psi_{i}^{*}\cdot\langle\hbox to0.0pt{\hbox to5.20486pt{\hfil\text@underline{\phantom{\hbox{$d$}}}\hglue 1.0pt\hfil}\hss}d\,\hbox to0.0pt{\hbox to5.20486pt{\hfil\text@underline{\phantom{\hbox{$d$}}}\hglue 1.0pt\hfil}\hss}d^{*}\rangle\cdot\hbox to0.0pt{\hbox to6.51392pt{\hfil\text@underline{\phantom{\hbox{$\psi$}}}\hglue 1.0pt\hfil}\hss}\psi_{j}=\lambda_{i}\delta_{ij}, (21)
where the dyadic d d d ddelimited-โจโฉ d ๐ d superscript๐\langle\hbox to0.0pt{\hbox to5.20486pt{\hfil\text@underline{\phantom{\hbox{$d$}}}\hglue 1.0pt\hfil}\hss}d\,\hbox to0.0pt{\hbox to5.20486pt{\hfil\text@underline{\phantom{\hbox{$d$}}}\hglue 1.0pt\hfil}\hss}d^{*}\rangle is CC C C CC C ๐ถ\hbox to0.0pt{\hbox to7.86249pt{\hfil\text@underline{\phantom{\hbox{$\hbox to0.0pt{\hbox to7.86249pt{\hfil\text@underline{\phantom{\hbox{$C$}}}\hglue 1.0pt\hfil}\hss}C$}}}\hglue 1.0pt\hfil}\hss}\hbox to0.0pt{\hbox to7.86249pt{\hfil\text@underline{\phantom{\hbox{$C$}}}\hglue 1.0pt\hfil}\hss}C, the correlation matrix of the data vector: ( d d d d)ijdidjsubscript d ๐ d superscript๐๐๐subscript๐๐subscriptsuperscript๐๐(\hbox to0.0pt{\hbox to5.20486pt{\hfil\text@underline{\phantom{\hbox{$d$}}}\hglue 1.0pt\hfil}\hss}d\,\hbox to0.0pt{\hbox to5.20486pt{\hfil\text@underline{\phantom{\hbox{$d$}}}\hglue 1.0pt\hfil}\hss}d^{*})_{ij}\equiv d_{i}d^{*}_{j}. Now, the effect of operating this matrix on one of the ฯ ฯi ฯ subscript๐๐\hbox to0.0pt{\hbox to6.51392pt{\hfil\text@underline{\phantom{\hbox{$\psi$}}}\hglue 1.0pt\hfil}\hss}\psi_{i} must be expandable in terms of the complete set, which shows that the ฯ ฯj ฯ subscript๐๐\smash{\hbox to0.0pt{\hbox to6.51392pt{\hfil\text@underline{\phantom{\hbox{$\psi$}}}\hglue 1.0pt\hfil}\hss}\psi_{j}} must be the eigenvectors of the correlation matrix:
d d d d ฯ ฯj=ฮปj ฯ ฯj.delimited-โจโฉ d ๐ d superscript๐ ฯ subscript๐๐subscript๐๐ ฯ subscript๐๐\displaystyle\langle\hbox to0.0pt{\hbox to5.20486pt{\hfil\text@underline{\phantom{\hbox{$d$}}}\hglue 1.0pt\hfil}\hss}d\,\hbox to0.0pt{\hbox to5.20486pt{\hfil\text@underline{\phantom{\hbox{$d$}}}\hglue 1.0pt\hfil}\hss}d^{*}\rangle\cdot\hbox to0.0pt{\hbox to6.51392pt{\hfil\text@underline{\phantom{\hbox{$\psi$}}}\hglue 1.0pt\hfil}\hss}\psi_{j}=\lambda_{j}\hbox to0.0pt{\hbox to6.51392pt{\hfil\text@underline{\phantom{\hbox{$\psi$}}}\hglue 1.0pt\hfil}\hss}\psi_{j}. (22)
Vogeley & Szalay further show that these uncorrelated modes are optimal for representing the data: if the modes are arranged in order of decreasing $`\lambda `$, and the series expansion truncated after $`n`$ terms, the rms truncation error is minimized for this choice of eigenmodes. To prove this, consider the truncation error
ฯต ฯต= d di=1nai ฯ ฯi=i=n+1ai ฯ ฯi. ฯต italic-ฯต d ๐superscriptsubscript๐1๐subscript๐๐ ฯ subscript๐๐superscriptsubscript๐๐1subscript๐๐ ฯ subscript๐๐\displaystyle\hbox to0.0pt{\hbox to4.05904pt{\hfil\text@underline{\phantom{\hbox{$\epsilon$}}}\hglue 1.0pt\hfil}\hss}\epsilon=\hbox to0.0pt{\hbox to5.20486pt{\hfil\text@underline{\phantom{\hbox{$d$}}}\hglue 1.0pt\hfil}\hss}d-\sum_{i=1}^{n}a_{i}\hbox to0.0pt{\hbox to6.51392pt{\hfil\text@underline{\phantom{\hbox{$\psi$}}}\hglue 1.0pt\hfil}\hss}\psi_{i}=\sum_{i=n+1}^{\infty}a_{i}\hbox to0.0pt{\hbox to6.51392pt{\hfil\text@underline{\phantom{\hbox{$\psi$}}}\hglue 1.0pt\hfil}\hss}\psi_{i}. (23)
The square of this is
$`ฯต^2={\displaystyle \underset{i=n+1}{\overset{\mathrm{}}{}}}|a_i|^2,`$ (24)
where |ai|2= ฯ ฯi CC C C ฯ ฯidelimited-โจโฉsuperscriptsubscript๐๐2 ฯ superscriptsubscript๐๐ CC C ๐ถ ฯ subscript๐๐\langle|a_{i}|^{2}\rangle=\hbox to0.0pt{\hbox to6.51392pt{\hfil\text@underline{\phantom{\hbox{$\psi$}}}\hglue 1.0pt\hfil}\hss}\psi_{i}^{*}\cdot\hbox to0.0pt{\hbox to7.86249pt{\hfil\text@underline{\phantom{\hbox{$\hbox to0.0pt{\hbox to7.86249pt{\hfil\text@underline{\phantom{\hbox{$C$}}}\hglue 1.0pt\hfil}\hss}C$}}}\hglue 1.0pt\hfil}\hss}\hbox to0.0pt{\hbox to7.86249pt{\hfil\text@underline{\phantom{\hbox{$C$}}}\hglue 1.0pt\hfil}\hss}C\cdot\hbox to0.0pt{\hbox to6.51392pt{\hfil\text@underline{\phantom{\hbox{$\psi$}}}\hglue 1.0pt\hfil}\hss}\psi_{i}, as before. We want to minimize $`ฯต^2`$ by varying the ฯ ฯi ฯ subscript๐๐\hbox to0.0pt{\hbox to6.51392pt{\hfil\text@underline{\phantom{\hbox{$\psi$}}}\hglue 1.0pt\hfil}\hss}\psi_{i}, but we need to do this in a way that preserves normalization. This is achieved by introducing a Lagrange multiplier, and minimizing
ฯ ฯi CC C C ฯ ฯi+ฮป(1 ฯ ฯi ฯ ฯi). ฯ superscriptsubscript๐๐ CC C ๐ถ ฯ subscript๐๐๐1 ฯ superscriptsubscript๐๐ ฯ subscript๐๐\displaystyle\sum\hbox to0.0pt{\hbox to6.51392pt{\hfil\text@underline{\phantom{\hbox{$\psi$}}}\hglue 1.0pt\hfil}\hss}\psi_{i}^{*}\cdot\hbox to0.0pt{\hbox to7.86249pt{\hfil\text@underline{\phantom{\hbox{$\hbox to0.0pt{\hbox to7.86249pt{\hfil\text@underline{\phantom{\hbox{$C$}}}\hglue 1.0pt\hfil}\hss}C$}}}\hglue 1.0pt\hfil}\hss}\hbox to0.0pt{\hbox to7.86249pt{\hfil\text@underline{\phantom{\hbox{$C$}}}\hglue 1.0pt\hfil}\hss}C\cdot\hbox to0.0pt{\hbox to6.51392pt{\hfil\text@underline{\phantom{\hbox{$\psi$}}}\hglue 1.0pt\hfil}\hss}\psi_{i}+\lambda(1-\hbox to0.0pt{\hbox to6.51392pt{\hfil\text@underline{\phantom{\hbox{$\psi$}}}\hglue 1.0pt\hfil}\hss}\psi_{i}^{*}\cdot\hbox to0.0pt{\hbox to6.51392pt{\hfil\text@underline{\phantom{\hbox{$\psi$}}}\hglue 1.0pt\hfil}\hss}\psi_{i}). (25)
This is easily solved if we consider the more general problem where ฯ ฯi ฯ superscriptsubscript๐๐\hbox to0.0pt{\hbox to6.51392pt{\hfil\text@underline{\phantom{\hbox{$\psi$}}}\hglue 1.0pt\hfil}\hss}\psi_{i}^{*} and ฯ ฯi ฯ subscript๐๐\hbox to0.0pt{\hbox to6.51392pt{\hfil\text@underline{\phantom{\hbox{$\psi$}}}\hglue 1.0pt\hfil}\hss}\psi_{i} are independent vectors:
CC C C ฯ ฯi=ฮปฯi. CC C ๐ถ ฯ subscript๐๐๐subscript๐๐\displaystyle\hbox to0.0pt{\hbox to7.86249pt{\hfil\text@underline{\phantom{\hbox{$\hbox to0.0pt{\hbox to7.86249pt{\hfil\text@underline{\phantom{\hbox{$C$}}}\hglue 1.0pt\hfil}\hss}C$}}}\hglue 1.0pt\hfil}\hss}\hbox to0.0pt{\hbox to7.86249pt{\hfil\text@underline{\phantom{\hbox{$C$}}}\hglue 1.0pt\hfil}\hss}C\cdot\hbox to0.0pt{\hbox to6.51392pt{\hfil\text@underline{\phantom{\hbox{$\psi$}}}\hglue 1.0pt\hfil}\hss}\psi_{i}=\lambda\psi_{i}. (26)
In short, the eigenvectors of CC C C CC C ๐ถ\hbox to0.0pt{\hbox to7.86249pt{\hfil\text@underline{\phantom{\hbox{$\hbox to0.0pt{\hbox to7.86249pt{\hfil\text@underline{\phantom{\hbox{$C$}}}\hglue 1.0pt\hfil}\hss}C$}}}\hglue 1.0pt\hfil}\hss}\hbox to0.0pt{\hbox to7.86249pt{\hfil\text@underline{\phantom{\hbox{$C$}}}\hglue 1.0pt\hfil}\hss}C are optimal in a least-squares sense for expanding the data. The process of truncating the expansion is a form of lossy data compression, since the size of the data vector can be greatly reduced without significantly affecting the fidelity of the resulting representation of the universe.
The process of diagonalizing the covariance matrix of a set of data also goes by the more familiar name of principal components analysis, so what is the difference between the KL approach and PCA? In the above discussion, they are identical, but the idea of choosing an optimal eigenbasis is more general than PCA. Consider the case where the covariance matrix can be decomposed into a โsignalโ and a โnoiseโ term:
CC C C= SS S S+ NN N N, CC C ๐ถ SS S ๐ NN N ๐\displaystyle\hbox to0.0pt{\hbox to7.86249pt{\hfil\text@underline{\phantom{\hbox{$\hbox to0.0pt{\hbox to7.86249pt{\hfil\text@underline{\phantom{\hbox{$C$}}}\hglue 1.0pt\hfil}\hss}C$}}}\hglue 1.0pt\hfil}\hss}\hbox to0.0pt{\hbox to7.86249pt{\hfil\text@underline{\phantom{\hbox{$C$}}}\hglue 1.0pt\hfil}\hss}C=\hbox to0.0pt{\hbox to6.70831pt{\hfil\text@underline{\phantom{\hbox{$\hbox to0.0pt{\hbox to6.70831pt{\hfil\text@underline{\phantom{\hbox{$S$}}}\hglue 1.0pt\hfil}\hss}S$}}}\hglue 1.0pt\hfil}\hss}\hbox to0.0pt{\hbox to6.70831pt{\hfil\text@underline{\phantom{\hbox{$S$}}}\hglue 1.0pt\hfil}\hss}S+\hbox to0.0pt{\hbox to9.12497pt{\hfil\text@underline{\phantom{\hbox{$\hbox to0.0pt{\hbox to9.12497pt{\hfil\text@underline{\phantom{\hbox{$N$}}}\hglue 1.0pt\hfil}\hss}N$}}}\hglue 1.0pt\hfil}\hss}\hbox to0.0pt{\hbox to9.12497pt{\hfil\text@underline{\phantom{\hbox{$N$}}}\hglue 1.0pt\hfil}\hss}N, (27)
where SS S S SS S ๐\hbox to0.0pt{\hbox to6.70831pt{\hfil\text@underline{\phantom{\hbox{$\hbox to0.0pt{\hbox to6.70831pt{\hfil\text@underline{\phantom{\hbox{$S$}}}\hglue 1.0pt\hfil}\hss}S$}}}\hglue 1.0pt\hfil}\hss}\hbox to0.0pt{\hbox to6.70831pt{\hfil\text@underline{\phantom{\hbox{$S$}}}\hglue 1.0pt\hfil}\hss}S depends on cosmological parameters that we might wish to estimate, whereas NN N N NN N ๐\hbox to0.0pt{\hbox to9.12497pt{\hfil\text@underline{\phantom{\hbox{$\hbox to0.0pt{\hbox to9.12497pt{\hfil\text@underline{\phantom{\hbox{$N$}}}\hglue 1.0pt\hfil}\hss}N$}}}\hglue 1.0pt\hfil}\hss}\hbox to0.0pt{\hbox to9.12497pt{\hfil\text@underline{\phantom{\hbox{$N$}}}\hglue 1.0pt\hfil}\hss}N is some fixed property of the experiment under consideration. In the simplest imaginable case, NN N N NN N ๐\hbox to0.0pt{\hbox to9.12497pt{\hfil\text@underline{\phantom{\hbox{$\hbox to0.0pt{\hbox to9.12497pt{\hfil\text@underline{\phantom{\hbox{$N$}}}\hglue 1.0pt\hfil}\hss}N$}}}\hglue 1.0pt\hfil}\hss}\hbox to0.0pt{\hbox to9.12497pt{\hfil\text@underline{\phantom{\hbox{$N$}}}\hglue 1.0pt\hfil}\hss}N might be a diagonal matrix, so PCA diagonalizes both SS S S SS S ๐\hbox to0.0pt{\hbox to6.70831pt{\hfil\text@underline{\phantom{\hbox{$\hbox to0.0pt{\hbox to6.70831pt{\hfil\text@underline{\phantom{\hbox{$S$}}}\hglue 1.0pt\hfil}\hss}S$}}}\hglue 1.0pt\hfil}\hss}\hbox to0.0pt{\hbox to6.70831pt{\hfil\text@underline{\phantom{\hbox{$S$}}}\hglue 1.0pt\hfil}\hss}S and NN N N NN N ๐\hbox to0.0pt{\hbox to9.12497pt{\hfil\text@underline{\phantom{\hbox{$\hbox to0.0pt{\hbox to9.12497pt{\hfil\text@underline{\phantom{\hbox{$N$}}}\hglue 1.0pt\hfil}\hss}N$}}}\hglue 1.0pt\hfil}\hss}\hbox to0.0pt{\hbox to9.12497pt{\hfil\text@underline{\phantom{\hbox{$N$}}}\hglue 1.0pt\hfil}\hss}N. In this case, ranking the PCA modes by eigenvalue would correspond to ordering the modes according to signal-to-noise ratio. Data compression by truncating the mode expansion then does the sensible thing: it rejects all modes of low signal-to-noise ratio.
However, in general these matrices will not commute, and there will not be a single set of eigenfunctions that are common to the SS S S SS S ๐\hbox to0.0pt{\hbox to6.70831pt{\hfil\text@underline{\phantom{\hbox{$\hbox to0.0pt{\hbox to6.70831pt{\hfil\text@underline{\phantom{\hbox{$S$}}}\hglue 1.0pt\hfil}\hss}S$}}}\hglue 1.0pt\hfil}\hss}\hbox to0.0pt{\hbox to6.70831pt{\hfil\text@underline{\phantom{\hbox{$S$}}}\hglue 1.0pt\hfil}\hss}S and NN N N NN N ๐\hbox to0.0pt{\hbox to9.12497pt{\hfil\text@underline{\phantom{\hbox{$\hbox to0.0pt{\hbox to9.12497pt{\hfil\text@underline{\phantom{\hbox{$N$}}}\hglue 1.0pt\hfil}\hss}N$}}}\hglue 1.0pt\hfil}\hss}\hbox to0.0pt{\hbox to9.12497pt{\hfil\text@underline{\phantom{\hbox{$N$}}}\hglue 1.0pt\hfil}\hss}N matrices. Normally, this would be taken to mean that it is impossible to find a set of coordinates in which both are diagonal. This conclusion can however be evaded, as follows. When considering the effect of coordinate transformations on vectors and matrices, we are normally forced to consider only rotation-like transformations that preserve the norm of a vector (e.g. in quantum mechanics, so that states stay normalized). Thus, we write d d= RR R R d d d superscript๐ RR R ๐
d ๐\hbox to0.0pt{\hbox to5.20486pt{\hfil\text@underline{\phantom{\hbox{$d$}}}\hglue 1.0pt\hfil}\hss}d^{\prime}=\hbox to0.0pt{\hbox to7.67015pt{\hfil\text@underline{\phantom{\hbox{$\hbox to0.0pt{\hbox to7.67015pt{\hfil\text@underline{\phantom{\hbox{$R$}}}\hglue 1.0pt\hfil}\hss}R$}}}\hglue 1.0pt\hfil}\hss}\hbox to0.0pt{\hbox to7.67015pt{\hfil\text@underline{\phantom{\hbox{$R$}}}\hglue 1.0pt\hfil}\hss}R\cdot\hbox to0.0pt{\hbox to5.20486pt{\hfil\text@underline{\phantom{\hbox{$d$}}}\hglue 1.0pt\hfil}\hss}d, where RR R R RR R ๐
\hbox to0.0pt{\hbox to7.67015pt{\hfil\text@underline{\phantom{\hbox{$\hbox to0.0pt{\hbox to7.67015pt{\hfil\text@underline{\phantom{\hbox{$R$}}}\hglue 1.0pt\hfil}\hss}R$}}}\hglue 1.0pt\hfil}\hss}\hbox to0.0pt{\hbox to7.67015pt{\hfil\text@underline{\phantom{\hbox{$R$}}}\hglue 1.0pt\hfil}\hss}R is unitary, so that RR R R RR R R= II I I RR R ๐
RR R superscript๐
II I ๐ผ\hbox to0.0pt{\hbox to7.67015pt{\hfil\text@underline{\phantom{\hbox{$\hbox to0.0pt{\hbox to7.67015pt{\hfil\text@underline{\phantom{\hbox{$R$}}}\hglue 1.0pt\hfil}\hss}R$}}}\hglue 1.0pt\hfil}\hss}\hbox to0.0pt{\hbox to7.67015pt{\hfil\text@underline{\phantom{\hbox{$R$}}}\hglue 1.0pt\hfil}\hss}R\cdot\hbox to0.0pt{\hbox to7.67015pt{\hfil\text@underline{\phantom{\hbox{$\hbox to0.0pt{\hbox to7.67015pt{\hfil\text@underline{\phantom{\hbox{$R$}}}\hglue 1.0pt\hfil}\hss}R$}}}\hglue 1.0pt\hfil}\hss}\hbox to0.0pt{\hbox to7.67015pt{\hfil\text@underline{\phantom{\hbox{$R$}}}\hglue 1.0pt\hfil}\hss}R^{\dagger}=\hbox to0.0pt{\hbox to5.18054pt{\hfil\text@underline{\phantom{\hbox{$\hbox to0.0pt{\hbox to5.18054pt{\hfil\text@underline{\phantom{\hbox{$I$}}}\hglue 1.0pt\hfil}\hss}I$}}}\hglue 1.0pt\hfil}\hss}\hbox to0.0pt{\hbox to5.18054pt{\hfil\text@underline{\phantom{\hbox{$I$}}}\hglue 1.0pt\hfil}\hss}I. If RR R R RR R ๐
\hbox to0.0pt{\hbox to7.67015pt{\hfil\text@underline{\phantom{\hbox{$\hbox to0.0pt{\hbox to7.67015pt{\hfil\text@underline{\phantom{\hbox{$R$}}}\hglue 1.0pt\hfil}\hss}R$}}}\hglue 1.0pt\hfil}\hss}\hbox to0.0pt{\hbox to7.67015pt{\hfil\text@underline{\phantom{\hbox{$R$}}}\hglue 1.0pt\hfil}\hss}R is chosen so that its columns are the eigenvalues of NN N N NN N ๐\hbox to0.0pt{\hbox to9.12497pt{\hfil\text@underline{\phantom{\hbox{$\hbox to0.0pt{\hbox to9.12497pt{\hfil\text@underline{\phantom{\hbox{$N$}}}\hglue 1.0pt\hfil}\hss}N$}}}\hglue 1.0pt\hfil}\hss}\hbox to0.0pt{\hbox to9.12497pt{\hfil\text@underline{\phantom{\hbox{$N$}}}\hglue 1.0pt\hfil}\hss}N, then the transformed noise matrix, RR R R NN N N RR R R RR R ๐
NN N ๐ RR R superscript๐
\smash{\hbox to0.0pt{\hbox to7.67015pt{\hfil\text@underline{\phantom{\hbox{$\hbox to0.0pt{\hbox to7.67015pt{\hfil\text@underline{\phantom{\hbox{$R$}}}\hglue 1.0pt\hfil}\hss}R$}}}\hglue 1.0pt\hfil}\hss}\hbox to0.0pt{\hbox to7.67015pt{\hfil\text@underline{\phantom{\hbox{$R$}}}\hglue 1.0pt\hfil}\hss}R\cdot\hbox to0.0pt{\hbox to9.12497pt{\hfil\text@underline{\phantom{\hbox{$\hbox to0.0pt{\hbox to9.12497pt{\hfil\text@underline{\phantom{\hbox{$N$}}}\hglue 1.0pt\hfil}\hss}N$}}}\hglue 1.0pt\hfil}\hss}\hbox to0.0pt{\hbox to9.12497pt{\hfil\text@underline{\phantom{\hbox{$N$}}}\hglue 1.0pt\hfil}\hss}N\cdot\hbox to0.0pt{\hbox to7.67015pt{\hfil\text@underline{\phantom{\hbox{$\hbox to0.0pt{\hbox to7.67015pt{\hfil\text@underline{\phantom{\hbox{$R$}}}\hglue 1.0pt\hfil}\hss}R$}}}\hglue 1.0pt\hfil}\hss}\hbox to0.0pt{\hbox to7.67015pt{\hfil\text@underline{\phantom{\hbox{$R$}}}\hglue 1.0pt\hfil}\hss}R^{\dagger}}, is diagonal. Nevertheless, if the transformed SS S S SS S ๐\hbox to0.0pt{\hbox to6.70831pt{\hfil\text@underline{\phantom{\hbox{$\hbox to0.0pt{\hbox to6.70831pt{\hfil\text@underline{\phantom{\hbox{$S$}}}\hglue 1.0pt\hfil}\hss}S$}}}\hglue 1.0pt\hfil}\hss}\hbox to0.0pt{\hbox to6.70831pt{\hfil\text@underline{\phantom{\hbox{$S$}}}\hglue 1.0pt\hfil}\hss}S is not diagonal, the two will not commute. This apparently insuperable problem can be solved by using the fact that the data vectors are entirely abstract at this stage. There is therefore no reason not to consider the further transformation of scaling the data, so that NN N N NN N ๐\hbox to0.0pt{\hbox to9.12497pt{\hfil\text@underline{\phantom{\hbox{$\hbox to0.0pt{\hbox to9.12497pt{\hfil\text@underline{\phantom{\hbox{$N$}}}\hglue 1.0pt\hfil}\hss}N$}}}\hglue 1.0pt\hfil}\hss}\hbox to0.0pt{\hbox to9.12497pt{\hfil\text@underline{\phantom{\hbox{$N$}}}\hglue 1.0pt\hfil}\hss}N becomes proportional to the identity matrix. This means that the transformation is no longer unitary โ but there is no physical reason to object to a change in the normalization of the data vectors.
Suppose we therefore make a further transformation
d dโฒโฒ= WW W W d d d superscript๐โฒโฒ WW W ๐ d superscript๐\displaystyle\hbox to0.0pt{\hbox to5.20486pt{\hfil\text@underline{\phantom{\hbox{$d$}}}\hglue 1.0pt\hfil}\hss}d^{\prime\prime}=\hbox to0.0pt{\hbox to10.83334pt{\hfil\text@underline{\phantom{\hbox{$\hbox to0.0pt{\hbox to10.83334pt{\hfil\text@underline{\phantom{\hbox{$W$}}}\hglue 1.0pt\hfil}\hss}W$}}}\hglue 1.0pt\hfil}\hss}\hbox to0.0pt{\hbox to10.83334pt{\hfil\text@underline{\phantom{\hbox{$W$}}}\hglue 1.0pt\hfil}\hss}W\cdot\hbox to0.0pt{\hbox to5.20486pt{\hfil\text@underline{\phantom{\hbox{$d$}}}\hglue 1.0pt\hfil}\hss}d^{\prime} (28)
The matrix WW W W WW W ๐\hbox to0.0pt{\hbox to10.83334pt{\hfil\text@underline{\phantom{\hbox{$\hbox to0.0pt{\hbox to10.83334pt{\hfil\text@underline{\phantom{\hbox{$W$}}}\hglue 1.0pt\hfil}\hss}W$}}}\hglue 1.0pt\hfil}\hss}\hbox to0.0pt{\hbox to10.83334pt{\hfil\text@underline{\phantom{\hbox{$W$}}}\hglue 1.0pt\hfil}\hss}W is related to the rotated noise matrix:
NN N N=diag(n1,n2,) WW W W=diag(1/n1,1/n2,).formulae-sequence NN N superscript๐diagsubscript๐1subscript๐2 WW W ๐diag1subscript๐11subscript๐2\displaystyle\hbox to0.0pt{\hbox to9.12497pt{\hfil\text@underline{\phantom{\hbox{$\hbox to0.0pt{\hbox to9.12497pt{\hfil\text@underline{\phantom{\hbox{$N$}}}\hglue 1.0pt\hfil}\hss}N$}}}\hglue 1.0pt\hfil}\hss}\hbox to0.0pt{\hbox to9.12497pt{\hfil\text@underline{\phantom{\hbox{$N$}}}\hglue 1.0pt\hfil}\hss}N^{\prime}={\rm diag}\,(n_{1},n_{2},\dots)\quad\Rightarrow\quad\hbox to0.0pt{\hbox to10.83334pt{\hfil\text@underline{\phantom{\hbox{$\hbox to0.0pt{\hbox to10.83334pt{\hfil\text@underline{\phantom{\hbox{$W$}}}\hglue 1.0pt\hfil}\hss}W$}}}\hglue 1.0pt\hfil}\hss}\hbox to0.0pt{\hbox to10.83334pt{\hfil\text@underline{\phantom{\hbox{$W$}}}\hglue 1.0pt\hfil}\hss}W={\rm diag}\,(1/\sqrt{n_{1}},1/\sqrt{n_{2}},\dots). (29)
This transformation is termed prewhitening by Vogeley & Szalay (1996), since it converts the noise matrix to white noise, in which each pixel has a unit noise that is uncorrelated with other pixels. The effect of this transformation on the full covariance matrix is
Cijโฒโฒdiโฒโฒdjโฒโฒ CC C Cโฒโฒ=( WW W W RR R R) CC C C( WW W W RR R R)\displaystyle C_{ij}^{\prime\prime}\equiv\langle d_{i}^{\prime\prime}d_{j}^{\prime\prime}{}^{*}\rangle\quad\Rightarrow\quad\hbox to0.0pt{\hbox to7.86249pt{\hfil\text@underline{\phantom{\hbox{$\hbox to0.0pt{\hbox to7.86249pt{\hfil\text@underline{\phantom{\hbox{$C$}}}\hglue 1.0pt\hfil}\hss}C$}}}\hglue 1.0pt\hfil}\hss}\hbox to0.0pt{\hbox to7.86249pt{\hfil\text@underline{\phantom{\hbox{$C$}}}\hglue 1.0pt\hfil}\hss}C^{\prime\prime}=(\hbox to0.0pt{\hbox to10.83334pt{\hfil\text@underline{\phantom{\hbox{$\hbox to0.0pt{\hbox to10.83334pt{\hfil\text@underline{\phantom{\hbox{$W$}}}\hglue 1.0pt\hfil}\hss}W$}}}\hglue 1.0pt\hfil}\hss}\hbox to0.0pt{\hbox to10.83334pt{\hfil\text@underline{\phantom{\hbox{$W$}}}\hglue 1.0pt\hfil}\hss}W\cdot\hbox to0.0pt{\hbox to7.67015pt{\hfil\text@underline{\phantom{\hbox{$\hbox to0.0pt{\hbox to7.67015pt{\hfil\text@underline{\phantom{\hbox{$R$}}}\hglue 1.0pt\hfil}\hss}R$}}}\hglue 1.0pt\hfil}\hss}\hbox to0.0pt{\hbox to7.67015pt{\hfil\text@underline{\phantom{\hbox{$R$}}}\hglue 1.0pt\hfil}\hss}R)\cdot\hbox to0.0pt{\hbox to7.86249pt{\hfil\text@underline{\phantom{\hbox{$\hbox to0.0pt{\hbox to7.86249pt{\hfil\text@underline{\phantom{\hbox{$C$}}}\hglue 1.0pt\hfil}\hss}C$}}}\hglue 1.0pt\hfil}\hss}\hbox to0.0pt{\hbox to7.86249pt{\hfil\text@underline{\phantom{\hbox{$C$}}}\hglue 1.0pt\hfil}\hss}C\cdot(\hbox to0.0pt{\hbox to10.83334pt{\hfil\text@underline{\phantom{\hbox{$\hbox to0.0pt{\hbox to10.83334pt{\hfil\text@underline{\phantom{\hbox{$W$}}}\hglue 1.0pt\hfil}\hss}W$}}}\hglue 1.0pt\hfil}\hss}\hbox to0.0pt{\hbox to10.83334pt{\hfil\text@underline{\phantom{\hbox{$W$}}}\hglue 1.0pt\hfil}\hss}W\cdot\hbox to0.0pt{\hbox to7.67015pt{\hfil\text@underline{\phantom{\hbox{$\hbox to0.0pt{\hbox to7.67015pt{\hfil\text@underline{\phantom{\hbox{$R$}}}\hglue 1.0pt\hfil}\hss}R$}}}\hglue 1.0pt\hfil}\hss}\hbox to0.0pt{\hbox to7.67015pt{\hfil\text@underline{\phantom{\hbox{$R$}}}\hglue 1.0pt\hfil}\hss}R)^{\dagger} (30)
After this transformation, the noise and signal matrices certainly do commute, and the optimal modes for expanding the new data are once again the PCA eigenmodes in the new coordinates:
CC C Cโฒโฒ ฯ ฯiโฒโฒ=ฮป ฯ ฯiโฒโฒ. CC C superscript๐ถโฒโฒ ฯ superscriptsubscript๐๐โฒโฒ๐ ฯ superscriptsubscript๐๐โฒโฒ\displaystyle\hbox to0.0pt{\hbox to7.86249pt{\hfil\text@underline{\phantom{\hbox{$\hbox to0.0pt{\hbox to7.86249pt{\hfil\text@underline{\phantom{\hbox{$C$}}}\hglue 1.0pt\hfil}\hss}C$}}}\hglue 1.0pt\hfil}\hss}\hbox to0.0pt{\hbox to7.86249pt{\hfil\text@underline{\phantom{\hbox{$C$}}}\hglue 1.0pt\hfil}\hss}C^{\prime\prime}\cdot\hbox to0.0pt{\hbox to6.51392pt{\hfil\text@underline{\phantom{\hbox{$\psi$}}}\hglue 1.0pt\hfil}\hss}\psi_{i}^{\prime\prime}=\lambda\hbox to0.0pt{\hbox to6.51392pt{\hfil\text@underline{\phantom{\hbox{$\psi$}}}\hglue 1.0pt\hfil}\hss}\psi_{i}^{\prime\prime}. (31)
These eigenmodes must be expressible in terms of some modes in the original coordinates, e ei e subscript๐๐\hbox to0.0pt{\hbox to4.65627pt{\hfil\text@underline{\phantom{\hbox{$e$}}}\hglue 1.0pt\hfil}\hss}e_{i}:
ฯ ฯiโฒโฒ=( WW W W RR R R) e ei. ฯ superscriptsubscript๐๐โฒโฒ WW W ๐ RR R ๐
e subscript๐๐\displaystyle\hbox to0.0pt{\hbox to6.51392pt{\hfil\text@underline{\phantom{\hbox{$\psi$}}}\hglue 1.0pt\hfil}\hss}\psi_{i}^{\prime\prime}=(\hbox to0.0pt{\hbox to10.83334pt{\hfil\text@underline{\phantom{\hbox{$\hbox to0.0pt{\hbox to10.83334pt{\hfil\text@underline{\phantom{\hbox{$W$}}}\hglue 1.0pt\hfil}\hss}W$}}}\hglue 1.0pt\hfil}\hss}\hbox to0.0pt{\hbox to10.83334pt{\hfil\text@underline{\phantom{\hbox{$W$}}}\hglue 1.0pt\hfil}\hss}W\cdot\hbox to0.0pt{\hbox to7.67015pt{\hfil\text@underline{\phantom{\hbox{$\hbox to0.0pt{\hbox to7.67015pt{\hfil\text@underline{\phantom{\hbox{$R$}}}\hglue 1.0pt\hfil}\hss}R$}}}\hglue 1.0pt\hfil}\hss}\hbox to0.0pt{\hbox to7.67015pt{\hfil\text@underline{\phantom{\hbox{$R$}}}\hglue 1.0pt\hfil}\hss}R)\cdot\hbox to0.0pt{\hbox to4.65627pt{\hfil\text@underline{\phantom{\hbox{$e$}}}\hglue 1.0pt\hfil}\hss}e_{i}. (32)
In these terms, the eigenproblem is
( WW W W RR R R) CC C C( WW W W RR R R)( WW W W RR R R) e ei=ฮป( WW W W RR R R) e ei. WW W ๐ RR R ๐
CC C ๐ถsuperscript WW W ๐ RR R ๐
WW W ๐ RR R ๐
e subscript๐๐๐ WW W ๐ RR R ๐
e subscript๐๐\displaystyle(\hbox to0.0pt{\hbox to10.83334pt{\hfil\text@underline{\phantom{\hbox{$\hbox to0.0pt{\hbox to10.83334pt{\hfil\text@underline{\phantom{\hbox{$W$}}}\hglue 1.0pt\hfil}\hss}W$}}}\hglue 1.0pt\hfil}\hss}\hbox to0.0pt{\hbox to10.83334pt{\hfil\text@underline{\phantom{\hbox{$W$}}}\hglue 1.0pt\hfil}\hss}W\cdot\hbox to0.0pt{\hbox to7.67015pt{\hfil\text@underline{\phantom{\hbox{$\hbox to0.0pt{\hbox to7.67015pt{\hfil\text@underline{\phantom{\hbox{$R$}}}\hglue 1.0pt\hfil}\hss}R$}}}\hglue 1.0pt\hfil}\hss}\hbox to0.0pt{\hbox to7.67015pt{\hfil\text@underline{\phantom{\hbox{$R$}}}\hglue 1.0pt\hfil}\hss}R)\cdot\hbox to0.0pt{\hbox to7.86249pt{\hfil\text@underline{\phantom{\hbox{$\hbox to0.0pt{\hbox to7.86249pt{\hfil\text@underline{\phantom{\hbox{$C$}}}\hglue 1.0pt\hfil}\hss}C$}}}\hglue 1.0pt\hfil}\hss}\hbox to0.0pt{\hbox to7.86249pt{\hfil\text@underline{\phantom{\hbox{$C$}}}\hglue 1.0pt\hfil}\hss}C\cdot(\hbox to0.0pt{\hbox to10.83334pt{\hfil\text@underline{\phantom{\hbox{$\hbox to0.0pt{\hbox to10.83334pt{\hfil\text@underline{\phantom{\hbox{$W$}}}\hglue 1.0pt\hfil}\hss}W$}}}\hglue 1.0pt\hfil}\hss}\hbox to0.0pt{\hbox to10.83334pt{\hfil\text@underline{\phantom{\hbox{$W$}}}\hglue 1.0pt\hfil}\hss}W\cdot\hbox to0.0pt{\hbox to7.67015pt{\hfil\text@underline{\phantom{\hbox{$\hbox to0.0pt{\hbox to7.67015pt{\hfil\text@underline{\phantom{\hbox{$R$}}}\hglue 1.0pt\hfil}\hss}R$}}}\hglue 1.0pt\hfil}\hss}\hbox to0.0pt{\hbox to7.67015pt{\hfil\text@underline{\phantom{\hbox{$R$}}}\hglue 1.0pt\hfil}\hss}R)^{\dagger}\cdot(\hbox to0.0pt{\hbox to10.83334pt{\hfil\text@underline{\phantom{\hbox{$\hbox to0.0pt{\hbox to10.83334pt{\hfil\text@underline{\phantom{\hbox{$W$}}}\hglue 1.0pt\hfil}\hss}W$}}}\hglue 1.0pt\hfil}\hss}\hbox to0.0pt{\hbox to10.83334pt{\hfil\text@underline{\phantom{\hbox{$W$}}}\hglue 1.0pt\hfil}\hss}W\cdot\hbox to0.0pt{\hbox to7.67015pt{\hfil\text@underline{\phantom{\hbox{$\hbox to0.0pt{\hbox to7.67015pt{\hfil\text@underline{\phantom{\hbox{$R$}}}\hglue 1.0pt\hfil}\hss}R$}}}\hglue 1.0pt\hfil}\hss}\hbox to0.0pt{\hbox to7.67015pt{\hfil\text@underline{\phantom{\hbox{$R$}}}\hglue 1.0pt\hfil}\hss}R)\cdot\hbox to0.0pt{\hbox to4.65627pt{\hfil\text@underline{\phantom{\hbox{$e$}}}\hglue 1.0pt\hfil}\hss}e_{i}=\lambda(\hbox to0.0pt{\hbox to10.83334pt{\hfil\text@underline{\phantom{\hbox{$\hbox to0.0pt{\hbox to10.83334pt{\hfil\text@underline{\phantom{\hbox{$W$}}}\hglue 1.0pt\hfil}\hss}W$}}}\hglue 1.0pt\hfil}\hss}\hbox to0.0pt{\hbox to10.83334pt{\hfil\text@underline{\phantom{\hbox{$W$}}}\hglue 1.0pt\hfil}\hss}W\cdot\hbox to0.0pt{\hbox to7.67015pt{\hfil\text@underline{\phantom{\hbox{$\hbox to0.0pt{\hbox to7.67015pt{\hfil\text@underline{\phantom{\hbox{$R$}}}\hglue 1.0pt\hfil}\hss}R$}}}\hglue 1.0pt\hfil}\hss}\hbox to0.0pt{\hbox to7.67015pt{\hfil\text@underline{\phantom{\hbox{$R$}}}\hglue 1.0pt\hfil}\hss}R)\cdot\hbox to0.0pt{\hbox to4.65627pt{\hfil\text@underline{\phantom{\hbox{$e$}}}\hglue 1.0pt\hfil}\hss}e_{i}. (33)
This can be simplified using WW W W WW W W= NN N N1\hbox to0.0pt{\hbox to10.83334pt{\hfil\text@underline{\phantom{\hbox{$\hbox to0.0pt{\hbox to10.83334pt{\hfil\text@underline{\phantom{\hbox{$W$}}}\hglue 1.0pt\hfil}\hss}W$}}}\hglue 1.0pt\hfil}\hss}\hbox to0.0pt{\hbox to10.83334pt{\hfil\text@underline{\phantom{\hbox{$W$}}}\hglue 1.0pt\hfil}\hss}W^{\dagger}\cdot\hbox to0.0pt{\hbox to10.83334pt{\hfil\text@underline{\phantom{\hbox{$\hbox to0.0pt{\hbox to10.83334pt{\hfil\text@underline{\phantom{\hbox{$W$}}}\hglue 1.0pt\hfil}\hss}W$}}}\hglue 1.0pt\hfil}\hss}\hbox to0.0pt{\hbox to10.83334pt{\hfil\text@underline{\phantom{\hbox{$W$}}}\hglue 1.0pt\hfil}\hss}W=\hbox to0.0pt{\hbox to9.12497pt{\hfil\text@underline{\phantom{\hbox{$\hbox to0.0pt{\hbox to9.12497pt{\hfil\text@underline{\phantom{\hbox{$N$}}}\hglue 1.0pt\hfil}\hss}N$}}}\hglue 1.0pt\hfil}\hss}\hbox to0.0pt{\hbox to9.12497pt{\hfil\text@underline{\phantom{\hbox{$N$}}}\hglue 1.0pt\hfil}\hss}N^{\prime}{}^{-1} and NN N N=1 RR R R NN N N1 RR R R\hbox to0.0pt{\hbox to9.12497pt{\hfil\text@underline{\phantom{\hbox{$\hbox to0.0pt{\hbox to9.12497pt{\hfil\text@underline{\phantom{\hbox{$N$}}}\hglue 1.0pt\hfil}\hss}N$}}}\hglue 1.0pt\hfil}\hss}\hbox to0.0pt{\hbox to9.12497pt{\hfil\text@underline{\phantom{\hbox{$N$}}}\hglue 1.0pt\hfil}\hss}N^{\prime}{}^{-1}=\hbox to0.0pt{\hbox to7.67015pt{\hfil\text@underline{\phantom{\hbox{$\hbox to0.0pt{\hbox to7.67015pt{\hfil\text@underline{\phantom{\hbox{$R$}}}\hglue 1.0pt\hfil}\hss}R$}}}\hglue 1.0pt\hfil}\hss}\hbox to0.0pt{\hbox to7.67015pt{\hfil\text@underline{\phantom{\hbox{$R$}}}\hglue 1.0pt\hfil}\hss}R\cdot\hbox to0.0pt{\hbox to9.12497pt{\hfil\text@underline{\phantom{\hbox{$\hbox to0.0pt{\hbox to9.12497pt{\hfil\text@underline{\phantom{\hbox{$N$}}}\hglue 1.0pt\hfil}\hss}N$}}}\hglue 1.0pt\hfil}\hss}\hbox to0.0pt{\hbox to9.12497pt{\hfil\text@underline{\phantom{\hbox{$N$}}}\hglue 1.0pt\hfil}\hss}N^{-1}\hbox to0.0pt{\hbox to7.67015pt{\hfil\text@underline{\phantom{\hbox{$\hbox to0.0pt{\hbox to7.67015pt{\hfil\text@underline{\phantom{\hbox{$R$}}}\hglue 1.0pt\hfil}\hss}R$}}}\hglue 1.0pt\hfil}\hss}\hbox to0.0pt{\hbox to7.67015pt{\hfil\text@underline{\phantom{\hbox{$R$}}}\hglue 1.0pt\hfil}\hss}R^{\dagger}, to give
CC C C NN N N1 e ei=ฮป e ei, CC C ๐ถ NN N superscript๐1 e subscript๐๐๐ e subscript๐๐\displaystyle\hbox to0.0pt{\hbox to7.86249pt{\hfil\text@underline{\phantom{\hbox{$\hbox to0.0pt{\hbox to7.86249pt{\hfil\text@underline{\phantom{\hbox{$C$}}}\hglue 1.0pt\hfil}\hss}C$}}}\hglue 1.0pt\hfil}\hss}\hbox to0.0pt{\hbox to7.86249pt{\hfil\text@underline{\phantom{\hbox{$C$}}}\hglue 1.0pt\hfil}\hss}C\cdot\hbox to0.0pt{\hbox to9.12497pt{\hfil\text@underline{\phantom{\hbox{$\hbox to0.0pt{\hbox to9.12497pt{\hfil\text@underline{\phantom{\hbox{$N$}}}\hglue 1.0pt\hfil}\hss}N$}}}\hglue 1.0pt\hfil}\hss}\hbox to0.0pt{\hbox to9.12497pt{\hfil\text@underline{\phantom{\hbox{$N$}}}\hglue 1.0pt\hfil}\hss}N^{-1}\cdot\hbox to0.0pt{\hbox to4.65627pt{\hfil\text@underline{\phantom{\hbox{$e$}}}\hglue 1.0pt\hfil}\hss}e_{i}=\lambda\hbox to0.0pt{\hbox to4.65627pt{\hfil\text@underline{\phantom{\hbox{$e$}}}\hglue 1.0pt\hfil}\hss}e_{i}, (34)
so the required modes are eigenmodes of CC C C NN N N1 CC C ๐ถ NN N superscript๐1\hbox to0.0pt{\hbox to7.86249pt{\hfil\text@underline{\phantom{\hbox{$\hbox to0.0pt{\hbox to7.86249pt{\hfil\text@underline{\phantom{\hbox{$C$}}}\hglue 1.0pt\hfil}\hss}C$}}}\hglue 1.0pt\hfil}\hss}\hbox to0.0pt{\hbox to7.86249pt{\hfil\text@underline{\phantom{\hbox{$C$}}}\hglue 1.0pt\hfil}\hss}C\cdot\hbox to0.0pt{\hbox to9.12497pt{\hfil\text@underline{\phantom{\hbox{$\hbox to0.0pt{\hbox to9.12497pt{\hfil\text@underline{\phantom{\hbox{$N$}}}\hglue 1.0pt\hfil}\hss}N$}}}\hglue 1.0pt\hfil}\hss}\hbox to0.0pt{\hbox to9.12497pt{\hfil\text@underline{\phantom{\hbox{$N$}}}\hglue 1.0pt\hfil}\hss}N^{-1}. However, care is required when considering the orthonormality of the e ei e subscript๐๐\hbox to0.0pt{\hbox to4.65627pt{\hfil\text@underline{\phantom{\hbox{$e$}}}\hglue 1.0pt\hfil}\hss}e_{i}: ฯ ฯi ฯ ฯj= e ei NN N N1 e ej ฯ superscriptsubscript๐๐ ฯ subscript๐๐ e superscriptsubscript๐๐ NN N superscript๐1 e subscript๐๐\smash{\hbox to0.0pt{\hbox to6.51392pt{\hfil\text@underline{\phantom{\hbox{$\psi$}}}\hglue 1.0pt\hfil}\hss}\psi_{i}^{\dagger}\cdot\hbox to0.0pt{\hbox to6.51392pt{\hfil\text@underline{\phantom{\hbox{$\psi$}}}\hglue 1.0pt\hfil}\hss}\psi_{j}=\hbox to0.0pt{\hbox to4.65627pt{\hfil\text@underline{\phantom{\hbox{$e$}}}\hglue 1.0pt\hfil}\hss}e_{i}^{\dagger}\cdot\hbox to0.0pt{\hbox to9.12497pt{\hfil\text@underline{\phantom{\hbox{$\hbox to0.0pt{\hbox to9.12497pt{\hfil\text@underline{\phantom{\hbox{$N$}}}\hglue 1.0pt\hfil}\hss}N$}}}\hglue 1.0pt\hfil}\hss}\hbox to0.0pt{\hbox to9.12497pt{\hfil\text@underline{\phantom{\hbox{$N$}}}\hglue 1.0pt\hfil}\hss}N^{-1}\cdot\hbox to0.0pt{\hbox to4.65627pt{\hfil\text@underline{\phantom{\hbox{$e$}}}\hglue 1.0pt\hfil}\hss}e_{j}}, so the e ei e subscript๐๐\hbox to0.0pt{\hbox to4.65627pt{\hfil\text@underline{\phantom{\hbox{$e$}}}\hglue 1.0pt\hfil}\hss}e_{i} are not orthonormal. If we write d d=iai e ei d ๐subscript๐subscript๐๐ e subscript๐๐\hbox to0.0pt{\hbox to5.20486pt{\hfil\text@underline{\phantom{\hbox{$d$}}}\hglue 1.0pt\hfil}\hss}d=\sum_{i}a_{i}\hbox to0.0pt{\hbox to4.65627pt{\hfil\text@underline{\phantom{\hbox{$e$}}}\hglue 1.0pt\hfil}\hss}e_{i}, then
ai=( NN N N1 e ei) d d ฯ ฯi d d.subscript๐๐superscript NN N superscript๐1 e subscript๐๐ d ๐ ฯ superscriptsubscript๐๐ d ๐\displaystyle a_{i}=(\hbox to0.0pt{\hbox to9.12497pt{\hfil\text@underline{\phantom{\hbox{$\hbox to0.0pt{\hbox to9.12497pt{\hfil\text@underline{\phantom{\hbox{$N$}}}\hglue 1.0pt\hfil}\hss}N$}}}\hglue 1.0pt\hfil}\hss}\hbox to0.0pt{\hbox to9.12497pt{\hfil\text@underline{\phantom{\hbox{$N$}}}\hglue 1.0pt\hfil}\hss}N^{-1}\cdot\hbox to0.0pt{\hbox to4.65627pt{\hfil\text@underline{\phantom{\hbox{$e$}}}\hglue 1.0pt\hfil}\hss}e_{i})^{\dagger}\cdot\hbox to0.0pt{\hbox to5.20486pt{\hfil\text@underline{\phantom{\hbox{$d$}}}\hglue 1.0pt\hfil}\hss}d\equiv\hbox to0.0pt{\hbox to6.51392pt{\hfil\text@underline{\phantom{\hbox{$\psi$}}}\hglue 1.0pt\hfil}\hss}\psi_{i}^{\dagger}\cdot\hbox to0.0pt{\hbox to5.20486pt{\hfil\text@underline{\phantom{\hbox{$d$}}}\hglue 1.0pt\hfil}\hss}d. (35)
Thus, the modes used to extract the compressed data by dot product satisfy CC C C ฯ ฯ=ฮป NN N N ฯ ฯ CC C ๐ถ ฯ ๐๐ NN N ๐ ฯ ๐\hbox to0.0pt{\hbox to7.86249pt{\hfil\text@underline{\phantom{\hbox{$\hbox to0.0pt{\hbox to7.86249pt{\hfil\text@underline{\phantom{\hbox{$C$}}}\hglue 1.0pt\hfil}\hss}C$}}}\hglue 1.0pt\hfil}\hss}\hbox to0.0pt{\hbox to7.86249pt{\hfil\text@underline{\phantom{\hbox{$C$}}}\hglue 1.0pt\hfil}\hss}C\cdot\hbox to0.0pt{\hbox to6.51392pt{\hfil\text@underline{\phantom{\hbox{$\psi$}}}\hglue 1.0pt\hfil}\hss}\psi=\lambda\hbox to0.0pt{\hbox to9.12497pt{\hfil\text@underline{\phantom{\hbox{$\hbox to0.0pt{\hbox to9.12497pt{\hfil\text@underline{\phantom{\hbox{$N$}}}\hglue 1.0pt\hfil}\hss}N$}}}\hglue 1.0pt\hfil}\hss}\hbox to0.0pt{\hbox to9.12497pt{\hfil\text@underline{\phantom{\hbox{$N$}}}\hglue 1.0pt\hfil}\hss}N\cdot\hbox to0.0pt{\hbox to6.51392pt{\hfil\text@underline{\phantom{\hbox{$\psi$}}}\hglue 1.0pt\hfil}\hss}\psi, or finally
SS S S ฯ ฯ=ฮป NN N N ฯ ฯ, SS S ๐ ฯ ๐๐ NN N ๐ ฯ ๐\displaystyle\hbox to0.0pt{\hbox to6.70831pt{\hfil\text@underline{\phantom{\hbox{$\hbox to0.0pt{\hbox to6.70831pt{\hfil\text@underline{\phantom{\hbox{$S$}}}\hglue 1.0pt\hfil}\hss}S$}}}\hglue 1.0pt\hfil}\hss}\hbox to0.0pt{\hbox to6.70831pt{\hfil\text@underline{\phantom{\hbox{$S$}}}\hglue 1.0pt\hfil}\hss}S\cdot\hbox to0.0pt{\hbox to6.51392pt{\hfil\text@underline{\phantom{\hbox{$\psi$}}}\hglue 1.0pt\hfil}\hss}\psi=\lambda\,\hbox to0.0pt{\hbox to9.12497pt{\hfil\text@underline{\phantom{\hbox{$\hbox to0.0pt{\hbox to9.12497pt{\hfil\text@underline{\phantom{\hbox{$N$}}}\hglue 1.0pt\hfil}\hss}N$}}}\hglue 1.0pt\hfil}\hss}\hbox to0.0pt{\hbox to9.12497pt{\hfil\text@underline{\phantom{\hbox{$N$}}}\hglue 1.0pt\hfil}\hss}N\cdot\hbox to0.0pt{\hbox to6.51392pt{\hfil\text@underline{\phantom{\hbox{$\psi$}}}\hglue 1.0pt\hfil}\hss}\psi, (36)
given a redefinition of $`\lambda `$. The optimal modes are thus eigenmodes of NN N N1 SS S S NN N superscript๐1 SS S ๐\hbox to0.0pt{\hbox to9.12497pt{\hfil\text@underline{\phantom{\hbox{$\hbox to0.0pt{\hbox to9.12497pt{\hfil\text@underline{\phantom{\hbox{$N$}}}\hglue 1.0pt\hfil}\hss}N$}}}\hglue 1.0pt\hfil}\hss}\hbox to0.0pt{\hbox to9.12497pt{\hfil\text@underline{\phantom{\hbox{$N$}}}\hglue 1.0pt\hfil}\hss}N^{-1}\cdot\hbox to0.0pt{\hbox to6.70831pt{\hfil\text@underline{\phantom{\hbox{$\hbox to0.0pt{\hbox to6.70831pt{\hfil\text@underline{\phantom{\hbox{$S$}}}\hglue 1.0pt\hfil}\hss}S$}}}\hglue 1.0pt\hfil}\hss}\hbox to0.0pt{\hbox to6.70831pt{\hfil\text@underline{\phantom{\hbox{$S$}}}\hglue 1.0pt\hfil}\hss}S, hence the name signal-to-noise eigenmodes (Bond 1995; Bunn 1996).
It is interesting to appreciate that the set of KL modes just discussed is also the โbestโ set of modes to choose from a completely different point of view: they are the modes that are optimal for estimation of a parameter via maximum likelihood. Suppose we write the compressed data vector, x x x ๐ฅ\hbox to0.0pt{\hbox to5.71527pt{\hfil\text@underline{\phantom{\hbox{$x$}}}\hglue 1.0pt\hfil}\hss}x, in terms of a non-square matrix AA A A AA A ๐ด\hbox to0.0pt{\hbox to7.50002pt{\hfil\text@underline{\phantom{\hbox{$\hbox to0.0pt{\hbox to7.50002pt{\hfil\text@underline{\phantom{\hbox{$A$}}}\hglue 1.0pt\hfil}\hss}A$}}}\hglue 1.0pt\hfil}\hss}\hbox to0.0pt{\hbox to7.50002pt{\hfil\text@underline{\phantom{\hbox{$A$}}}\hglue 1.0pt\hfil}\hss}A (whose rows are the basis vectors ฯ ฯi ฯ superscriptsubscript๐๐\hbox to0.0pt{\hbox to6.51392pt{\hfil\text@underline{\phantom{\hbox{$\psi$}}}\hglue 1.0pt\hfil}\hss}\psi_{i}^{*}):
x x= AA A A d d. x ๐ฅ AA A ๐ด d ๐\displaystyle\hbox to0.0pt{\hbox to5.71527pt{\hfil\text@underline{\phantom{\hbox{$x$}}}\hglue 1.0pt\hfil}\hss}x=\hbox to0.0pt{\hbox to7.50002pt{\hfil\text@underline{\phantom{\hbox{$\hbox to0.0pt{\hbox to7.50002pt{\hfil\text@underline{\phantom{\hbox{$A$}}}\hglue 1.0pt\hfil}\hss}A$}}}\hglue 1.0pt\hfil}\hss}\hbox to0.0pt{\hbox to7.50002pt{\hfil\text@underline{\phantom{\hbox{$A$}}}\hglue 1.0pt\hfil}\hss}A\cdot\hbox to0.0pt{\hbox to5.20486pt{\hfil\text@underline{\phantom{\hbox{$d$}}}\hglue 1.0pt\hfil}\hss}d. (37)
The transformed covariance matrix is
DD D D x x x x= AA A A CC C C AA A A. DD D ๐ทdelimited-โจโฉ x ๐ฅ x superscript๐ฅ AA A ๐ด CC C ๐ถ AA A superscript๐ด\displaystyle\hbox to0.0pt{\hbox to8.55695pt{\hfil\text@underline{\phantom{\hbox{$\hbox to0.0pt{\hbox to8.55695pt{\hfil\text@underline{\phantom{\hbox{$D$}}}\hglue 1.0pt\hfil}\hss}D$}}}\hglue 1.0pt\hfil}\hss}\hbox to0.0pt{\hbox to8.55695pt{\hfil\text@underline{\phantom{\hbox{$D$}}}\hglue 1.0pt\hfil}\hss}D\equiv\langle\hbox to0.0pt{\hbox to5.71527pt{\hfil\text@underline{\phantom{\hbox{$x$}}}\hglue 1.0pt\hfil}\hss}x\hbox to0.0pt{\hbox to5.71527pt{\hfil\text@underline{\phantom{\hbox{$x$}}}\hglue 1.0pt\hfil}\hss}x^{\dagger}\rangle=\hbox to0.0pt{\hbox to7.50002pt{\hfil\text@underline{\phantom{\hbox{$\hbox to0.0pt{\hbox to7.50002pt{\hfil\text@underline{\phantom{\hbox{$A$}}}\hglue 1.0pt\hfil}\hss}A$}}}\hglue 1.0pt\hfil}\hss}\hbox to0.0pt{\hbox to7.50002pt{\hfil\text@underline{\phantom{\hbox{$A$}}}\hglue 1.0pt\hfil}\hss}A\cdot\hbox to0.0pt{\hbox to7.86249pt{\hfil\text@underline{\phantom{\hbox{$\hbox to0.0pt{\hbox to7.86249pt{\hfil\text@underline{\phantom{\hbox{$C$}}}\hglue 1.0pt\hfil}\hss}C$}}}\hglue 1.0pt\hfil}\hss}\hbox to0.0pt{\hbox to7.86249pt{\hfil\text@underline{\phantom{\hbox{$C$}}}\hglue 1.0pt\hfil}\hss}C\cdot\hbox to0.0pt{\hbox to7.50002pt{\hfil\text@underline{\phantom{\hbox{$\hbox to0.0pt{\hbox to7.50002pt{\hfil\text@underline{\phantom{\hbox{$A$}}}\hglue 1.0pt\hfil}\hss}A$}}}\hglue 1.0pt\hfil}\hss}\hbox to0.0pt{\hbox to7.50002pt{\hfil\text@underline{\phantom{\hbox{$A$}}}\hglue 1.0pt\hfil}\hss}A^{\dagger}. (38)
For the case where the original data obeyed Gaussian statistics, this is true for the compressed data also, so the likelihood is
2ln=lndet DD D D+ x x DD D D1 x x+constant2det DD D ๐ท x superscript๐ฅ DD D superscript๐ท1 x ๐ฅconstant\displaystyle-2\ln{\cal L}=\ln{\rm det}\,\hbox to0.0pt{\hbox to8.55695pt{\hfil\text@underline{\phantom{\hbox{$\hbox to0.0pt{\hbox to8.55695pt{\hfil\text@underline{\phantom{\hbox{$D$}}}\hglue 1.0pt\hfil}\hss}D$}}}\hglue 1.0pt\hfil}\hss}\hbox to0.0pt{\hbox to8.55695pt{\hfil\text@underline{\phantom{\hbox{$D$}}}\hglue 1.0pt\hfil}\hss}D+\hbox to0.0pt{\hbox to5.71527pt{\hfil\text@underline{\phantom{\hbox{$x$}}}\hglue 1.0pt\hfil}\hss}x^{*}\cdot\hbox to0.0pt{\hbox to8.55695pt{\hfil\text@underline{\phantom{\hbox{$\hbox to0.0pt{\hbox to8.55695pt{\hfil\text@underline{\phantom{\hbox{$D$}}}\hglue 1.0pt\hfil}\hss}D$}}}\hglue 1.0pt\hfil}\hss}\hbox to0.0pt{\hbox to8.55695pt{\hfil\text@underline{\phantom{\hbox{$D$}}}\hglue 1.0pt\hfil}\hss}D^{-1}\cdot\hbox to0.0pt{\hbox to5.71527pt{\hfil\text@underline{\phantom{\hbox{$x$}}}\hglue 1.0pt\hfil}\hss}x+{\rm constant} (39)
The normal variance on some parameter $`p`$ (on which the covariance matrix depends) is
$`{\displaystyle \frac{1}{\sigma _p^2}}={\displaystyle \frac{d^2[2\mathrm{ln}]}{dq^2}}.`$ (40)
Without data, we donโt know this, so it is common to use the expectation value of the rhs as an estimate (recently, there has been a tendency to dub this the โFisher matrixโ).
We desire to optimize $`\sigma _p`$ by an appropriate choice of data-compression vectors, ฯ ฯi ฯ subscript๐๐\hbox to0.0pt{\hbox to6.51392pt{\hfil\text@underline{\phantom{\hbox{$\psi$}}}\hglue 1.0pt\hfil}\hss}\psi_{i}. By writing $`\sigma _p`$ in terms of AA A A AA A ๐ด\hbox to0.0pt{\hbox to7.50002pt{\hfil\text@underline{\phantom{\hbox{$\hbox to0.0pt{\hbox to7.50002pt{\hfil\text@underline{\phantom{\hbox{$A$}}}\hglue 1.0pt\hfil}\hss}A$}}}\hglue 1.0pt\hfil}\hss}\hbox to0.0pt{\hbox to7.50002pt{\hfil\text@underline{\phantom{\hbox{$A$}}}\hglue 1.0pt\hfil}\hss}A, CC C C CC C ๐ถ\hbox to0.0pt{\hbox to7.86249pt{\hfil\text@underline{\phantom{\hbox{$\hbox to0.0pt{\hbox to7.86249pt{\hfil\text@underline{\phantom{\hbox{$C$}}}\hglue 1.0pt\hfil}\hss}C$}}}\hglue 1.0pt\hfil}\hss}\hbox to0.0pt{\hbox to7.86249pt{\hfil\text@underline{\phantom{\hbox{$C$}}}\hglue 1.0pt\hfil}\hss}C and d d d ๐\hbox to0.0pt{\hbox to5.20486pt{\hfil\text@underline{\phantom{\hbox{$d$}}}\hglue 1.0pt\hfil}\hss}d, it may eventually be shown that the desired optimal modes satisfy
(ddp CC C C) ฯ ฯ=ฮป CC C C ฯ ฯ.๐๐๐ CC C ๐ถ ฯ ๐๐ CC C ๐ถ ฯ ๐\displaystyle\left({d\over dp}\,\hbox to0.0pt{\hbox to7.86249pt{\hfil\text@underline{\phantom{\hbox{$\hbox to0.0pt{\hbox to7.86249pt{\hfil\text@underline{\phantom{\hbox{$C$}}}\hglue 1.0pt\hfil}\hss}C$}}}\hglue 1.0pt\hfil}\hss}\hbox to0.0pt{\hbox to7.86249pt{\hfil\text@underline{\phantom{\hbox{$C$}}}\hglue 1.0pt\hfil}\hss}C\right)\cdot\hbox to0.0pt{\hbox to6.51392pt{\hfil\text@underline{\phantom{\hbox{$\psi$}}}\hglue 1.0pt\hfil}\hss}\psi=\lambda\,\hbox to0.0pt{\hbox to7.86249pt{\hfil\text@underline{\phantom{\hbox{$\hbox to0.0pt{\hbox to7.86249pt{\hfil\text@underline{\phantom{\hbox{$C$}}}\hglue 1.0pt\hfil}\hss}C$}}}\hglue 1.0pt\hfil}\hss}\hbox to0.0pt{\hbox to7.86249pt{\hfil\text@underline{\phantom{\hbox{$C$}}}\hglue 1.0pt\hfil}\hss}C\cdot\hbox to0.0pt{\hbox to6.51392pt{\hfil\text@underline{\phantom{\hbox{$\psi$}}}\hglue 1.0pt\hfil}\hss}\psi. (41)
For the case where the parameter of interest is the cosmological power, the matrix on the lhs is just proportional to SS S S SS S ๐\hbox to0.0pt{\hbox to6.70831pt{\hfil\text@underline{\phantom{\hbox{$\hbox to0.0pt{\hbox to6.70831pt{\hfil\text@underline{\phantom{\hbox{$S$}}}\hglue 1.0pt\hfil}\hss}S$}}}\hglue 1.0pt\hfil}\hss}\hbox to0.0pt{\hbox to6.70831pt{\hfil\text@underline{\phantom{\hbox{$S$}}}\hglue 1.0pt\hfil}\hss}S, so we have to solve the eigenproblem
SS S S ฯ ฯ=ฮป CC C C ฯ ฯ. SS S ๐ ฯ ๐๐ CC C ๐ถ ฯ ๐\displaystyle\hbox to0.0pt{\hbox to6.70831pt{\hfil\text@underline{\phantom{\hbox{$\hbox to0.0pt{\hbox to6.70831pt{\hfil\text@underline{\phantom{\hbox{$S$}}}\hglue 1.0pt\hfil}\hss}S$}}}\hglue 1.0pt\hfil}\hss}\hbox to0.0pt{\hbox to6.70831pt{\hfil\text@underline{\phantom{\hbox{$S$}}}\hglue 1.0pt\hfil}\hss}S\cdot\hbox to0.0pt{\hbox to6.51392pt{\hfil\text@underline{\phantom{\hbox{$\psi$}}}\hglue 1.0pt\hfil}\hss}\psi=\lambda\,\hbox to0.0pt{\hbox to7.86249pt{\hfil\text@underline{\phantom{\hbox{$\hbox to0.0pt{\hbox to7.86249pt{\hfil\text@underline{\phantom{\hbox{$C$}}}\hglue 1.0pt\hfil}\hss}C$}}}\hglue 1.0pt\hfil}\hss}\hbox to0.0pt{\hbox to7.86249pt{\hfil\text@underline{\phantom{\hbox{$C$}}}\hglue 1.0pt\hfil}\hss}C\cdot\hbox to0.0pt{\hbox to6.51392pt{\hfil\text@underline{\phantom{\hbox{$\psi$}}}\hglue 1.0pt\hfil}\hss}\psi. (42)
With a redefinition of $`\lambda `$, this becomes
SS S S ฯ ฯ=ฮป NN N N ฯ ฯ. SS S ๐ ฯ ๐๐ NN N ๐ ฯ ๐\displaystyle\hbox to0.0pt{\hbox to6.70831pt{\hfil\text@underline{\phantom{\hbox{$\hbox to0.0pt{\hbox to6.70831pt{\hfil\text@underline{\phantom{\hbox{$S$}}}\hglue 1.0pt\hfil}\hss}S$}}}\hglue 1.0pt\hfil}\hss}\hbox to0.0pt{\hbox to6.70831pt{\hfil\text@underline{\phantom{\hbox{$S$}}}\hglue 1.0pt\hfil}\hss}S\cdot\hbox to0.0pt{\hbox to6.51392pt{\hfil\text@underline{\phantom{\hbox{$\psi$}}}\hglue 1.0pt\hfil}\hss}\psi=\lambda\,\hbox to0.0pt{\hbox to9.12497pt{\hfil\text@underline{\phantom{\hbox{$\hbox to0.0pt{\hbox to9.12497pt{\hfil\text@underline{\phantom{\hbox{$N$}}}\hglue 1.0pt\hfil}\hss}N$}}}\hglue 1.0pt\hfil}\hss}\hbox to0.0pt{\hbox to9.12497pt{\hfil\text@underline{\phantom{\hbox{$N$}}}\hglue 1.0pt\hfil}\hss}N\cdot\hbox to0.0pt{\hbox to6.51392pt{\hfil\text@underline{\phantom{\hbox{$\psi$}}}\hglue 1.0pt\hfil}\hss}\psi. (43)
The optimal modes for parameter estimation in the linear case are thus identical to the PCA modes of the prewhitened data discussed above. The more general expression was given by Tegmark, Taylor & Heavens (1997), and it is only in this case, where the covariance matrix is not necessarily linear in the parameter of interest, that the KL method actually differs from PCA.
The reason for going to all this trouble is that the likelihood can now be evaluated much more rapidly, using the compressed data. This allows extensive model searches over large parameter spaces that would be unfeasible with the original data (since inversion of an $`N\times N`$ covariance matrix takes a time proportional to $`N^3`$). Note however that the price paid for this efficiency is that a different set of modes need to be chosen depending on the model of interest, and that these modes will not in general be optimal for expanding the dataset itself. Nevertheless, it may be expected that application of these methods will inevitably grow as datasets increase in size. Present applications mainly prove that the techniques work: see Matsubara, Szalay & Landy (1999) for application to the LCRS, or Padmanabhan, Tegmark & Hamilton (1999) for the UZC survey. The next generation of experiments will probably be forced to resort to data compression of this sort, rather than using it as an elegant alternative method of analysis.
## 4 Redshift-space effects
Peculiar velocity fields are responsible for the distortion of the clustering pattern in redshift space, as first clearly articulated by Kaiser (1987). For a survey that subtends a small angle (i.e. in the distant-observer approximation), a good approximation to the anisotropic redshift-space Fourier spectrum is given by the Kaiser function together with a damping term from nonlinear effects:
$`\delta _k^s=\delta _k^r(1+\beta \mu ^2)D(k\sigma \mu ),`$ (44)
where $`\beta =\mathrm{\Omega }_m^{0.6}/b`$, $`b`$ being the linear bias parameter of the galaxies under study, and $`\mu =\widehat{๐ค}\widehat{๐ซ}`$. For an exponential distribution of relative small-scale peculiar velocities (as seen empirically), the damping function is $`D(y)(1+y^2/2)^{1/2}`$, and $`\sigma 400\mathrm{km}\mathrm{s}^1`$ is a reasonable estimate for the pairwise velocity dispersion of galaxies (e.g. Ballinger, Peacock & Heavens 1996).
In principle, this distortion should be a robust way to determine $`\mathrm{\Omega }`$ (or at least $`\beta `$). In practice, the effect has not been easy to see with past datasets. This is mainly a question of depth: a large survey is needed in order to beat down the shot noise, but this tends to favour bright spectroscopic limits. This limits the result both because relatively few modes in the linear regime are sampled, and also because local survey volumes will tend to violate the small-angle approximation. Strauss & Willick (1995) and Hamilton (1997a) review the practical application of redshift-space distortions. In the next section, preliminary results are presented from the 2dF redshift survey, which shows the distortion effect clearly for the first time.
## 5 The state of the art in LSS
### 5.1 The APM survey
In the past few years, much attention has been attracted by the estimate of the galaxy power spectrum from the APM survey (Baugh & Efstathiou 1993, 1994; Maddox et al. 1996). The APM result was generated from a catalogue of $`10^6`$ galaxies derived from UK Schmidt Telescope photographic plates scanned with the Cambridge Automatic Plate Measuring machine; because it is based on a deprojection of angular clustering, it is immune to the complicating effects of redshift-space distortions. The difficulty, of course, is in ensuring that any low-level systematics from e.g. spatial variations in magnitude zero point are sufficiently well controlled that they do not mask the cosmological signal, which is of order $`w(\theta )\stackrel{<}{}0.01`$ at separations of a few degrees.
The best evidence that the APM survey has the desired uniformity is the scaling test, where the correlations in fainter magnitude slices are expected to move to smaller scales and be reduced in amplitude. If we increase the depth of the survey by some factor $`D`$, the new angular correlation function will be
$`w^{}(\theta )={\displaystyle \frac{1}{D}}w(D\theta ).`$ (45)
The APM survey passes this test well; once the overall redshift distribution is known, it is possible to obtain the spatial power spectrum by inverting a convolution integral:
$`w(\theta )={\displaystyle _0^{\mathrm{}}}y^4\varphi ^2๐y{\displaystyle _0^{\mathrm{}}}\pi \mathrm{\Delta }^2(k)J_0(ky\theta )๐k/k^2`$ (46)
(where zero spatial curvature is assumed). Here, $`\varphi (y)`$ is the comoving density at comoving distance $`y`$, normalized so that $`y^2\varphi (y)๐y=1`$.
This integral was inverted numerically by Baugh & Efstathiou (1993), and gives an impressively accurate determination of the power spectrum. The error estimates are derived empirically from the scatter between independent regions of the sky, and so should be realistic. If there are no undetected systematics, these error bars say that the power is very accurately determined. The APM result has been investigated in detail by a number of authors (e.g. Gaztaรฑaga & Baugh 1998; Eisenstein & Zaldarriaga 1999) and found to be robust; this has significant implications if true.
### 5.2 Past redshift surveys
Because of the sheer number of galaxies, plus the large volume surveyed, the APM survey outperforms redshift surveys of the past, at least for the purpose of determining the power spectrum. The largest surveys of recent years (CfA: Huchra et al. 1990; LCRS: Shectman et al. 1996; PSCz: Saunders et al. 1999) contain of order $`10^4`$ galaxy redshifts, and their statistical errors are considerably larger than those of the APM. On the other hand, it is of great importance to compare the results of deprojection with clustering measured directly in 3D.
This comparison was carried out by Peacock & Dodds (1994; PD94). The exercise is not straightforward, because the 3D results are affected by redshift-space distortions; also, different galaxy tracers can be biased to different extents. The approach taken was to use each dataset to reconstruct an estimate of the linear spectrum, allowing the relative bias factors to float in order to make these estimates agree as well as possible (figure 5). To within a scatter of perhaps a factor 1.5 in power, the results were consistent with a $`\mathrm{\Gamma }0.25`$ CDM model. Even though the subsequent sections will discuss some possible disagreements with the CDM models at a higher level of precision, the general existence of CDM-like curvature in the spectrum is likely to be an important clue to the nature of the dark matter.
### 5.3 The 2dF survey
The proper resolution of many of the observational questions regarding the large-scale distribution of galaxies requires new generations of redshift survey that push beyond the $`N=10^5`$ barrier. Two groups are pursuing this goal. The Sloan survey (e.g. Margon 1999) is using a dedicated 2.5-m telescope to measure redshifts for approximately 700,000 galaxies to $`r=18.2`$ in the North Galactic Cap. The 2dF survey (e.g. Colless 1999) is using a fraction of the time on the 3.9-m Anglo-Australian Telescope plus Two-Degree Field spectrograph to measure 250,000 galaxies from the APM survey to $`B_J=19.45`$ in the South Galactic Cap. At the time of writing, the Sloan spectroscopic survey has yet to commence. However, the 2dF project has measured 77,000 redshifts, and some preliminary clustering results are given below. For more details of the survey, particularly the team members whose hard work has made all this possible, see http://www.mso.anu.edu.au/2dFGRS/.
One of the advantages of 2dF is that it is a fully sampled survey, so that the space density out to the depth imposed by the magnitude limit (median $`z=0.12`$) is as high as nature allows: apart from a tail of low surface brightness galaxies (inevitably omitted from any spectroscopic survey), the 2dF measure all the galaxies that exist over a cosmologically representative volume. It is the first to achieve this goal. The fidelity of the resulting map of the galaxy distribution can be seen in figure 6, which shows a small subset of the data: a slice of thickness 4 degrees, centred at declination $`27^{}`$.
An issue with using the 2dF data in their current form is that the sky has to be divided into circular โtilesโ each two degrees in diameter (โ2dFโ = โtwo-degree fieldโ, within which the AAT is able to measure 400 spectra simultaneously; see http://www.aao.gov.au/2df/ for details of the instrument). The tiles are positioned adaptively, so that larger overlaps occur in regions of high galaxy density. It this way, it is possible to place a fibre on $`>95\%`$ of all galaxies. However, while the survey is in progress, there exist parts of the sky where the overlapping tiles have not yet been observed, and so the effective sampling fraction is only $`50\%`$. These effects can be allowed for in two different ways. In clustering analyses, we compare the counts of pairs (or $`n`$-tuplets) of galaxies in the data to the corresponding counts involving an unclustered random catalogue. The effects of variable sampling can therefore be dealt with either by making the density of random points fluctuate according to the sampling, or by weighting observed galaxies by the reciprocal of the sampling factor for the zone in which they lie. The former approach is better from the point of view of shot noise, but the latter may be safer if there is any suspicion that the sampling fluctuations are correlated with real structure on the sky. In practice, both strategies give identical answers for the results below.
At the two-point level, the most direct quantity to compute is the redshift-space correlation function. This is an anisotropic function of the orientation of a galaxy pair, owing to peculiar velocities. We therefore evaluate $`\xi `$ as a function of 2D separation in terms of coordinates both parallel and perpendicular to the line of sight. If the comoving radii of two galaxies are $`y_1`$ and $`y_2`$ and their total separation is $`r`$, then we define coordinates
$`\pi |y_1y_2|;\sigma =\sqrt{r^2\pi ^2}.`$ (47)
The correlation function measured in these coordinates is shown in figure 7. In evaluating $`\xi (\sigma ,\pi )`$, the optimal radial weight discussed above has been applied, so that the noise at large $`r`$ should be representative of true cosmic scatter.
The correlation-function results display very clearly the two signatures of redshift-space distortions discussed above. The fingers of God from small-scale random velocities are very clear, as indeed has been the case from the first redshift surveys (e.g. Davis & Peebles 1983). However, this is arguably the first time that the large-scale flattening from coherent infall has been really obvious in the data.
A good way to quantify the flattening is to analyze the clustering as a function of angle into Legendre polynomials:
$`\xi _{\mathrm{}}(r)={\displaystyle \frac{2\mathrm{}+1}{2}}{\displaystyle _1^1}\xi (\sigma =r\mathrm{sin}\theta ,\pi =r\mathrm{cos}\theta )P_{\mathrm{}}(\mathrm{cos}\theta )d\mathrm{cos}\theta .`$ (48)
The quadrupole-to-monopole ratio should be a clear indicator of coherent infall. In linear theory, it is given by
$`{\displaystyle \frac{\xi _2}{\xi _0}}=f(n){\displaystyle \frac{4\beta /3+4\beta ^2/7}{1+2\beta /3+\beta ^2/5}},`$ (49)
where $`f(n)=(3+n)/n`$ (Hamilton 1992). On small and intermediate scales, the effective spectral index is negative, so the quadrupole-to-monopole ratio should be negative, as observed.
However, it is clear that the results on the largest scales are still significantly affected by finger-of-God smearing. The best way to interpret the observed effects is to calculate the same quantities for a model. To achieve this, we use the observed APM 3D power spectrum, plus the distortion model discussed above. This gives the plots shown in figure 8. The free parameter is $`\beta `$, and this is set at a value of 0.5, approximately consistent with other arguments for a universe with $`\mathrm{\Omega }=0.3`$ and little large-scale bias (e.g. Peacock 1997). Although a quantitative comparison has not yet been carried out, it is clear that this plot closely resembles the observed data.
By the end of 2001, the size of the 2dF survey should have expanded by a factor 3, increasing the pair counts tenfold. It should then be possible to trace the correlations well beyond the present limit, and follow the redshift-space distortion well into the linear regime. However, the biggest advantage of a survey of this size and uniformity is the ability to subdivide it. All analyses to date have lumped together very different kinds of galaxies, whereas we know from morphological segregation that different classes of galaxy have spatial distributions that differ from each other. The homogeneous 2dF data allow classification into different galaxy types (representing, physically, a sequence of star-formation rates), from the spectra alone (Folkes et al. 1999). It will be a critical test to see if the distortion signature can be picked up in each type individually. Although the large-scale behaviour of each galaxy type will probably be quite similar, differences in the clustering properties will inevitably arise on smaller scales, giving important information about the sequence of galaxy formation.
## 6 Small-scale clustering
### 6.1 History
One of the earliest models to be used to interpret the galaxy correlation function was to consider a density field composed of randomly-placed independent clumps with some universal density profile (Neyman, Scott & Shane 1953; Peebles 1974). Since the clumps are placed at random, the only correlations arise from points in the same clump. The correlations are easily deduced by using statistical isotropy: calculate the excess number of pairs separated by a distance $`r`$ in the $`z`$ direction (chosen as some arbitrary polar axis in a spherically-symmetric clump). For power-law clumps, with $`\rho =nBr^ฯต`$, truncated at $`r=R`$, this model gives $`\xi r^{32ฯต}`$ in the limit $`rR`$, provided $`3/2<ฯต<3`$. Values $`ฯต>3`$ are unphysical, and require a small-scale cutoff to the profile. There is no such objection to $`ฯต<3/2`$, and the expression for $`\xi `$ tends to a constant for small $`r`$ in this case (see Yano & Gouda 1999).
A long-standing problem for this model is that the correlation function in this case is much flatter than is observed for galaxies: $`\xi r^{1.8}`$ is the canonical slope, requiring $`ฯต=2.4`$. The first reaction may be to say that the model is incredibly naive by comparison with our sophisticated present understanding of the nonlinear evolution of CDM density fields. However, as will be shown below, it may after all contain more than a grain of truth.
### 6.2 The CDM clustering problems
A number of authors have pointed out that the detailed spectral shape inferred from galaxy data appears to be inconsistent with that of nonlinear evolution from CDM initial conditions. (e.g. Efstathiou, Sutherland & Maddox 1990; Klypin, Primack & Holtzman 1996; Peacock 1997). Perhaps the most detailed work was carried out by the VIRGO consortium, who carried out $`N=256^3`$ simulations of a number of CDM models (Jenkins et al. 1998). Their results are shown in figure 9, which gives the nonlinear power spectrum at various times (cluster normalization is chosen for $`z=0`$) and contrasts this with the APM data. The lower small panels are the scale-dependent bias that would required if the model did in fact describe the real universe, defined as
$`b(k)\left({\displaystyle \frac{\mathrm{\Delta }_{\mathrm{gals}}^2(k)}{\mathrm{\Delta }_{\mathrm{mass}}^2}}\right)^{1/2}.`$ (50)
In all cases, the required bias is non-monotonic; it rises at $`k\stackrel{>}{}5h^1\mathrm{Mpc}`$, but also displays a bump around $`k0.1h^1\mathrm{Mpc}`$. If real, this feature seems impossible to understand as a genuine feature of the mass power spectrum; certainly, it is not at a scale where the effects of even a large baryon fraction would be expected to act (Eisenstein et al. 1998; Meiksin, White & Peacock 1999).
## 7 Bias
The conclusions from the above discussion are either that the physics of dark matter and structure formation are more complex than in CDM models, or that the relation between galaxies and the overall matter distribution is sufficiently complicated that the effective bias is not a simple slowly-varying monotonic function of position.
### 7.1 Simple bias models
The simplest assumption is that all the complicated physical effects leading to galaxy formation depend in a causal (but nonlinear) way on the local mass density, so that we write
$`\rho _{\mathrm{light}}=f(\rho _{\mathrm{mass}}).`$ (51)
Coles (1993) showed that, under rather general assumptions, this equation would lead to an effective bias that was a monotonic function of scale. This issue was investigated in some detail by Mann, Peacock & Heavens (1998), who verified Colesโ conclusion in practice for simple few-parameter forms for $`f`$, and found in all cases that the effective bias varied rather weakly with scale. The APM results thus are either inconsistent with a CDM universe, or require non-local bias.
A puzzle with regard to this conclusion is provided by the work of Jing, Mo & Bรถrner (1998). They evaluated the projected real-space correlations for the LCRS survey (see figure 10). This statistic also fails to match the prediction of CDM models, but this can be amended by introducing a simple antibias scheme, in which galaxy formation is suppressed in the most massive haloes. This scheme should in practice be very similar to the Mann, Peacock & Heavens recipe of a simple weighting of particles as a function of the local density; indeed, the main effect is a change of amplitude, rather than shape of the correlations. The puzzle is this: if the APM power spectrum is used to predict the projected correlation function, the result agrees almost exactly with the LCRS. Either projected correlations are a rather insensitive statistic, or perhaps the Baugh & Efstathiou deconvolution procedure used to get $`P(k)`$ has exaggerated the significance of features in the spectrum. The LCRS results are one reason for treating the apparent conflict between APM and CDM with caution.
### 7.2 Halo correlations
In reality, bias is unlikely to be completely causal, and this has led some workers to explore stochastic bias models, in which
$`\rho _{\mathrm{light}}=f(\rho _{\mathrm{mass}})+ฯต,`$ (52)
where $`ฯต`$ is a random field that is uncorrelated with the mass density (Pen 1998; Dekel & Lahav 1999). Although truly stochastic effects are possible in galaxy formation, a relation of the above form is expected when the galaxy and mass densities are filtered on some scale (as they always are, in practice). Just averaging a galaxy density that is a nonlinear function of the mass will lead to some scatter when comparing with the averaged mass field; a scatter will also arise when the relation between mass and light is non-local, however, and this may be the dominant effect.
The simplest and most important example of non-locality in the galaxy-formation process is to recognize that galaxies will generally form where there are galaxy-scale haloes of dark matter. In the past, it was generally believed that dissipative processes were critically involved in galaxy formation, since pure collisionless evolution would lead to the destruction of galaxy-scale haloes when they are absorbed into the creation of a larger-scale nonlinear system such as a group or cluster. However, it turns out that this overmerging problem was only an artefact of inadequate resolution. When a simulation is carried out with $`10^6`$ particles in a rich cluster, the cores of galaxy-scale haloes can still be identified after many crossing times (Ghigna et al. 1997). Furthermore, if catalogues of these โsub-haloesโ are created within a cosmological-sized simulation, their correlation function is quite different from that of the mass, resembling the single power law seen in galaxies (e.g. Klypin et al. 1999; Ma 1999).
These are very important results, and they hold out the hope that many of the issues concerning where galaxies form in the cosmic density field can be settled within the domain of collisionless simulations. Dissipative physics will still be needed to understand in detail the star-formation history within a galaxy-scale halo. Nevertheless, the idea that there may be a one-to-one correspondence between galaxies and galaxy-scale dark-matter haloes is clearly an enormous simplification โ and one that increases the chance of making robust predictions of the statistical properties of the galaxy population.
### 7.3 Numerical galaxy formation
The formation of galaxies must be a non-local process to some extent. The modern paradigm was introduced by White & Rees (1978): galaxies form through the cooling of baryonic material in virialized haloes of dark matter. The virial radii of these systems are in excess of 0.1 Mpc, so there is the potential for large differences in the correlation properties of galaxies and dark matter on these scales.
A number of studies have indicated that the observed galaxy correlations may indeed be reproduced by CDM models. The most direct approach is a numerical simulation that includes gas, and relevant dissipative processes. This is challenging, but just starting to be feasible with current computing power (Pearce et al. 1999). The alternative is โsemianalyticโ modelling, in which the merging history of dark-matter haloes is treated via the extended Press-Schechter theory (Bond et al. 1991), and the location of galaxies within haloes is estimated using dynamical-friction arguments (e.g. Cole et al. 1996; Kauffmann et al. 1996; Somerville & Primack 1997). Both these approaches have yielded similar conclusions, and shown how CDM models can match the galaxy data: specifically, the low-density flat $`\mathrm{\Lambda }`$CDM model that is favoured on other grounds can yield a correlation function that is close to a single power law over $`1000\stackrel{>}{}\xi \stackrel{>}{}1`$, even though the mass correlations show a marked curvature over this range (Pearce et al. 1999; Benson et al. 1999; see figure 11). These results are impressive, yet it is frustrating to have a result of such fundamental importance emerge from a complicated calculational apparatus. There is thus some motivation for constructing a simpler heuristic model that captures the main processes at work in the full semianalytic models. The following section describes an approach of this sort (Peacock & Smith, in preparation).
### 7.4 Halo-ology and bias
We mentioned above the early model of Neyman, Scott & Shane (1953), in which the nonlinear density field was taken to be a superposition of randomly-placed clumps. With our present knowledge about the evolution of CDM universes, we can make this idealised model considerably more realistic: hierarchical models are expected to contain a distribution of masses of clumps, which have density profiles that are more complicated than isothermal spheres. These issues are well studied in $`N`$-body simulations, and highly accurate fitting formulae exist, both for the mass function and for the density profiles. Briefly, we use the mass function of Sheth & Tormen (1999; ST) and the halo profiles of Moore et al. (1999; M99).
$`f(\nu )`$ $`=0.21617[1+(\sqrt{2}/\nu ^2)^{0.3}]\mathrm{exp}[\nu ^2/(2\sqrt{2})]`$ $`F(>\nu )`$ $`=0.32218[1\mathrm{erf}(\nu /2^{3/4})]`$ $`+0.14765\mathrm{\Gamma }[0.2,\nu ^2/(2\sqrt{2})],`$ (53)
where $`\mathrm{\Gamma }`$ is the incomplete gamma function.
Recently, it has been claimed by Moore et al. (1999; M99) that the commonly-adopted density profile of Navarro, Frenk & White (1996; NFW) is in error at small $`r`$. M99 proposed the alternative form
$`\rho /\rho _b={\displaystyle \frac{\mathrm{\Delta }_c}{y^{3/2}(1+y^{3/2})}};(r<r_{\mathrm{vir}});yr/r_c.`$ (54)
Using this model, it is then possible to calculate the correlations of the nonlinear density field, neglecting only the large-scale correlations in halo positions. The power spectrum determined in this way is shown in figure 12, and turns out to agree very well with the exact nonlinear result on small and intermediate scales. The lesson here is that a good deal of the nonlinear correlations of the dark matter field can be understood as a distribution of random clumps, provided these are given the correct distribution of masses and mass-dependent density profiles.
How can we extend this model to understand how the clustering of galaxies can differ from that of the mass? There are two distinct ways in which a degree of bias is inevitable:
(1) Halo occupation numbers. For low-mass haloes, the probability of obtaining an $`L^{}`$ galaxy must fall to zero. For haloes with mass above this lower limit, the number of galaxies will in general not scale with halo mass.
(2) Nonlocality. Galaxies can orbit within their host haloes, so the probability of forming a galaxy depends on the overall halo properties, not just the density at a point. Also, the galaxies will end up at special places within the haloes: for a halo containing only one galaxy, the galaxy will clearly mark the halo centre. In general, we expect one central galaxy and a number of satellites.
The numbers of galaxies that form in a halo of a given mass is the prime quantity that numerical models of galaxy formation aim to calculate. However, for a given assumed background cosmology, the answer may be determined empirically. Galaxy redshift surveys have been analyzed via grouping algorithms similar to the โfriends-of-friendsโ method widely employed to find virialized clumps in $`N`$-body simulations. With an appropriate correction for the survey limiting magnitude, the observed number of galaxies in a group can be converted to an estimate of the total stellar luminosity in a group. This allows a determination of the All Galaxy System (AGS) luminosity function: the distribution of virialized clumps of galaxies as a function of their total luminosity, from small systems like the Local Group to rich Abell clusters.
The AGS function for the CfA survey was investigated by Moore, Frenk & White (1993), who found that the result in blue light was well described by
$`d\varphi =\varphi ^{}\left[(L/L^{})^\beta +(L/L^{})^\gamma \right]^1dL/L^{},`$ (55)
where $`\varphi ^{}=0.00126h^3\mathrm{Mpc}^3`$, $`\beta =1.34`$, $`\gamma =2.89`$; the characteristic luminosity is $`M^{}=21.42+5\mathrm{log}_{10}h`$ in Zwicky magnitudes, corresponding to $`M_B^{}=21.71+5\mathrm{log}_{10}h`$, or $`L^{}=7.6\times 10^{10}h^2L_{}`$, assuming $`M_B^{}=5.48`$. One notable feature of this function is that it is rather flat at low luminosities, in contrast to the mass function of dark-matter haloes (see Sheth & Tormen 1999). It is therefore clear that any fictitious galaxy catalogue generated by randomly sampling the mass is unlikely to be a good match to observation. The simplest cure for this deficiency is to assume that the stellar luminosity per virialized halo is a monotonic, but nonlinear, function of halo mass. The required luminosityโmass relation is then easily deduced by finding the luminosity at which the integrated AGS density $`\mathrm{\Phi }(>L)`$ matches the integrated number density of haloes with mass $`>M`$. The result is shown in figure 13.
We can now return to the halo-based galaxy power spectrum and use the correct occupation number, $`N`$, as a function of mass. This is needs a little care at small numbers, however, since the number of haloes with occupation number unity affects the correlation properties strongly. These haloes contribute no correlated pairs, so they simply dilute the signal from the haloes with $`N2`$. The existence of antibias on intermediate scales can probably be traced to the fact that a large fraction of galaxy groups contain only one $`>L_{}`$ galaxy. Finally, we need to put the galaxies in the correct location, as discussed above. If one galaxy always occupies the halo centre, with others acting as satellites, the small-scale correlations automatically follow the slope of the halo density profile, which keeps them steep. The results of this exercise are shown in figure 14.
Although it is encouraging that it is possible to find simple models in which it is possible to understand the observed correlation properties of galaxies, there are other longstanding puzzles concerning the galaxy distribution. Arguably the chief of these concerns the dynamical properties of galaxies, in particular the pairwise peculiar velocity dispersion. This statistic has been the subject of debate, and preferred values have crept up in recent years, to perhaps 450 or $`500\mathrm{km}\mathrm{s}^1`$ at projected separations around 1 Mpc (e.g. Jing, Mo & Bรถrner 1998), most simple models predict a higher figure. Clearly, the amplitude of peculiar velocities depends on the normalization of the fluctuation spectrum; however, if this is set from the abundance of rich clusters, then Jenkins et al. (1998) found that reasonable values were predicted for large-scale streaming velocities, independent of $`\mathrm{\Omega }`$. However, Jenkins et al. also found a robust prediction for the pairwise peculiar velocity dispersion around 1 Mpc of about $`800\mathrm{km}\mathrm{s}^1`$. The observed galaxy velocity field appears to have a higher โcosmic Mach numberโ than the predicted dark-matter distribution.
This difficulty is also solved by the simple bias model discussed here. Two factors contribute: the variation of occupation number with mass downweights the contribution of more massive groups, with larger velocity dispersions. Also, where one galaxy is centred on a halo, it gains a peculiar velocity which is that of the centre of mass of the halo, but does not reflect the internal velocity dispersion of the halo. Given a full $`N`$-body simulation, it is easy enough to predict what would be expected for a realistic bias model: one needs to construct a halo catalogue, calculating the peculiar velocities and internal velocity dispersions of each halo. Knowing the occupation number as a function of mass, a montecarlo catalogue of โgalaxiesโ complete with peculiar velocities can be generated. As shown in figure 15, the effect of the empirical bias recipe advocated here is sufficient to reduce the predicted dispersion into agreement with observation. The simple model outlined here thus gives a consistent picture, and it is tempting to believe that it may capture some of the main features of realistic models for galaxy bias.
## 8 Conclusions
It should be clear from these lectures that large-scale structure has advanced enormously as a field in the past two decades. Many of our long-standing ambitions have been realised; in some cases, much faster than we might have expected. Of course, solutions for old problems generate new difficulties. We now have good measurements of the clustering spectrum and its evolution, and it is arguable that the discussion of section 7.4 captures the main features of the placement of galaxies with respect to the mass. However, a fairly safe bet is that one of the major results from new large surveys such as 2dF and Sloan will be a heightened appreciation of the subtleties of this problem.
Nevertheless, we should not be depressed if problems remain. Observationally, we are moving from an era of 20% โ 50% accuracy in measures of large-scale structure to a future of pinpoint precision. This maturing of the subject will demand more careful analysis and rejection of some of our existing tools and habits of working. The prize for rising to this challenge will be the ability to claim a real understanding of the development of structure in the universe. We are not there yet, but there is a real prospect that the next 5โ10 years may see this remarkable goal achieved.
## Acknowledgements
I thank my colleagues in the 2dF Galaxy Redshift Survey for permission to reproduce our joint results in section 5.3, and Robert Smith for the joint work reported in section 7.4.
## References
Ballinger W.E., Peacock J.A., Heavens A.F., 1996, MNRAS, 282, 877
Baugh C.M., Efstathiou G., 1993, MNRAS, 265, 145
Baugh C.M., Efstathiou G., 1994, MNRAS, 267, 323
Benson A.J., Cole S., Frenk C.S., Baugh C.M., Lacey C.G., 1999, astro-ph/9903343
Bond J.R., Cole S., Efstathiou G., Kaiser N., 1991, Astrophys. J., 379, 440
Bond J.R., 1995, Phys. Rev. Lett., 74, 4369
Bunn E.F., 1995, PhD thesis, Univ. of California, Berkeley
Coles P., 1993, MNRAS, 262, 1065
Colless M., 1999, Phil. Trans. R. Soc. Lond. A, 357, 105
Davis M., Peebles P.J.E., 1983, Astrophys. J., 267, 465
Dekel A., Lahav O., 1999, Astrophys. J., 520, 24
Efstathiou G., Sutherland W., Maddox S.J., 1990, Nature, 348, 705
Eisenstein D.J., Hu W., 1998, ApJ, 496, 605
Eisenstein D.J., Zaldarriaga M., 1999, astro-ph/9912149
Eke V.R., Cole S., Frenk C.S., 1996, Mon. Not. R. Astr. Soc., 282, 263
Feldman H.A., Kaiser N., Peacock J.A., 1994, Astrophys. J., 426, 23
Folkes S., et al., 1999, Mon. Not. R. Astr. Soc., 308, 459
Gaztaรฑaga E., Baugh C.M. 1998, Mon. Not. R. Astr. Soc., 294, 229
Ghigna S., Moore B., Governato F., Lake G., Quinn T., Stadel J., 1998, Mon. Not. R. Astr. Soc., 300, 146
Goldberg D.M., Strauss M.A., 1998, ApJ, 495, 29
Hamilton A.J.S., Kumar P., Lu E., Matthews A., 1991, Astrophys. J., 374, L1
Hamilton A.J.S., 1992, Astrophys. J., 385, L5
Hamilton A.J.S., 1997a, astro-ph/9708102
Hamilton A.J.S., 1997b, Mon. Not. R. Astr. Soc., 289, 285
Hamilton A.J.S., 1997c, Mon. Not. R. Astr. Soc., 289, 295
Huchra J.P., Geller M.J., de Lapparant V., Corwin H.G., 1990, Astrophys. J. Suppl., 72, 433
Jain B., Mo H.J., White S.D.M., 1995, Mon. Not. R. Astr. Soc., 276, L25
Jenkins A., Frenk C.S., Pearce F.R., Thomas P.A., Colberg J.M., White S.D.M., Couchman H.M.P., Peacock J.A., Efstathiou G., Nelson A.H., 1998, ApJ, 499, 20
Jing Y.P., Mo H.J., Bรถrner G., 1998, ApJ, 494, 1
Kaiser N., 1987, MNRAS, 227, 1
Klypin A., Primack J., Holtzman J., 1996, Astrophys. J., 466, 13
Klypin A., Gottloeber S., Kravtsov A.V., Khokhlov A.M., 1999, 516, 530
Ma C.-P., 1999, Astrophys. J., 510, 32
Maddox S. Efstathiou G., Sutherland W.J., 1996, MNRAS, 283, 1227
Mann R.G., Peacock J.A., Heavens A.F., 1998, MNRAS, 293, 209
Margon B., 1999, Phil. Trans. R. Soc. Lond. A, 357, 93
Matsubara T., Szalay A.S., Landy S.D., 1999, astro-ph/9911151
Meiksin A.A., White M., Peacock J.A., MNRAS, 1999, 304, 851
Moore B., Frenk C.S., White S.D.M., 1993, Mon. Not. R. Astr. Soc., 261, 827
Moore B., Quinn T., Governato F., Stadel J., Lake G., 1999, astro-ph/9903164
Navarro J.F., Frenk C.S., White S.D.M., 1996, ApJ, 462, 563
Neyman, Scott & Shane 1953, ApJ, 117, 92
Padmanabhan N., Tegmark M., Hamilton A.J.S., 1999, astro-ph/9911421
Peacock J.A., Dodds S.J., 1994, MNRAS, 267, 1020
Peacock J.A., Dodds S.J., 1996, MNRAS, 280, L19
Peacock J.A., 1997, Mon. Not. R. Astr. Soc., 284, 885
Pearce F.R., et al., 1999, astro-ph/9905160
Peebles P.J.E., 1974, A&A, 32, 197
Pen, W.-L., 1998, ApJ, 504, 601
Saunders W., et al., 1999, astro-ph/9909191
Shectman S.A., Landy S.D., Oemler A., Tucker D.L., Lin H., Kirshner R.P., Schechter P.L., 1996, Astrophys. J., 470, 172
Sheth R.K., Tormen G., 1999, astro-ph/9901122
Strauss M.A., Willick J.A., 1995, Physics Reports, 261, 271
Sugiyama N., 1995, Astrophys. J. Suppl., 100, 281
Tegmark M., 1996, Mon. Not. R. Astr. Soc., 280, 299
Tegmark M., Taylor A.N., Heavens A.F., 1997, Astrophys. J., 480, 22
Tegmark M., Hamilton A.J.S., Strauss M.A., Vogeley M.S., Szalay A.S., 1998, Astrophys. J., 499, 555
Viana P.T., Liddle A.R., 1996, MNRAS, 281, 323
Vogeley M.S., Szalay, A.S., 1996, Astrophys. J., 465, 34
White S.D.M., Rees M., 1978, Mon. Not. R. Astr. Soc., 183, 341
White S.D.M., Efstathiou G., Frenk C.S., 1993, Mon. Not. R. Astr. Soc., 262, 1023
Yano T., Gouda N., 1999, astro-ph/9906375 |
warning/0002/hep-ph0002188.html | ar5iv | text | # 1 Introduction
## 1 Introduction
Ever since the pioneering scattering experiments by Rutherford in the early 1900s, there has been the persistent question of how matter would behave during very hard collisions. Two examples are the Landau picture where matter would come to a complete stop before exploding due to the enormous pressure built-up or the Bjorken picture where target and projectile would pass through each other experiencing only partial deceleration in the process. Present day high energy nuclear collision experiments provide not only an arena for settling this age-old question but also to see if there is any energy dependence that would change the collision scenario from one picture to another.
Another very good reason to study this is related to the original central goal of these experiments, which is to recreate deconfined matter, the so-called quark-gluon plasma, believed to have existed only in the early universe. To confirm its existence in the laboratory at the up-coming Relativistic Heavy Ion Collider (RHIC) at Brookhaven or the Large Hadron Collider (LHC) at CERN, one relies heavily on finding evidence amongst the many produced particles. These will be affected by the environment in which they are produced. In particular the net baryon number left in the central collision region or the amount of baryon stopping affects, for example, photon and dilepton production. It is therefore of considerable interest to measure this in experiment. Indeed this has been done by NA49 at CERNโs SPS by measuring the net proton or baryon rapidity distribution of the final hadrons.
Recently it was proposed to measure this using the bremsstrahlung associated with slowing down of the baryons during the collisions. This method should be simpler than measuring hadron rapidities because of the extreme forward focus of the photons emitted from relativistically moving targets and projectiles. The measurement is therefore much more localized and only photons instead of different types of hadrons need to be detected. From a pragmatic point of view, it is better because most hadron detectors are not able to cover the full rapidity range. From a physical viewpoint, it can say something about the space-time evolution of the collisions . This information will definitely not be available from rapidity measurements alone. In addition, it can distinguish between the Landau and Bjorken picture described above.
In view of the two different ways of measuring baryon stopping, one would like a way to compare the two and be able to communicate between them. In any case, whenever one talks about baryon stopping today, one is usually referring to the shape of the rapidity distribution. As far as we are aware, there is no attempt to quantify it in any way so that the measurement can be put into a more concrete footing. In this paper, we will make just such an attempt. Whenever we refer to baryon stopping, we are referring to stopping in rapidity space as is traditionally the case and not in velocity space. It will be shown that stopping in rapidity and velocity space are entirely different matters and therefore it is very important that it is made clear in which it is being described. There can be no confusion between the two.
## 2 Quantifying Baryon Stopping or Transparency
In order to quantify baryon stopping or its inverse, baryon transparency, we aim to find a quantity, called $`๐ฎ`$, that satisfies the following requirements.
* It should be equal to unity if there is complete stopping. That is, all baryons end up having $`y=0`$.
* If there is full transparency, and all final baryons move with the original initial rapidity $`y_0`$, this quantity should be zero.
* One should be able to define its inverse, baryon transparency, $`๐ฏ`$ which is related to $`๐ฎ`$ by the simple relation $`๐ฎ=1๐ฏ`$ so that it has the opposite value in case (i) and (ii).
* $`๐ฎ`$ and $`๐ฏ`$ should both be equal to half or approximately so when it is clear that there is half stopping and half transparency in rapidity space. For example, $`๐ฎ=๐ฏ=1/2`$ when the rapidity distribution is totally flat.
* For different degrees of stopping, $`๐ฎ`$ should range between $`0`$ and $`1`$, signifying transparent to opaque in that order (or the inverse for $`๐ฏ`$).
* It should be obtainable from both hadron rapidity data and from bremsstrahlung measurements. This requirement is essential for bridging the two types of measurements.
A quantity that satisfies all these is
$$๐ฎ=1\frac{(1v_0^2\mathrm{cos}^2\theta _{1/2})}{2v_0^2\mathrm{cos}\theta _{1/2}}_{\mathrm{}}^+\mathrm{}๐y\frac{v(y)\rho (y)}{1v(y)\mathrm{cos}\theta _{1/2}}$$
(1)
and therefore
$$๐ฏ=\frac{(1v_0^2\mathrm{cos}^2\theta _{1/2})}{2v_0^2\mathrm{cos}\theta _{1/2}}_{\mathrm{}}^+\mathrm{}๐y\frac{v(y)\rho (y)}{1v(y)\mathrm{cos}\theta _{1/2}}.$$
(2)
Here $`v_0=v(y_0)`$ is the initial velocity in the center of mass frame which is related to the initial rapidity $`y_0`$ by the general relation between velocity and rapidity
$$v(y)=\mathrm{tanh}y.$$
(3)
The $`\rho (y)`$ is proportional to the final baryon rapidity distribution $`dN/dy`$. For symmetric target and projectile, it is defined by
$$_{\mathrm{}}^+\mathrm{}๐y\rho (y)=2.$$
(4)
To see that specifications (i), (ii) and (iii) are satisfied, one can consider two distributions
$$\rho (y)=2\delta (y)$$
(5)
and
$$\rho (y)=\delta (yy_0)+\delta (y+y_0)$$
(6)
which correspond to full stopping and full transparency, respectively. It is easily verified that one gets $`๐ฎ=1`$, $`๐ฏ=0`$ from the first distribution and $`๐ฎ=0`$, $`๐ฏ=1`$ from the second.
Now for specification (iv), the flat rapidity distribution that respects Eq. (4) would be
$`\rho (y)`$ $`=`$ $`1/y_0\mathrm{for}|y|y_0`$
$`=`$ $`0\mathrm{o}\mathrm{t}\mathrm{h}\mathrm{e}\mathrm{r}\mathrm{w}\mathrm{i}\mathrm{s}\mathrm{e}.`$
Using the formula
$$\frac{dy}{1\mathrm{cos}\theta \mathrm{tanh}y}=\frac{y}{\mathrm{sin}^2\theta }+\frac{\mathrm{cos}\theta }{\mathrm{sin}^2\theta }\left\{\mathrm{ln}\mathrm{cosh}y+\mathrm{ln}(1\mathrm{cos}\theta \mathrm{tanh}y)\right\}$$
(8)
and the requirement that $`๐ฎ=๐ฏ=1/2`$ for a flat distribution, we end up with the equation
$$(2v_0^2)\mathrm{cos}\theta _{1/2}v_0^2\mathrm{cos}^3\theta _{1/2}\frac{(1v_0^2\mathrm{cos}^2\theta _{1/2})}{y_0}\mathrm{ln}\left(\frac{1+v_0\mathrm{cos}\theta _{1/2}}{1v_0\mathrm{cos}\theta _{1/2}}\right)=0.$$
(9)
This equation defines the value of the angle $`\theta _{1/2}`$ which is the angle at which a flat rapidity distribution will yield $`๐ฎ=๐ฏ=1/2`$ for a given $`v_0`$ or $`y_0`$. At this stage, $`\theta _{1/2}`$ is no more than a numerical quantity but its meaning will be explained below. In Table 1, some values of $`\theta _{1/2}`$ have been computed at the various accelerators. The first case at SPS is for Pb+Pb and the second is for S+S collisions. The angle becomes smaller as we go to higher energies.
Our specification (iv) requires that rapidity distributions intuitively half way between stopping and transparency in rapidity space should be $`0.5`$. To show that this is indeed the case, we use a simple test distribution depicted in Fig. 1. In this figure, we have two blocks at the height of $`1/(L\lambda )y_0`$ because of Eq. (4) and symmetric about $`y=0`$. The parameters $`L`$ and $`\lambda `$ allows us the freedom of a range of distributions. For our purpose here, we set $`L=1\lambda `$ and vary the value of $`\lambda `$. For any value of $`\lambda `$, we have two blocks symmetric about $`\pm y_0/2`$, and thus they should all have $`๐ฎ๐ฏ0.5`$. We tabulated the value of $`๐ฎ`$ for a set of values of $`\lambda `$ at the various accelerators in Table 2. We recover the flat distribution when $`\lambda =0`$ and we have two delta functions sitting at $`\pm y_0/2`$ when $`\lambda =1/2`$. We see that (vi) is better satisfied as we go to higher and higher energies. Since this way of using bremsstrahlung will only be done at RHIC or at LHC, this is good enough and we consider this specification met. In any case, they are all fairly close to $`0.5`$.
One can introduce another special case which is without any doubt half way between opaque and transparent. That is
$$\rho (y)=\frac{1}{2}\left(\delta (yy_0)+\delta (y+y_0)+2\delta (y_0)\right).$$
(10)
This distribution is artificial but is ideal for our purpose here. This distribution describes half the baryons from the projectile and half from the target sitting at $`y=0`$ and half of them from each initial nucleus traveling with the original $`y_0`$. This is easily worked out to give the exact result $`๐ฎ=๐ฏ=1/2`$. The purpose of this last distribution is to show that we have a sensible definition in hand.
## 3 Difference Between Rapidity and Velocity Space
When one looks for a numerical definition for stopping or transparency, one encounters the question of whether this should be in rapidity or velocity space. If one was working within Newtonian mechanics, velocity space would have been the automatic choice. This is quite logical since one could easily associate stopping with the slowing down of the incoming clusters of nucleons. However, it is also traditional to speak of baryon stopping while referring implicitly to the shape of $`dN/dy`$ in rapidity space. So it is in rapidity space that we gave this definition in the previous section. It must be stressed that the difference is huge between the two spaces. For example, the flat distribution given in Eq. (2) at RHIC in rapidity space will appear as in Fig. 2 (a) in velocity space because
$$\rho (y)=\frac{1}{y_0}\rho (v)=\frac{1}{y_0(1v^2)}.$$
(11)
The halfway point of stopping in rapidity space at RHIC would appear to be much closer to transparent in velocity space! On the contrary, a flat distribution in velocity space
$$\rho (v)=\frac{1}{v_0}\rho (y)=\frac{1}{v_0\mathrm{cosh}^2y}$$
(12)
would be more opaque than transparent when one switches to rapidity. This is shown in Fig. 2 (b). Using this distribution, one finds $`๐ฎ0.97`$ or $`๐ฏ0.03`$.
For this reason, we have insisted that our definition be in rapidity space, conforming to the convention used in the heavy ion collision community.
## 4 How To Obtain $`๐ฎ`$ and $`๐ฏ`$ From Bremsstrahlung Measurements
In Sect. 2, the definition given for baryon stopping $`๐ฎ`$ and transparency $`๐ฏ`$ depended on the distribution $`\rho (y)`$ which, for symmetric collisions, is related to the net baryon rapidity distribution by $`dN_{B\overline{B}}/dy=A\rho (y)`$ or net proton rapidity distribution by $`dN_{p\overline{p}}/dyZ\rho (y)`$. Here $`Z`$ and $`A`$ are the atomic and mass number respectively of the incoming target or projectile. This assumes that the net proton rapidity distribution is representative of or approximately proportional to the net baryon distribution. This does not seem to be too bad an assumption if one examines data from SPS . Since one could easily calculate the value of $`๐ฎ`$ and $`๐ฏ`$ for some given data of $`dN/dy`$, the important question would be how to meet specification (vi) in Sect. 2. The main difficulty is how to obtain the quantity $`๐ฎ`$ or $`๐ฏ`$ from bremsstrahlung measurements. To solve this problem, we now state that $`๐ฎ`$ can be related to the intensity distribution of the bremsstrahlung emitted in nuclear collisions by
$$๐ฎ=\frac{1v_0^2\mathrm{cos}^2\theta _{1/2}}{v_0^2\mathrm{sin}2\theta _{1/2}}\left(\frac{4\pi ^2}{\alpha Z^2}\frac{d^2I}{d\omega d\mathrm{\Omega }}|_{\genfrac{}{}{0pt}{}{\omega 0}{\theta =\theta _{1/2}}}\right)^{1/2}.$$
(13)
The numerical quantity $`\theta _{1/2}`$ has now been given the physical meaning of the opening angle from the beam pipe in which direction the soft photons should be measured. To see exactly where this angle lies in relation to other directions, one can work out the angle of maximum intensity assuming, for example, the case of full stopping $`๐ฎ=1`$ or complete opacity $`๐ฏ=0`$ using the formulae given in ref. . In this case, the cosine of this angle $`\theta _{\mathrm{max}}`$ is related to the initial velocity $`v_0`$ by
$$\mathrm{cos}\theta _{\mathrm{max}}=(2v_0^2)^{1/2}.$$
(14)
It then works out at RHIC to be $`\theta _{\mathrm{max}}=0.538^o`$ and $`\theta _{\mathrm{max}}=0.072^o`$ at LHC. In view of the fact that the intensity distribution falls off with the opening angle $`\theta `$ from the beam pipe, these $`\theta _{\mathrm{max}}`$ are not too far from those $`\theta _{1/2}`$ in Table 1 at the respective accelerators so that there will be sufficient intensity at the $`\theta _{1/2}`$ to enable the photon measurements.
If the reader has not guessed it already, we will now disclose the physical meaning of the somewhat mysterious quantity $`๐ฎ`$ or $`๐ฏ`$. The inverse of the prefactor to the rapidity integral in Eq. (2) is in fact proportional to the radiation amplitude for full stopping . So one can now see why the rapidity integral itself will always be less than or equal to the inverse of this prefactor (but see the next paragraph concerning the soft photon requirement). Because of $`๐ฎ`$ and $`๐ฏ`$ have an origin in the bremsstrahlung intensity distribution, they depend on the (charge) hadron rapidity distribution which automatically allows them to bridge the two different methods of determining baryon stopping. So far we have not mentioned the contribution from charged mesons to photon emissions. They could potentially ruin Eq. (13). However, pions are by far the most abundant meson type and they can be positively as well as negatively charged. Therefore contribution to bremsstrahlung from mesons cancel out to a large extent .
Although it has been expressed in Eq. (13) that the intensity distribution should be for low energy photons, in practice a few to tens of MeV should be good enough. The reason for the soft photon requirement is to remove all nuclear structural dependence as well as any potential interference effects. In ref. it was shown that if there were more than one component in the acceleration of the nuclear clusters during the collisions, this would result in enhancement in and oscillations of the intensity distribution $`dI/d\omega d\mathrm{\Omega }`$ with $`\omega `$. This is a direct result of the interference between the various components in the acceleration. For the purpose of our definition, interference would unfortunately just taint any value of $`๐ฎ`$ obtained from Eq. (13). Only in the soft $`\omega `$ limit is it free from this type of effect. As seen in Fig. 7 of ref. , the intensity at small $`\omega `$ is invariant under this.
## 5 Examples
We will now try some example distributions and work out their $`๐ฎ`$ and $`๐ฏ`$ values. For the test distribution in Fig. 1 we vary the two parameters $`L`$ and $`\lambda `$ to obtain the stopping values for the different cases. The first five entries in Table 3 are for central single-block distributions centering around $`y=0`$. As the distribution is widened, $`๐ฎ`$ decreases towards the flat $`0.5`$ value, as expected. The subsequent groups of entries are for symmetric two-block distributions shifting progressively away from the center to either side towards $`|y|=1`$. Thus in each group there is the tendency $`๐ฎ0`$ and $`๐ฏ1`$ which is how a sensible definition should behave.
Admittedly these distributions are only test cases designed to show how the numerical definition works. However, more realistic distributions can always be approximately reconstructed from thin blocks (strips) of varying heights so the simple distributions used do not affect in any way how the definition meets the specifications stipulated in Sec. 2. It may be that in practice our definition would need to be refined but here we have laid the groundwork for a simple but sensible numerical definition for baryon stopping. For the actual applications of this to real data, and for other ways of using bremsstrahlung from more realistic collision models than those used in , we refer the reader to .
## Acknowledgments
The author would like to thank Joe Kapusta for a critical reading of the manuscript and for discussions. This work was supported by the U.S. Department of Energy under grant DE-FG02-87ER40328. |
warning/0002/math-ph0002042.html | ar5iv | text | # The Klein-Gordonโs field. A counter-example of the classical limit
## 1. Introduction
In this work we will study the classical limit of Klein-Gordonโs field, with an homogeneous potential which does not depend on Planckโs constant.
First we will see that, in this case, the Klein-Gordonโs equation is equivalent to a hamiltonian system, composed by an infinite number of harmonic oscillators with frequencies which depend on time. Once we have seen this equivalence, we will quantize these oscillators and we will obtain the energy and the electric charge operators. With the energy operator, we will obtain the quantum equation of Klein-Gordonโs field. We will also see that we can find all the eigenfunctions of the energy and the electric charge operators. Consequently, with all those eigenfunctions we can construct the Fockโs space.
After that, we will study the quantum dynamic of vacuum state. We will see that, if the space dimension is $`1`$, when $`\mathrm{}0`$, the probability that does not exist any particle-antiparticle pair, converges to $`1`$. However, in dimension $`2`$ or $`3`$, we will prove that, when $`\mathrm{}0`$, this probability does not converge to $`1`$. Consequently, in dimension $`2`$ or $`3`$, the classical limit is not true.
The notation that we are going to use, is the following
$`<,>`$ euclidean scalar product.
$`<,>_2`$ scalar product of $`^2`$.
$`||.||_2`$ norm $`^2`$.
$`||.||_2`$ norm $`\mathrm{}`$.
## 2. The Klein-Gordonโs field with an homogeneous potential
To simplify, we will take $`m=c=e=1`$.
If we apply the Correspondence Principle $`Ei\mathrm{}_t`$, $`\stackrel{}{p}i\mathrm{}\stackrel{}{}`$ to the relativistic relation $`E^2=|\stackrel{}{p}+\stackrel{}{f}(t)|^2+1`$, we obtain the Klein-Gordonโs equation,
$$\mathrm{}^2_t^2\psi =|i\mathrm{}\stackrel{}{}+\stackrel{}{f}(t)|^2\psi +\psi .$$
One important property of the K-Gโs equation is the electric charge conservation $`\dot{\rho }(t)=0`$, where $`\rho (t)=<i\mathrm{}_t\psi ,\psi >_2+<\psi ,i\mathrm{}_t\psi >_2`$. However there is no norme square conservation $`\dot{\psi (t)}_2^20`$. Then, to be able to speak about probabilities, we have to consider that the K-Gโs field describes an infinity of harmonic oscillators. After that, we have to quantize these oscillators to arrive to an equation of this type, $`i\mathrm{}_t|\mathrm{\Phi }>=H|\mathrm{\Phi }>`$, where $`H`$ is an self-adjoint operator.
### 2.1. The Quantization of Klein-Gordonโs field
Suppose that the domain is finite. To simplify we take the n-dimensional interval $`[\pi ,\pi ]^n`$.
The lagrangian and the energy of the system at $`t`$ time are:
$$L(t)=\mathrm{}^2_t\psi _2^2(i\mathrm{}\stackrel{}{}+\stackrel{}{f}(t))\psi _2^2\psi _2^2$$
$$E(t)=\mathrm{}^2_t\psi _2^2+(i\mathrm{}\stackrel{}{}+\stackrel{}{f}(t))\psi _2^2+\psi _2^2.$$
We expand $`\psi `$ in Fourierโs serie, $`\psi (\stackrel{}{x},t)=_{\stackrel{}{k}^n}A_\stackrel{}{k}(t)\psi _\stackrel{}{k}(\stackrel{}{x})`$, with $`\psi _\stackrel{}{k}(\stackrel{}{x})=\frac{e^{i<\stackrel{}{k},\stackrel{}{x}>}}{(2\pi )^{\frac{n}{2}}}`$. Then
$$L(t)=\underset{\stackrel{}{k}^n}{}\mathrm{}^2|\dot{A}_\stackrel{}{k}|^2ฯต_\stackrel{}{k}^2(t)|A_\stackrel{}{k}|^2,\text{ where }ฯต_\stackrel{}{k}(t)=\sqrt{|\mathrm{}\stackrel{}{k}+\stackrel{}{f}(t)|^2+1}.$$
With the momenta $`B_\stackrel{}{k}=\mathrm{}^2\dot{A}_\stackrel{}{k}`$, we obtain
$$E(t)=\underset{\stackrel{}{k}^n}{}\frac{|B_\stackrel{}{k}|^2}{\mathrm{}^2}+ฯต_\stackrel{}{k}^2(t)|A_\stackrel{}{k}|^2,\rho (t)=\underset{\stackrel{}{k}^n}{}\frac{i}{\mathrm{}}(A_\stackrel{}{k}^{}B_\stackrel{}{k}A_\stackrel{}{k}B_\stackrel{}{k}^{}).$$
We make the real canonical change
$$B_\stackrel{}{k}=\frac{\mathrm{}}{\sqrt{2}}(P_\stackrel{}{k}+i\overline{P}_\stackrel{}{k});A_\stackrel{}{k}=\frac{1}{\mathrm{}\sqrt{2}}(Q_\stackrel{}{k}+i\overline{Q}_\stackrel{}{k}),$$
and let $`\omega _\stackrel{}{k}(t)=\frac{ฯต_\stackrel{}{k}(t)}{\mathrm{}}`$ be the corresponding frequency, then $`E(t)`$ and $`\rho (t)`$ take the form
$$E(t)=\frac{1}{2}\underset{\stackrel{}{k}^n}{}(P_\stackrel{}{k}^2+\omega _\stackrel{}{k}^2(t)Q_\stackrel{}{k}^2)+(\overline{P}_\stackrel{}{k}^2+\omega _\stackrel{}{k}^2(t)\overline{Q}_\stackrel{}{k}^2)$$
$$\rho (t)=\frac{1}{\mathrm{}}\underset{\stackrel{}{k}^n}{}(\overline{Q}_\stackrel{}{k}P_\stackrel{}{k}Q_\stackrel{}{k}\overline{P}_\stackrel{}{k}).$$
That is the energy decomposition in oscillators. Notice, the K-Gโs equation is equivalent to the hamiltonian system
$`\{\begin{array}{ccc}\dot{Q}_\stackrel{}{k}& =& P_\stackrel{}{k}\\ \dot{P}_\stackrel{}{k}& =& \omega _\stackrel{}{k}^2(t)Q_\stackrel{}{k}\end{array}\{\begin{array}{ccc}\dot{\overline{Q}}_\stackrel{}{k}& =& \overline{P}_\stackrel{}{k}\\ \dot{\overline{P}}_\stackrel{}{k}& =& \omega _\stackrel{}{k}^2(t)\overline{Q}_\stackrel{}{k}\end{array}`$
Now, to obtain the quantum theory, what we have to do is to quantize these oscillators, i.e. $`P_\stackrel{}{k}i\mathrm{}_{Q_\stackrel{}{k}}`$, $`\overline{P}_\stackrel{}{k}i\mathrm{}_{\overline{Q}_\stackrel{}{k}}`$, and the equation will be
$$i\mathrm{}_t\mathrm{\Phi }=\frac{1}{2}\underset{\stackrel{}{k}^n}{}[(\mathrm{}^2_{Q_\stackrel{}{k}}^2+\omega _\stackrel{}{k}^2(t)Q_\stackrel{}{k}^2)+(\mathrm{}^2_{\overline{Q}_\stackrel{}{k}}^2+\omega _\stackrel{}{k}^2(t)\overline{Q}_\stackrel{}{k}^2)]\mathrm{\Phi }\underset{\stackrel{}{k}^n}{}\omega _\stackrel{}{k}(t)\mathrm{\Phi }.$$
Now, we will look for the eigenfunctions of the energy and of the electric charge operators. First, we have to introduce the creation and anihilation operators for particles and antiparticles
$$a_\stackrel{}{k}(t)=\frac{1}{2\sqrt{ฯต_\stackrel{}{k}(t)}}[(\mathrm{}_{Q_\stackrel{}{k}}+\omega _\stackrel{}{k}(t)Q_\stackrel{}{k})+i(\mathrm{}_{\overline{Q}_\stackrel{}{k}}+\omega _\stackrel{}{k}(t)\overline{Q}_\stackrel{}{k})]$$
$$a_\stackrel{}{k}^+(t)=\frac{1}{2\sqrt{ฯต_\stackrel{}{k}(t)}}[(\mathrm{}_{Q_\stackrel{}{k}}+\omega _\stackrel{}{k}(t)Q_\stackrel{}{k})i(\mathrm{}_{\overline{Q}_\stackrel{}{k}}+\omega _\stackrel{}{k}(t)\overline{Q}_\stackrel{}{k})]$$
$$b_\stackrel{}{k}(t)=\frac{1}{2\sqrt{ฯต_\stackrel{}{k}(t)}}[(\mathrm{}_{Q_\stackrel{}{k}}+\omega _\stackrel{}{k}(t)Q_\stackrel{}{k})i(\mathrm{}_{\overline{Q}_\stackrel{}{k}}+\omega _\stackrel{}{k}(t)\overline{Q}_\stackrel{}{k})]$$
$$b_\stackrel{}{k}^+(t)=\frac{1}{2\sqrt{ฯต_\stackrel{}{k}(t)}}[(\mathrm{}_{Q_\stackrel{}{k}}+\omega _\stackrel{}{k}(t)Q_\stackrel{}{k})+i(\mathrm{}_{\overline{Q}_\stackrel{}{k}}+\omega _\stackrel{}{k}(t)\overline{Q}_\stackrel{}{k})].$$
Then
$$E(t)=\underset{\stackrel{}{k}^n}{}ฯต_\stackrel{}{k}(t)(a_\stackrel{}{k}^+(t)a_\stackrel{}{k}(t)+b_\stackrel{}{k}^+(t)b_\stackrel{}{k}(t))$$
$$\rho (t)=\underset{\stackrel{}{k}^n}{}(a_\stackrel{}{k}^+(t)a_\stackrel{}{k}(t)b_\stackrel{}{k}^+(t)b_\stackrel{}{k}(t)).$$
We construct the vacuum state at $`t`$ time.
Letโs consider $`\varphi _\stackrel{}{k}^{0,0}(Q_\stackrel{}{k},\overline{Q}_\stackrel{}{k},t)=\sqrt{\frac{\omega _\stackrel{}{k}(t)}{\pi \mathrm{}}}e^{\frac{\omega _\stackrel{}{k}(t)}{2\mathrm{}}(Q_\stackrel{}{k}^2+\overline{Q}_\stackrel{}{k}^2)}`$, then the vacumm state at $`t`$ time, $`|0>(t)`$, is
$$|0>(t)=\underset{\stackrel{}{k}^n}{}\varphi _\stackrel{}{k}^{0,0}(Q_\stackrel{}{k},\overline{Q}_\stackrel{}{k},t),$$
because
$$E(t)|0>(t)=0\rho (t)|0>(t)=0.$$
Starting from this state we will define all the others. In fact,
the state $`|1_\stackrel{}{k}^+>(t)=a_\stackrel{}{k}^+(t)|0>(t)`$, verifies
$$E(t)|1_\stackrel{}{k}^+>(t)=ฯต_\stackrel{}{k}(t)|1_\stackrel{}{k}^+>(t)\rho (t)|1_\stackrel{}{k}^+>(t)=|1_\stackrel{}{k}^+>(t),$$
consequently, $`|1_\stackrel{}{k}^+>(t)`$ is the state of a particle with energy $`ฯต_\stackrel{}{k}(t)`$ at $`t`$ time.
The state $`|1_\stackrel{}{k}^{}>(t)=b_\stackrel{}{k}^+(t)|0>(t)`$, verifies
$$E(t)|1_\stackrel{}{k}^{}>(t)=ฯต_\stackrel{}{k}(t)|1_\stackrel{}{k}^{}>(t)\rho (t)|1_\stackrel{}{k}^{}>(t)=|1_\stackrel{}{k}^{}>(t),$$
consequently, $`|1_\stackrel{}{k}^{}>(t)`$ is the state of an antiparticle with energy $`ฯต_\stackrel{}{k}(t)`$ at $`t`$ time.
In general, we consider series
$`\{n_\stackrel{}{k}\}:\begin{array}{ccc}^n& & \\ \stackrel{}{k}& & n_\stackrel{}{k}\end{array}`$
and let
$$|\{n_\stackrel{}{k}\};\{m_\stackrel{}{k}\}>(t)=\underset{\stackrel{}{k}^n}{}\frac{(a_\stackrel{}{k}^+(t))^{n_\stackrel{}{k}}}{\sqrt{n_\stackrel{}{k}!}}\frac{(b_\stackrel{}{k}^+(t))^{m_\stackrel{}{k}}}{\sqrt{m_\stackrel{}{k}!}}|0>(t).$$
Then $`|\{n_\stackrel{}{k}\};\{m_\stackrel{}{k}\}>(t)`$, verifies
$$E(t)|\{n_\stackrel{}{k}\};\{m_\stackrel{}{k}\}>(t)=\underset{\stackrel{}{l}^n}{}ฯต_\stackrel{}{l}(t)(n_\stackrel{}{l}+m_\stackrel{}{l})|\{n_\stackrel{}{k}\};\{m_\stackrel{}{k}\}>(t)$$
$$\rho (t)|\{n_\stackrel{}{k}\};\{m_\stackrel{}{k}\}>(t)=\underset{\stackrel{}{l}^n}{}(n_\stackrel{}{l}m_\stackrel{}{l})|\{n_\stackrel{}{k}\};\{m_\stackrel{}{k}\}>(t).$$
Consequently, the state $`|\{n_\stackrel{}{k}\};\{m_\stackrel{}{k}\}>(t)`$ contains, at $`t`$ time, $`n_\stackrel{}{k}`$ particles and $`m_\stackrel{}{k}`$ antiparticles with energy $`ฯต_\stackrel{}{k}(t)`$, for each $`\stackrel{}{k}^n`$.
## 3. The counter-example
### 3.1. Quantum dynamic
First, we study the case $`\stackrel{}{f}(t)\stackrel{}{0}`$, then
$$E=\frac{1}{2}\left[\underset{\stackrel{}{k}^n}{}(\mathrm{}^2_{Q_\stackrel{}{k}}^2+\omega _\stackrel{}{k}^2Q_\stackrel{}{k}^2)+(\mathrm{}^2_{\overline{Q}_\stackrel{}{k}}^2+\omega _\stackrel{}{k}^2\overline{Q}_\stackrel{}{k}^2)\right]\underset{\stackrel{}{k}^n}{}\omega _\stackrel{}{k},$$
where $`\omega _\stackrel{}{k}=\frac{\sqrt{|\mathrm{}\stackrel{}{k}|^2+1}}{\mathrm{}}`$. Notice that the energy does not depend on time, then the eigenvalues $`|\{n_\stackrel{}{k}\};\{m_\stackrel{}{k}\}>(t)|\{n_\stackrel{}{k}\};\{m_\stackrel{}{k}\}>`$ do not depend on time. Therefore, the solution of the problem
$`\{\begin{array}{ccc}i\mathrm{}_t|\mathrm{\Psi }>& =& E|\mathrm{\Psi }>\\ |\mathrm{\Psi }>(0)& =& |\{n_\stackrel{}{k}\};\{m_\stackrel{}{k}\}>,\end{array}`$
is
$$T_q^t|\{n_\stackrel{}{k}\};\{m_\stackrel{}{k}\}>=e^{\frac{i}{\mathrm{}}_{\stackrel{}{l}^n}ฯต_\stackrel{}{l}(n_\stackrel{}{l}+m_\stackrel{}{l})t}|\{n_\stackrel{}{k}\};\{m_\stackrel{}{k}\}>.$$
In particular, $`T_q^t|0>=|0>`$, i.e., when $`\stackrel{}{f}(t)0`$, the vacuum state is invariant for the quantum dynamic, and there is no creation and anihilation particle-antiparticle pairs.
We now study the vacuum dynamic when $`\stackrel{}{f}(t)\stackrel{}{0}`$. Let $`T_q^t|0>(0)`$ be the solution of the problem
$`\{\begin{array}{ccc}i\mathrm{}_t|\mathrm{\Psi }>& =& E(t)|\mathrm{\Psi }>\\ |\mathrm{\Psi }>(0)& =& |0>(0),\end{array}`$
then $`T_q^t|0>(0)=_{\stackrel{}{k}^n}T_{\mathrm{}}^t\varphi _\stackrel{}{k}^{0,0}(Q_\stackrel{}{k},\overline{Q}_\stackrel{}{k},0)`$, where $`T_{\mathrm{}}^t\varphi _\stackrel{}{k}^{0,0}(Q_\stackrel{}{k},\overline{Q}_\stackrel{}{k},0)`$ is the solution of problem
(7) $`\{\begin{array}{ccc}i\mathrm{}_t\varphi & =& [\frac{1}{2}(\mathrm{}^2_{Q_\stackrel{}{k}}^2+\omega _\stackrel{}{k}^2(t)Q_\stackrel{}{k}^2\mathrm{}^2_{\overline{Q}_\stackrel{}{k}}^2+\omega _\stackrel{}{k}^2(t)\overline{Q}_\stackrel{}{k}^2)\omega _\stackrel{}{k}(t)]\varphi \\ \varphi (0)& =& \varphi _\stackrel{}{k}^{0,0}(Q_\stackrel{}{k},\overline{Q}_\stackrel{}{k},0).\end{array}`$
Denote by $`P_{\mathrm{}}^0(t)=|(t)<0|T_q^t|0>(0)|^2`$, the probability that it does not exist any particle-antiparticle pair at $`t`$ time.
Then, we have the
###### Theorem 3.1.
Let $`n`$ be the dimension of the space and suppose that $`\stackrel{}{f}๐_0^{\mathrm{}}(0,T)`$, then:
If $`n=1`$ we have
$$\underset{\mathrm{}0}{lim}P_{\mathrm{}}^0(t)=1t.$$
If $`n=2`$ or $`3`$, at $`t`$ time such that $`\dot{\stackrel{}{f}}(t)\stackrel{}{0}`$, we have
$$\underset{\mathrm{}0}{lim}P_{\mathrm{}}^0(t)1.$$
Consequently, in the case $`n=2`$ or $`3`$, at $`t`$ time such that $`\dot{\stackrel{}{f}}(t)\stackrel{}{0}`$, we do not obtain the classical limit.
###### Remark.
In dimension $`1`$ the result is valid for periodic potentials, i.e., for $`f(x,t)=_{k=0}^N[f_k(t)\mathrm{sin}(kx)+g_k(t)\mathrm{cos}(kx)]`$, we have $`lim_\mathrm{}0P_{\mathrm{}}^0(t)=1`$.
###### Remark.
Le Theorem 3.1 is valid for the Diracโs field.
## 4. Proofs
To make the proof of theorem we need the following
###### Lemma 4.1.
The solution of the problem (7) is
$$T_{\mathrm{}}^t\varphi _\stackrel{}{k}^{0,0}(0)=A_\stackrel{}{k}(t)\varphi _\stackrel{}{k}^{0,0}(t)+\left(\frac{i\mathrm{}\dot{ฯต}_\stackrel{}{k}(t)}{4ฯต_\stackrel{}{k}^2(t)}+\mathrm{}^2B_\stackrel{}{k}(t)\right)\varphi _\stackrel{}{k}^{1,1}(t)+\mathrm{}^2\gamma _\stackrel{}{k}(t),$$
with
$`|1|A_\stackrel{}{k}(t)|^2|\frac{\mathrm{}^2K}{ฯต_\stackrel{}{k}^4}`$.
$`|B_\stackrel{}{k}(t)|^2,\gamma _\stackrel{}{k}(t)_2^2\frac{K}{ฯต_\stackrel{}{k}^4}`$; $`\gamma _\stackrel{}{k}(t)\varphi _\stackrel{}{k}^{0,0}(t),\varphi _\stackrel{}{k}^{1,1}(t)`$.
Where,
$`\varphi _\stackrel{}{k}^{1,1}(t)=a_\stackrel{}{k}^+(t)b_\stackrel{}{k}^+(t)\varphi _\stackrel{}{k}^{0,0}(t)`$.
$`K`$ is a constant independent on $`\stackrel{}{k}`$, $`\mathrm{}`$ and $`t`$.
$`ฯต_\stackrel{}{k}=\sqrt{|\mathrm{}\stackrel{}{k}|^2+1}`$.
With this lemma we can make the
Proof of Theorem 3.1:
If $`n=1`$, $`P_{\mathrm{}}^0(t)=_k|A_k(t)|^2`$. We write, $`|A_k(t)|^2=1+\overline{A}_k(t)`$, then
$`P_{\mathrm{}}^0(t)=1+{\displaystyle \frac{1}{1!}}{\displaystyle \underset{k}{}}\overline{A}_k(t)+{\displaystyle \frac{1}{2!}}{\displaystyle \underset{\begin{array}{c}k_1,k_2\\ k_1k_2\end{array}}{}}\overline{A}_{k_1}(t)\overline{A}_{k_2}(t)+{\displaystyle \frac{1}{3!}}{\displaystyle \underset{\begin{array}{c}k_1,k_2,k_3\\ k_ik_j;ij\\ i,j=1,2,3\end{array}}{}}\overline{A}_{k_1}(t)\overline{A}_{k_2}(t)\overline{A}_{k_3}(t)+\mathrm{}`$
We bound
$`|P_{\mathrm{}}^0(t)1|{\displaystyle \frac{1}{1!}}{\displaystyle \underset{k}{}}|\overline{A}_k(t)|+{\displaystyle \frac{1}{2!}}{\displaystyle \underset{\begin{array}{c}k_1,k_2\\ k_1k_2\end{array}}{}}|\overline{A}_{k_1}(t)||\overline{A}_{k_2}(t)|+\mathrm{}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{n!}}\left({\displaystyle \underset{k}{}}|\overline{A}_k(t)|\right)^n`$
$`{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{n!}}\left(K\mathrm{}{\displaystyle \underset{k}{}}{\displaystyle \frac{\mathrm{}}{ฯต_k^4}}\right)^n{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{n!}}\left(K\mathrm{}{\displaystyle \underset{k}{}}{\displaystyle \frac{\mathrm{}}{ฯต_k^2}}\right)^n,`$
since $`_k\frac{\mathrm{}}{ฯต_k^2}_{}\frac{dx}{x^2+1}+\mathrm{}=\pi +\mathrm{}`$, we have
$`|P_{\mathrm{}}^0(t)1|{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{n!}}(K\mathrm{}(\pi +\mathrm{}))^n=e^{K\mathrm{}(\pi +\mathrm{})}1,`$
therefore
$$\underset{\mathrm{}0}{lim}P_{\mathrm{}}^0(t)=1.$$
We now study the case $`n=2`$. The case $`n=3`$ is analogous.
Denote by $`P_{\mathrm{}}^1(t)=_{\stackrel{}{k}^2}|(t)<1_\stackrel{}{k}^+1_\stackrel{}{k}^{}|T_q^t|0>(0)|^2`$, the probability that at $`t`$ time, does exist a particle-antiparticle pair, then
$$P_{\mathrm{}}^1(t)=\underset{\stackrel{}{k}^2}{}\left|\frac{i\mathrm{}\dot{ฯต}_\stackrel{}{k}(t)}{4ฯต_\stackrel{}{k}^2(t)}+\mathrm{}^2B_\stackrel{}{k}(t)\right|^2\underset{\begin{array}{c}\stackrel{}{l}^2\\ \stackrel{}{l}\stackrel{}{k}\end{array}}{}|A_\stackrel{}{l}(t)|^2.$$
We calcule
$`\underset{\mathrm{}0}{lim}\left|P_{\mathrm{}}^1(t){\displaystyle \underset{\stackrel{}{k}^2}{}}{\displaystyle \frac{\mathrm{}^2}{16}}{\displaystyle \frac{\dot{ฯต}_\stackrel{}{k}^2(t)}{ฯต_\stackrel{}{k}^4(t)}}P_{\mathrm{}}^0(t)\right|\underset{\mathrm{}0}{lim}\left|P_{\mathrm{}}^1(t){\displaystyle \underset{\stackrel{}{k}^2}{}}{\displaystyle \frac{\mathrm{}^2}{16}}{\displaystyle \frac{\dot{ฯต}_\stackrel{}{k}^2(t)}{ฯต_\stackrel{}{k}^4(t)}}{\displaystyle \underset{\begin{array}{c}\stackrel{}{l}^2\\ \stackrel{}{l}\stackrel{}{k}\end{array}}{}}|A_\stackrel{}{l}(t)|^2\right|`$
$`+\underset{\mathrm{}0}{lim}\left|{\displaystyle \underset{\stackrel{}{k}^2}{}}{\displaystyle \frac{\mathrm{}^2}{16}}{\displaystyle \frac{\dot{ฯต}_\stackrel{}{k}^2(t)}{ฯต_\stackrel{}{k}^4(t)}}(1|A_\stackrel{}{k}(t)|^2){\displaystyle \underset{\begin{array}{c}\stackrel{}{l}^2\\ \stackrel{}{l}\stackrel{}{k}\end{array}}{}}|A_\stackrel{}{l}(t)|^2\right|`$
$`\underset{\mathrm{}0}{lim}{\displaystyle \underset{\stackrel{}{k}^2}{}}\left({\displaystyle \frac{\mathrm{}^3}{2}}{\displaystyle \frac{|\dot{ฯต}_\stackrel{}{k}(t)|}{ฯต_\stackrel{}{k}^2(t)}}|B_\stackrel{}{k}(t)|+\mathrm{}^4|B_\stackrel{}{k}(t)|^2+{\displaystyle \frac{\mathrm{}^2}{16}}{\displaystyle \frac{\dot{ฯต}_\stackrel{}{k}^2(t)}{ฯต_\stackrel{}{k}^4(t)}}|1|A_\stackrel{}{k}(t)|^2|\right).`$
Because of the lemma (4.1) and the relation $`ฯต_\stackrel{}{k}^2Cฯต_\stackrel{}{k}^2(t)`$ where $`C=2(1+\stackrel{}{f}_{\mathrm{}}^2)`$, we obtain
$`\underset{\mathrm{}0}{lim}\left|P_{\mathrm{}}^1(t){\displaystyle \underset{\stackrel{}{k}^2}{}}{\displaystyle \frac{\mathrm{}^2}{16}}{\displaystyle \frac{\dot{ฯต}_\stackrel{}{k}^2(t)}{ฯต_\stackrel{}{k}^4(t)}}P_{\mathrm{}}(t)\right|`$
$`\underset{\mathrm{}0}{lim}{\displaystyle \underset{\stackrel{}{k}^2}{}}{\displaystyle \frac{\mathrm{}^2}{ฯต_\stackrel{}{k}^4}}\underset{\mathrm{}0}{lim}\mathrm{}\left(\sqrt{K}C\dot{\stackrel{}{f}}_{\mathrm{}}+K\mathrm{}(\dot{\stackrel{}{f}}_{\mathrm{}}^2+1)\right)=0,`$
because $`lim_\mathrm{}0_{\stackrel{}{k}^2}\frac{\mathrm{}^2}{ฯต_\stackrel{}{k}^4}=_^2\frac{d\stackrel{}{x}^2}{(|\stackrel{}{x}|^2+1)^2}=\pi `$.
Therefore, we have proved that
$$\underset{\mathrm{}0}{lim}P_{\mathrm{}}^1(t)=\underset{\mathrm{}0}{lim}\underset{\stackrel{}{k}^2}{}\frac{\mathrm{}^2}{16}\frac{\dot{ฯต}_\stackrel{}{k}^2(t)}{ฯต_\stackrel{}{k}^4(t)}P_{\mathrm{}}^0(t).$$
With this result, we can prove that for $`n=2`$, if $`\dot{\stackrel{}{f}}(t)\stackrel{}{0}`$, then $`lim_\mathrm{}0P_{\mathrm{}}^0(t)1`$. In fact, we take $`t_0`$ such that $`\dot{\stackrel{}{f}}(t_0)\stackrel{}{0}`$ and we assume that $`lim_\mathrm{}0P_{\mathrm{}}^0(t_0)=1`$. Thus, $`lim_\mathrm{}0P_{\mathrm{}}^1(t_0)=0`$.
However
$`\underset{\mathrm{}0}{lim}P_{\mathrm{}}^1(t_0)=\underset{\mathrm{}0}{lim}{\displaystyle \underset{\stackrel{}{k}^2}{}}{\displaystyle \frac{\mathrm{}^2}{16}}{\displaystyle \frac{\dot{ฯต}_\stackrel{}{k}^2(t_0)}{ฯต_\stackrel{}{k}^4(t_0)}}\underset{\mathrm{}0}{lim}P_{\mathrm{}}^0(t_0)=(\text{for hypothesis})=`$
$`\underset{\mathrm{}0}{lim}{\displaystyle \underset{\stackrel{}{k}^2}{}}{\displaystyle \frac{\mathrm{}^2}{16}}{\displaystyle \frac{\dot{ฯต}_\stackrel{}{k}^2(t_0)}{ฯต_\stackrel{}{k}^4(t_0)}}={\displaystyle \frac{1}{16}}{\displaystyle _^2}{\displaystyle \frac{<\dot{\stackrel{}{f}}(t_0),\stackrel{}{x}>^2}{(|\stackrel{}{x}|^2+1)^3}}๐\stackrel{}{x}^20,`$
because $`\dot{\stackrel{}{f}}(t_0)\stackrel{}{0}`$. Therefore, we have a contradiction and in consequence, $`lim_\mathrm{}0P_{\mathrm{}}^0(t_0)1`$. โ
Now, we make the
Proof of lemma 4.1:
First, we will construct a semi-classical solution of the problem (7). To search a semi-classical solution, we have to consider the functions $`\varphi _\stackrel{}{k}^{s,s}(t)=\frac{(a_\stackrel{}{k}^+(t))^s(b_\stackrel{}{k}^+(t))^s}{s!}\varphi _\stackrel{}{k}^{0,0}(t)`$ with $`s`$.
We write the problem (7) in the form
$`\{\begin{array}{ccc}i\mathrm{}_t\varphi & =& H_\stackrel{}{k}(t)\varphi \\ \varphi (0)& =& \varphi _\stackrel{}{k}^{0,0}(0),\end{array}`$
where $`H_\stackrel{}{k}(t)=ฯต_\stackrel{}{k}(t)(a_\stackrel{}{k}^+(t)a_\stackrel{}{k}(t)+b_\stackrel{}{k}^+(t)b_\stackrel{}{k}(t))`$. We expand the solution in powers serie of $`\mathrm{}`$, in the following form, $`T_{\mathrm{}}^t\varphi _\stackrel{}{k}^{0,0}(0)=_{j,s}\mathrm{}^{s+j}A_{s,\stackrel{}{k}}^j(t)\varphi _\stackrel{}{k}^{s,s}(t)`$. Then, because of following
###### Lemma 4.2.
$$\dot{\varphi }_\stackrel{}{k}^{s,s}(t)=\frac{\dot{ฯต}_\stackrel{}{k}(t)}{2ฯต_\stackrel{}{k}(t)}(s\varphi _\stackrel{}{k}^{s1,s1}(t)(s+1)\varphi _\stackrel{}{k}^{s+1,s+1}(t)),$$
we obtain, after having equalized the powers of $`\mathrm{}`$, the system:
If $`s=0`$
$$\dot{A}_{0,\stackrel{}{k}}^0=0;\dot{A}_{0,\stackrel{}{k}}^j+\frac{\dot{ฯต}_\stackrel{}{k}(t)}{2ฯต_\stackrel{}{k}(t)}A_{1,\stackrel{}{k}}^{j1}=0,\text{ for }j>0.$$
If $`s>0`$
$$i\frac{\dot{ฯต}_\stackrel{}{k}(t)}{2ฯต_\stackrel{}{k}(t)}A_{s1,\stackrel{}{k}}^02ฯต_\stackrel{}{k}(t)A_{s,\stackrel{}{k}}^0=0.$$
$$i\dot{A}_{s,\stackrel{}{k}}^0i\frac{\dot{ฯต}_\stackrel{}{k}(t)}{2ฯต_\stackrel{}{k}(t)}sA_{s1,\stackrel{}{k}}^12sฯต_\stackrel{}{k}(t)A_{s,\stackrel{}{k}}^1=0.$$
$$i\dot{A}_{s,\stackrel{}{k}}^{j1}+i\frac{\dot{ฯต}_\stackrel{}{k}(t)}{2ฯต_\stackrel{}{k}(t)}\left((s+1)A_{s+1,\stackrel{}{k}}^{j2}sA_{s1,\stackrel{}{k}}^j\right)2sฯต_\stackrel{}{k}(t)A_{s,\stackrel{}{k}}^j=0,\text{ for }j>1.$$
We obtain the solution of the system by recurrence. In fact
$$A_{0,\stackrel{}{k}}^0(t)1;A_{1,\stackrel{}{k}}^0(t)=i\frac{\dot{ฯต}_\stackrel{}{k}(t)}{4ฯต_\stackrel{}{k}^2(t)};A_{0,\stackrel{}{k}}^1(t)=_0^ti\frac{\dot{ฯต}_\stackrel{}{k}^2(\tau )}{8ฯต_\stackrel{}{k}^3(\tau )}๐\tau .$$
$$A_{1,\stackrel{}{k}}^1(t)=\frac{1}{2ฯต_\stackrel{}{k}(t)}(i\dot{A}_{1,\stackrel{}{k}}^0i\frac{\dot{ฯต}_\stackrel{}{k}(t)}{2ฯต_\stackrel{}{k}(t)}A_{0,\stackrel{}{k}}^1);A_{2,\stackrel{}{k}}^0(t)=i\frac{\dot{ฯต}_\stackrel{}{k}(t)}{4ฯต_\stackrel{}{k}^2(t)}A_{1,\stackrel{}{k}}^0(t).$$
$$A_{0,\stackrel{}{k}}^2(t)=_0^t\frac{\dot{ฯต}_\stackrel{}{k}(\tau )}{2ฯต_\stackrel{}{k}(\tau )}A_{1,\stackrel{}{k}}^1(\tau )๐\tau ;A_{1,\stackrel{}{k}}^2(t)=\frac{1}{2ฯต_\stackrel{}{k}(t)}(i\dot{A}_{1,\stackrel{}{k}}^1+i\frac{\dot{ฯต}_\stackrel{}{k}(t)}{2ฯต_\stackrel{}{k}(t)}(2A_{2,\stackrel{}{k}}^0A_{0,\stackrel{}{k}}^2))$$
$$A_{3,\stackrel{}{k}}^0(t)=i\frac{\dot{ฯต}_\stackrel{}{k}(t)}{4ฯต_\stackrel{}{k}^2(t)}A_{2,\stackrel{}{k}}^0(t);A_{0,\stackrel{}{k}}^3(t)=_0^t\frac{\dot{ฯต}_\stackrel{}{k}(\tau )}{2ฯต_\stackrel{}{k}(\tau )}A_{1,\stackrel{}{k}}^2(\tau )๐\tau $$
$$A_{2,\stackrel{}{k}}^1(t)=\frac{1}{4ฯต_\stackrel{}{k}(t)}(i\dot{A}_{2,\stackrel{}{k}}^0i\frac{\dot{ฯต}_\stackrel{}{k}(t)}{ฯต_\stackrel{}{k}(t)}A_{1,\stackrel{}{k}}^1);etc.$$
With these solutions, and the relation $`ฯต_\stackrel{}{k}^2Cฯต_\stackrel{}{k}^2(t)`$, we obtain the
###### Lemma 4.3.
If $`s,j3`$ we have
$$|A_{s,\stackrel{}{k}}^j(t)|\frac{\overline{C}}{ฯต_\stackrel{}{k}^{2s+j}}\text{ for }s>0;|A_{0,\stackrel{}{k}}^j(t)|\frac{\overline{C}}{ฯต_\stackrel{}{k}^{2+j}}\text{ for }j>0$$
$$|\dot{A}_{s,\stackrel{}{k}}^j(t)|\frac{g(t)}{ฯต_\stackrel{}{k}^{2s+j}}\text{ for }s>0;|\dot{A}_{0,\stackrel{}{k}}^j(t)|\frac{g(t)}{ฯต_\stackrel{}{k}^{2+j}}\text{ for }j>0,$$
where $`\overline{C}`$ is a constant independent on $`\stackrel{}{k}`$, and $`g(t)๐_0^{\mathrm{}}(0,T)`$ is a function independent on $`\stackrel{}{k}`$.
Now, we show that the function
$`\overline{\varphi }_\stackrel{}{k}(t)=(A_{0,\stackrel{}{k}}^0+\mathrm{}A_{0,\stackrel{}{k}}^1+\mathrm{}^2A_{0,\stackrel{}{k}}^2+\mathrm{}^3A_{0,\stackrel{}{k}}^3)\varphi _\stackrel{}{k}^{0,0}(t)+(\mathrm{}A_{1,\stackrel{}{k}}^0+\mathrm{}^2A_{1,\stackrel{}{k}}^1`$
$`+\mathrm{}^3A_{1,\stackrel{}{k}}^2)\varphi _\stackrel{}{k}^{1,1}(t)+(\mathrm{}^2A_{2,\stackrel{}{k}}^0+\mathrm{}^3A_{2,\stackrel{}{k}}^1)\varphi _\stackrel{}{k}^{2,2}(t)+\mathrm{}^3A_{3,\stackrel{}{k}}^0\varphi _\stackrel{}{k}^{3,3}(t)`$
is a semi-classical solution. In fact, we calcule
$`(i\mathrm{}_tH_\stackrel{}{k})\overline{\varphi }_\stackrel{}{k}(t)`$ $`=`$ $`2i\mathrm{}^4{\displaystyle \frac{\dot{ฯต}_\stackrel{}{k}(t)}{ฯต_\stackrel{}{k}(t)}}A_{3,\stackrel{}{k}}^0\varphi _\stackrel{}{k}^{4,4}(t)+\mathrm{}^4(i\dot{A}_{3,\stackrel{}{k}}^0i{\displaystyle \frac{3\dot{ฯต}_\stackrel{}{k}(t)}{2ฯต_\stackrel{}{k}(t)}}A_{2,\stackrel{}{k}}^1)\varphi _\stackrel{}{k}^{3,3}(t)`$
$`+\mathrm{}^4(i\dot{A}_{2,\stackrel{}{k}}^1+i{\displaystyle \frac{\dot{ฯต}_\stackrel{}{k}(t)}{2ฯต_\stackrel{}{k}(t)}}(3A_{3,\stackrel{}{k}}^02A_{1,\stackrel{}{k}}^2))\varphi _\stackrel{}{k}^{2,2}(t)`$
$`+\mathrm{}^4(i\dot{A}_{1,\stackrel{}{k}}^2+i{\displaystyle \frac{\dot{ฯต}_\stackrel{}{k}(t)}{2ฯต_\stackrel{}{k}(t)}}(2A_{2,\stackrel{}{k}}^1A_{0,\stackrel{}{k}}^3))\varphi _\stackrel{}{k}^{1,1}(t).`$
We deduce from the lemma (4.3), that
$$(i\mathrm{}_tH_\stackrel{}{k})\overline{\varphi }_\stackrel{}{k}(t)_2^2\frac{2\mathrm{}^8}{ฯต_\stackrel{}{k}^8}(3g^2(t)+14C\overline{C}^2|\dot{\stackrel{}{f}}(t)|^2).$$
Furthermore, if using that
$$T_{\mathrm{}}^t\varphi _\stackrel{}{k}^{0,0}(0)\overline{\varphi }_\stackrel{}{k}(t)_2\frac{1}{\mathrm{}}_0^t(i\mathrm{}_\tau H_\stackrel{}{k}(\tau ))\overline{\varphi }_\stackrel{}{k}(\tau )_2๐\tau ,$$
we obtain
$`T_{\mathrm{}}^t\varphi _\stackrel{}{k}^{0,0}(0)\overline{\varphi }_\stackrel{}{k}(t)_2{\displaystyle \frac{\sqrt{2}\mathrm{}^3}{ฯต_\stackrel{}{k}^4}}{\displaystyle _0^t}\sqrt{3g^2(\tau )+14C\overline{C}^2|\dot{\stackrel{}{f}}(\tau )|^2}๐\tau `$
$`{\displaystyle \frac{\sqrt{2}\mathrm{}^3}{ฯต_\stackrel{}{k}^4}}{\displaystyle _0^T}\sqrt{3g^2(\tau )+14C\overline{C}^2|\dot{\stackrel{}{f}}(\tau )|^2}๐\tau {\displaystyle \frac{\mathrm{}^3\stackrel{~}{C}}{ฯต_\stackrel{}{k}^4}}.`$
Therefore, $`T_{\mathrm{}}^t\varphi _\stackrel{}{k}^{0,0}(0)`$ has the form
$`T_{\mathrm{}}^t\varphi _\stackrel{}{k}^{0,0}(0)=(A_{0,\stackrel{}{k}}^0+\mathrm{}A_{0,\stackrel{}{k}}^1+\mathrm{}^2A_{0,\stackrel{}{k}}^2+\mathrm{}^3A_{0,\stackrel{}{k}}^3+\mathrm{}^3F_\stackrel{}{k})\varphi _\stackrel{}{k}^{0,0}(t)+(\mathrm{}A_{1,\stackrel{}{k}}^0+\mathrm{}^2A_{1,\stackrel{}{k}}^1`$
$`+\mathrm{}^3A_{1,\stackrel{}{k}}^2+\mathrm{}^3G_\stackrel{}{k})\varphi _\stackrel{}{k}^{1,1}(t)+(\mathrm{}^2A_{2,\stackrel{}{k}}^0+\mathrm{}^3A_{2,\stackrel{}{k}}^1+\mathrm{}^3I_\stackrel{}{k})\varphi _\stackrel{}{k}^{2,2}(t)+\mathrm{}^3\beta _\stackrel{}{k}(t),`$
with $`|F_\stackrel{}{k}(t)|,|G_\stackrel{}{k}(t)|,|I_\stackrel{}{k}(t)|\frac{\stackrel{~}{C}}{ฯต_\stackrel{}{k}^4}`$, $`\beta _\stackrel{}{k}(t)_2\frac{\stackrel{~}{C}+\overline{C}}{ฯต_\stackrel{}{k}^4}`$ and $`\beta _\stackrel{}{k}(t)\varphi _\stackrel{}{k}^{0,0}(t),\varphi _\stackrel{}{k}^{1,1}(t),\varphi _\stackrel{}{k}^{2,2}(t).`$
Finally, if we take
$$A_\stackrel{}{k}(t)=A_{0,\stackrel{}{k}}^0(t)+\mathrm{}A_{0,\stackrel{}{k}}^1(t)+\mathrm{}^2A_{0,\stackrel{}{k}}^2(t)+\mathrm{}^3A_{0,\stackrel{}{k}}^3(t)+\mathrm{}^3F_\stackrel{}{k}(t)$$
$$B_\stackrel{}{k}(t)=A_{1,\stackrel{}{k}}^1(t)+\mathrm{}A_{1,\stackrel{}{k}}^2(t)+\mathrm{}G_\stackrel{}{k}(t)$$
$$\gamma _\stackrel{}{k}(t)=(A_{2,\stackrel{}{k}}^0(t)+\mathrm{}A_{2,\stackrel{}{k}}^1(t)+\mathrm{}I_\stackrel{}{k}(t))\varphi _\stackrel{}{k}^{2,2}(t)+\mathrm{}\beta _\stackrel{}{k}(t),$$
and $`K=4(1+\overline{C}+\stackrel{~}{C})^2`$, we obtain the proof of lemma 4.1. โ
To finish the work we will make
Another proof of theorem 3.1:
First, we study the case of dimension $`2`$. Since
$`A_{0,\stackrel{}{k}}^2(t)={\displaystyle _0^t}{\displaystyle \frac{\dot{ฯต}_\stackrel{}{k}(\tau )}{4ฯต_\stackrel{}{k}^2(\tau )}}(i\dot{A}_{1,\stackrel{}{k}}^0i{\displaystyle \frac{\dot{ฯต}_\stackrel{}{k}(t)}{2ฯต_\stackrel{}{k}(t)}}A_{0,\stackrel{}{k}}^1)๐\tau ={\displaystyle \frac{\dot{ฯต}_\stackrel{}{k}^2(t)}{32ฯต_\stackrel{}{k}^4(t)}}{\displaystyle \frac{1}{2}}\left({\displaystyle _0^t}{\displaystyle \frac{\dot{ฯต}_\stackrel{}{k}^2(\tau )}{8ฯต_\stackrel{}{k}^3(\tau )}}๐\tau \right)^2,`$
and $`A_{0,\stackrel{}{k}}^3(t)`$ is imaginary, we have
$$|A_\stackrel{}{k}(t)|^2=1\mathrm{}^2\frac{\dot{ฯต}_\stackrel{}{k}^2(t)}{16ฯต_\stackrel{}{k}^4(t)}+h^4J_\stackrel{}{k}(t),$$
with $`|J_\stackrel{}{k}(t)|\frac{\overline{K}}{ฯต_\stackrel{}{k}^4}`$, where $`\overline{K}`$ is a constant independent on $`\stackrel{}{k}`$ and $`\mathrm{}`$.
Starting form this relation, we have
$$\underset{\mathrm{}0}{lim}P_{\mathrm{}}^0(t)=\underset{\mathrm{}0}{lim}\underset{\stackrel{}{k}^2}{}|A_\stackrel{}{k}(t)|^2=\underset{\mathrm{}0}{lim}\underset{\stackrel{}{k}^2}{}\left(1\mathrm{}^2\frac{\dot{ฯต}_\stackrel{}{k}^2(t)}{16ฯต_\stackrel{}{k}^4(t)}\right).$$
We now calcule
$`{\displaystyle \underset{\stackrel{}{k}^2}{}}\left(1\mathrm{}^2{\displaystyle \frac{\dot{ฯต}_\stackrel{}{k}^2(t)}{16ฯต_\stackrel{}{k}^4(t)}}\right)=1{\displaystyle \frac{\mathrm{}^2}{1!}}{\displaystyle \underset{\stackrel{}{k}^2}{}}{\displaystyle \frac{\dot{ฯต}_\stackrel{}{k}^2(t)}{16ฯต_\stackrel{}{k}^4(t)}}+{\displaystyle \frac{\mathrm{}^4}{2!}}{\displaystyle \underset{\begin{array}{c}\stackrel{}{k}_1,\stackrel{}{k}_2^2\\ \stackrel{}{k}_1\stackrel{}{k}_2\end{array}}{}}{\displaystyle \frac{\dot{ฯต}_{\stackrel{}{k}_1}^2(t)}{16ฯต_{\stackrel{}{k}_1}^4(t)}}{\displaystyle \frac{\dot{ฯต}_{\stackrel{}{k}_2}^2(t)}{16ฯต_{\stackrel{}{k}_2}^4(t)}}\mathrm{}`$
To make this calcul we will use the following
###### Lemma 4.4.
If $`n2`$ and $`f_\stackrel{}{k}0\stackrel{}{k}^n`$, then
$`\left({\displaystyle \underset{\stackrel{}{k}}{}}f_\stackrel{}{k}\right)^n{\displaystyle \underset{\begin{array}{c}\stackrel{}{k}_1,\mathrm{},\stackrel{}{k}_n\\ \stackrel{}{k}_i\stackrel{}{k}_j,\text{if }ij\end{array}}{}}f_{\stackrel{}{k}_1}\mathrm{}f_{\stackrel{}{k}_n}{\displaystyle \frac{n(n1)}{2}}\left({\displaystyle \underset{\stackrel{}{k}}{}}f_\stackrel{}{k}\right)^{n2}{\displaystyle \underset{\stackrel{}{k}}{}}f_\stackrel{}{k}^2,`$
consequently,
$`\left|{\displaystyle \underset{\stackrel{}{k}^2}{}}\left(1\mathrm{}^2{\displaystyle \frac{\dot{ฯต}_\stackrel{}{k}^2(t)}{16ฯต_\stackrel{}{k}^4(t)}}\right){\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{n!}}\left(\mathrm{}^2{\displaystyle \underset{\stackrel{}{k}^2}{}}{\displaystyle \frac{\dot{ฯต}_\stackrel{}{k}^2(t)}{16ฯต_\stackrel{}{k}^4(t)}}\right)^n\right|`$
$`{\displaystyle \underset{n=2}{\overset{\mathrm{}}{}}}{\displaystyle \frac{n(n1)}{2}}\left({\displaystyle \underset{\stackrel{}{k}^2}{}}{\displaystyle \frac{\mathrm{}^2\dot{ฯต}_\stackrel{}{k}^2(t)}{16ฯต_\stackrel{}{k}^4(t)}}\right)^{n2}{\displaystyle \underset{\stackrel{}{k}^2}{}}{\displaystyle \frac{\mathrm{}^4\dot{ฯต}_\stackrel{}{k}^4(t)}{16^2ฯต_\stackrel{}{k}^8(t)}}{\displaystyle \frac{1}{n!}}`$
$`{\displaystyle \frac{\mathrm{}^2\dot{\stackrel{}{f}}_{\mathrm{}}^2}{32}}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{(n1)!}}\left({\displaystyle \underset{\stackrel{}{k}^2}{}}{\displaystyle \frac{\mathrm{}^2\dot{ฯต}_\stackrel{}{k}^2(t)}{16ฯต_\stackrel{}{k}^4(t)}}\right)^n.`$
We use that, $`lim_\mathrm{}0_{\stackrel{}{k}^2}\frac{\mathrm{}^2\dot{ฯต}_\stackrel{}{k}^2(t)}{16ฯต_\stackrel{}{k}^4(t)}=\frac{1}{16}_^2\frac{<\dot{\stackrel{}{f}}(t),\stackrel{}{x}>^2}{(|\stackrel{}{x}|^2+1)^3}๐\stackrel{}{x}^2`$ and $`_{n=1}^{\mathrm{}}\frac{x^n}{(n1)!}=xe^x`$, then we obtain
$`\underset{\mathrm{}0}{lim}\left|{\displaystyle \underset{\stackrel{}{k}^2}{}}\left(1\mathrm{}^2{\displaystyle \frac{\dot{ฯต}_\stackrel{}{k}^2(t)}{16ฯต_\stackrel{}{k}^4(t)}}\right){\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}\left(\mathrm{}^2{\displaystyle \underset{\stackrel{}{k}^2}{}}{\displaystyle \frac{\dot{ฯต}_\stackrel{}{k}^2(t)}{16ฯต_\stackrel{}{k}^4(t)}}\right)^n\right|`$
$`\underset{\mathrm{}0}{lim}{\displaystyle \frac{\mathrm{}^2\dot{\stackrel{}{f}}_{\mathrm{}}^2}{32}}{\displaystyle \frac{1}{16}}{\displaystyle _^2}{\displaystyle \frac{<\dot{\stackrel{}{f}}(t),\stackrel{}{x}>^2}{(|\stackrel{}{x}|^2+1)^3}}๐\stackrel{}{x}^2e^{\frac{1}{16}_^2\frac{<\dot{\stackrel{}{f}}(t),\stackrel{}{x}>^2}{(|\stackrel{}{x}|^2+1)^3}๐\stackrel{}{x}^2}=0.`$
By virtue of this result, we have
$`\underset{\mathrm{}0}{lim}P_{\mathrm{}}^0(t)=\underset{\mathrm{}0}{lim}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{n!}}\left(\mathrm{}^2{\displaystyle \underset{\stackrel{}{k}^2}{}}{\displaystyle \frac{\dot{ฯต}_\stackrel{}{k}^2(t)}{16ฯต_\stackrel{}{k}^4(t)}}\right)^n=e^{\frac{1}{16}_^2\frac{<\dot{\stackrel{}{f}}(t),\stackrel{}{x}>^2}{(|\stackrel{}{x}|^2+1)^3}๐\stackrel{}{x}^2}.`$
Therefore, $`lim_\mathrm{}0P_{\mathrm{}}(t)<1`$ if $`\dot{\stackrel{}{f}}(t)\stackrel{}{0}`$.
Now, it is easy to calcule $`lim_\mathrm{}0P_{\mathrm{}}^1(t)`$. In fact, in the first proof of theorem 3.1, we have obtained
$$\underset{\mathrm{}0}{lim}P_{\mathrm{}}^1(t)=\underset{\mathrm{}0}{lim}\mathrm{}^2\underset{\stackrel{}{k}^2}{}\frac{\dot{ฯต}_\stackrel{}{k}^2(t)}{16ฯต_\stackrel{}{k}^4(t)}P_{\mathrm{}}^0(t),$$
then
$$\underset{\mathrm{}0}{lim}P_{\mathrm{}}^1(t)=\frac{1}{16}_^2\frac{<\dot{\stackrel{}{f}}(t),\stackrel{}{x}>^2}{(|\stackrel{}{x}|^2+1)^3}๐\stackrel{}{x}^2e^{\frac{1}{16}_^2\frac{<\dot{\stackrel{}{f}}(t),\stackrel{}{x}>^2}{(|\stackrel{}{x}|^2+1)^3}๐\stackrel{}{x}^2}.$$
In general, let
$$P_{\mathrm{}}^n(t)=\frac{1}{n!}\underset{\begin{array}{c}\stackrel{}{k}_1,\mathrm{},\stackrel{}{k}_n^2\\ \stackrel{}{k}_i\stackrel{}{k}_j,\text{if }ij\end{array}}{}|(t)<1_{\stackrel{}{k}_1}^+1_{\stackrel{}{k}_1}^{}\mathrm{}1_{\stackrel{}{k}_n}^+1_{\stackrel{}{k}_n}^{}|T_q^t|0>(0)|^2,$$
be the probability, that at $`t`$ time, does exist $`n`$ particle-antiparticle pairs. Then we have
$$\underset{\mathrm{}0}{lim}P_{\mathrm{}}^n(t)=\frac{1}{n!}\left(\frac{1}{16}_^2\frac{<\dot{\stackrel{}{f}}(t),\stackrel{}{x}>^2}{(|\stackrel{}{x}|^2+1)^3}๐\stackrel{}{x}^2\right)^ne^{\frac{1}{16}_^2\frac{<\dot{\stackrel{}{f}}(t),\stackrel{}{x}>^2}{(|\stackrel{}{x}|^2+1)^3}๐\stackrel{}{x}^2}.$$
To finish the proof, we have to consider the case of dimension $`3`$. We can prove, proceeding as the case of dimension $`2`$, that
$$\underset{\mathrm{}0}{lim}P_{\mathrm{}}^1(t)=\underset{\mathrm{}0}{lim}\mathrm{}^2\underset{\stackrel{}{k}^3}{}\frac{\dot{ฯต}_\stackrel{}{k}^2(t)}{16ฯต_\stackrel{}{k}^4(t)}P_{\mathrm{}}^0(t).$$
However, $`lim_\mathrm{}0\mathrm{}^2_{\stackrel{}{k}^3}\frac{\dot{ฯต}_\stackrel{}{k}^2(t)}{16ฯต_\stackrel{}{k}^4(t)}=\mathrm{}`$ if $`\dot{\stackrel{}{f}}(t)\stackrel{}{0}`$, whence we conclude that, $`lim_\mathrm{}0P_{\mathrm{}}^0(t)=0`$ if $`\dot{\stackrel{}{f}}(t)\stackrel{}{0}`$.
Another proof of last result, is the following
$$\underset{\mathrm{}0}{lim}P_{\mathrm{}}^0(t)=\underset{\mathrm{}0}{lim}\underset{\stackrel{}{k}^3}{}\left(1\mathrm{}^2\frac{\dot{ฯต}_\stackrel{}{k}^2(t)}{16ฯต_\stackrel{}{k}^4(t)}\right).$$
Since $`_{\stackrel{}{k}^3}\left(1\mathrm{}^2\frac{\dot{ฯต}_\stackrel{}{k}^2(t)}{16ฯต_\stackrel{}{k}^4(t)}\right)_{\stackrel{}{k}^3}\left(1\frac{L\mathrm{}^3\dot{ฯต}_\stackrel{}{k}^2(t)}{16ฯต_\stackrel{}{k}^4(t)}\right)`$ if $`L\mathrm{}1`$, we obtain
$`\underset{\mathrm{}0}{lim}P_{\mathrm{}}^0(t)\underset{L\mathrm{}}{lim}\underset{\mathrm{}0}{lim}{\displaystyle \underset{\stackrel{}{k}^3}{}}\left(1{\displaystyle \frac{L\mathrm{}^3\dot{ฯต}_\stackrel{}{k}^2(t)}{16ฯต_\stackrel{}{k}^4(t)}}\right)=\underset{L\mathrm{}}{lim}e^{\frac{L}{16}_^3\frac{<\dot{\stackrel{}{f}}(t),\stackrel{}{x}>^2}{(|\stackrel{}{x}|^2+1)^3}๐\stackrel{}{x}^2}`$
$`=\{\begin{array}{ccc}0& \text{if}& \dot{\stackrel{}{f}}(t)\stackrel{}{0}\\ 1& \text{if}& \dot{\stackrel{}{f}}(t)=\stackrel{}{0}.\end{array}`$
## 5. References
1. J.D.BJORKEN and S.D.DRELL, Relativistic Quantum Fields; McGraw-Hill Book Co., New York (1965).
2. P.A.M.DIRAC, The principles of quantum mechanics; Oxford University Press (1958).
3. S.A.FULLING, Aspects of Quantum Field Theory in Curved Spacetime; Cambridge University Press (1989).
4. W.GREINER, B.MรLLER, J.RAFELSKI, Quantum Electrodynamics of Strong Fields; Springer-Verlang (1985).
5. G.A.HAGEDORN, Semiclassical quantum mechanics I. The $`\mathrm{}0`$ limit for coherent states; Comm. Math. Phys. 71, no. 1, pag. 77-93, (1980).
6. J.HARO, El lรญmit clร ssic de la mecร nica quร ntica; Tesi Doctoral, U.A.B. (1997).
7. J.HARO, รtude classique de lโรฉquation de Dirac; Ann. Fond. Louis de Broglie 23, no. 3-4, pag. 166-172, (1998).
8. J.HARTHONG, รtudes sur la mรฉcanique quantique; Asterisque, 111, (1984).
9. V.P.MASLOV and M.V.FEDORIUK, Semi-classical aproximation in quantum mechanics; D.Riedel Publishing Compay, Dordrecht, Holland (1981). |
warning/0002/gr-qc0002073.html | ar5iv | text | # (๐ฟฬ_๐,๐)-spaces. Length of a vector and angle between two vectors
## 1 Introduction
In previous papers , the notions of contravariant and covariant affine connections are considered for contravariant and covariant tensor fields over differentiable manifolds. It was shown that these two different (not only by sign) connections can be introduced by means of changing the canonical definition of the bases of dual vector spaces (respectively of dual vector fields). The deviation operator and its applications to deviation equations over differentiable manifolds with contravariant and covariant affine connections \[the s.c. ($`\overline{L}_n,g`$)-spaces\] are investigated.
In the Einstein theory of gravitation (ETG) kinematic notions related to the notion relative velocity such as shear velocity tensor (shear velocity, shear) $`\sigma `$, rotation velocity tensor (rotation velocity, rotation) $`\omega `$ and expansion velocity (expansion) $`\theta `$, are used in finding solutions of special types of Einsteinโs field equations and in the description of the properties of the (pseudo) Riemannian spaces without torsion ($`V_n`$-spaces). By means of these notions a classification of $`V_n`$-spaces, admitting special types of geodesic vector fields, has been proposed . The same kinematic characteristics are also necessary for description of the projections of the Riemannian (curvature) tensor and the Ricci tensor along a non-isotropic (non-null) vector field , and in obtaining and using the Raychaudhuri identity , in $`V_n`$-spaces.
The kinematic characteristics connected with the notion relative velocity can be generalized for vector fields over differentiable manifolds with contravariant and covariant affine connections and metrics \[($`\overline{L}_n,g`$)-spaces\] so that in the case of ($`L_n,g`$)- and $`V_n`$-spaces \[as a special case of ($`\overline{L}_n,g`$)-spaces\], and for normalized non-isotropic vector fields these characteristics are the same as those introduced in the ETG. In an analogous way as in the case of the kinematic characteristics, related to the notion of relative velocity, it is possible to introduce kinematic characteristics, related to the notion of relative acceleration such as shear acceleration tensor (shear acceleration), rotation acceleration tensor (rotation acceleration) and expansion acceleration .
The corresponding for ($`\overline{L}_n,g`$)-spaces notions of relative velocity and relative acceleration are considered in . By means of these kinematic characteristics, several other types of notions such as shear velocity, shear acceleration, rotation velocity, rotation acceleration, expansion velocity and expansion acceleration are investigated. The connections between the kinematic characteristics related to the relative acceleration and these related to the relative velocity are also found. The auto-parallel vector fields in ($`\overline{L}_n,g`$)-spaces are classified on the basis of the kinematic characteristics. The generalizations compared with those in ($`L_n,g`$)-spaces (differentiable manifold with affine connection and metric) appear only in the explicit forms of the expressions, written in a corresponding basis (or in other words - only in index forms).
In this paper the basic notions of the length of a vector field and the cosine of the angle between two vector fields are considered over $`(\overline{L}_n,g)`$-spaces as well as their changes along an other vector field. These notions are necessary for introduction and consideration of different types of transports in $`(\overline{L}_n,g)`$-spaces such as Fermi-Walker transports and conformal transports .
## 2 Length of a contravariant vector field
The square of the length of a contravariant vector field $`u`$ is determined by means of the covariant metric tensor $`g`$ as
$$u^2=\pm u^2=g(u,u)\text{ , }uT(M)\text{ , }u0\text{ .}$$
(1)
###### Definition 1
The length of a contravariant vector field $`u`$ is the positive square root of the absolute value of the square of the length of this field, i.e.
$$l_u=u=g(u,u)^{\frac{1}{2}}\text{ , }l_u(x)=u_x=u(x)\text{ , }xM\text{ ,}$$
(2)
where $`l_u`$ is the length of the contravariant vector field $`u`$ and $`l_u(x)`$ is the length of the contravariant vector $`u(x)`$ in a point $`xM`$.
With respect to their lengths, the contravariant vector fields can be divided in two classes: null- or isotropic vector fields ($`l_u=0`$) and non-null or non-isotropic vector fields ($`l_u0`$). In the cases of a positive definite covariant metric $`g`$ (Sgn $`g=n`$, $`dimM=n`$) the isotropic vector field is identically equal to zero ($`u0`$, $`u^\alpha 0`$). In the cases of an indefinite covariant metric $`g`$ (Sgn $`g<n`$ or Sgn $`g>n`$, $`dimM=n`$) the isotropic vector field with equal to zero length can have different from zero components in an arbitrary basis, i.e. it is not identically equal to zero in the points, where it has been defined.
The changes of the length of a contravariant vector field under the action of the covariant differential operator is determined on one side, by the action of the covariant operator on a function (here the length $`l_u`$) over the manifold and on the other - by the structure of the length $`l_u`$ itself and by the commutation relation between the covariant operator and the contraction operator
$$\begin{array}{c}_\xi u^2=\pm _\xi (l_u^2)=_\xi [g(u,u)]=\xi [g(u,u)]=\pm \xi (l_u^2)=\\ =\pm 2l_u(\xi l_u)=(_\xi g)(u,u)+2g(_\xi u,u)\text{ ,}\end{array}\text{ }$$
(3)
from where it follows
$$\xi l_u=\pm \frac{1}{2l_u}[(_\xi g)(u,u)+2g(_\xi u,u)]\text{ , }l_u0\text{ .}$$
(4)
In the case of an null (isotropic) contravariant vector field $`u`$, $`g(u,u)=0`$, there is a relation between the covariant derivative $`_\xi g`$ of the covariant metric tensor $`g`$ and the covariant derivative $`_\xi u`$ of the vector field $`u`$
$$(_\xi g)(u,u)=2g(_\xi u,u)\text{ .}$$
(5)
If the contravariant vector field $`u`$ is transported parallel along the contravariant vector field $`\xi `$, the change of the length of $`u`$ obeys the condition
$$\xi l_u=\pm \frac{1}{2l_u}(_\xi g)(u,u)\text{ , }_\xi u=0\text{ , }l_u0\text{ ,}$$
(6)
and if $`l_u=0`$ and $`_\xi u=0`$, then the condition $`(_\xi g)(u,u)=0`$ follows for the covariant metric $`g`$ or for the vector field $`\xi `$.
One of the essential characteristics of the different types of transport is their influence on the change of the length of a contravariant vector field due to the action of the covariant differential operator. In the case of manifolds with affine connections and metric which allow different type of transports of the covariant metric, there are transports under which the length of a contravariant vector field does not change.
The change of the length of a contravariant vector field under the action of the Lie differential operator, i. e. the change of length under draggings-along, can be described in an analogous way as in the cases of transports
$$\begin{array}{c}\mathrm{\pounds }_\xi u^2=\pm \mathrm{\pounds }_\xi u^2=\mathrm{\pounds }_\xi [g(u,u)]=\pm 2l_u(\xi l_u)=\\ =(\mathrm{\pounds }_\xi g)(u,u)+2g(\mathrm{\pounds }_\xi u,u)\text{ ,}\end{array}$$
(7)
from where it follows
$$\begin{array}{c}l_u0:\xi l_u=\pm \frac{1}{2l_u}[(\mathrm{\pounds }_\xi g)(u,u)+2g(\mathrm{\pounds }_\xi u,u)]\text{ ,}\\ l_u=0:(\mathrm{\pounds }_\xi g)(u,u)=2g(\mathrm{\pounds }_\xi u,u)\text{ .}\end{array}$$
(8)
For parallel dragging of $`u`$ along $`\xi `$ ($`\mathrm{\pounds }_\xi u=0`$), the change of $`l_u`$ can be written in the form
$$\begin{array}{c}l_u0:\xi l_u=\pm \frac{1}{2l_u}(\mathrm{\pounds }_\xi g)(u,u)\text{ ,}\\ l_u=0:(\mathrm{\pounds }_\xi g)(u,u)=0\text{ .}\end{array}$$
(9)
One of the main characteristics of the different draggings-along of a covariant metric $`g`$ is the change of length of a contravariant vector field under draggings-along. The different types of draggings-along of $`g`$ induce different changes of the length of a given contravariant vector field.
The length change of contravariant vector fields under different types of โtransportโ and โdraggings-alongโ of the covariant metric $`g`$ has been used in mathematical models of physical systems, described by means of differentiable manifolds with affine connections and metric. The change of length has also been used for giving certain properties and characteristics with physical interpretation of these systems.
## 3 Cosine of the angle between two contravariant vector fields
The cosine of the angle between two contravariant vector fields is determined by means of the scalar product of both vector fields and their lengths. From the relation
$$g(u,v)=u.v.\mathrm{cos}(u,v)=l_u.l_v.\mathrm{cos}(u,v)\text{ ,}$$
(10)
the definition for the cosine between the vector fields $`u`$ and $`v`$ can be written as
$$\mathrm{cos}(u,v)=\frac{1}{l_u.l_v}.g(u,v)\text{ , }u,vT(M)\text{ , }l_u0\text{ , }l_v0\text{ .}$$
(11)
From the definition it follows, that two non-isotropic contravariant vector fields $`u`$ and $`v`$ are orthogonal to each other, if their scalar product is equal to zero, i. e.
$$uv:g(u,v)=0\text{ : }\mathrm{cos}(u,v)=0\text{ , }u,vT(M)\text{ .}$$
(12)
The notion cosine between two isotropic contravariant vector fields or between one non-isotropic and one isotropic contravariant vector field cannot be introduced by means of the above definition.
###### Remark 1
If one assumes that (10) is also fulfilled for the scalar product of two isotropic (null) vector fields (or for one non-isotropic and one isotropic vector field), then from (10) and (12) it could be fixed that: (a) two isotropic contravariant vector fields are always orthogonal to each other; (b) every isotropic contravariant vector field is orthogonal to each non-isotropic contravariant vector field. The last statement leads to the condition $`g(u,v)=0`$, ($`l_u=0`$, $`l_v0`$), which is not fulfilled in general. In the case (a) the value of $`\mathrm{cos}(u,v)`$ as a value of a free determined limited function can be fixed to zero, i. e. $`\mathrm{cos}(u,v)=0`$, $`l_u=0`$, $`l_v=0`$. This value corresponds to the notion orthogonal contravariant non-isotropic (non-null) vector fields. In this way, by definition, the cosine between two isotropic vector fields is equal to zero. On the other side, every isotropic contravariant vector field can be presented as a co-linear to other given isotropic vector field, i.e. if $`uT(M)`$, $`l_u=0`$, $`vT(M)`$, $`l_v=0`$ with $`u=\kappa .v`$, $`\kappa C^r(M)`$. Therefore, under the introduced notion of orthogonality of two isotropic contravariant vector fields the following statement is valid: every two co-linear isotropic contravariant vector fields are orthogonal to each other.
The cosine change of the angle between two non-isotropic vector fields under the action of the covariant differential operator is determined on the grounds of the definition of the notion cosine and the commutation relation between the covariant and contraction operator
$$\begin{array}{c}_\xi \mathrm{cos}(u,v)=\xi [\mathrm{cos}(u,v)]=\frac{1}{l_u.l_v}[(_\xi g)(u,v)+g(_\xi u,v)+g(u,_\xi v)]\\ [\frac{1}{l_u}(\xi l_u)+\frac{1}{l_v}(\xi l_v)].\mathrm{cos}(u,v)=\\ =[\xi (\mathrm{log}l_u)+\xi (\mathrm{log}l_v)].\mathrm{cos}(u,v)+\\ +\frac{1}{l_u.l_v}[(_\xi g)(u,v)+g(_\xi u,v)+g(u,_\xi v)]=\\ =[\xi (\mathrm{log}(l_u.l_v))].\mathrm{cos}(u,v)+\\ +\frac{1}{l_u.l_v}[(_\xi g)(u,v)+g(_\xi u,v)+g(u,_\xi v)]\text{ .}\end{array}$$
(13)
The last expression can also be written in the form
$$\begin{array}{c}\xi [\mathrm{cos}(u,v)]=\frac{1}{l_u.l_v}[(_\xi g)(u,v)+g(_\xi u,v)+g(u,_\xi v)]\\ \frac{1}{2}\{\pm \frac{1}{l_u^2}[(_\xi g)(u,u)+2g(_\xi u,u)]\pm \frac{1}{l_v^2}[(_\xi g)(v,v)+2g(_\xi v,v)]\}.\mathrm{cos}(u,v)\text{ .}\end{array}$$
(14)
The conditions for transports of the covariant metric $`g`$ and the conditions for transports of the contravariant vector fields as well determined the change of the cosine of the angle between two contravariant vector fields.
Under the action of the Lie differential operator, the change of the angle between two contravariant vector fields is determined by the action of this operator on a function over manifold and by its commutation relation with the contraction operator
$$\begin{array}{c}\xi [\mathrm{cos}(u,v)]=\frac{1}{l_u.l_v}[(\mathrm{\pounds }_\xi g)(u,v)+g(\mathrm{\pounds }_\xi u,v)+g(u,\mathrm{\pounds }_\xi v)]\\ [\frac{1}{l_u}(\xi l_u)+\frac{1}{l_v}(\xi l_v)].\mathrm{cos}(u,v)=\\ =\frac{1}{l_ul_v}[(\mathrm{\pounds }_\xi g)(u,v)+g(\mathrm{\pounds }_\xi u,v)+g(u,\mathrm{\pounds }_\xi v)]\\ [\xi (\mathrm{log}l_u)+\xi (\mathrm{log}l_v)].\mathrm{cos}(u,v)=\\ =\frac{1}{l_u.l_v}[(\mathrm{\pounds }_\xi g)(u,v)+g(\mathrm{\pounds }_\xi u,v)+g(u,\mathrm{\pounds }_\xi v)]\\ \frac{1}{2}\{\pm \frac{1}{l_u^2}[(\mathrm{\pounds }_\xi g)(u,u)+2g(\mathrm{\pounds }_\xi u,u)]\pm \frac{1}{l_v^2}[(\mathrm{\pounds }_\xi g)(v,v)+2g(\mathrm{\pounds }_\xi v,v)]\}.\mathrm{cos}(u,v)\text{ .}\end{array}$$
(15)
Different draggings-along determine the change of the cosine of the angle between two contravariant non-isotropic vector fields dragged along other contravariant vector field.
The change of the cosine of the angle between two contravariant non-isotropic vector fields under transport or dragging- along is a characteristic of great importance for mathematical models, describing physical systems by means of vector fields over differentiable manifolds with contravariant and covariant affine connections. The geometric characteristics of vector fields (length, angle between two vector fields) have been connected with the characteristics of the physical system and in this way the kinematic structure can be determined for describing physical processes.
## 4 Length of a covariant vector field
The square of the length of a covariant vector field is determined as
$$\overline{g}(p,p)=p^2=\pm p^2=\pm l_p^2\text{ , }pT^{}(M)\text{ , }l_p0\text{ .}$$
(16)
By means of this definition the notion length of a covariant vector field is introduced.
###### Definition 2
The length of a covariant vector field $`p`$ is the positive square root of the absolute value of the square of the length of this field, i.e.
$$l_p=p=\overline{g}(p,p)^{\frac{1}{2}}\text{ , }l_p(x)=p_x=p(x)\text{ , }xM\text{ ,}$$
(17)
where $`l_p`$ is the length of the covariant vector field $`p`$ and $`l_p(x)`$ is the length of the covariant vector $`p(x)`$ in a point $`xM`$.
With respect to their lengths the covariant vector fields are also divided in two types: isotropic (null), $`(l_p=0`$, $`l_p(x)=0`$, $`xM)`$, and non-isotropic (non-null), $`(l_p0`$, $`l_p(x)0`$, $`xM)`$, covariant vector fields and vectors. In the cases of definite contravariant metric $`\overline{g}`$, ($`Sgn\overline{g}=n`$, $`dimM=n`$), the isotropic covariant vector filed is identically equal to zero ($`p0`$, $`p_\alpha 0`$). In the cases of indefinite contravariant metric $`\overline{g}`$, ($`Sgn\overline{g}<n`$ or $`Sgn\overline{g}>n`$, $`dimM=n`$) an isotropic (null) covariant vector field would have different from zero components in an arbitrary basis, i. e. it is not equal to zero in the region where it is defined.
The change of the length of a covariant vector field $`p`$ under the action of the covariant differential operator $`_\xi `$ can be found in an analogous way as in the case of contravariant vector field
$$\xi l_p=\pm \frac{1}{2l_p}[(_\xi \overline{g})(p,p)+2\overline{g}(_\xi p,p)]\text{ , }l_p0\text{ .}$$
(18)
The change of the length of a covariant vector field $`p`$ under the action of the Lie differential operator $`\mathrm{\pounds }_\xi `$ is determined by the expression
$$\xi l_p=\pm \frac{1}{2l_p}[(\mathrm{\pounds }_\xi \overline{g})(p,p)+2\overline{g}(_\xi p,p)]\text{ .}$$
(19)
For isotropic (null) covariant vector fields the following relations are fulfilled
$$\begin{array}{c}(_\xi \overline{g})(p,p)=2\overline{g}(_\xi p,p)\text{ , }l_p=0\text{ ,}\\ (\mathrm{\pounds }_\xi \overline{g})(p,p)=2\overline{g}(\mathrm{\pounds }_\xi p,p)\text{ , }l_p=0\text{ .}\end{array}$$
(20)
Different transports and draggings along of the contravariant metric tensor $`\overline{g}`$ induce analogous changes of the covariant vector fields as in the case of contravariant vector fields.
## 5 Cosine of the angle between two covariant vector fields
By means of the scalar product of two covariant vector fields and the length of a covariant vector field, the cosine of the angle between two covariant vector fields is determined by the relation
$$\overline{g}(p,q)=p.q.\mathrm{cos}(p,q)=l_p.l_q.\mathrm{cos}(p,q)\text{ , }p,qT^{}(M)\text{ .}$$
(21)
###### Definition 3
The cosine of the angle between two non-isotropic covariant vector fields $`p`$ and $`q`$ is by definition
$$\mathrm{cos}(p,q)=\frac{1}{l_pl_q}.\overline{g}(p,q)\text{ , }p,qT^{}(M)\text{ , }l_p0\text{ , }l_q0\text{ .}$$
(22)
###### Definition 4
Two covariant vector fields $`p`$ and $`q`$ are orthogonal to each other, when
$$\overline{g}(p,q)=0\text{ .}$$
(23)
If the two covariant vector fields $`p`$ and $`q`$ are null (isotropic) vector fields, then from the last definition (21) it follows that they are orthogonal to each other.
The change of the cosine of the angle between two covariant vector fields under the action of the covariant differential operator $`_\xi `$ can be presented in the form
$$\begin{array}{c}\xi [\mathrm{cos}(p,q)]=\frac{1}{l_pl_q}[(_\xi \overline{g})(p,q)+\overline{g}(_\xi p,q)+\overline{g}(p,_\xi q)]\\ [\xi (\mathrm{log}l_p)+\xi (\mathrm{log}l_q)].\mathrm{cos}(p,q)=\\ =\frac{1}{l_pl_q}[(_\xi \overline{g})(p,q)+\overline{g}(_\xi p,q)+\overline{g}(p,_\xi q)]\\ \frac{1}{2}\{\pm \frac{1}{l_p^2}[(_\xi \overline{g})(p,p)+2\overline{g}(_\xi p,p)]\pm \\ \pm \frac{1}{l_q^2}[(_\xi \overline{g})(q,q)+2\overline{g}(_\xi q,q)]\}.\mathrm{cos}(p,q)\text{ .}\end{array}$$
(24)
The change of the cosine of the angle between two covariant vector fields under the action of the Lie differential operator $`\mathrm{\pounds }_\xi `$ can be found in the form
$$\begin{array}{c}\xi [\mathrm{cos}(p,q)]=\frac{1}{l_pl_q}[(\mathrm{\pounds }_\xi \overline{g})(p,q)+\overline{g}(\mathrm{\pounds }_\xi p,q)+\overline{g}(p,\mathrm{\pounds }_\xi q)]\\ [\xi (\mathrm{log}l_p)+\xi (\mathrm{log}l_q)].\mathrm{cos}(p,q)=\\ =\frac{1}{l_pl_q}[(\mathrm{\pounds }_\xi \overline{g})(p,q)+\overline{g}(\mathrm{\pounds }_\xi p,q)+\overline{g}(p,\mathrm{\pounds }_\xi q)]\\ \frac{1}{2}\{\pm \frac{1}{l_p^2}[(\mathrm{\pounds }_\xi \overline{g})(p,p)+2\overline{g}(\mathrm{\pounds }_\xi p,p)]\pm \\ \pm \frac{1}{l_q^2}[(\mathrm{\pounds }_\xi \overline{g})(q,q)+2\overline{g}(\mathrm{\pounds }_\xi q,q)]\}.\mathrm{cos}(p,q)\text{ .}\end{array}$$
(25)
The cosine changes of the angle between two non-isotropic covariant vector fields under transports or draggings-along of the contravariant metric can be found in analogous way as in the case of the covariant metric.
## 6 Relative velocity and change of the length of a contravariant vector field
Let we now consider the influence of the kinematic characteristics related to the relative velocity upon the change of the length of a contravariant vector field.
Let $`l_\xi =g(\xi ,\xi )^{\frac{1}{2}}`$ be the length of a contravariant vector field $`\xi `$. The rate of change $`ul_\xi `$ of $`l_\xi `$ along a contravariant vector field $`u`$ can be expressed in the form $`\pm \mathrm{\hspace{0.17em}2}.l_\xi .(ul_\xi )=(_ug)(\xi ,\xi )+2g(_u\xi ,\xi )`$. By the use of the projections of $`\xi `$ and $`_u\xi `$ along and orthogonal to $`u`$ (see the chapter about kinematic characteristics and relative velocity) we can find the relations
$$\begin{array}{c}2g(_u\xi ,\xi )=2.\frac{l}{e}.g(_u\xi ,u)+2g(_{rel}v,\xi _{})\text{ ,}\\ (_ug)(\xi ,\xi )=(_ug)(\xi _{},\xi _{})+2.\frac{l}{e}.(_ug)(\xi _{},u)+\frac{l^2}{e^2}.(_ug)(u,u)\text{ .}\end{array}$$
Then, it follows for $`\pm \mathrm{\hspace{0.17em}2}.l_\xi .(ul_\xi )`$ the expression
$$\begin{array}{c}\pm \mathrm{\hspace{0.17em}2}.l_\xi .(ul_\xi )=(_ug)(\xi _{},\xi _{})+2.\frac{l}{e}.[(_ug)(\xi _{},u)+g(_u\xi ,u)]+\\ +\frac{l^2}{e^2}.(_ug)(u,u)+2g(_{rel}v,\xi _{})\text{ ,}\end{array}$$
(26)
where
$$g(_{rel}v,\xi _{})=\frac{l}{e}.h_u(a,\xi _{})+h_u(\mathrm{\pounds }_u\xi ,\xi _{})+d(\xi _{},\xi _{})\text{ ,}$$
(27)
$$d(\xi _{},\xi _{})=\sigma (\xi _{},\xi _{})+\frac{1}{n1}.\theta .l_\xi _{}^2\text{ .}$$
(28)
For finding out the last two expressions the following relations have been used:
$$g(\overline{g}(h_u)a,\xi _{})=h_u(a,\xi _{})\text{ , }g(\overline{g}(h_u)(\mathrm{\pounds }_u\xi ),\xi _{})=h_u(\mathrm{\pounds }_u\xi ,\xi _{})\text{ ,}$$
(29)
$$g(\overline{g}[d(\xi )],\xi _{})=d(\xi _{},\xi _{})\text{ , }d(\xi )=d(\xi _{})\text{ .}$$
(30)
Special case: $`g(u,\xi )=l:=0:\xi =\xi _{}`$.
$$\pm \mathrm{\hspace{0.17em}2}.l_\xi _{}.(ul_\xi _{})=(_ug)(\xi _{},\xi _{})+2g(_{rel}v,\xi _{})\text{ .}$$
(31)
Special case: $`V_n`$-spaces: $`_\eta g=0`$ for $`\eta T(M)`$ ($`g_{ij;k}=0`$), $`g(u,\xi )=l:=0:\xi =\xi _{}`$.
$$\pm l_\xi _{}.(ul_\xi _{})=g(_{rel}v,\xi _{})\text{ .}$$
(32)
In $`(\overline{L}_n,g)`$-spaces as well as in $`(L_n,g)`$-spaces the covariant derivative $`_ug`$ of the metric tensor field $`g`$ along $`u`$ can be decomposed in its trace free part $`{}_{}{}^{s}_{u}^{}g`$ and its trace part $`\frac{1}{n}.Q_u.g`$ as
$$_ug=^s_ug+\frac{1}{n}.Q_u.g\text{ , }dimM=n\text{ ,}$$
where
$$\overline{g}[^s_ug]=0\text{ , }Q_u=\overline{g}[_ug]=g^{\overline{k}\overline{l}}.g_{kl;j}.u^j=Q_j.u^j\text{ , }Q_j=g^{\overline{k}\overline{l}}.g_{kl;j}\text{ .}$$
The covariant vector $`\overline{Q}=\frac{1}{n}.Q=\frac{1}{n}.Q_j.dx^j=\frac{1}{n}.Q_\alpha .e^\alpha `$ is called Weylโs covector field. The operator $`_u=^s_u+\frac{1}{n}.Q_u`$ is called trace free covariant operator.
If we use now the decomposition of $`_ug`$ in the expression for $`\pm \mathrm{\hspace{0.17em}2}.l_\xi .(ul_\xi )`$ we find the relation
$$\begin{array}{c}\pm \mathrm{\hspace{0.17em}2}.l_\xi .(ul_\xi )=(^s_ug)(\xi ,\xi )+\frac{1}{n}.Q_u.l_\xi ^2+2g(_u\xi ,\xi )=\\ =(^s_ug)(\xi _{},\xi _{})+\\ +\frac{l}{e}.[2.(^s_ug)(\xi _{},u)+2.g(_u\xi ,u)+\frac{l}{e}.(^s_ug)(u,u)]+\\ +\frac{1}{n}.Q_u.(l_\xi _{}^2+\frac{l^2}{e})+2.g(_{rel}v,\xi _{})\text{ ,}\end{array}$$
(33)
where $`l_\xi _{}^2=g(\xi _{},\xi _{})`$, $`l=g(\xi ,u)`$.
For $`l_\xi 0:`$
$$ul_\xi =\pm \frac{1}{2.l_\xi }(^s_ug)(\xi ,\xi )\pm \frac{1}{2.n}.Q_u.l_\xi \pm \frac{1}{l_\xi }.g(_u\xi ,\xi )\text{ .}$$
(34)
In the case of a parallel transport ($`_u\xi =0`$) of $`\xi `$ along $`u`$ the change $`ul_\xi `$ of the length $`l_\xi `$ is
$$ul_\xi =\pm \frac{1}{2.l_\xi }(^s_ug)(\xi ,\xi )\pm \frac{1}{2.n}.Q_u.l_\xi \text{ . }$$
(35)
Special case: $`_u\xi =0`$ and $`{}_{}{}^{s}_{u}^{}g=0`$.
$$ul_\xi =\pm \frac{1}{2.n}.Q_u.l_\xi \text{ . }$$
(36)
If $`u=\frac{d}{ds}=u^i._i=(dx^i/ds)._i`$, then
$`l_\xi (s+ds)`$ $``$ $`l_\xi (s)+{\displaystyle \frac{dl_\xi }{ds}}.ds=l_\xi (s)\pm {\displaystyle \frac{1}{2.n}}.Q_u(s).l_\xi (s).ds=`$
$`=`$ $`(1\pm {\displaystyle \frac{1}{2.n}}.Q_u(s).ds).l_\xi (s)=\mathrm{}_u(s).l_\xi (s)\text{ ,}`$
$`\mathrm{}_u(s)`$ $`=`$ $`1\pm {\displaystyle \frac{1}{2.n}}.Q_u(s).ds\text{ .}`$ (37)
Therefore, the rate of change of $`l_\xi `$ along $`u`$ is linear to $`l_\xi `$.
Special case: $`g(u,\xi )=l:=0:\xi =\xi _{}`$.
$$\pm 2.l_\xi _{}.(ul_\xi _{})=(^s_ug)(\xi _{},\xi _{})+\frac{1}{n}.Q_u.l_\xi _{}^2+2.g(_{rel}v,\xi _{})\text{ .}$$
$$ul_\xi _{}=\pm \frac{1}{2.l_\xi _{}}.(^s_ug)(\xi _{},\xi _{})\pm \frac{1}{2n}.Q_u.l_\xi _{}\pm \frac{1}{l_\xi _{}}.g(_{rel}v,\xi _{})\text{ , }l_\xi _{}0\text{ .}$$
(38)
Special case: Quasi-metric transports: $`_ug:=2.g(u,\eta ).g`$, $`u`$, $`\eta T(M)`$.
$$\pm 2.l_\xi .(ul_\xi )=2.g(u,\eta ).(l_\xi _{}^2+\frac{l^2}{e})+2.[\frac{l}{e}.g(_u\xi ,u)+g(_{rel}v,\xi _{})]\text{ .}$$
(39)
## 7 Change of the cosine between two contravariant vector fields and the relative velocity
The cosine between two contravariant vector fields $`\xi `$ and $`\eta `$ has been defined as $`g(\xi ,\eta )=l_\xi .l_\eta .\mathrm{cos}(\xi ,\eta )`$. The rate of change of the cosine along a contravariant vector field $`u`$ can be found in the form
$$\begin{array}{c}l_\xi .l_\eta .\{u[\mathrm{cos}(\xi ,\eta )]\}=(_ug)(\xi ,\eta )+g(_u\xi ,\eta )+g(\xi ,_u\eta )\\ [l_\eta .(ul_\xi )+l_\xi .(ul_\eta )].\mathrm{cos}(\xi ,\eta )\text{ .}\end{array}$$
Special case: $`_u\xi =0`$, $`_u\eta =0`$, $`{}_{}{}^{s}_{u}^{}g=0`$.
$$l_\xi .l_\eta .\{u[\mathrm{cos}(\xi ,\eta )]\}=\frac{1}{n}.Q_u.g(\xi ,\eta )[l_\eta .(ul_\xi )+l_\xi .(ul_\eta )].\mathrm{cos}(\xi ,\eta )\text{ .}$$
Since $`g(\xi ,\eta )=l_\xi .l_\eta .\mathrm{cos}(\xi ,\eta )`$, it follows from the last relation
$$l_\xi .l_\eta .\{u[\mathrm{cos}(\xi ,\eta )]\}=\{\frac{1}{n}.Q_u.l_\xi .l_\eta [l_\eta .(ul_\xi )+l_\xi .(ul_\eta )]\}.\mathrm{cos}(\xi ,\eta )\text{ .}$$
Therefore, if $`\mathrm{cos}(\xi ,\eta )=0`$ between two parallel transported along $`u`$ vector fields $`\xi `$ and $`\eta `$, then the right angle between them \[determined by the condition $`\mathrm{cos}(\xi ,\eta )=0`$\] does not change along the contravariant vector field $`u`$. In the cases, when $`\mathrm{cos}(\xi ,\eta )0`$, the rate of change of the cosine of the angle between two vector fields $`\xi `$ and $`\eta `$ is linear to $`\mathrm{cos}(\xi ,\eta )`$.
By the use of the definitions and the relations:
$${}_{rel}{}^{}v_{\xi }^{}:=\overline{g}[h_u(_u\xi )]=_{rel}v\text{ , }_{rel}v_\eta :=\overline{g}[h_u(_u\eta )]\text{ ,}$$
(40)
$$\begin{array}{c}g(_u\xi ,\eta )=\frac{1}{e}.g(u,\eta ).g(_u\xi ,u)+g(_{rel}v_\xi ,\eta )\text{ ,}\\ g(_u\eta ,\xi )=\frac{1}{e}.g(u,\xi ).g(_u\eta ,u)+g(_{rel}v_\eta ,\xi )\text{ ,}\end{array}$$
(41)
$$(_ug)(\xi ,\eta )=(^s_ug)(\xi ,\eta )+\frac{1}{n}.Q_u.g(\xi ,\eta )\text{ ,}$$
(42)
$$\begin{array}{c}(^s_ug)(\xi ,\eta )=(^s_ug)(\xi _{},\eta _{})+\frac{l}{e}.(^s_ug)(u,\eta _{})+\frac{\overline{l}}{e}.(^s_ug)(\xi _{},u)+\\ +\frac{l}{e}.\frac{\overline{l}}{e}.(^s_ug)(u,u)\text{ , }\overline{l}=g(u,\eta )\text{ , }\eta _{}=\overline{g}[h_u(\eta )]\text{ , }l=g(u,\xi )\text{ ,}\end{array}$$
(43)
$$\begin{array}{c}(_ug)(\xi ,\eta )=(^s_ug)(\xi ,\eta )+\frac{1}{n}.Q_u.g(\xi ,\eta )=\\ =(^s_ug)(\xi _{},\eta _{})+\frac{l}{e}.(^s_ug)(u,\eta _{})+\frac{\overline{l}}{e}.(^s_ug)(\xi _{},u)+\\ +\frac{l}{e}.\frac{\overline{l}}{e}.(^s_ug)(u,u)+\frac{1}{n}.Q_u.[\frac{l.\overline{l}}{e}+g(\xi _{},\eta _{})]\text{ ,}\end{array}$$
(44)
the expression of $`l_\xi .l_\eta .\{u[\mathrm{cos}(\xi ,\eta )]\}`$ follows in the form
$$\begin{array}{c}l_\xi .l_\eta .\{u[\mathrm{cos}(\xi ,\eta )]\}=(^s_ug)(\xi _{},\eta _{})+\frac{l}{e}.[(^s_ug)(u,\eta _{})+g(_u\eta ,u)]+\\ +\frac{\overline{l}}{e}.[(^s_ug)(\xi _{},u)+g(_u\xi ,u)]+\frac{l.\overline{l}}{e^2}.(^s_ug)(u,u)+\\ +\frac{1}{n}.Q_u.[\frac{l.\overline{l}}{e}+g(\xi _{},\eta _{})]+g(_{rel}v_\xi ,\eta )+g(_{rel}v_\eta ,\xi )\\ [l_\eta .(ul_\xi )+l_\xi .(ul_\eta )].\mathrm{cos}(\xi ,\eta )\text{ .}\end{array}$$
(45)
Special case: $`g(u,\xi )=l:=0`$, $`g(u,\eta )=\overline{l}:=0:\xi =\xi _{}`$, $`\eta =\eta _{}`$.
$$\begin{array}{c}l_\xi _{}.l_\eta _{}.\{u[\mathrm{cos}(\xi _{},\eta _{})]\}=(^s_ug)(\xi _{},\eta _{})+\frac{1}{n}.Q_u.l_\xi _{}.l_\eta _{}.\mathrm{cos}(\xi _{},\eta _{})+\\ +g(_{rel}v_\xi _{},\eta _{})+g(_{rel}v_\eta _{},\xi _{})[l_\eta _{}.(ul_\xi _{})+l_\xi _{}.(ul_\eta _{})].\mathrm{cos}(\xi _{},\eta _{})\text{ ,}\end{array}$$
(46)
where $`g(\xi _{},\eta _{})=l_\xi _{}.l_\eta _{}.\mathrm{cos}(\xi _{},\eta _{})`$.
The kinematic characteristics related to the relative velocity and used in considerations of the rate of change of the length of contravariant vector fields as well as the change of the angle between two contravariant vector fields could also be useful for description of the motion of physical systems in $`(\overline{L}_n,g)`$-spaces.
Acknowledgments
This work is supported in part by the National Science Foundation of Bulgaria under Grant No. F-642. |
warning/0002/cond-mat0002107.html | ar5iv | text | # Development of Magnetism in Strongly Correlated Cerium Systems: Non-Kondo Mechanism for Moment Collapse
\[
## Abstract
We present an ab initio based method which gives clear insight into the interplay between the hybridization, the coulomb exchange, and the crystal-field interactions, as the degree of 4f localization is varied across a series of strongly correlated cerium systems. The results for the ordered magnetic moments, magnetic structure, and ordering temperatures are in excellent agreement with experiment, including the occurence of a moment collapse of non-Kondo origin. In contrast, standard ab initio density functional calculations fail to predict, even qualitatively, the trend of the unusual magentic properties.
PACS: 71.27+a, 71.28+d, 71.10Fd, 75.30Mb, 75.20Hr, 75.10Lp
\]
The difficulties and interest in treating strongly correlated electron systems, and the consequences of correlation effects on magnetic behavior in the transitional 4f or 5f localization regime, provide one of the central problems of condensed matter physics.<sup>1-3</sup> The transitional regime behavior is neither atomiclike nor itinerant. This gives rise to an extremely interesting range of phenomena, but also causes very great difficulties in treating the theory of these phenomena adequately, especially in a way providing the ability to predict the behavior of specific materials.<sup>1-3</sup> An adequate treatment requires treating the interelectronic coulomb interaction, i.e. the correlation effects, as constrained by exchange symmetry.<sup>4-6</sup> In this letter, we demonstrate an approach for treating these difficulties in predicting the interesting and complex behavior of an important series of cerium compounds.
The isostructural (rock-salt structure) series of the cerium monopnictides CeX (X = P, As, Sb, Bi) and monochalcogenides (X=S, Se, Te) have become prototype model systems for study, because of their unusual magnetic properties.<sup>7-13</sup> This series of strongly correlated electron systems offers the opportunity to vary systematically, through chemical pressure, the lattice constant and the cerium-cerium separation on going down the pnictogen or chalcogen column, and hence tailor the degree of 4f localization from the strongly correlated limit in the heavier systems to the weakly correlated limit in the lighter systems.<sup>7-13</sup> The calculated single-impurity Kondo temperature, T<sub>K</sub>, presented below, is much smaller than the magnetic ordering temperature in these systems, and hence this series lies in the magnetic regime of the Kondo phase diagram.<sup>14</sup> Nevertheless, in this work we demonstrate that the sensitivity of the hybridization, coulomb exchange, and crystal-field interactions with the chemical environment gives rise to a variety of unusual and interesting magnetic properties across the series, in agreement with experiment, including the occurence of a non-Kondo magnetic moment collapse.
This class of cerium systems exhibits large magnetic anisotropy which changes from the $`<001>`$ direction in the pnictides to the $`<111>`$ direction in the chalcogenides. The low-temperature ordered magnetic moment increases with increasing lattice constant for the pnictides from 0.80$`\mu _B`$ in CeP to 2.1$`\mu _B`$ in CeSb and CeBi,<sup>7-8</sup> while it decreases with increasing lattice constant for the chalcogenides from 0.57$`\mu _B`$ in CeS to 0.3$`\mu _B`$ in CeTe.<sup>7-9</sup> The magnetic moment collapse from CeSb to CeTe, with both systems having about the same lattice constant, is indicative of the sensitivity of the magnetic interactions to chemical environment. The experimentally observed low-temperature structure in CeBi and CeSb is the $`<001>`$ antiferromagnetic type IA ($`)`$, whereas in CeAs and CeP the structure is the $`<001>`$ antiferromagnetic type I ($``$).<sup>7,15</sup> The ordering temperature increases from 8K in CeP to 26K in CeBi for the pnictides, whereas it decreases from 8.4K in CeS to an unusually low 2.2K in CeTe.<sup>7-11</sup> Another unusual feature of this series of cerium compounds is the large suppression of the crystal field (CF) splitting of the Ce<sup>3+</sup> free-ion 4$`f_{5/2}`$ multiplet from values expected from the behavior of the heavier isostructural rare-earth pnictides or chalcogenides.<sup>16</sup> This can be understood<sup>17</sup> as arising from band-f hybridization effects. In both the cerium monopnictides and monochalcogenides, the CF splitting between the $`\mathrm{\Gamma }_7`$ doublet and the $`\mathrm{\Gamma }_8`$ quartet decreases with increasing anion size, from 150 K for CeP to 10 K in CeBi and from 130 K for CeS to 30 K for CeTe, and it is about the same for the same row in both series, a rather surprising result in view of the additional valence electron on the chalcogen ion.<sup>18</sup> Neutron scattering experiments have shown<sup>19</sup> that the $`\mathrm{\Gamma }_7`$-doublet is the CF ground state in all the cerium pnictides and chalcogenides.
In this paper we present material-predictive results from two ab initio based methods to study the change of magnetic properties across this series of cerium systems. The first, ab initio based, method gives clear insight into the role of the three pertinent interactions: 1) The band-f hybridization-induced inter-cerium magnetic coupling; 2) the corresponding effects of band-f coulomb exchange; and 3) the crystal-field interaction. This approach allows us also to understand the interplay between these interactions as the degree of 4f localization is varied across the series. The predictive calculations give results for the magnetic moments, magnetic structure, and ordering temperatures in excellent agreement with experiment. Thus, this approach allows to understand and predict a number of key features of observed behavior. First, is the very low moment and low ordering temperature of the antiferromagnetism observed in CeTe, an incipient heavy Fermion system. (For a review of theory and experimental behavior of heavy Fermion systems see references 2,3,20,21.) This ab initio-based method, described below, predicts the magnetic moment and ordering temperature collapse from CeSb to CeTe, both systems having about the same lattice constant but CeTe having an additional p electron. The origin of the moment collapse is of non-Kondo origin. The earlier work of Sheng and Cooper<sup>5</sup> showed that this magnetic ordering reduction is accurately predicted without including any crystal-field effects. An erroneous statement appears in the recent review article by Santini et al.<sup>22</sup> stating that crystal-field effects played an important role in the calculated results of Sheng and Cooper.<sup>5</sup> This is incorrect, since crystal-field effects were not included in these calculations. We show in this paper that including the crystal-field effects modifies this behavior only quantitatively. Second, our results demonstrate that, while the band-f coulomb exchange mediated interatomic 4f-4f interactions dominate the magnetic behavior for the heavier systems, which are more localized because of the larger Ce-Ce separation, the opposite is true for the lighter, more delocalized systems, where the hybridization-mediated coupling dominates the magnetic behavior. This reflects the great sensitivity of the relative importance of hybridization and coulomb exchange effects on magnetic ordering depending on the degree of 4f localization. Third, we show that for the lighter more delocalized systems the crystal-field interactions are much larger than the inter-cerium interactions and hence dominate the magnetic behavior. Finally, we predict the experimentally observed change of the ground-state magnetic structure from the $`<001>`$ antiferromagnetic type IA ($`)`$ in CeBi and CeSb to the $`<001>`$ antiferromagnetic type I ($``$) in CeAs and CeP. On the other hand, the second ab initio method, based on density functional theory within the local density approximation (LDA),<sup>23,24</sup> fails to predict, even qualitatively, the trend of magnetic properties in this series of strongly correlated electron systems.
The first, ab initio based, method employs the degenerate Anderson lattice model which incorporates explicitly the hybridization and the coulomb exchange interactions on an equal footing<sup>4,5</sup>
$`H`$ $`=`$ $`{\displaystyle \underset{k}{}}ฯต_kc_k^+c_k+{\displaystyle \underset{Rm}{}}ฯต_mf_m^+(R)f_m(R)`$
$`+{\displaystyle \frac{U}{2}}{\displaystyle \underset{R,mm^{}}{}}n_m(R)n_m^{}(R)`$
$`+{\displaystyle \underset{kmR}{}}[V_{km}e^{i๐ค๐}c_k^+f_m(R)+H.C.]`$
$`{\displaystyle \underset{kk^{}}{}}{\displaystyle \underset{mm^{}R}{}}J_{mm^{}}(๐ค,๐ค^{})e^{i(๐ค๐ค^{})๐}c_k^+f_m^+(R)c_k^{}f_m^{}(R).`$
(1)
The parameters entering the model Hamiltonian, i.e., the band energies $`ฯต_k`$, the f-state energy $`ฯต_m`$, the on-site coulomb repulsion U, the hybridization matrix elements, V<sub>km</sub>, and the band-f coulomb exchange $`J_{mm^{}}(๐ค,๐ค^{})=\varphi _k^{}(r_1)\psi _m^{}(r_2)\left|\frac{1}{r_{12}}\right|\psi _m^{}(r_1)\varphi _k^{}(r_2)`$ are evaluated on a wholly ab initio basis from non-spin polarized full potential linear muffin tin orbital<sup>23</sup> (FPLMTO) calculations. Here, $`r_{12}`$ stands for $`|r_1r_2|`$; $`\varphi _k`$ are the non-f basis states of the FPLMTO, and $`\psi _m`$ are the localized f states. Because the size of both the hybridization and coulomb exchange matrix elements are much smaller ($``$ 0.1 eV) than the intraatomic coulomb interaction U (6eV), one can apply perturbation theory and evaluate the anisotropic two-ion 6X6 interaction matrices<sup>4,5</sup>, $`E_{m_1m_1^{}}^{m_2m_2^{}}(๐_\mathrm{๐}๐_\mathrm{๐})`$, which couple the two f-ions. The exchange interactions have three contributions: the wholly band-f coulomb exchange mediated interaction proportional to $`J_{mm^{}}^2(๐ค,๐ค^{})`$, the wholly hybridization-mediated exchange interaction proportional to $`V_{km}^4`$, and the cross term proportional to $`V_{km}^2J_{mm^{}}(๐ค,๐ค^{})`$. With the two-ion interactions having been determined, the low-temperature magnetic moment and the ordering temperature can be determined by use of a mean field calculation.<sup>4,5,8</sup> We have previously applied this ab initio based method to investigate the effect of hybridization-induced cerium-cerium interactions<sup>4,17</sup> and the combined effect of both the hybridization and coulomb induced interactions<sup>5</sup> on the magnetic properties of the heavier cerium pnictides and chalcogenides (CeBi, CeSb, and CeTe). However, these calculations did not take into account the crystal field interaction and employed a warped muffin-tin LMTO calculation for the parameters entering the model. The excellent agreement found<sup>5</sup> with experiment for the low-temperature magnetic moment and ordering temperature is relatively unaffected by the CF interaction, because the CF interaction in the heavier cerium systems is smaller than the two-ion exchange interactions.
The second method employs ab initio spin polarized electronic structure calculations based on the FPLMTO method<sup>23</sup> using 1) only spin polarization, with the orbital polarization included only through the spin-orbit coupling, and 2) both the spin and orbital polarization polarization.<sup>24</sup> In these calculations the 4f states are treated as band states. The orbital polarization is taken into account by means of an eigenvalue shift<sup>24</sup>, $`\mathrm{\Delta }V_m=E^3L_zm_l`$, for the 4f atom. Here, L<sub>z</sub> is the z-component of the cerium total orbital moment, m<sub>l</sub> is the magnetic quantum number, and E<sup>3</sup> is the Racah parameter evaluated self-consistently at each iteration.
The crystalline field, which was neglected in the previous calculations,<sup>4,5</sup> is expected to affect the magnetic behavior considerably, if it is large. It is important to emphasize that since in the first method the 4f states are treated as core states, they interact only with the spherical component of the effective one-electron potential. Thus, the interaction of the atomic-like 4f state with the non-spherical components of the potential, giving rise to the CF splitting, $`\mathrm{\Delta }_{CF}=ฯต_{\mathrm{\Gamma }_8}ฯต_{\mathrm{\Gamma }_7}`$, is not included in the calculation of the model Hamiltonian parameters. In this paper, we generalize the first, ab initio based, method to include both the interatomic 4f-4f coupling and the crystal-field interactions on an equal footing and to employ a full potential LMTO evaluation of the model Hamiltonian parameters. While the effect of the full potential on both the hybridization and coulomb exchange interactions is small, including the CF interaction will be shown to play a role as important as the interatomic 4f-4f interactions for understanding and predicting the overall trend in the unusual magnetic properties, as as one chemically tunes the degree of 4f localization across this series of strongly correlated electron systems. The resultant Hamiltonian is<sup>4,5</sup>
$`H`$ $`=`$ $`{\displaystyle \underset{i,j}{}}{\displaystyle \underset{_{ฯต,\sigma }^{\mu ,\nu }}{}}\xi _{\mu \nu }^{ฯต\sigma }(\theta _{ij})e^{i(\mu \nu +ฯต\sigma )\varphi _{ij}}c_ฯต^{}(j)c_\sigma (j)c_\mu ^{}(i)c_\nu (i)`$
$`+B_4{\displaystyle \underset{i}{}}\left(O_4^0(i)+5O_4^4(i)\right),`$
(2)
where the $`\xi _{\mu \nu }^{ฯต\sigma }(\theta _{ij})`$ are the two-ion 4f-4f interaction matrices rotated to a common crystal-lattice axis, and the $`O_4^0`$ and $`O_4^4`$ are the Stevens operators equivalents acting on the Ce<sup>3+</sup> free-ion 4$`f_{5/2}`$ multiplet.<sup>25</sup> The CF splitting is $`\mathrm{\Delta }_{CF}`$ = 360$`B_4`$; a positive $`B_4`$ value gives the $`\mathrm{\Gamma }_7`$ ground state, which is experimentally observed.<sup>19</sup> While our work in progress is aimed at evaluating the CF splitting on a wholly ab initio basis, in the absence of an ab initio value of the CF interaction in this class of strongly correlated cerium systems, the $`\mathrm{\Delta }_{CF}`$ is set to the experimental values listed in Table 3.<sup>10,19</sup>
In Table I, we present the calculated values of the zero-temperature cerium magnetic moment from the FPLMTO electronic structure calculations. Listed in the table are values both with and without the orbital polarization correction taken into account. Note, the importance of including the orbital polarization in these 4f correlated electron systems. As expected, in all cases, the orbital polarization is found to be opposite to the spin polarization. Comparison of the total energies predicts that the magnetic anisotropy changes from the $`<001>`$ direction in the pnictides to the $`<111>`$ in the chalcogenides, in agreement with experiment. On the other hand, except perhaps for the lighter chalcogenides (CeS and CeSe), comparison of the ab initio and experimental values for the magnetic moment indicates the failure of the LDA calculations to treat properly the correlation effects of the 4f states (treated as valence states) within the LDA as the degree of 4f correlations increases in the heavier pnictide systems. Furthermore, the ab inito calculations fail to predict the large moment collapse from CeSb to CeTe, the latter being described as an incipient heavy Fermion system.<sup>2,3,20,21</sup>
In Table II, we list the values of the $`m`$ = $`m^{}`$ =1/2 matrix elements (characteristic matrix elements of the 6X6 exchange interaction matrix) for the first three nearest-neighbor shells for the light (CeP and CeS) and the heavier compounds (CeSb and CeTe). Listed separately in this table are the three contributions to the interatomic 4f-4f interactions arising from band-f hybridization (V<sup>4</sup>), band-f coulomb exchange (J<sup>2</sup>), and the cross term. It is important to note that while the coulomb exchange mediated interactions dominate the magnetic behavior for the heavier, more localized, 4f systems, the opposite is true for the lighter, more delocalized, systems where the hybridization mediated interactions dominate the magnetic behavior. This change of behavior of the interatomic 4f-4f interactions is a result of the sensitivity of the hybridization and coulomb exchange to the degree of 4f localization. Equally important, is that while both first and second nearest-neighbor 4f-4f interactions are ferromagnetic for CeSb, there is an interplay between ferromagnetic first nearest-neighbor and antiferromagnetic second nearest-neighbor interactions for CeTe.(These interactions are mediated via scattering of conduction electrons). This results in a saturated ordered moment for CeSb and in the ordered magnetic moment collapse for CeTe (see Table III).
In order to determine whether the magnetic moment collapse might be of Kondo origin, we have evaluated the single-impurity Kondo temperature,<sup>26</sup> k<sub>B</sub>T<sub>K</sub> = De$`^{\frac{1}{2\rho (E_F)|J(E_F)|}}`$, across the series. Here, D is the bandwidth of the conduction electron states, $`\rho (E_F)`$ is the density of states of the conduction electrons at the Fermi energy, and $`J(E_F)`$ is the conduction electron-f exchange interaction at the Fermi energy, which has contributions both from the coulomb exchange interaction in Eq. (2), provided that it is negative, and the hybridization-induced exchange interaction $`|J_{hyb}(E_F)|=\frac{V^2(E_F)U}{(|E_fE_F|)(|E_fE_F|+U)}`$, where $`J_{hyb}(E_F)<0`$. We find that the coulomb exchange interaction in Eq. (2) evaluated at the Fermi energy is positive across the entire series and hence cannot give rise to the Kondo effect. Thus, only the hybridization-induced exchange interaction, $`J_{hyb}(E_F)`$, can give rise to the Kondo effect.<sup>26</sup> Using the ab initio values of the parameters entering the expression for T<sub>K</sub>, we find that T$`{}_{K}{}^{}`$ T<sub>ord</sub> across the entire series(T$`{}_{K}{}^{}<`$ 10<sup>-4</sup>K). The Kondo temperatures in the monopnictide series is smaller than that in the chaclogenides, due to the fact that in the pnictides the Fermi energy lies in the pseudogap, resulting in low $`\rho (E_F)`$. These results suggest that the moment collapse from CeSb to CeTe is of non-Kondo origin. Rather, it results from an interplay of ferromagnetic and antiferromagnetic interatomic 4f-4f interactions which arises purely from differences in the underlying electronic structure.
Listed in Table III are the calculated zero-temperature ordered moment and ordering temperature, T<sub>N</sub>, from the first, ab initio based, method, with and without the CF interaction. It is clear that for the heavier systems (CeBi, CeSb, CeTe) the effect of the CF interaction on the magnetic moments is small, and it is slightly more pronounced on the ordering temperatures. This is due to the fact that for the more localized systems the CF interaction is smaller than the two-ion interactions. This is the reason that the previous calculations,<sup>5</sup> neglecting the CF interaction, gave results in very good agreement with experiment. On the other hand, for the lighter more delocalized systems the CF interactions are much larger than the interatomic 4f-4f interactions, and hence dominate the magnetic behavior. The overall decrease of the magnetic moments in the presence of the CF interaction in all systems, arises from the mixing of the off-diagonal angular momentum states $`|\pm 5/2>`$ and $`|3/2>`$ states from the CF interaction with $`\mathrm{\Gamma }_7`$ ground state. Overall, we find that the first, ab initio based, approach which takes into account all three pertinent interactions (hybridization, coulomb exchange, and CF interactions) on an equal footing, yields results for both the zero-temperature moment and the ordered temperature (a more stringent test for the theory) in excellent agreement with experiment.
A final corroboration of the success of the first ab initio based method is that it predicts the experimentally observed change of the ground-state magnetic structure from the $`<001>`$ antiferromagnetic type IA ($`)`$ in CeBi and CeSb to the $`<001>`$ antiferromagnetic type I ($``$) in CeAs and CeP.<sup>7,15</sup> More specifically, the sign \[ferromagnetic (F) or antiferromagnetic (AF)\] of the $`|\pm 5/2>`$ matrix elements of the 6X6 exchange matrix determines the interplanar interaction between successive (001) Ce planes. We find, that for the heavier compounds (CeBi and CeSb) the $`|\pm 5/2>`$ matrix elements of the coulomb exchange matrix are FM and hence favor the $``$ type, while in the lighter systems (CeAs and CeP) the $`|\pm 5/2>`$ matrix elements of the hybridization-induced two-ion matrix are AF, and hence they favor the $``$ type.
In conclusion, we have applied two different, ab inito based and ab initio LDA, methods to study the dramatic change of magnetic properties across a series of strongly correlated electron systems which offer the opportunity to chemically tailor the different type of interactions (band-f hybridization, band-f coulomb exchange, and CF interactions), pertinent to the unusual magnetic behavior. The first, ab initio based, approach which explicitly takes into account the interplay of the three pertinent interactions, gives results in excellent agreement with experiment for all compounds in the series, including the moment collapse from CeSb to CeTe and the trend of moments and ordering temperatures across the series. The remaining problem of determining on a wholly ab initio basis the suppressed crystal-field interactions in this class of systems poses a theoretical challenge for future theoretical work. On the other hand, the second, fully ab initio LDA, method gives good results for the lighter chalcogenide systems, but it entirely fails to give, even qualitatively, the trend of the unusual magnetic behavior.
The research at California State University Northridge (CSUN) was supported by the National Science Foundation under Grant No. DMR-9531005, by the US Army Grant No. DAAH04-95, and the CSUN Office of Research and Sponsored Projects, and at West Virginia University by the NSF under Grant No. DMR-9120333. |
warning/0002/astro-ph0002058.html | ar5iv | text | # *
## 1 Introduction
The nature of the dark matter in the haloes of galaxies is an outstanding problem in astrophysics. Over the last several decades there has been great debate about whether this matter is baryonic or must be exotic. Many astronomers believed that a stellar or substellar solution to this problem might be the most simple and therefore most plausible explanation. However, in the last few years, these candidates have been ruled out as significant components of the Galactic Halo. I will discuss limits on these stellar candidates, and argue for my personal conviction that: Most of the dark matter in the Galactic Halo must be nonbaryonic.
Until recently, stellar candidates for the dark matter, including faint stars, brown dwarfs, white dwarfs, and neutron stars, were extremely popular. However, recent analysis of various data sets has shown that faint stars and brown dwarfs probably constitute no more than a few percent of the mass of our Galaxy . Specifically, using Hubble Space Telescope and parallax data, we showed that faint stars and brown dwarfs contribute no more than 1% of the mass density of the Galaxy. Microlensing experiments (the MACHO , and EROS ) experiments), which were designed to look for Massive Compact Halo Objects (MACHOs), also failed to find these light stellar objects and ruled out substellar dark matter candidates in the $`(10^710^2)M_{}`$ mass range.
Recently white dwarfs have received attention as possible dark matter candidates. Interest in white dwarfs has been motivated by microlensing events interpreted as being in the Halo, with a best fit mass of $`0.5M_{}`$. However, I will show that stellar remnants including white dwarfs and neutron stars are extremely problematic as dark matter candidates. A combination of excessive infrared radiation, mass budget issues and chemical abundances constrains the abundance of stellar remnants in the Halo quite severely, as shown below. Hence, white dwarfs, brown dwarfs, faint stars, and neutron stars are either ruled out or extremely problematic as dark matter candidates. Thus the puzzle remains, What are the 14 MACHO events that have been interpreted as being in the Halo of the Galaxy? Are some of them actually located elsewhere, such as in the LMC itself? These questions are currently unanswered. As regards the dark matter in the Halo of our Galaxy, one is driven to nonbaryonic constituents as the bulk of the matter. Possibilities include supersymmetric particles, axions, primordial black holes, or other exotic candidates.
In this talk I will focus on the arguments against stellar remnants as candidates for a substantial fraction of the dark matter, as white dwarfs in particular have been the focus of attention as potential explanations of microlensing data. For a discussion of limits on faint stars and brown dwarfs, see earlier conference proceedings by Freese, Fields, and Graff ( and ).
## 2 White Dwarfs
Stellar remnants (white dwarfs and neutron stars) face a number of problems and issues as dark matter candidates: 1) infrared radiation; 2) IMF (initial mass function); 3) baryonic mass budget; 4) element abundances.
We find that none of the expected signatures in the above list of a significant white dwarf component in the Galactic Halo are seen to exist.
### 2.1 Constraints from multi-TeV $`\gamma `$-rays seen by HEGRA
The mere existence of multi-TeV $`\gamma `$-rays seen in the HEGRA experiment places a powerful constraint on the allowed abundance of white dwarfs. This arises because the progenitors of the white dwarfs would produce infrared radiation that would prevent the $`\gamma `$-rays from getting here. The $`\gamma `$-rays and infrared photons would interact via $`\gamma \gamma e^+e^{}`$.
Multi-TeV $`\gamma `$-rays from the blazar Mkn 501 at a redshift z=0.034 are seen in the HEGRA detector. The cross section for (1-10)TeV $`\gamma `$-rays peaks at infrared photon energies of (0.03-3)eV. Photons in this energy range would be produced in abundance by the progenitor stars to white dwarfs and neutron stars. By requiring that the optical depth due to $`\gamma \gamma e^+e^{}`$ be less than one for a source at $`z=0.034`$ we limit the cosmological density of stellar remnants to $`\mathrm{\Omega }_{\mathrm{WD}}(13)\times 10^3h^1`$. This constraint is quite robust and model independent, as it applies to a variety of models for stellar physics, star formation rate and redshift, mass function, and clustering.
### 2.2 Mass Budget Issues
#### Contribution of Machos to the Mass Density of the Universe:
(based on work by Fields, Freese, and Graff ) There is a potential problem in that too many baryons are tied up in Machos and their progenitors (Fields, Freese, and Graff). We begin by estimating the contribution of Machos to the mass density of the universe. Microlensing results predict that the total mass of Machos in the Galactic Halo out to 50 kpc is $`M_{\mathrm{Macho}}=(1.33.2)\times 10^{11}M_{}.`$ Now one can obtain a โMacho-to-lightโ ratio for the Halo by dividing by the luminosity of the Milky Way (in the B-band), $`L_{MW}(1.32.5)\times 10^{10}L_{},`$ to obtain $`(M/L)_{\mathrm{Macho}}=(5.225)M_{}/L_{}.`$ ยฟFrom the ESO Slice Project Redshift survey , the luminosity density of the Universe in the $`B`$ band is $`_B=1.9\times 10^8hL_{}\mathrm{Mpc}^3.`$ If we assume that the $`M/L`$ which we defined for the Milky Way is typical of the Universe as a whole, then the universal mass density of Machos is
$$\mathrm{\Omega }_{\mathrm{Macho}}\rho _{\mathrm{Macho}}/\rho _c=(0.00360.017)h^1$$
(1)
where the critical density $`\rho _c3H_0^2/8\pi G=2.71\times 10^{11}h^2M_{}\mathrm{Mpc}^3`$.
We will now proceed to compare our $`\mathrm{\Omega }_{\mathrm{Macho}}`$ derived in Eq. 1 with the baryonic density in the universe, $`\mathrm{\Omega }_\mathrm{B}`$, as determined by primordial nucleosynthesis. To conservatively allow for the full range of possibilities, we will adopt $`\mathrm{\Omega }_\mathrm{B}=(0.0050.022)h^2.`$ Thus, if the Galactic halo Macho interpretation of the microlensing results is correct, Machos make up an important fraction of the baryonic matter of the Universe. Specifically, the central values give
$$\mathrm{\Omega }_{\mathrm{Macho}}/\mathrm{\Omega }_\mathrm{B}0.7.$$
(2)
However, the lower limit on this fraction is considerably less restrictive,
$$\frac{\mathrm{\Omega }_{\mathrm{Macho}}}{\mathrm{\Omega }_\mathrm{B}}\frac{1}{6}h\frac{1}{12}.$$
(3)
#### Mass Budget constraints from Machos as Stellar Remnants: White Dwarfs or Neutron Stars
In general, white dwarfs, neutron stars, or black holes all came from significantly heavier progenitors. Hence, the excess mass left over from the progenitors must be added to the calculation of $`\mathrm{\Omega }_{\mathrm{Macho}}`$; the excess mass then leads to stronger constraints. Typically we find the contribution of Macho progenitors to the mass density of the universe to be $`\mathrm{\Omega }_{\mathrm{prog}}=4\mathrm{\Omega }_{\mathrm{Macho}}=(0.0160.08)h^1`$. The central values of all the numbers now imply $`\mathrm{\Omega }_{\mathrm{prog}}3\mathrm{\Omega }_B`$, which is obviously unacceptable. One is driven to the lowest values of $`\mathrm{\Omega }_{\mathrm{Macho}}`$ and highest value of $`\mathrm{\Omega }_B`$ to avoid this problem.
### 2.3 On Carbon and Nitrogen
The overproduction of carbon and/or nitrogen produced by white dwarf progenitors is one of the greatest difficulties faced by a white dwarf dark matter scenario, as first noted by Gibson and Mould . Stellar carbon yields for zero metallicity stars are quite uncertain. Still, according to the yields by , a star of mass 2.5$`M_{}`$ will produce about twice the solar enrichment of carbon. However, stars in our galactic halo have carbon abundance in the range $`10^410^2`$ solar. Hence the ejecta of a large population of white dwarfs would have to be removed from the galaxy via a galactic wind.
However, carbon abundances in intermediate redshift Ly$`\alpha `$ forest lines have recently been measured to be quite low, at the $`10^2`$ solar level , for Ly$`\alpha `$ systems at $`z3`$ with column densities $`N3\times 10^{15}\mathrm{cm}^2`$ (for lower column densities, the mean C/H drops to $`10^{3.5}`$ solar .
In order to maintain carbon abundances as low as $`10^2`$ solar, only about $`10^2`$ of all baryons can have passed through the intermediate mass stars that were the predecessors of Machos . Such a fraction can barely be accommodated for the remnant density predicted from our extrapolation of the Macho group results, and would be in conflict with $`\mathrm{\Omega }_{\mathrm{prog}}`$ in the case of a single burst of star formation. Note that stars heavier than 4$`M_{}`$ may replace the carbon overproduction problem with nitrogen overproduction .
Using the yields described above, we calculated the C and N that would result from the stellar processing for a variety of initial mass functions for the white dwarf progenitors. We used a chemical evolution model based on a code described in Fields & Olive to obtain our numerical results. Our results are presented in the figure.
In the figure, we make the parameter choices that are in agreement with D and He<sup>4</sup> measurements (see the discussion below) and are the least restrictive when comparing with the Ly$`\alpha `$ measurements. We take an initial mass function (IMF) sharply peaked at 2$`M_{}`$, so that there are very few progenitor stars heavier than 3$`M_{}`$ (this IMF is required by D and He<sup>4</sup> measurements). In addition (see the figures in Fields, Freese, and Graff ) we have considered a variety of other parameter choices. By comparing with the observations, we obtain the limit, $`\mathrm{\Omega }_{\mathrm{WD}}h2\times 10^4`$. As a caveat, note that it is possible that carbon never leaves the (zero metallicity) white dwarf progenitors, so that carbon overproduction is not a problem .
### 2.4 Deuterium and Helium
Because of the uncertainty in the C and N yields from low-metallicity stars, we have also calculated the D and He<sup>4</sup> abundances that would be produced by white dwarf progenitors. These are far less uncertain as they are produced farther out from the center of the star and do not have to be dredged up from the core. Panel a) in the figure displays our results. Also shown are the initial values from big bang nucleosynthesis and the (very generous) range of primordial values of D and He<sup>4</sup> from observations. ยฟFrom D and He alone, we can see that the white dwarf progenitor IMF must be peaked at low masses, $`2M_{}`$. We obtain $`\mathrm{\Omega }_{\mathrm{WD}}0.003`$.
## 3 Conclusions
#### A Zero Macho Halo?
The possibility exists that the 14 microlensing events that have been interpreted as being in the Halo of the Galaxy are in fact due to some other lensing population. One of the most difficult aspects of microlensing is the degeneracy of the interpretation of the data, so that it is currently impossible to determine whether the lenses lie in the Galactic Halo, or in the Disk of the Milky Way, or in the LMC. In particular, it is possible that the LMC is thicker than previously thought so that the observed events are due to self-lensing of the LMC. All these possibilities are being investigated. More data are required in order to identify where the lenses are.
Microlensing experiments have ruled out baryonic dark matter objects in the mass range $`10^7M_{}`$ all the way up to $`10^2M_{}`$. In this talk I discussed the heavier possibilities in the range $`10^2M_{}`$ to a few $`M_{}`$. Brown dwarfs and faint stars are ruled out as significant dark matter components; they contribute no more than 1% of the Halo mass density. Stellar remnants are not able to explain the dark matter of the Galaxy either; none of the expected signatures of stellar remnants, i.e., infrared radiation, large baryonic mass budget, and C,N, and He<sup>4</sup> abundances, are found observationally.
Hence, in conclusion, 1) Nonbaryonic dark matter in our Galaxy seems to be required, and 2) The nature of the Machos seen in microlensing experiments and interpreted as the dark matter in the Halo of our Galaxy remains a mystery. Are we driven to primordial black holes , nonbaryonic Machos (Machismos?), mirror matter Machos () or perhaps a no-Macho Halo? |
warning/0002/hep-ph0002239.html | ar5iv | text | # 1 Introduction
## 1 Introduction
Often in applications of quantum field theories, it is desirable to treat a part of the problem semiclassically and compute the quantum corrections in a controlled expansion around the classical solutions. This applies to a broad spectrum of problems, including for example the Fermi liquid theory, as well as the study of the expanding and cooling plasma in the heavy ion collisions . In these, and in many other cases, it is often sufficient to assume that the background is very slowly varying and to treat the problem at leading order in gradients. Such an adiabatic (quasiparticle) approximation may make sense if the physical effect of the background one is interested in is already present at the trivial order. However, there are situations where a semiclassical approximation is justified, but the system displays an important physical effect only at a nontrivial order in gradients. An important example of this occurs at the electroweak transition, where the problem is to study particle propagation in the presence of a spatially varying, CP-violating Higgs condensate. The CP-violating effects of the condensate on the plasma dynamics arise only at the first nontrivial order in gradient expansion . The induced perturbations in the particle and antiparticle densities, coupled with the weak anomaly, lead to a local baryon production mechanism, capable of creating the observed matter-antimatter asymmetry of the Universe .
Unfortunately, no consistent derivation of the semiclassical quantum Boltzman equation for a fermionic field beyond the trivial order in gradients exist in the literature. Moreover, there exists a well known no-go theorem , according to which first nontrivial corrections to fermion propagator arise only at second order in derivative expansion. This is in apparent contrast with the classical force mechanism (CFM) for baryogenesis , according to which the nontrivial effects on fermion propagation are felt already at the first order in gradients. In this letter we show that the no-go theorem does not hold universally, but instead, whenever the fermion self energy contains a pseudoscalar term, there are corrections to the propagator already at the first order in gradients. A simple example of such a term is provided by a Dirac field with a complex spatially (or temporally) varying mass term of the form
$$\widehat{m}(x)=m_R(x)+i\gamma _5m_I(x)|m|e^{i\gamma ^5\theta }.$$
(1)
This term could be of course replaced by generic scalar and pseudoscalar operators; we retain the present notation however, because of its clear physical implications. Namely, this type of term has precisely the form relevant for baryogenesis, where such an effective mass term is induced by a coupling of the fermion with a spatially varying complex scalar field condensate during a first order phase transition. Here we shall consider a simple toy model with the lagrangian
$$=i\overline{\psi }/\psi \overline{\psi }_Lm\psi _R\overline{\psi }_Rm^{}\psi _L+_{\mathrm{int}},$$
(2)
where $`m=m_R+im_I`$ is a complex mass. The detailed form of the interaction lagrangian $`_{\mathrm{int}}`$ is not relevant for us here. Realistic cases studied in literature in connection with the CFM baryogenesis include the two doublet extension of the standard model and the chargino sector of the minimal supersymmetric standard model .
This letter is a part of a series dedicated to derivation of semiclassical transport equations for quantum fields in varying backgrounds beyond the trivial order in gradients. We first explicitly construct the fermion propagator for the theory (2) to the first nontrivial order in gradients. Corrections appear as higher order poles multiplied by derivatives of the complex phase of the mass. The position of the poles is not shifted however, and the physical content of the correction terms only becomes apparent after consideration of spectral integrals over test functions by a technique developed in . It will be shown that the spectral function, while not expressible as a sum of ordinary delta-functions even in the on-shell limit, projects a test function to a sharp, but shifted energy shell. We then show that this shell corresponds to the physical semiclassical dispersion relation of a standard WKB wave function, and describe how the current computed using the spectral function representation can be related to motion of WKB wave packets.
## 2 Propagator
Our first task is to solve the retarded and advanced propagators for the fermionic field $`\psi `$ in the weak coupling approximation and in an expansion in gradients of the slowly varying mass terms. We will work in the Keldysh closed time contour (CTC) formalism , and the basic quantity of interest is the path-ordered propagator
$$iG_๐^\psi (x,y)=T_๐[\psi (y)\overline{\psi }(x)].$$
(3)
Writing the Schwinger-Dyson equations for $`G_๐`$ in the real time formulation implies the following formally exact equations for the retarded and advanced propagators $`G^{r,a}`$ in the Wigner representation
$$e^i\mathrm{}\left\{(iG_0)^1\right\}\left\{iG^{r,a}(k;X)\right\}=1,$$
(4)
where $`\mathrm{}\left\{f\right\}\left\{g\right\}\frac{1}{2}(_Xf_kg_kf_Xg)`$ and the tree level propagator is given by
$$(iG_0^{r,a})^1k/m_R(X)i\gamma _5m_I(X)\mathrm{\Sigma }^{r,a}\pm is_\omega ฯต,$$
(5)
where $`s_\omega =\mathrm{sign}(\omega )`$, $`k`$ denotes the canonical $`4`$-momentum and $`X(x+y)/2`$ the average position. From now on we shall assume a weak coupling limit and neglect the self-energy $`\mathrm{\Sigma }^{r,a}=\mathrm{\Sigma }_Ri\omega \mathrm{\Gamma }`$; in particular we thus consider the on-shell limit $`\mathrm{\Gamma }0`$. We will also suppress the $`iฯต`$-prescription which differentiates between $`G^{r,a}`$ even when $`\mathrm{\Gamma }0`$, which is introduced as a mnemonic tool to maintain information about the integration path required to fully specify how $`G`$ acts operationally. A formal solution to equation (5) can easily be found in the gradient approximation. Indeed, to first order in gradients we have:
$$iG=iG_0iG_0i\mathrm{}\left\{(iG_0)^1\right\}iG_0\left\{(iG_0)^1\right\}iG_0,$$
(6)
where the free propagator is given by
$$iG_0=\frac{1}{k^2|m|^2}(k/+m_R(X)i\gamma _5m_I(X)).$$
(7)
After some algebra one can write the expression (6) into the following covariant form:
$$iG=iG_0+\frac{1}{(k^2|m|^2)^2}[i(k_Xk//_X)|m|e^{i\theta \gamma _5}+\gamma ^5|m|^2/_X\theta ].$$
(8)
There is clearly a nonvanishing correction to the propagator at the leading nontrivial order. This is not in disagreement with the no-go theorem by , because, as we shall see, all nontrivial consequences of (8) are proportional to the complex pseudoscalar mass (or more generally to a pseudoscalar self-energy function), which was not considered in Ref. . An important point to observe about (8) is that the pole of the propagator is not shifted, but instead a new second order pole, proportional to gradients of $`m`$ has appeared. This behaviour is analogous to one previously observed in the case of a scalar field .
For simplicity we shall from now on restrict ourselves to the case of a planar symmetry with only spatially varying mass: $`m=m(x_3)`$. In this case one can Lorentz boost to the frame in which $`๐ค_{}=\mathrm{๐}`$, so that $`\omega \stackrel{~}{\omega }s_\omega \left(\omega ^2๐ค_{}^{\mathrm{\hspace{0.33em}2}}\right)^{1/2}`$. With these definitions Eq. (8) simplifies to
$`i\gamma ^0G`$ $`=`$ $`{\displaystyle \frac{1}{z}}(\stackrel{~}{\omega }k_3\gamma ^5S^3+\gamma ^0m_Ri\gamma ^0\gamma ^5m_I)`$ (9)
$``$ $`{\displaystyle \frac{1}{z^2}}\left[i\stackrel{~}{\omega }(m_R^{}\gamma ^0\gamma ^5im_I^{}\gamma ^0)+|m|^2\theta ^{}\right]S^3,`$
where we combined various $`\gamma `$-matrices into the spin operator $`S^3=\gamma ^0\gamma ^3\gamma ^5`$, which measures the spin in the $`\widehat{๐ค}_3`$-direction and in the frame in which $`๐ค_{}=\mathrm{๐}`$, and we used the definition $`zk^2|m|^2`$. Allowing $`z`$ a small complex value, one sees that (9) yields $`G^r`$ ($`G^a`$) for $`\mathrm{Im}(z)<0`$ ($`>0`$). A further very important simplification arises from the fact that pseudoscalar interactions conserve spin (one can show that the spin-operator $`S^3`$ commutes with the propagator $`G`$). The problem then becomes diagonal in spin and, by replacement $`S^3s`$, the $`4\times 4`$ structure of (9) becomes block-diagonal, so that the problem can be reduced to a $`2\times 2`$ problem. In the $`2\times 2`$ formalism, it is more convenient to use the Pauli matrices for the spinor algebra, which can be effected by the following identifications:
$$\gamma ^0\sigma ^1,i\gamma ^0\gamma ^5\sigma ^2\gamma ^5\sigma ^3.$$
(10)
where we have taken the chiral representation for the $`\gamma `$ matrices. In this way (9) can be recast in a particularly simple form:
$`i\sigma ^1G_s`$ $`=`$ $`{\displaystyle \frac{1}{z}}(\stackrel{~}{\omega }+sk_3\sigma ^3+m_R\sigma ^1+m_I\sigma ^2)`$ (11)
$``$ $`{\displaystyle \frac{s}{z^2}}\left[\stackrel{~}{\omega }(m_I^{}\sigma ^1m_R^{}\sigma ^2)+|m|^2\theta ^{}\right]`$
$``$ $`b_0+๐\sigma .`$
From the last expression in particular it is evident that there exists a basis in which the propagator $`G`$ is diagonal $`G_sG_{d,s}=b_0+s\mathrm{sign}(\stackrel{~}{\omega })|๐|\sigma ^3`$. In the Dirac notation the propagator then reads
$$i\gamma ^0G_d=s_{\stackrel{~}{\omega }}\left(\frac{|\stackrel{~}{\omega }|}{z}\frac{ss_{\stackrel{~}{\omega }}|m|^2\theta ^{}}{z^2}\right)\left(1s_{\stackrel{~}{\omega }}S^3\gamma ^5\right),$$
(12)
where $`s_{\stackrel{~}{\omega }}\mathrm{sign}(\stackrel{~}{\omega })`$. Physically the equation (12) means that the states with equal $`ss_{\stackrel{~}{\omega }}`$ propagate identically. Formally, this is expressed by the projector $`(1s_{\stackrel{~}{\omega }}S^3\gamma ^5)/2`$, which identifies the states with equal entries on the diagonal. In particular this means that particles and antiparticles with opposite spin have identical spectral functions. However, for particles and antiparticles of same spin (helicity) the spectral functions differ by an amount proportional to the CP-violating angle $`\theta ^{}`$. This is how in the context of baryogenesis CP-violating backgrounds induce bias in the particle and antiparticle distributions. In order to understand how this bias propagates in a plasma, it is then necessary to study the relevant quantum trasport equations with the effect of collisions included .
## 3 Spectral integrals
To study the physical consequences of the propagator (8), we now consider spectral integrals over test functions , representing some generic observables, such as the generalized particle distribution function. We define:
$$G[๐ฏ]\mathrm{Tr}\frac{i}{2\pi }_0^{\mathrm{}}๐\omega \gamma ^0๐๐ฏ,$$
(13)
where $`๐\frac{i}{2}(G^rG^a)`$ is the spectral function, $`๐ฏ`$ is a test function, and the trace is taken over the Dirac indices. Because scalar and pseudoscalar backgrounds considered here conserve spin, we can take $`๐ฏ`$ to be diagonal in the spin space and a scalar function in a spin $`2\times 2`$ block. For backgrounds that violate spin this analysis would have to be generalized to include spin mixing. The $`\gamma ^0`$-factor in Eq. (13) is added to make the projection operator explicitly hermitean. The spectral function $`๐`$ carries information on the physical spectrum of excitations in the system. In particular to zeroth order in gradients it becomes a simple on-shell projector <sup>1</sup><sup>1</sup>1Note that we are assuming the on-shell limit and neglecting possible thermal and/or vacuum corrections, so that $`\mathrm{\Sigma }^{r,a}=0`$. (cf. Eq. (7)):
$$๐\frac{k/+m_Ri\gamma _5m_I}{2\omega _0}\left[\delta (\omega \omega _0)\delta (\omega +\omega _0)\right],.$$
(14)
where $`\omega _0=(k_3^2+|m|^2)^{1/2}`$. In a spatially varying background the spectral function can no longer be expressed in terms of $`\delta `$-functions projecting onto the physical energy shells. As discussed in , it can however be written as a sum of projectors onto complex shells, but this is not the approach we will pursue here. Instead, by extending the momentum to the complex plane, the integral (13) can be turned into a contour integral over the variable $`zk^2|m|^2`$ :
$$G[๐ฏ]=\mathrm{Tr}\frac{i}{4\pi }_{๐_0}๐z\frac{i\gamma ^0G๐ฏ}{\sqrt{\omega _0^2+z}},$$
(15)
where $`G`$ is the propagator (8), and $`๐_0`$ is a contour encircling the pole at $`z=0`$ ($`\omega =\omega _0`$) once clockwise. We can then compute (15) by making use of the residue theorem:
$$G[๐ฏ]=\mathrm{Tr}\mathrm{Res}_{z=0}\left[\frac{i\gamma ^0G๐ฏ}{2\sqrt{\omega _0^2+z}}\right].$$
(16)
Because the highest pole in the expression (8) is of the second order, only the terms up to the linear order in $`z`$ in the expansions of various terms appearing in (16) contribute. In particular the gradient expansion in $`_{k_z}๐ฏ`$ then terminates at first order. That is to say that all $`_{k_z}^l๐ฏ`$-terms for $`l>1`$ are of at least second order in gradients. An important consequence of this truncation is that a complete description of a fermionic plasma, consistent to first order in gradients, requires at most two independent distribution functions, in contrast to the four distribution functions required for a full description of the bosonic case at the lowest nontrivial order in gradients .
It is now a simple matter to compute the residue in Eq. (16) of the propagator (11) or (12). (For notational simplicity from now on we shall denote $`\stackrel{~}{\omega }`$ by $`\omega `$.) To the order we are working
$`G[๐ฏ]`$ $`=`$ $`๐ฏ_0+{\displaystyle \frac{ss_\omega |m|^2\theta ^{}}{2\omega _0^2}}\left(๐ฏ_0\omega _0๐ฏ_0^{}\right)`$ (17)
$`=`$ $`{\displaystyle \frac{\omega _0}{\omega _{\mathrm{sc}}}}๐ฏ_0(\omega _{\mathrm{sc}}),`$
where $`๐ฏ_0๐ฏ(\omega _0)`$, $`๐ฏ_0^{}(_\omega ๐ฏ)(\omega _0)`$, and we have defined the semiclassical energy as
$$\omega _{\mathrm{sc}}\omega _0\frac{ss_\omega |m|^2\theta ^{}}{2\omega _0^2}.$$
(18)
The test function thus still gets projected to a sharply defined energy shell, which is however shifted with respect to the lowest order shell $`\omega _0`$.
## 4 Current
As a physically motivated application of the above, we now consider a current of spin states. In the present formalism this can be written as follows:
$$j_s^\mu (X)=\mathrm{Tr}\frac{d^4k}{(2\pi )^4}i\gamma ^\mu G^<(k;X)P_s,$$
(19)
where $`P_s=(1+sS^3)/2`$ is the spin projector and $`G^<(k;X)`$ is the quantum Wigner function. Assuming the decomposition $`iG^<(k;X)P_s=๐P_sn_s`$, where $`n_s`$ are generalized particle distribution functions of spin states, and inserting this expression into (19) it is a simple matter to show that
$`j_s^0(X)=2{\displaystyle \frac{d^3k}{(2\pi )^3}\frac{\omega _0}{\omega _{\mathrm{sc}}}f_s^{\mathrm{sc}}(๐ค;X)}`$
$`j_s^3(X)=2{\displaystyle \frac{d^3k}{(2\pi )^3}\frac{k_3}{\omega _0}f_s^0(๐ค;X)},`$ (20)
where the (momentum space) semiclassical distribution function $`f_s^{\mathrm{sc}}(๐ค;X)n_s(\omega _{\mathrm{sc}},๐ค;X)`$ is the projection of $`n_s`$ onto the semiclassical shell $`\omega =\omega _{\mathrm{sc}}`$, while the distribution function $`f_s^0(๐ค;X)n_s(\omega _0,๐ค;X)`$ is the projection of $`n_s`$ onto the unshifted shell $`\omega =\omega _0`$. Further, in Eq. (20) we have reintroduced the dependencies on the trivial perpendicular directions $`๐ค_{}`$. Finally, changing the variables $`k_3k_{\mathrm{sc}}`$, $`\omega _{\mathrm{sc}}\omega `$ ($`dk_3d\omega (\omega _{\mathrm{sc}}/dk_3)^1`$) in the first integral, and $`k_3k_0(\omega ^2|m|^2)^{1/2}`$, $`\omega _0\omega `$ ($`dk_3d\omega \omega _0/k_3`$) in the second integral, the current reduces to
$$j_s^\mu =2\frac{d\omega d^2k_{||}}{(2\pi )^3}(\frac{\omega }{k_{\mathrm{sc}}}f_s^{\mathrm{sc}},\widehat{๐ค}f_s^0),$$
(21)
where the distribution functions $`f_s^{\mathrm{sc}}=f_s^{\mathrm{sc}}(\omega ,๐ค_{};X)n_s(\omega ,k_{\mathrm{sc}},๐ค_{};X)`$ and $`f_s^0=f_s^0(\omega ,๐ค_{};X)`$ $`n_s(\omega ,k_0,๐ค_{};X)`$ are now defined on the energy phase space $`(\omega ,๐ค_{};X)`$. The semiclassical momentum $`k_{\mathrm{sc}}=k_{\mathrm{sc}}(\omega ,๐ค_{};X)`$ in Eq. (21) is the one obtained by inverting the semiclassical energy-momentum relation (18):
$$k_{\mathrm{sc}}=k_0\left(1+\frac{ss_\omega |m|^2\theta ^{}}{2k_0\omega ^2}\right).$$
(22)
Of course, in equilibrium the spatial currents cancel. However, if the preferred coordinate frame is boosted with respect to the plasma rest frame, as is the case for a propagating phase transition front, then there are net currents flowing through the space. One can then determine how these induced currents are affected by the presence of a background mass term (1). Since particles and antiparticles of a given spin (helicity) $`s`$ and energy $`|\omega |`$ have different effective momenta $`k_{\mathrm{sc}}`$, the corresponding particle and antiparticle densities $`j_s^0`$ are different in the region of a spatially varying background. This is how local particle-antiparticle asymmetric density and current profiles get created by a spatially varying and CP-violating background, which is the phenomenon at the heart of most electroweak baryogenesis mechanisms.
## 5 Field theoretical vs WKB approach
Let us now make a connection between the above results and the dispersion relation obtained by using the WKB method. Through a standard calculation , one finds
$$k_{\mathrm{wkb}}=k_0+\frac{ss_\omega \theta ^{}}{2k_0}(\omega \pm sk_0)+\alpha ^{},k_0=\sqrt{\omega ^2|m|^2},$$
(23)
where the arbitrary function $`\alpha ^{}`$ follows from the reparametrization invariance of the theory given by the lagrangian (2) invariant under (global) $`U(1)`$-transformations $`\psi e^{i\alpha }\psi `$. It is then obvious that $`k_{\mathrm{wkb}}`$ cannot represent a physical quantity. The group velocity
$$v_g(_\omega k_{\mathrm{wkb}})^1,$$
(24)
corresponding to the stationary phase of the WKB wave packet, is a well defined physical quantity however. The physical momentum of a WKB state is then given by $`k_{\mathrm{phys}}=\omega v_g(\omega )`$. Computing the derivative (24) one immediately finds the physical reparametrization invariant WKB dispersion relation
$$k_{\mathrm{phys}}\omega v_g=k_0\left(1+\frac{ss_\omega |m|^2\theta ^{}}{2k_0\omega ^2}\right),$$
(25)
which is identical to Eq. (22). With this the semiclassical momentum (22) can be reinterpreted as the physical momentum of a WKB state. In this sense we have established the equivalence between the physical WKB dispersion relation and the field theoretical energy-momentum shell corresponding to the propagator (8), or equivalently (12).
We now return to discuss the current in Eq. (21). The difference between the WKB approach and the field theoretical one presented here is the fact that there are two mutually independent distribution functions $`f_s^{\mathrm{sc}}`$ and $`f_s^0`$. In the same spirit as it was done for the scalar field in Ref. , we may introduce a coherent quantum density $`f_s^{\mathrm{qc}}=f_s^{\mathrm{sc}}f_s^0`$, which is a dynamical measure of quantum coherence on phase space between the semiclassical and the unshifted shell. At the leading order in gradients $`f_s^{\mathrm{qc}}`$ of course vanishes. In analogy to the scalar field case we then expect that in the frequent scattering limit $`f_s^{\mathrm{qc}}`$ gets suppressed and can be consistently neglected. Assuming this is the case, we may then set $`f_s^{\mathrm{qc}}f_s`$, $`f_s^0f_s`$, so that finally the current (21) reduces to the following form
$$j_s^\mu 2\frac{d\omega d^2k_{||}}{(2\pi )^3}(\frac{1}{v_g};\widehat{๐ค})f_s.$$
(26)
This limit is in perfect agreement with the naive WKB result , implying the following simple physical interpretation: when quasiparticles arrive into a region of a spatially varying background, they either speed up or slow down. If they slow down for example, the local particle density increasess proportionally to the factor $`1/v_g`$. At the same time the flux, given by the spatial part of the current, remains unchanged, as it should. To prove (26) rigorously requires derivation and solution of the relevant quantum transport equations to determine selfconsistently the distribution functions $`f_s^{\mathrm{sc}}`$ and $`f_s^0`$. This problem will be further considered in a forthcoming work .
While it was to be expected that the field theoretical calculation, unlike the WKB-computation, is automatically invariant under the field redefinition $`\psi e^{i\alpha }\psi `$, the fact that the two methods lead to identical dispersion relations and currents (in the frequent scattering limit) is highly nontrivial. Indeed, it should be stressed that the physical WKB-shell has emerged from the field theoretical calculation indirectly, not as a simple shift of the propagator poles, but rather by the operation of a modified spectral projector on test functions, which is effected via the technique of spectral integrals.
## 6 Conclusions
We have studied the fermion propagator coupled with a spatially varying classical background field. We represented this coupling by a spatially varying complex mass term, and proved that the propagator receives nontrivial corrections already at the first order in the gradients of the mass term $`m`$. For this result it was crucial that $`m`$ contains a nonvanishing complex pseudoscalar piece. We then explicitly constructed the fermion propagator and the associated spectral function to the first order in gradients. The gradient corrections were shown to arise, not as shifts in the poles, but rather as additional higher order poles multiplied by the gradients of the complex phase of $`m`$. We then considered a spectral integral of a simple test function, and showed that it nevertheless gets projected to a sharp energy shell, which acquires a spin dependent shift with respect to the lowest order shell. The new shell was then shown to coincide with the physical dispersion relation of a WKB-state. We also computed the current in the field theoretical way, and showed that it can be interpreted in terms of moving WKB wave packets. We note that it is straightforward to generalize our results to a generic pseudoscalar self-energy term and to temporally varying fields.
Here we have concentrated on the propagation of fermions in a nontrivial backgrounds. The next logical step is to study the transport properties in a plasma containing fermions. Drawing on the insights obtained from the study of the scalar case , the results presented here suggest that a consistent description of transport at nontrivial order in gradient expansion will require a set of coupled equations for number density and an additional function describing quantum coherence in phase space. We also expect that in the frequent scattering limit the quantum coherence effects are suppressed, and an effective quasiparticle picture is recovered .
## Acknowledgements
We wish to thank Felipe Freire for collaboration at the early stages of this project, and Dietrich Bรถdecker for illuminating discussions. KK thanks CERN for hospitality during the completion of this work. |
warning/0002/hep-ph0002291.html | ar5iv | text | # SELF-CONSISTENT APPROXIMATIONS: APPLICATION TO A QUASIPARTICLE DESCRIPTION OF THE THERMODYNAMIC PROPERTIES OF RELATIVISTIC PLASMAS
## 1 INTRODUCTION
Nuclear matter at very high densities is expected, because QCD is asymptotically free, to be in the form of a plasma of deconfined quarks and gluons. $`^\mathrm{?}`$ Such plasmas were present in the early universe and may exist in the cores of neutrons stars; current experiments aim to produce and study them in collisions of ultrarelativistic nuclei. $`^\mathrm{?}`$ In this talk, we address the question of providing a microscopic description of the thermodynamic properties of relativistic plasmas. The elementary excitations of weakly coupled quark-gluon plasmas are in fact, up to color factors, very similar in character to those of relativistic electromagnetic plasmas. Here, for simplicity we confine ourselves to the thermodynamic properties of QED plasmas. Unless otherwise stated, the properties of the electron and photon degrees of freedom which we consider can easily be translated into those of quark and gluon degrees of freedom in QCD plasmas.
Relativistic plasmas exhibit many of the familiar many-body effects encountered in non-relativistic condensed matter systems, including screening and the existence of collective modes. $`^\mathrm{?}`$ Debye screening shields the static electric component of the interaction but does not affect static magnetic fields. Screening of finite frequency fields arises from Landau damping mechanisms in a manner similar to the anomalous skin effect in metals. $`^\mathrm{?}`$ Early studies of relativistic plasmas have demonstrated the existence of bosonic and fermionic collective modes. $`^\mathrm{?}`$ The bosonic excitations are the familiar longitudinal and transverse plasmon oscillations and appear as poles in the photon propagator. The fermionic collective excitations are unique to relativistic systems and develop as poles in the low momentum region of the electron propagator. Investigations of the fermionic spectrum at one-loop order have revealed the existence of a rich structure, including a gap at zero momentum and two distinct excitation branches at small momenta. $`^\mathrm{?}`$
In the limit of small coupling constant $`g`$, these medium effects enter the theory through one-loop order corrections, which are systematically implemented by the โHard-Thermal-Loopโ (HTL) perturbation scheme of Braaten and Pisarski. $`^\mathrm{?}`$ In a plasma of massless particles at temperature $`T`$, medium effects generically develop over the long wavelength scale $`1/(gT)`$. (By comparison, the interparticle spacing is $`1/T`$.) Difficulties, whose solution lies beyond the HTL scheme, arise from the lack of static screening of magnetic interactions. The long-ranged nature of static magnetic fields causes the fermion quasiparticle damping rate to be logarithmically divergent at finite temperature in perturbation theory. Blaizot and Iancu have shown that in high temperature QED inclusion of multiple scattering ร la Bloch-Nordsieck leads to well-defined, divergence-free quasiparticle modes. $`^\mathrm{?}`$ In high temperature QCD, magnetic screening may actually appear at the scale $`1/g^2T`$ due to gluon self-interactions. The physics of the strong interactions at this scale is however non-perturbative and is thus again beyond the reach of HTL scheme. Related infrared difficulties due to long-ranged gauge fields arise in current issues of condensed matter theory, for instance in gauge field models of high $`T_c`$ superconductors and of the fractional quantum Hall effect. $`^\mathrm{?}`$
In this talk, we present a general framework for analyzing the effects of the gap in the fermion spectrum and of the long-ranged gauge fields on the thermodynamic properties of relativistic plasmas. Since we are dealing with interacting degrees of freedom, care is needed to avoid overcounting. We generalize to relativistic systems the $`\mathrm{\Phi }`$-derivable conserving approximations developed and Baym $`^\mathrm{?}`$ (see also $`^\mathrm{?}`$) in the context of quantum transport theories. As we shall see, this formulation permits the resolution of the entropy of a relativistic plasma into components relative to its elementary excitations. We will also show that the self-consistency of the approach implies subtle relationships between the contributions to the entropy from matter and interaction fields; these relationships are the manifestation of a proper counting of the degrees of freedom.
## 2 FROM THE THERMODYNAMICAL POTENTIAL TO THE ENTROPY
We illustrate the derivation for a hot relativistic QED plasma, a gas of electrons and positrons in equilibrium with photons at temperature $`Tm`$, the electron mass. For simplicity, we set $`m=0`$. At finite temperature, the system is not Lorentz invariant; only rotational invariance is obeyed. There is thus a preferred frame โ the rest frame of the system โ in which medium effects act separately on the longitudinal and transverse components of the interaction. It is thus natural to choose the Coulomb gauge. We work in the imaginary time formalism.
As we have shown $`^\mathrm{?}`$ the free energy $`\mathrm{\Omega }`$ can be written in terms of the fully dressed electron and photon propagators $`G`$ and $`D`$ as
$`\mathrm{\Omega }/T`$ $`=`$ $`\mathrm{\Phi }[G,D]\mathrm{Tr}\mathrm{\Sigma }G+\mathrm{Tr}\mathrm{log}(\gamma ^0G)+{\displaystyle \frac{1}{2}}\mathrm{Tr}\mathrm{\Pi }D{\displaystyle \frac{1}{2}}\mathrm{Tr}\mathrm{log}(D).`$ (1)
The traces are over the four momenta and over spin and polarization states, so for instance
$`\mathrm{Tr}\mathrm{\Sigma }G\mathrm{tr}{\displaystyle \underset{p,n}{}}\mathrm{\Sigma }(\omega _n,p)G(\omega _n,p),\mathrm{Tr}\mathrm{\Pi }D\mathrm{tr}{\displaystyle \underset{q,n}{}}\mathrm{\Pi }(\omega _n,q)D(\omega _n,q),`$ (2)
where $`\omega _n=(2n+1)i\pi T`$ and $`\omega _n=2ni\pi T`$ are the electron and photon Matsubara frequencies, and $`\mathrm{tr}`$ represents the traces over spins or polarizations. The second and third terms in Eq. (1) represent a contribution from the electron degrees of freedom and have a form similar to those encountered in non-relativistic version of conserving approximations. $`^\mathrm{?}`$ The generalization to relativistic systems is implemented here by the last two terms, which correspond to treating the electromagnetic interaction as a dynamical degree of freedom. The coupling between electron and photon modes is described by the functional $`\mathrm{\Phi }[G,D]`$, the sum of all two-particle-irreducible skeleton diagrams, expressed in terms of fully dressed $`G`$ and $`D`$ instead of the bare electron and photon propagators $`G_0`$ and $`D_0`$. Under a simultaneous variation of the propagators $`G`$ and $`D`$,
$`\delta \mathrm{\Phi }[G,D]`$ $`=`$ $`\mathrm{Tr}\mathrm{\Sigma }\delta G{\displaystyle \frac{1}{2}}\mathrm{Tr}\mathrm{\Pi }\delta D,`$ (3)
as one easily sees by removing an electron or photon line in a given diagram contributing to $`\mathrm{\Phi }`$. The electron self-energy $`\mathrm{\Sigma }[G,D]`$ and the polarization operator $`\mathrm{\Pi }[G,D]`$ are here functionals of $`G`$ and $`D`$ and satisfy Dysonโs equations
$`G^1=G_0^1\mathrm{\Sigma },D^1=D_0^1\mathrm{\Pi }.`$ (4)
We emphasize that the representation of the thermodynamical potential in Eq. (1) is exact. This result is based upon two essential properties of the $`\mathrm{\Phi }`$ functional. The first given in Eq. (3), implies that the potential $`\mathrm{\Omega }`$ is stationary under variations of $`G`$ and $`D`$ that do not modify the free propagators $`G_0`$ and $`D_0`$. The second property follows from the topology of the diagrams contributing to $`\mathrm{\Phi }`$. As each vertex of a given $`\mathrm{\Phi }`$-diagram is connected to two electron lines and one photon line, the functional $`\mathrm{\Phi }`$ remains constant under the scaling transformations
$`\mathrm{\Phi }[G,D;g]`$ $`=`$ $`\mathrm{\Phi }[s^fG,s^{2f2v}D;s^vg],`$ (5)
where $`g`$ is the coupling constant at a vertex. To gain insight into the meaning of this property, we note that when combined with Eq. (3) it implies $`\mathrm{Tr}\mathrm{\Sigma }G=\mathrm{Tr}\mathrm{\Pi }D`$. This is one example of the interdependencies that arise between the electron and photon contributions to thermodynamic quantities when treated self-consistently. The traces $`\mathrm{Tr}\mathrm{\Sigma }G`$ and $`\mathrm{Tr}\mathrm{\Pi }D`$ both represent the interaction energy between electrons and photons. These expressions are equal if $`\mathrm{\Phi }`$ satisfies the scaling invariance of Eq. (5).
The principle behind the conserving approximations consists of choosing a physically motivated subset of diagrams contributing to $`\mathrm{\Phi }`$, and deriving the corresponding self-energy functionals via Eq. (3). Dysonโs equations then provide self-consistent equations for the electron and photon propagators. These approximations preserve the functional and topological properties of $`\mathrm{\Phi }`$, and a proper counting of the degrees of freedom is ensured. In particular, it follows from the stationarity property of the thermodynamical potential $`\mathrm{\Omega }`$ that the approximations obtained for the propagators $`G`$ and $`D`$ lead to current densities and an energy momentum tensor that obey the continuity equations for the conservation of charge, momentum, and energy $`^\mathrm{?}`$.
As we shall see shortly, the entropy $`S=\mathrm{\Omega }/T|_{\mu ,V}`$ is a quantity well-suited to be analyzed in terms of its quasiparticle content. Converting the frequency sums into integrals in the usual way<sup>a</sup><sup>a</sup>aDue to the zero Matsubara mode, the sum over photon frequencies becomes a principal value integral., we write $`S`$ as
$`S={\displaystyle \frac{\mathrm{\Omega }}{T}}|_{\mu ,V}`$ $`=`$ $`S_f+S_b+S^{},`$ (6)
where
$`S_f`$ $``$ $`{\displaystyle \underset{p\pm }{}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \frac{d\omega _p}{\pi }}{\displaystyle \frac{f}{T}}\left(\mathrm{Im}\mathrm{\Sigma }_\pm \mathrm{Re}G_\pm +\mathrm{Im}\mathrm{log}(G_\pm ^1)\right),`$ (7)
$`S_b`$ $``$ $`{\displaystyle \underset{q,l}{}}\text{ }\text{ }{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{d\omega _q}{\pi }}{\displaystyle \frac{n}{T}}\left(\mathrm{Im}\mathrm{\Pi }_l\mathrm{Re}D_l+\mathrm{Im}\mathrm{log}(D_l^1)\right),`$ (8)
$`S^{}`$ $``$ $`{\displaystyle \frac{(T\mathrm{\Phi })}{T}}{\displaystyle \underset{p\pm }{}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \frac{d\omega _p}{\pi }}{\displaystyle \frac{f}{T}}\mathrm{Im}G_\pm \mathrm{Re}\mathrm{\Sigma }_\pm {\displaystyle \underset{q,l}{}}\text{ }\text{ }{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{d\omega _q}{\pi }}{\displaystyle \frac{n}{T}}\mathrm{Im}D_l\mathrm{Re}\mathrm{\Pi }_l`$
and where $`f=(\mathrm{exp}\beta (\omega _p\mu )+1)^1`$ and $`n(\omega _q)=(\mathrm{exp}\beta \omega _q1)^1`$ are Fermi and Bose occupation factors. The sum over internal degrees of freedom includes in $`S_f`$ a sum over states with helicity equal ($`+`$) and opposite to ($``$) their chirality, while it involves in $`S_b`$ a sum over the two transverse and the longitudinal polarizations. We note that in the expressions above only the temperature derivatives acting on occupation factors need to be retained. All other temperature dependences cancel out as a consequence of the stationarity of $`\mathrm{\Omega }`$$`^\mathrm{?}`$ The terms $`S_f`$ and $`S_b`$ represent the contributions from the electron and photon elementary modes. The term $`S^{}`$, which describes residual interactions between quasiparticles, has the particular property that it vanishes at one-loop order, $`^\mathrm{?}`$ corresponding to the $`\mathrm{\Phi }`$-diagram and the self-energies in Fig. 1. In the following, we will stay at this level of approximation and neglect $`S^{}`$. An example of a higher-order interaction term $`S^{}`$ appears in the description of the thermodynamic properties of <sup>3</sup>He, where it describes the effects of repeated scattering of particle-hole pairs and gives rise to the $`T^3\mathrm{log}T`$ term in the low temperature specific heat. $`^\mathrm{?}`$
## 3 QUASIPARTICLE ANALYSIS OF THE ENTROPY OF AN ELECTROMAGNETIC PLASMA
### 3.1 A quasiparticle description of the thermodynamics of a QED plasma
The reasons for choosing the entropy as the thermodynamic quantity to be analyzed in terms of its quasiparticle content are twofold. First, a simple inspection of Eqs. (7)-(LABEL:sp) shows that, because of the presence of occupation factors in the integrand, the entropy does not suffer from ultraviolet divergences. Potential renormalization problems are thus avoided as all other thermodynamic quantities can be obtained from the entropy by simple integration. Second, a spectral representation can be introduced by recognizing that
$`{\displaystyle \frac{f}{T}}`$ $`=`$ $`{\displaystyle \frac{\sigma _f}{\omega _p}}(\omega _p),{\displaystyle \frac{n}{T}}={\displaystyle \frac{\sigma _n}{\omega _q}}(\omega _q),`$ (10)
where $`\sigma _ff\mathrm{log}f(1f)\mathrm{log}(1f)`$ and $`\sigma _nn\mathrm{log}n+(1+n)\mathrm{log}(1+n)`$ are the entropy contributions from a single electron mode with energy $`\omega _p`$ and from a single photon mode with frequency $`\omega _q`$. Integrating Eqs. (7) and (8) by parts, we obtain the spectral representations
$`S_f`$ $`=`$ $`{\displaystyle \underset{p}{}}{\displaystyle \frac{d\omega _p}{2\pi }\sigma _f(\omega _p)A_s(\omega _p,p)},S_b={\displaystyle \underset{q}{}}\text{ }\text{ }{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{d\omega _q}{2\pi }}\sigma _b(\omega _q)B_s(\omega _q,q)`$ (11)
where the spectral density $`A_s`$ is defined as
$`A_s(\omega _p,p)`$ $`=`$ $`{\displaystyle \underset{\pm }{}}{\displaystyle \frac{}{\omega _p}}\left\{\mathrm{Re}G_\pm \mathrm{Im}\mathrm{\Sigma }_\pm +2\mathrm{I}\mathrm{m}\mathrm{log}(G_\pm )\right\},`$ (12)
and $`B_s`$ has a similar structure in terms of $`D_l`$ and $`\mathrm{\Pi }_l`$.
The physical content of the spectral densities can be elucidated as follows. For brevity, we only consider $`A_s`$. If the system develops well-defined excitations, the fermion propagator takes the following form near a quasiparticle pole:
$`G_\pm (\omega _p+i0^+,p){\displaystyle \frac{Z_\pm }{\omega _p\epsilon _p^\pm +i\mathrm{\Gamma }_p^\pm /2}},`$ (13)
where $`Z_\pm `$ is the residue of the quasiparticle with energy $`\epsilon _\pm `$ and inverse lifetime $`\mathrm{\Gamma }_p^\pm =2\mathrm{I}\mathrm{m}\mathrm{\Sigma }_\pm `$. When the variation of $`\mathrm{\Gamma }_\pm `$ around the quasiparticle pole is negligible, the density $`A_s`$ near the pole takes the form
$`A_{s\pm }(\omega _p,p)`$ $``$ $`{\displaystyle \frac{(Z_\pm \mathrm{\Gamma }_\pm )^3/2}{((\omega _p\epsilon _p^\pm )^2+(Z_\pm \mathrm{\Gamma }_\pm /2)^2)^2}}.`$ (14)
Thus, in the non-interacting limit $`\mathrm{\Gamma }_\pm 0`$, $`A_s`$ reduces to a delta function and $`S_f`$ in Eq. (11) coincides with the entropy of a gas of free electrons with energies $`\epsilon _p^\pm `$. When $`\mathrm{\Gamma }_\pm `$ is different from zero but small, so that well-defined quasiparticle modes exist, $`A_s`$ is a sharply peaked function which falls off faster than a Lorentzian. The excited quasiparticles as well as the degrees of freedom associated with their finite lifetimes thus contribute to the entropy $`S_f`$. The form of $`B_s`$ close to photon quasiparticles is qualitatively similar to that of $`A_s`$. Both $`A_s`$ and $`B_s`$ also have support over wide bands of continuum states, which describe the contributions from Landau damping effects. In higher order approximations, the continuum bands also include multiparticle states.
### 3.2 Comparison to perturbative expansions
The present $`\mathrm{\Phi }`$ conserving approximations also provide a general framework for organizing calculations in perturbation theory. To illustrate we comment on the effects of the infrared structure and the presence of a gap in the fermion spectrum on the plasma thermodynamics. As mentioned earlier, the lack of static screening of magnetic fields is responsible for a logarithmic divergence of the damping rate evaluated in perturbation theory, $`\mathrm{\Gamma }_\pm (\omega _p)g^2T\mathrm{log}(gT/|\omega _p\epsilon _p^\pm |)`$. This behavior implies that $`A_s`$ vanishes at the quasiparticle pole, since $`A_s1/\mathrm{\Gamma }1/\mathrm{log}(gT/|\omega _p\epsilon _p^\pm |)`$ as $`\omega _p\epsilon _p^\pm `$. There is thus no well-defined quasiparticles at this level of approximation; a correct quasiparticle description of the plasma thermodynamics must include a proper treatment of the long-ranged gauge fields for frequencies close to the poles $`\epsilon _p^\pm `$.
The effects of the gap in the fermion spectrum on thermodynamic quantities can also readily be estimated. Inspection of $`S_f`$ in Eq. (7) shows that low momentum modes with $`\omega _ppgT`$ contribute a factor $`g^5T^3`$ to the entropy per unit volume, in contrast with the contribution from the long-wavelength plasmon modes, $`S_{\mathrm{plasmons}}g^3T^3`$$`^\mathrm{?}`$ The difference in the order of magnitude can be traced back to statistics. For low frequencies $`\omega _p,\omega _qgT`$, the temperature derivative $`f/T\omega _p/T^2`$ in Eq. (7) is two powers of g smaller than the derivative $`n/T1/\omega _q`$ in Eq. (8). Hence $`S_fg^2S_{\mathrm{plasmons}}`$.
Further comparisons between our expressions and the standard results of perturbation theory reveal relations that follow from the self-consistent character of the derivation. For the formalism to reproduce all terms in the plasmon contribution to the entropy, $`g^3T^3`$, one needs to include contributions from both long wavelength photons, via $`S_b`$, and electron modes, via $`S_f`$. These last terms describe the effects of the plasmons on the single particle energy of the electrons, via the real part of the self-energy, and are automatically included in a self-consistent treatment. The approach presented in this talk has recently been applied to purely gluonic QCD by Blaizot et al., $`^\mathrm{?}`$ who have succeeded in obtaining an accurate quasiparticle description of lattice data for the equation of state down to temperatures twice $`T_c`$. There again, it is important to take into account the effects of soft modes on the dispersion relation of the hard modes to reproduce the entropy correctly up the order $`g^3`$. A description of these effects in terms of microscopic processes can be found in Ref. $`^\mathrm{?}`$.
## 4 CONCLUSIONS
In this talk, we generalized the concept of conserving approximations to relativistic theories. We have illustrated the technique for a hot electromagnetic plasma and shown how the approach allows one to resolve the entropy into contributions from the elementary degrees of freedom. We have also pointed out the advantages of implementing $`\mathrm{\Phi }`$-conserving techniques as a basis for perturbation expansions.
## Acknowledgements
This research was supported in part by NSF Grants PHY98-00978 and PHY94-21309.
## References |
warning/0002/astro-ph0002133.html | ar5iv | text | # Models of Disk Evolution: Confrontation with Observations
## 1 Introduction
Over the last several years, there has been a steady increase in the number and quality of observations available for disk galaxies from $`z=0`$ and $`z=1`$. Schade et al. (1995,1996), using early ground and space based images of galaxies from the Canada-France Redshift Survey (CFRS), found a net increase in the surface brightness of galaxies to $`z1`$. Along the same lines, Roche et al. (1998), compiling 347 galaxies from the Medium Deep Survey and other surveys, concluded that disk galaxies had undergone a net evolution in surface brightness and a net devolution in size. Lilly et al. (1998), using structural parameters extracted from HST images of the combined CFRS and LDSS2 sample, concluded that there has been essentially no evolution in large disks out to $`z1`$. As a preliminary effort as part of the DEEP survey, Vogt et al. (1996,1997) found little evolution in the Tully-Fisher relationship ($`<0.3^m`$) out to $`z1`$. More recently, these observations have been augmented by the DEEP sample with 197 galaxies from the Groth strip to $`I<23.5`$, 1.5 magnitudes deeper than the LDSS2-CFRS sample. In a first paper, Simard et al. (1998) concluded that there had been little evolution in the disk surface brightness distribution to $`z1`$ contrary to previous claims.
A number of different approaches have been proposed for making specific predictions about disk evolution. Mo, Mao, & White (1998a) showed how the standard paradigm for hierarchical growth of structure combined with simple assumptions about angular momentum conservation led to simple scaling relationships for the change in disk properties as a function of redshift. Other authors (Ferrini et al. 1994; Prantzos & Aubert 1995; Prantzos & Silk 1998; Boissier & Prantzos 1999; Chiappini, Matteucci & Gratton 1997), taking more of a backwards approach to the problem, used detailed studies of the profiles of the Milky Way and other nearby galaxies to propose radially dependent models of star formation in disk galaxies, models which could be used to make detailed predictions about high-redshift disk evolution.
Already there have been a number of elegant studies in which both the backwards approach (Cayรณn, Silk, & Charlot 1996; Bouwens, Cayรณn, & Silk 1997; Roche et al. 1998) and the forwards approach (Mao, Mo, & White 1998; Steinmetz & Navarro 1999; Contardo, Steinmetz, & Fritze-von Alvensleben 1998; van den Bosch 1998; Mo, Mao, & White 1998b) have been used to interpret the observations available for disk galaxies, mostly to $`z1`$.
Unfortunately, none of these studies considered the important effect that a large spread in surface brightness could have on the interpretation of these observations, particularly the potentially large fraction of low surface brightness galaxies. In some studies, the surface brightness selection effects at low and high redshift were simply ignored, and in others, e.g., Roche et al. (1998), the spread was limited to $`0.3\text{mag/arcsec}^2`$ about Freemanโs law (Freeman 1970). Clearly, given the apparent large numbers of low surface brightness galaxies seen locally, it is quite logical to wonder if these galaxies are detectable in current high redshift surveys. Indeed, one might wonder whether these galaxies or the observed correlation between luminosity and surface brightness may have already affected the interpretation of high redshift observations. In light of the recent claim by Simard et al. (1999) that the apparent surface brightness evolution thus far inferred to $`z1`$ is completely due to surface brightness selection effects, such a study would seem to be especially timely. Secondly, none of these studies directly compared the predictions of the forward and backward approaches using the same observations. Simple comparisons of the scaling expected in surface brightness, size, luminosity, and number are useful for interpreting the high redshift observations.
To address these shortcomings, we shall therefore consider implementations of both approaches, normalize them to the observed $`z0`$ size-luminosity relationship, compare their predictions, and consider how each of them fares at explaining the observed disk evolution out to $`z1`$ incorporating all the selection effects as they are best understood. We commence by presenting our models (ยง2) and the observational samples with which we compare (ยง3). We present the results (ยง4), discuss them (ยง5), and then summarize our conclusions (ยง6). Throughout this study, we use $`H_0=50\text{km/s/Mpc}`$ unless otherwise noted.
## 2 Models
We begin by sketching the base $`z=0`$ model to be used for both the models which follow. We use a set of gaussian LFs based on those presented in Binggeli, Sandage, and Tammann (1988):
$$\varphi (M)dM=\frac{\varphi _0}{2\pi \sigma _M}\mathrm{exp}((\frac{MM_{}}{\sigma _M})^2)dM$$
(1)
We adjusted the bulge-to-total ($`B/T`$) distributions of these galaxy types to obtain fair agreement with the de Jong & van der Kruit (1994) sample. Finally, we adjusted the luminosity function so there was rough agreement with the combined Sabc and Sdm luminosity functions presented in Pozzetti, Bruzual, & Zamorani (1996). We present our parameterized populations in Table 1.
For the above luminosity functions, we convert $`z=0`$ $`B`$-band luminosity to mass using a constant mass-to-light ratio, where the mass of a $`M_{b_J}=21.1`$ galaxy is $`1.110^{12}M_{}`$. We assume a log-normal scatter of 0.3 dex to reproduce the observed Tully-Fisher scatter though variation in the formation times (van den Bosch 1998) and concentration indexes (Avila-Reese, Firmani, & Hernandez 1998) certainly play a role.
To translate this mass into a circular velocity and size, we calculate the time at which the ambient halo formed. Since halos are always accreting more mass and merging with larger halos, there is some ambiguity in defining this, so for simplicity we take it to equal the redshift at which half of the mass in a halo has been assembled. We determine the distribution of formation times using the procedure outlined in Section 2.5.2 of Lacey & Cole (1993). We take the circular velocity of the halo to be that corresponding to the halo at its formation time using Eq. (14).
Then, given the mass, circular velocity, and luminosity of the $`z=0`$ disk, we randomly draw the sizes $`r_e`$ from the following distribution:
$$\varphi (r_e)d\mathrm{log}r_e=\frac{1}{\sigma _\lambda \sqrt{2\pi }}\mathrm{exp}(\frac{1}{2}[\frac{\mathrm{log}r_e/r_e^{}0.4(MM_{,s})(1/3)}{\sigma _\lambda /\mathrm{ln}(10)}]^2)d\mathrm{log}r_e$$
(2)
where $`r_e^{}=6.9\text{kpc}`$, $`\sigma _\lambda =0.37`$, and $`M_{,s}=21.1`$. We provide a basic observational and theoretical motivation for this scaling in ยง2.1. Note that here the surface brightness is proportional to $`L^{1/3}`$, that the spread in the size distribution is proportional to the spread in the distribution of $`\lambda `$, and that the scale length (surface brightness) of the average $`L_{}`$ galaxy is exactly equal to that predicted by Freemanโs law.
We assume that the SED of disks and bulges is identical to that of a 10 Gyr-old stellar population with an e-folding time of 4.5 Gyr and a $`\tau _B=0.3`$ foreground dust screen, the extinction curve being that of Calzetti (1997). For simplicity, we assume this SED is constant independent of time. We use the Bruzual & Charlot tables from the Leitherer et al. (1996) compilation for this calculation. In the following models, we evolve the size, number, and luminosity of all galaxy types using simple single-valued functions of redshift:
$$R(z)=R(0)E_R(z)$$
(3)
$$N(z)=N(0)E_N(z)$$
(4)
$$L(z)=L(0)E_L(z)$$
(5)
For simplicity, we assume similar scalings in the properties of bulges as a function of time.
Since the color and luminosity of disks has been shown to be correlated with inclination, it is reasonable to suppose that disks are not transparent. Unfortunately, there is much controversy concerning the degree to which disks are or are not transparent. For better or worse, we will side-step this controversy and simply adopt the Tully & Fouquรฉ (1985) prescription for extinction in the $`B`$ band:
$$A_B=2.5\mathrm{log}(f(1\mathrm{exp}(\tau \mathrm{sec}i))+(12f)(\frac{1\mathrm{exp}(\tau \mathrm{sec}i)}{\tau \mathrm{sec}i}))$$
(6)
where $`\tau =0.55`$, $`f=0.25`$, and $`i`$ is the inclination of the disk, 0 corresponding to a face-on disk. We shall assume our model galaxies are always observed at an inclination of $`70\mathrm{deg}`$ and therefore always correct the observations to this inclination for comparison with the models. In the $`B`$ band, this corresponds to an extinction correction of $`0.67^m`$.
### 2.1 Hierarchical Model (Forwards Approach)
The use of simple scaling relationships between the properties of disks and the halos in which they live has provided a relatively successful way of explaining both the internal correlations between disk properties and their evolution to high redshift. In this picture developed by Fall & Estathiou (1980) and revived more recently by Dalcanton et al. (1997) and Mo et al. (1998) among others, the bivariate mass and angular momentum distribution nicely translates into a luminosity and surface brightness relationship for disk galaxies, mass translating directly into luminosity and the dimensionless angular momentum translating directly into surface brightness.
This picture has had much success in explaining the internal correlations between size, circular velocity, and mass. For example, De Jong & Lacey (1999) recently showed that the observed local bivariate luminosity-size distribution is nicely fit by this picture, albeit with a slightly smaller scatter in surface brightness than might otherwise be expected. For constant mass-to-light ratios, the rough $`MV_c^3`$ relationship for halos provides a relatively natural explanation for the luminosity-circular velocity (Tully-Fisher) relationship (Dalcanton, Spergel, & Summers 1997a; Mo et al. 1998a; Steinmetz & Navarro 1999; Contardo et al. 1998). Finally, the $`RV_c`$ relationship found in galaxy halos is similarly observed in the disk population (Courteau 1997).
On the other hand, this picture says little, if anything, about how the gas disk evolved over time and therefore what its local properties (i.e., spatial variations in the metallicity, color, stellar ages, gas density, etc.) or global properties (total gas mass) are, and therefore comparisons of this sort will depend upon the model adopted, whether it be one of the popular semi-analytic approaches (Cole et al. 1999, Somerville & Primack 1998) or a full N-body hydrodynamical simulation (e.g., Contardo et al. 1998). Moreover, it is now apparent that the simple scaling model is fundamentally flawed with regard to the implementation of galaxy formation theory as revealed by high resolution numerical simulations (Moore et al. 1999; Navarro & Steinmetz 1999; Steinmetz & Navarro 1999). Since it however is the only detailed model available, it is imperative to fully explore comparisons with data, properly incorporating observational selection effects, in order to establish the correct basis for ultimately refining the model.
For a detailed discussion of this picture, i.e., the idea that simple scalings in the properties of halos lead to simple scalings in the properties of disks, the reader is referred to Mo et al. (1998a) and later papers by the same authors. For the sake of clarity, we shall review some of this material. In the standard spherical collapse model for an Einstein-de Sitter universe, the density of the collapsed halo is $`18\pi ^2178`$ times the critical density of the universe at collapse time (see also Gunn & Gott 1972; Bertschinger 1985; Cole & Lacey 1996), but depends on the density of universe through the parameter $`x=1\mathrm{\Omega }(z)`$. Expressing the result in terms of the super critical density parameter $`\mathrm{\Delta }_c`$
$$\frac{M}{\frac{4}{3}\pi r_{vir}^3}=\mathrm{\Delta }_c\rho _c$$
(7)
Bryan & Norman (1998) found that for $`\mathrm{\Omega }+\mathrm{\Omega }_\mathrm{\Lambda }=1`$,
$$\mathrm{\Delta }_c18\pi ^2+82x39x^2$$
(8)
for $`\mathrm{\Omega }_\mathrm{\Lambda }=0`$,
$$\mathrm{\Delta }_c18\pi ^2+60x32x^2$$
(9)
where $`x=\mathrm{\Omega }(z)1`$.
Using the virial theorem, it is possible to write equations to relate the mass, radius, and circular velocity of each halo. As in Somerville & Primack (1998), it can be shown that
$$V_{vir}^2=\frac{GM}{r_{vir}}\frac{\mathrm{\Omega }_\mathrm{\Lambda }}{3}H(z)^2r_{vir}^2$$
(10)
where $`r_{vir}`$ is the halo size, $`V_{vir}`$ is the circular velocity of the halo at $`r_{vir}`$, $`G`$ is Newtonโs constant, and
$$H(z_f)=H_0\sqrt{\mathrm{\Omega }_{\mathrm{\Lambda },0}+(1\mathrm{\Omega }_0\mathrm{\Omega }_{\mathrm{\Lambda },0})(1+z_f)^2+\mathrm{\Omega }_0(1+z_f)^3}.$$
(11)
Using the fact that $`M=\mathrm{\Delta }_c\rho _c\frac{4}{3}\pi r_{vir}^3=\mathrm{\Delta }_cH(z)^2\frac{1}{2G}r_{vir}^3`$, we can rewrite this as
$$V_{vir}^2=\frac{1}{2}(\mathrm{\Delta }_c\mathrm{\Omega }_\mathrm{\Lambda })H(z)^2r_{vir}^2$$
(12)
or
$$r_{vir}=\frac{V_{vir}}{\sqrt{\frac{1}{2}(\mathrm{\Delta }_c\mathrm{\Omega }_\mathrm{\Lambda })}H(z)}$$
(13)
Similarly, we can now rewrite the halo mass as
$$M=\frac{V_{vir}^2r_{vir}}{G}=\frac{V_{vir}^3}{GH(z)\sqrt{\frac{1}{2}(\mathrm{\Delta }_c\mathrm{\Omega }_\mathrm{\Lambda })}},$$
(14)
Assuming the matter which settles in the disk to be some fraction $`m_d`$ of the halo mass and the angular momentum of this settling matter to be some fraction $`j_d`$ of the haloโs angular momentum, a straightforward derivation (e.g., Mo et al. 1998a) allows one to obtain
$$M_d=\frac{m_dV_{vir}^3}{GH(z_f)\sqrt{\frac{1}{2}(\mathrm{\Delta }_c\mathrm{\Omega }_\mathrm{\Lambda })}}$$
(15)
for the mass of the disk and
$$R_d=\frac{1}{\sqrt{2}}\left(\frac{j_d}{m_d}\right)\lambda r_{vir}$$
(16)
for the radius of the disk. The dimensionless angular momentum parameter $`\lambda `$ is defined as
$$\lambda =J|E|^{1/2}G^1M^{5/2}$$
(17)
where $`J`$ is angular momentum, $`M`$ is the mass, and $`E`$ is the total energy of the bound system.
There are three relatively simple reasons to go beyond this simple approach. First, the adiabatic contraction of the halo due to dissipation of baryons towards the halo center will modify the halo profile. Second, numerical simulations show that model halos actually have a Navarro, Frenk, & White (1997) profile for the the dark halo rather than the isothermal profile used above. Finally, in order to make comparisons back to the observations, it is important to consider the observationally-measured rotational velocities of the disk rather than the rotational velocities of the halo proper. Mo et al. (1998a) have found approximate fitting formulas for the consequent corrections made to the disk radius $`R_d`$ and the circular velocity at $`3R_d`$:
$$R_d=\frac{1}{\sqrt{2}}\left(\frac{j_d}{m_d}\right)\lambda r_{vir}f_c^{1/2}f_R$$
(18)
$$V_c(3R_d)=V_{vir}f_V$$
(19)
where approximate fitting functions for $`f_c`$, $`f_R`$, and $`f_V`$ are given by
$$f_c\frac{2}{3}+\left(\frac{c}{21.5}\right)^{0.7}$$
(20)
$$f_R\left(\frac{\lambda }{0.1}\right)^{0.06+2.71m_d+0.0047/\lambda }(13m_d+5.2m_d^2)(10.019c+0.00025c^2+0.52/c)$$
(21)
$$f_V\left(\frac{\lambda }{0.1}\right)^{2.67m_d0.0038/\lambda +0.2\lambda }(1+4.35m_d3.76m_d^2)\frac{1+0.057c0.00034c^21.54/c}{\left[c/(1+c)+\mathrm{ln}(1+c)\right]^{1/2}}$$
(22)
where $`c`$ is the standard halo concentration parameter for the Navarro et al. (1997) profile. For simplicity, we use $`m_d=0.05`$ and $`c=10`$ to convert $`V_{vir}`$ to $`V_c`$ in order to compare with the observations. We do not use Eqs. (15)-(16) for these comparisons. Note that larger values of $`m_d`$ render disks unstable at relatively faint surface brightnesses and thus have difficulty accounting for Freeman Law-type surface brightnesses (Freeman 1970).
On the basis of these simple halo scaling relations, the size and mass of disks at any redshift simply scales as $`1/H(z_f)\sqrt{\frac{1}{2}\mathrm{\Delta }_c(z_f)\mathrm{\Lambda }(z_f)}`$, $`z_f`$ being the redshift at which these high-redshift disks formed. Consequently, the size and luminosity scale as
$$r(z)=\frac{V_{vir}}{\sqrt{\frac{1}{2}(\mathrm{\Delta }_c(z_f(z))\mathrm{\Omega }_\mathrm{\Lambda }(z_f(z)))}H(z_f(z))}$$
(23)
$$L(z)=\frac{L}{M}(z)M(z)=\frac{L}{M}(z)\frac{V_{vir}^3}{GH(z_f(z))\sqrt{\frac{1}{2}(\mathrm{\Delta }_c(z_f(z))\mathrm{\Omega }_\mathrm{\Lambda }(z_f(z))}}f_V$$
(24)
Using Eq. (3) and (4), we now have our functions $`E_R(z)`$ and $`E_L(z)`$:
$$E_R(z)=\frac{H(z_f(0))\sqrt{\frac{1}{2}(\mathrm{\Delta }_c(z_f(0))\mathrm{\Omega }_\mathrm{\Lambda }(z_f(0)))}}{H(z_f(z))\sqrt{\frac{1}{2}(\mathrm{\Delta }_c(z_f(z))\mathrm{\Omega }_\mathrm{\Lambda }(z_f(z)))}}$$
(25)
$$E_L(z)=\frac{\gamma (0)H(z_f(0))\sqrt{\frac{1}{2}(\mathrm{\Delta }_c(z_f(0))\mathrm{\Omega }_\mathrm{\Lambda }(z_f(0)))}}{\gamma (z)H(z_f(z))\sqrt{\frac{1}{2}(\mathrm{\Delta }_c(z_f(z))\mathrm{\Omega }_\mathrm{\Lambda }(z_f(z)))}}$$
(26)
where $`\gamma (z)`$ is the mass-to-light ratio at redshift $`z`$. Note that this is quite different from scaling these disks simply in terms of the redshift at which these disks were observed, particularly in the case of low $`\mathrm{\Omega }`$ where little evolution in the size or baryonic mass of the disk population is expected.
We assume that the $`z=0`$ luminosity function scales in number as a function of $`z`$ in an analogous way to how the $`10^{12}M_{}`$ halos scale in number. Using the Press-Schechter (Press & Schechter 1974) mass function
$$N(M,z)dM=\sqrt{\frac{2}{\pi }}\frac{\rho _0}{M}\frac{\delta _c}{\sigma (M)D(z)}\mathrm{exp}\left(\frac{\delta _c^2}{2\sigma ^2(M)D(z)^2}\right)\frac{d\sigma (M)}{dM}dM$$
(27)
we see that the halo number density scales as
$$n[D(z)]^1$$
(28)
since the exponential factor remains approximately unity for the $`10^{12}M_{}`$ mass scale. In terms of the formalism of Eqs. (3-5),
$$E_N(z)=\frac{D(0)}{D(z)}$$
(29)
Here, $`D(z)`$, the growth factor, was computed using the formula tabulated in Carroll, Press, & Turner (1992).
We now provide a theoretical and observational justification for our size-luminosity distribution. Theoretically, in the Fall & Estathiou (1980) picture, the spread in surface brightnesses derives from the spread in dimensionless angular momenta for halos. An approximate parameterization of the dimensionless angular momentum distribution is
$$p(\lambda )=\frac{1}{\sqrt{2\pi \sigma _\lambda }}\mathrm{exp}\left[\frac{\mathrm{ln}(\lambda /\overline{\lambda })^2}{2\sigma _\lambda ^2}\right]\frac{d\lambda }{\lambda }$$
(30)
For $`\overline{\lambda }=0.05`$ and $`\sigma _\lambda =0.5`$, the above expression closely approximates the distribution obtained from N-body simulations (Warren et al. 1992; Cole & Lacey 1996; Catelan & Theuns 1996) and analytical treatments (Steinmetz & Bartelmann 1995). The above spread in dimensionless angular momentum directly translates into the following distribution of sizes:
$$\varphi (r)dr=\frac{1}{\sqrt{2\pi \sigma _\lambda }}\mathrm{exp}\left[\frac{\mathrm{ln}(r/\overline{r_e})^2}{2\sigma _\lambda ^2}\right]d\mathrm{log}r$$
(31)
For disks with a constant mass-to-light ratio, it follows from Eqs. (13-16) that $`r_dr_{vir}V_{vir}H(z_f)^1M_d^{1/3}H(z_f)^{2/3}L_d^{1/3}H(z_f)^{2/3}`$. Ignoring the dependence of $`H(z_f)^{2/3}`$ on the luminosity, it follows that the surface brightness ($`L_d/r_d^2`$) scales as $`L_d^{1/3}`$.
In fact, de Jong & Lacey (1998) found that the Mathewson, Ford, & Buchhorn (1992) data set gave a good fit to the following bivariate size-luminosity distribution with similar properties to those predicted above:
$`\mathrm{\Phi }(r_e,M)d\mathrm{log}r_edM={\displaystyle \frac{\mathrm{\Phi }_0}{\sigma _\lambda \sqrt{2\pi }}}\mathrm{exp}({\displaystyle \frac{1}{2}}[{\displaystyle \frac{\mathrm{log}r_e/r_e^{}0.4(MM_{})(2/\beta 1)}{\sigma _\lambda /\mathrm{ln}(10)}}]^2)`$ (32)
$`10^{0.4(MM_{})(\alpha +1)}\mathrm{exp}(10^{0.4(MM_{})})d\mathrm{log}r_edM`$
where $`\mathrm{\Phi }_0=0.0033\text{Mpc}^3`$, $`\alpha =1.04`$, $`\beta =3`$, $`M_{}=22.8`$, $`r_e^{}=7.9\text{kpc}`$, and $`\sigma _\lambda =0.37`$ (converting their sizes and luminosities from $`h_0=0.65`$ to the $`h_0=0.50`$ used here). Implicit in the above bivariate distribution is a distribution in sizes analogous to Eq. (31), a Schechter distribution in luminosity, and a $`SBL^{1/3}`$ correlation between luminosity and surface brightness. Similar scalings are apparent in the McGaugh & de Blok (1997) sample.
Now let us compare a typical $`L_{}`$ galaxy in this model with the observations. Using our stated assumption, a $`L_{}`$ galaxy has a mass of $`1.110^{12}M_{}`$. A typical formation time occurs at $`z=0.3`$. Using Eqs. (13-14), the circular velocity and size of the halo is 132 km/s and 272 kpc (compared to the 140 km/s and 241 kpc predicted assuming a constant $`200\rho _c`$ for the collapse density as in Mo et al. 1998). Then, using Eq. (18), the size of the disk is $`6.0`$ kpc.
This is smaller than the empirical findings of de Jong & Lacey (1998) (7.9 kpc), our own comparisons to local observations (ยง4.4) (6.9 kpc), and Freemanโs Law, which gives 6.9 kpc. Supposing this to be due to a slight cut-off at low values of the dimensionless angular momentum due to disk instabilities (Efstathiou, Lake & Negroponte 1982; Dalcanton et al. 1997; Mo et al. 1998a; van den Bosch 1998), we scale up the size of a typical $`L_{}`$ disk galaxy to 6.9 kpc and reduce the spread in dimensionless angular momenta to $`\sigma _\lambda =0.37`$, as found by de Jong & Lacey (1999) in the analysis of the Mathewson et al. (1992) sample.
Throughout our analysis, we shall take the $`\mathrm{\Omega }=0.3`$, $`\mathrm{\Omega }_\mathrm{\Lambda }=0.7`$ model as our preferred fiducial hierarchical model because of its better correspondence with the evolution in the number of small disks observed up to $`z1`$ (Mao et al. 1998). We evolve the mass-to-light ratio $`\gamma (z)`$ as $`(1+z)^{0.5}`$ to reproduce the observed evolution in the Tully-Fisher relationship (see ยง4.2).
### 2.2 Infall Model (Backwards Approach)
Instead of trying to determine how the global structural properties of disks evolve based on the corresponding properties of their ambient halos, it is also possible to examine a number of local disk galaxies in great detail and to use detailed models of their observed properties (gas profiles, stellar profiles, metallicity profiles, current SFR profiles, age-metallicity relationships) to determine how galaxies might have evolved to high redshift.
There is no consideration of how individual halos might evolve backwards in time in these models, both for simplicity and because of large uncertainties in the local distribution of dark matter. Consequently, while for the forwards approach, the entire evolution of disk properties derives from an evolution of the halo properties, the infall models considered here completely ignore these effects. Conversely, while the forwards approach presented here ignores issues related to the manner in which halo gas is converted into stars, for the infall model, such issues are important.
Naturally, given that one always adopts the observed local universe in this approach as known, this approach does not explain in and of itself why the local disk population is as it is. Indeed, it cannot since there is no link to the initial conditions. In this view, for the hierarchical approach, we adopted the local universe we did because it was a natural prediction of the model, and for the infall approach, we adopted it because it agrees with the observations.
We examine such a model for the evolution of local galaxies based upon the Prantzos & Aubert (1995) model for the star formation rates, metallicites, stars, and gas content for the Milky Way disk. We previously presented this infall model elsewhere (Bouwens et al. 1997; Cayรณn et al. 1996), and we shall revisit it here. This model ignores radial inflows for simplicity and takes the star formation rate to be proportional to both the gas surface density ($`\mathrm{\Sigma }_g`$) and the reciprocal of the radius $`r`$, which for a flat rotation curve is proportional to the epicyclic frequency:
$$\frac{d\mathrm{\Sigma }_{}(r,t)}{dt}=\frac{\mathrm{\Sigma }_g(r,t)}{\tau _g(r)}$$
(33)
where $`\tau _g(r)=[0.3(r/r_{})^1\text{Gyr}]^1`$. Physically, such a star formation rate results if the star formation rate is proportional to the rate at which molecular clouds collide (Wang & Silk 1994) or the periodic compression rate (Wyse & Silk 1989).
For simplicity, the accretion time scale $`\tau _{ff}`$ was taken to be independent of radius since a variation in this time scale is not strongly constrained by the observations (Prantzos & Aubert 1995). The spread in Hubble types was then naturally taken to arise from a spread in this time scale (Cayรณn et al. 1996). The equation for the evolution of the gas density is then
$$\frac{d\mathrm{\Sigma }_g(r)}{dt}=\frac{\mathrm{\Sigma }_g(r,T)+\mathrm{\Sigma }_{}(r,T)}{1e^{T/\tau _{ff}}}\frac{e^{t/\tau _{ff}}}{\tau _{ff}}\frac{\mathrm{\Sigma }_g(r)}{\tau _g}$$
(34)
where $`T`$ is the time from the formation of the disk to the present. Integrating these equations yields the result
$$\mathrm{\Sigma }_g(r,t)=\frac{\mathrm{\Sigma }_g(r,T)+\mathrm{\Sigma }_{}(r,T)}{1e^{T/\tau _{ff}}}\frac{e^{t/\tau _{ff}}e^{t/\tau _g}}{\tau _{ff}\tau _g}\tau _g$$
(35)
Given the fact that this model derived from only one galaxy, it is difficult to know how to extend its evolutionary predictions to galaxies with different luminosities and surface brightnesses. One possible means of extending this model to galaxies beyond the Milky Way involves simply scaling the star formation rates by the differential rotation rate. The change in the star formation rate would then be proportional to $`V_c/R`$ or the typical dynamical time for the disk. Since $`R`$ is roughly proportional to $`V_c`$, there would be no large change in the time scales as a function of disk mass or rotational velocity. Accordingly, Bouwens et al. (1997) simply elected to scale everything in size to reproduce all luminosities while conserving surface brightness, which is precisely what we have done here.
To determine the scaling relations for galaxies using the infall approach, we performed the calculation for each galaxy on a series of 30 different rings varying logarithmically in size, where the smallest is a circle with radius 0.2 kpc and the largest is a ring of radius 60 kpc with width 12 kpc as done in Bouwens et al. (1997) and Roche et al. (1998). We calculate the evolution in rest-frame $`B`$ band magnitudes by evolving each ring separately to keep track of its gas mass, stellar composition, and metallicity, and we output its colors using the Bruzual & Charlot instantaneous-burst metallicity-dependent spectral synthesis tables compiled in Leitherer et al. (1996).
In the rest-frame $`B`$ band, we found the following scaling relationships:
$$E_L(z)=10^{0.4(0.6z)}$$
(36)
$$E_R(z)=10.27z$$
(37)
Clearly, without number evolution, we take $`E_N(z)=1`$.
For the sake of clarity, we note that the present model differs from the one presented in Bouwens et al. (1997) in terms of both the luminosity functions used and the bulge-to-total distribution assumed. Furthermore, in the Bouwens et al. (1997) study, the preferred values of the age $`T`$ and gas-infall time scale $`\tau _{ff}`$ used in the Prantzos & Aubert (1995) study were scaled to reproduce the number counts. No such scaling was attempted in the present model and we simply use the same $`\tau _{ff}`$ for all disk types.
This model is similar in spirit to the size-luminosity evolution model presented by Roche et al. (1998) based on the infall models of Chiappini et al. (1997), which models the infall, star formation, and chemical evolution of both the thin and thick disk components. For the purposes of illustration, we shall compare the Chiappini et al. (1997) infall model to the one just described in ยง4.1, after which we will restrict our consideration to the redshift scalings given by the Prantzos & Aubert prescription. In this model, the star formation time scale is equal to
$$\tau =\{\begin{array}{cc}1\text{Gyr},\hfill & r<2\text{kpc}\hfill \\ (0.875r0.75)\text{Gyr},\hfill & r2\text{kpc}\hfill \end{array}$$
(38)
The star formation commences at $`t_{form}=(16\text{Gyr}0.35\tau )`$. Note that the time scale for star formation here depends on the radius to the first power as in our model, the preferred Prantzos & Aubert (1995) model, and the recent work by Boissier & Prantzos (1999).
## 3 Observations
### 3.1 Low-Redshift Samples
We make comparisons against local samples with information on the size, luminosity, and circular velocity of local galaxies. Though in principle we could have just used the Courteau (1997) sample, we follow Mao et al. (1998) in using a compilation of three different samples for the comparisons which follow to examine the three two-dimensional relationships.
de Jong & van der Kruit Sample:
The de Jong & van der Kruit (1994) sample provides a nice sample for examining the local size/magnitude relationship. It is selected from the Uppsala Catalogue of Galaxies (Nilson 1973, hereinafter UGC), over only $`12.5\%`$ of the sky and uses only relatively face-on ($`b>0.625`$) galaxies (37.5% of all orientations). Following de Jong (1996), we also take it to be diameter-limited in $`R`$ to galaxies larger than $`2^{}`$ at 24.7 $`R`$-band $`\text{mag/arcsec}^2`$. Whereas de Jong & van der Kruit (1994) in their treatment of their sample assume transparent disks, we correct observed magnitudes to an inclination of $`70\mathrm{deg}`$ using the Tully & Fouquรฉ (1985) inclination corrections.
Courteau (1997) Sample:
We use the Courteau (1997) sample to calibrate the local $`z=0`$ $`V_c\text{size}`$ relationship. The Courteau (1997) sample contains 304 Sb-Sc galaxies from the UGC with Zwicky magnitudes $`m_B<14.5`$, $`R`$-band angular diameters larger than $`1^{}`$, and $`B`$-band major axis $`<4^{}`$. We take their $`v_{opt}=V_c(3.2R_d)`$ as the circular velocity and the $`25r\text{mag/arcsec}^2`$ isophote as the radius.
Pierce & Tully (1988) Sample:
We use the Pierce & Tully (1988) sample to calibrate the local $`z=0`$ $`V_c\text{luminosity}`$ relationship. The Pierce & Tully sample was taken from galaxies in the area of the Ursa Major cluster and is complete up to $`B_T<13.3`$. It includes all galaxies which are not elliptical or S0, not more face on than $`30\mathrm{deg}`$, and not possessing confused H I profiles. Note that in this study and in the Vogt et al. (1996,1997) studies to be discussed, the observed absolute $`B`$-band magnitudes were corrected to intrinsic (unextincted) values using the Tully & Fouquรฉ (1985) inclination corrections.
### 3.2 High-Redshift Samples
Simard et al. (1999) Sample:
For the magnitude-radius relationship, we use the data presented by Simard et al. (1999). This data set contains structural information for $`200`$ galaxies to $`I<23.5`$ from 6 different WFPC2 pointings in the Groth strip ($`30\text{arcmin}^2`$). Spectra were obtained for only a fraction of the faint galaxies, but for galaxies with spectra, there was nearly 100% redshift identification. Following Simard et al. (1999), we can quantify the selection effects of this sample. The probability that a galaxy with apparent magnitude $`I_{814}`$ and radius $`r_d`$ would fall in the photometric sample is $`S_{UP}(I_{814},r_d)`$. From Figure 4 of Simard et al. (1999), we have approximated this as
$$\{\begin{array}{cc}1,\hfill & I_{814}+5\mathrm{log}r_d()<21,\hfill \\ 1\frac{2}{3}(I_{814}+5\mathrm{log}r_d()),\hfill & 21<I_{814}+5\mathrm{log}r_d()<22.5,\hfill \\ 0,\hfill & 22.5<I_{814}+5\mathrm{log}r_d().\hfill \end{array}$$
(Actually, this provides a steeper surface brightness cut-off than the selection function given in Simard et al. 1999.) The probability that a galaxy with apparent magnitude $`I_{814}`$ and radius $`r_d`$ would be selected from the photometric sample for spectroscopic follow-up is given by $`S_{PS}(I_{814},r_d)`$, which we have approximated as
$$\{\begin{array}{cc}1,\hfill & I_{814}<19.3,\hfill \\ 1\frac{0.8}{4}(I_{814}19.3),\hfill & 19.3<I_{814}<23.5,\hfill \\ 0,\hfill & 23.5<I_{814}\hfill \end{array}$$
(again by eyeballing Figure 4 of Simard et al. 1999). Putting these two selection effects together, the probability of selecting a galaxy with apparent magnitude $`I_{814}`$ and radius $`r_d`$ is simply the product of these quantities, namely, $`S_{UP}(I_{814},r_d)S_{PS}(I_{814},r_d)`$. Given our ignorance about the inclinations used in the Simard et al. (1999) study, we assume an average inclination of $`60\mathrm{deg}`$ in transforming the absolute magnitudes to an inclination of $`70\mathrm{deg}`$ using the Tully & Fouquรฉ (1985) law, which yields a correction of $`0.27^m`$ here.
Lilly et al. (1998) Sample:
We also use the data from the LDSS2-CFRS sample with HST WFPC2 follow-up to look at the magnitude-radius relationship. From Table 3 of Brinchmann et al. (1998), the effective area of the CFRS portion of this sample is $`0.01377\text{deg}^2`$ ($`49\text{arcmin}^2`$). The survey is magnitude limited to $`17.5<I<22.5`$, where the magnitudes are isophotal to $`28.0I_{AB}\text{mag/arcsec}^2`$. Though the surface brightness limit is quoted as $`24.5I_{AB}\text{mag/arcsec}^2`$, we have used the more conservative surface brightness limit $`23.5I_{AB}\text{mag/arcsec}^2`$. In principle, then, for any galaxy detected the isophotal magnitude should be approximately equal to the total magnitude. Lilly et al. (1998) chose to examine that subset of galaxies from this data-set which were disk-dominated and had disk scale lengths $`>4`$ kpc, a sample we shall henceforth refer to as the Lilly et al. (1998) large disk sample. Given that the central surface brightness is not strongly correlated with inclination angle, Lilly et al. (1998) concludes that disks are consistent with being opaque. We use the Tully & Fouquรฉ (1985) to transform the listed absolute magnitudes to an inclination of $`70\mathrm{deg}`$.
Vogt et al. (1997) Sample:
In contrast to high-redshift samples with both magnitude and radial information, high-redshift samples with circular velocity measurements are considerably smaller and possess less well-defined selection criteria. In fact, the Vogt et al. (1996,1997) sample with 16 galaxies is the largest such published sample. The selection criteria for this sample is still somewhat qualitative and patchy in nature. It considers galaxies with an inclination greater than $`30\mathrm{deg}`$, detectable line emission, undistorted disk morphology, and an extended profile. We assume an $`I_{814}<22.5`$ magnitude limit as used in the Vogt et al. (1997) sample.
## 4 Results
### 4.1 Basic Scalings
Before getting into a detailed comparison of the models with the observations, we begin by illustrating the manner in which the sizes of $`L_{}`$ ($`M1.210^{12}M_{}`$) galaxies typically evolve as a function of redshift for the different models in Figure 1. We present this evolution in terms of the rest-frame $`B`$ half-light radius, the specification of a band and a measure being necessary to the intrinsic band and profile dependence of evolution in the infall model.
The hierarchical models predict more size evolution than expected from both infall models considered here. Of course, the $`\mathrm{\Omega }=1`$ model possesses more size evolution than the $`\mathrm{\Omega }=0.3`$; $`\mathrm{\Omega }_\mathrm{\Lambda }=0.7`$ model, and the $`\mathrm{\Omega }=0.3`$; $`\mathrm{\Omega }_\mathrm{\Lambda }=0.7`$ model more size evolution than the $`\mathrm{\Omega }=0.1`$ model because of the steeper dependence of $`1/H(z)`$ at the typical formation redshifts of disks to $`z1`$. The predictions of both the Prantzos & Aubert (1995) and Chiappini et al. (1997) models are quite similar, encouragingly enough given that Prantzos & Aubert (1995) and Chiappini et al. (1997) models represent completely independent efforts to model the evolution of the Milky Way galaxy.
Taking $`R/V_c`$ as the measure of size for a given mass halo, we plot both the Courteau sample at $`z=0`$ and the higher-redshift data of Vogt et al. (1996,1997) scaled appropriately so that the Courteau data is centered on unevolved scale size of a galaxy at $`z=0`$. At face value, a comparison of the Vogt data with the Courteau data indicates that there has been size evolution from $`z=0`$ to $`z=1`$ as any of the models here would predict (Mao et al. 1998). Nevertheless, the lack of a strong trend in redshift across the Vogt data-set makes one suspicious that there may be selection effects at work or even systematic errors in the measurements of parameters which may produce the observed differences. In any case, strong conclusions must await the compilation of a larger high-redshift dataset, where the selection effects have been more carefully quantified.
We also present model scalings of the number, luminosity, and surface brightness expected for $`L_{}`$ galaxies in Figure 1. By construction, the hierarchical models produce more number evolution than the infall models, which involve no evolution in number. Except for the infall model based upon the Chiappini et al. (1997) prescription, our โInfallโ model produces similar evolution in luminosity as the hierarchical models to $`z1`$. The infall models also produce less surface brightness evolution than the hierarchical models.
### 4.2 Tully-Fisher Relationship
To assess hierarchical and infall models, we compare their predicted Tully-Fisher relationships with both low and high-redshift observations in the left panels of Figures 2-3. For the Monte-Carlo simulations, we use the same selection effects as already specified for the low (Pierce & Tully 1988) and high-redshift (Vogt et al. 1996, 1997) samples. We have added the Pierce & Tully (1992) fit to these plots for comparison. Naturally, for both our hierarchical and infall models, we obtained good agreement with the Pierce & Tully (1988) sample since we used that sample to adjust the mass-to-light ratio (we assumed that a $`M_{b_J}=21`$ galaxy had a mass $`1.210^{12}M_{}`$) and its assumed log-normal scatter (one sigma scatter of 0.3 dex). At higher redshift, we again obtain basic agreement with the Vogt et al. (1996,1997) sample for both our infall and hierarchical models.
### 4.3 Size-$`V_c`$ Relationship
We also compare our model predictions with the low and high-redshift observations for the size-$`V_c`$ relationship in the right panels of Figures 2-3. Again, we apply the selection criteria given in ยง3 to the low (Courteau 1997) and high-redshift (Vogt et al. 1996,1997) model results. At low redshift, we obtain basic agreement with the observations of Courteau (1997) though there seems to be a slight shift in our models toward larger sizes. No adjustment of our models was made to obtain agreement with the Courteau (1997) sample, and so this comparison can be considered a self-consistency check on our $`z=0`$ models.
At high redshift, sizes for both models are consistent, if not a little larger than the sizes in the Vogt et al. (1996,1997) sample. One of the most surprising thing about a comparison of the models with the observations is the significant size evolution observed in the lowest redshift bin ($`z<0.6`$) relative to the models. In fact, as discussed in relation to Figure 1, the low redshift ($`z<0.6`$) points seem to have undergone more size evolution than the high redshift ($`z>0.6`$) points. As the low redshift points are primarily low luminosity galaxies and the high redshift high luminosity galaxies, this could point to some luminosity-dependent evolutionary trend though the numbers are still too small to make any claims toward this end.
### 4.4 Size-Magnitude Relationship
As so often in making low-to-high redshift comparisons, freedom in the choice of the $`z=0`$ no-evolution model can be very important in interpreting the high-redshift results. In particular, due to the fact that each redshift bin in the Simard et al. (1999) sample contains galaxies of a particular luminosity, significant evolution in disk surface brightness would appear to be present, simply as a result of correlations between luminosity and surface brightness at $`z=0`$. There are also important surface brightness selection effects in constructing the Simard et al. (1999) sample. We illustrate the importance of these considerations in Figure 4 by presenting a no-evolution model, hereafter referred to as our fiducial no-evolution model, identical to our $`z=0`$ hierarchical and infall models except there is no-evolution in the disk size, number, surface brightness, or luminosity, i.e., $`E_L(z)=E_N(z)=E_R(z)=1`$. We also present the surface brightness distributions recovered from a similar no-evolution model differing only in its use of the $`M_{b_J}=21`$ surface brightness distribution for all luminosities. We also present the surface brightness distributions recovered both by including a less conservative surface brightness selection
$$S_{UP}(I_{814},r_d())=\{\begin{array}{cc}1,\hfill & I_{814}+5\mathrm{log}r_d()<21,\hfill \\ 1\frac{1}{3}(I_{814}+5\mathrm{log}r_d()),\hfill & 21<I_{814}+5\mathrm{log}r_d()<24,\hfill \\ 0,\hfill & 24<I_{814}+5\mathrm{log}r_d()\hfill \end{array}$$
(39)
(more resembling the one used by Simard et al. (1999)) and without including surface brightness ($`S_{UP}`$) selection at all. Notice the apparent increase in surface brightness for our fiducial model, where $`SBL^{1/3}`$, as a function of z relative to our constant surface brightness model. It is interesting to note that there is an absence of galaxies at surface brightnesses close to the threshold for detection in the Simard et al. (1999) sample.
With these caveats in mind, we compare our model $`B`$-band surface brightness distributions with both the low and high redshift observations in Figure 5. The models seem to be in rough agreement with the surface brightness distribution of the observations. Given that model populations increase in $`B`$-band surface brightness by $`1.5^m`$ to $`z1`$ (see Figure 1), this suggests a similar increase in the surface brightness of disks to $`z1`$. The models themselves show no large differences. As so often, the uncertainties in the $`z=0`$ modeling are large enough to preclude detailed discrimination among models, especially given the limited high redshift data sets.
In Figures 6-7, we plot the observed luminosity-size distributions at low (de Jong & van der Kruit 1994) and high (Simard et al. 1999) redshift and compare them with those obtained for the $`\mathrm{\Omega }=0.3/\mathrm{\Omega }_\mathrm{\Lambda }=0.7`$ hierarchical and infall models. We show a similar comparison of the hierarchical model with the high redshift Lilly et al. (1998) large disk data-set in Figure 8. We also plot the cumulative size and luminosity distributions along the vertical and horizontal axes, respectively, for both the observations (histogram) and the models (lines).
As in Figures 4-5, the models predict too many low surface brightness galaxies relative to the observations. This also results in too many large model galaxies and too many low luminosity model galaxies relative to the observed size and luminosity distributions at intermediate to low luminosities, especially for the Lilly et al. (1998) large disk sample even with our more conservative surface brightness limit. For large, luminous galaxies, there is no obvious change in numbers to high redshift.
In Figures 6-7, there are also relatively large differences in normalization between the observations and models. This is apparently the result of large-scale structure (there are small groups/clusters at $`z0.8`$ and $`z1.0`$ in the Groth Strip). It is therefore difficult to make extremely quantitative statements about the evolution in the number density of galaxies with specific surface brightnesses, sizes, and luminosities across the redshift intervals surveyed.
Nevertheless, the abundance of high surface brightness disk galaxies at high redshifts relative to the model predictions is surely conspicuous and suggests that there has been a significant increase in the number of high surface brightness disk galaxies to $`z1`$ as we have already argued in comparing the model and observed surface brightness distributions. Again, shifting the surface brightness distribution toward these high surface brightnesses (here by the $`1.5^m`$ predicted by the models) is an obvious way of accommodating this increase. Of course, any argument based on the normalization of specific galaxy populations is subject to considerable uncertainties important when such small contiguous areas are being probed.
To provide a visual comparison between our no-evolution and evolutionary models, we include a simulation of a patch of the HDF ($`I_{F814W}`$, $`B_{F450W}`$, and $`V_{F606W}`$) in Figure 9 using our fiducial no-evolution model (panel a), the $`\mathrm{\Omega }=1`$ hierarchical model (panel b), and the $`\mathrm{\Omega }=0.3`$,$`\mathrm{\Omega }_\mathrm{\Lambda }=0.7`$ infall model (panel c) for comparison with the actual HDF North and South (panel d). Clearly, the lack of high surface brightness galaxies is apparent in the no-evolution model relative to the HDF and even somewhat in the apparent in the $`\mathrm{\Omega }=1`$ hierarchical model relative to the HDF. Of course, our simulations do not include ellipticals or peculiars, so the actual HDF will include more bright objects than the simulations.
## 5 Discussion
There is a real question about a lack of low surface brightness galaxies relative to our predictions, especially as compared to the no-evolution model predictions. This conclusion is somewhat dependent on the assumed correlation between surface brightness and luminosity as is evident in Figure 4. This conclusion is also dependent on the selection biases against low surface brightness galaxies not being stronger than those considered here.
There is an extensive literature discussing surface brightness selection biases (Disney 1976; Allen & Shu 1979) and various attempts to derive the bivariate luminosity-surface brightness distribution of galaxies (McGaugh 1996; Dalcanton et al. 1997b; Sprayberry et al. 1997). Surface brightness has a particularly strong effect on isophotal magnitude determinations, especially for low surface brightness galaxies; and this can introduce significant errors in the magnitude determinations, so the effective volume probed for these galaxies is significantly smaller than it is for equivalent luminosity high surface brightness galaxies (McGaugh 1996).
Simard et al. (1999) in a detailed quantification of the selection effects of the DEEP sample do not consider the effect of surface brightness on the magnitudes and sizes recovered since typical errors were found to be $`0.2^m`$ (Simard 1999, private communication). Despite the relatively small size of this error, it is not entirely clear to the present authors that the errors would not become quite significant for the lowest surface brightness galaxies in the sample, particularly those just marginally detectable given the chosen object identification and photometric parameters. Secondly, Simard et al. (1999) considers disk galaxies to be optically thin whereas the observations of Lilly et al. (1998) are more consistent with disks being optically thick. Highly-inclined optically thin disks would be much more detectable than face-on or optically-thick disks. The upshot is that at many apparent magnitudes and radii, Simard et al. (1998) would suppose that at least some highly inclined galaxies would be detectable and therefore the selection function $`S_{UP}`$ there would be non-zero when in reality if disks were optically thick it would be zero. For these reasons, we used a slightly more conservative selection function in surface brightness than that given in Figure 4 of Simard et al. (1999) (see ยง3.2).
Another possibility, not considered here, is that low surface brightness galaxies might form relatively late, meaning that their mass-to-light ratios remain relatively large until relatively recent epochs. Of course, prima facie, this would seem unlikely given the apparently constant slope in the Tully-Fisher relationship to faint magnitudes.
In their own analysis of their sample of $`200`$ galaxies, Simard et al. (1999) concluded that there had been little evolution in the surface brightness distribution of disk galaxies when all selection effects had been carefully considered, quite in contrast to our estimated $`1.5^m`$ of $`B`$-band surface brightness evolution. Little consideration, however, was paid to the evolution in the total numbers of high surface brightness galaxies. Here, we find that the number of high surface brightness galaxies dramatically exceeds that predicted by the evolutionary models considered here, and we have argued that this provides evidence for an evolution in the surface brightness distribution of disk galaxies.
Our interpretation seems to be furthermore supported by the lack of low surface brightness galaxies relative to our models. For no-evolution in the disk surface brightness distribution really to be present as Simard et al. (1999) claims, high-redshift intervals should have similar numbers of low surface brightness galaxies to those found in local samples, and these galaxies seem to be deficient, even with respect to our models which show significant evolution in surface brightness.
While the conclusions of Simard et al. (1999) appear to have been carefully drawn, we would like to suggest that there are significant uncertainties in their determination of the mean surface brightnesses in the lowest redshift intervals and therefore the inferred evolution in surface brightness due to the small size of the low redshift samples considered. By applying the selection effects from the high-redshift bin identically to all redshift intervals, Simard et al. (1999) restricted their analysis to that fraction of disk galaxies exceeding the high-redshift surface brightness detection limit. Applying these selection criteria uniformly to all low redshift intervals severely pares down the low-redshift samples and significantly increases the uncertainty of their average surface brightness measure. Given the observed range in observed surface brightness ($`2\text{mag/arcsec}^2`$) and typical numbers ($`56`$) for the lowest redshift bins, there is a non-negligible uncertainty in the average surface brightness at low redshift, $`0.6^m`$.
Our estimates of $`1.5^m`$ of $`B`$-band surface brightness evolution are somewhat larger than that inferred by most authors. Roche et al. (1998) found $`0.9^m`$ of surface brightness evolution from $`z0.2`$ to $`z0.9`$, Lilly et al. (1998) found $`0.8^m`$ of surface brightness evolution in their large disk sample, and Schade et al. (1995,1996a) inferred $`1.2^m`$ and $`1.5^m`$ respectively to $`z0.8`$. Despite different differential measures of surface brightness evolution, most of these samples give similar values for the mean disk surface brightness near $`z1`$: $`20.79\pm 0.17`$ for the Roche et al. (1998) sample ($`0.65<z`$), $`19.9\pm 0.2`$ for the Simard et al. (1999) sample ($`0.9<z<1.1`$), $`20.7\pm 0.25`$ for the Lilly et al. (1998) large disk sample ($`0.5<z<0.75`$), $`20.2\pm 0.25`$ for the Schade et al. (1995) sample ($`0.5<z`$), and $`19.8\pm 0.1`$ for the Schade et al. (1996) sample ($`0.5<z<1.1`$). Consequently, differences in the surface brightness evolution inferred derive from differences in the $`z=0`$ surface brightness distributions assumed. We assume a distribution consistent with the local data of de Jong & van der Kruit (1994) and Mathewson et al. (1992) as a baseline for measuring evolution with a surface brightness peaking faintward of Freemanโs Law ($`21.7B\text{mag/arcsec}^2`$; Freeman 1970) while Schade et al. (1995,1996) simply makes reference to Freemanโs Law ($`21.65b_j\text{mag/arcsec}^2`$). Simard et al. (1999), Lilly et al. (1998), and Roche et al. (1998) measure surface brightness evolution differentially across their samples. Possible problems here are surface brightness selection effects and limited low-redshift samples.
## 6 Summary
In the present paper, we presented models based on two different approaches for predicting the evolution in disk properties: a hierarchical forwards approach, where the evolution in disk properties follows from corresponding changes in halo properties, and a backwards approach, where the evolution in disk properties follows from an infall model providing a close fit to numerous observables for the Milky Way. We normalized the models to the local $`z=0`$ observations, we made high-redshift predictions for the models, and we compared these predictions with high-redshift observations.
Our findings are as follows:
* The hierarchical and infall models predict relatively similar amounts of evolution in global properties (size, surface brightness, and luminosity) for disk galaxies to $`z1`$. Clearly, discriminating between the models will require a careful look at evolution in number (and therefore surveys over a much larger area) and/or measurements of certain internal properties, like color, star formation, or metallicity gradients of high redshift disks.
* There is an apparent lack of low surface brightness galaxies in the high-redshift observations of Simard et al. (1999) and Lilly et al. (1998) as compared to model predictions based on local observations (Mathewson et al. 1992; de Jong & van der Kruit 1994).
* Our model surface brightness distributions produce relatively good agreement with the observations, suggesting that the $`B`$-band surface brightness has evolved by $`1.5^m`$ from $`z=0`$ to $`z1`$ similar to that found in the models. This finding is supported by the fact that there is a significantly larger number of high surface brightness galaxies than in our model predictions, suggesting that there has been a significant evolution in number, most easily accommodated by shifting the mean surface brightness of the disk population to higher surface brightnesses. This is contrary to the conclusion reached by Simard et al. (1999) based on the same data.
Here the hierarchical and infall models were presented as competing models to describe the evolution in the properties of disk galaxies. If the hierarchical structural paradigm is roughly correct as is generally supposed, the infall model simply provides a modification of the basic hierarchical scalings to account for the fact that the gas infall rate or star formation efficiency is not the same at all radii. In this sort of scenario, if there is an appreciable formation of structure at low redshift, a consideration of hierarchical scalings is probably more appropriate and if there is not, a consideration of scalings following from infall models is probably more appropriate. Obviously, at the present time, a complete incorporation of radially dependent gas infall and star formation scenarios into a hierarchical paradigm is not merited given the lack of high-redshift data needed to constrain such hybrid models.
We acknowledge helpful discussions with David Schade and Luc Simard. We are especially grateful to Laura Cayรณn for helping compile some of the samples used here, for some useful discussions, and for providing a critical read of this document. We thank Stephane Courteau and Nicole Vogt for sending us their data in electronic form. This research has been supported in part by grants from NASA and the NSF. RJB would like to thank the Oxford astrophysics department for its hospitality while this work was being carried out. |
warning/0002/math0002054.html | ar5iv | text | # F-regular and F-pure rings vs. log terminal and log canonical singularities
## Introduction
The notions of F-regular and F-pure rings are defined by Hochster and Huneke \[HH1\] and Hochster and Roberts \[HR\], respectively, by using the Frobenius map in characteristic $`p>0`$. These notions have some similarity to the notion of rational singularities defined for singularities of characteristic zero. But if we look more closely, we find that F-regularity is strictly stronger than rational singularity and that there is no implication between F-purity and rational singularity. (There are several variants of the concept of โF-regularโ rings which are expected and in some cases proved to be equivalent to each other. In this paper, we always concern โstrongly F-regularโ rings as in Definition 1.1 (2).)
In dimension two, F-regular rings and F-pure rings are well investigated and we find very strong connection between F-regular rings and โquotient singularitiesโ \[W2\], \[Ha2\]. To generalize this result to higher dimension, we notice that there is a notion called log terminal singularities for rings of characteristic zero, which is equivalent to quotient singlarities in dimension two. Also, the notion of log canonical singularities is defined similarly as log terminal singularities. Looking more closely, we find that F-regular (resp. F-pure) rings and log terminal (resp. log-canonical) singularities have very similar properties.
The notions of F-regular and F-pure rings are also defined for rings of โcharacteristic zero.โ Namely, we say that a ring essentially of finite type over a field of characteristic zero is of F-regular (resp. F-pure) type if its reduction to characteristic $`p`$ is F-regular (resp. F-pure) for infinitely many $`p`$.
In this paper, we will show that a $``$-Gorenstein ring of F-regular (resp. F-pure) type has log terminal (resp. log canonical) singularities. Actually, we have the following result in characteristic $`p>0`$, as a special case of our main theorem (Theorem 3.3).
###### Theorem
Let $`(A,๐ช)`$ be a $``$-Gorenstein normal local ring of characteristic $`p>0`$ and $`f:XY=\mathrm{Spec}A`$ be a proper birational morphism with $`X`$ normal. Let $`E=_{i=1}^sE_i`$ be the exceptional divisor of $`f`$ and let
$$K_X=f^{}K_Y+\underset{i=1}{\overset{s}{}}a_iE_i$$
as in $`(\mathrm{3.1.1})`$. If $`A`$ is F-pure $`(`$resp. strongly F-regular$`)`$, then $`a_i1`$ $`(`$resp. $`a_i>1)`$ for every $`i`$.
If $`A`$ is a normal local ring essentially of finite type over a field of characteristic zero and if $`f:XY=\mathrm{Spec}A`$ is a resolution of singularities of $`A`$, then we can apply the above theorem to โreduction modulo $`p`$โ of this map. Also, this result is applicable to strongly F-regular (resp. normal F-pure) rings of fixed characteristic $`p`$. It is crucially used in \[Ha2\] to classify two-dimensional F-regular and normal F-pure rings in every characteristic $`p`$.
The technique we use in our proof is that of Frobenius splitting used by \[MR\]. If a ring has some splitting of Frobenius, we can lift it to an open set of its resolution. Then the splitting affects โdiscrepancyโ of the exceptional divisors.
The converse of the โF-regularโ part of the above theorem has been established as well \[Ha3\], and we have a Frobenius characterization of log terminal singularities. Namely, a ring in characteristic zero has log terminal singularities if and only if it is of F-regular type and $``$-Gorenstein. See also \[MS2\], \[S1, 3\].
The other new ingredient of this paper is an attempt to generalize the notions of F-regular and F-pure rings to those for โpairs.โ By a pair, we mean a pair $`(A,\mathrm{\Delta })`$ of a normal ring $`A`$ and a $``$-divisor $`\mathrm{\Delta }`$ on $`\mathrm{Spec}A`$. The notions of log terminal and log canonical singularities are defined not only for normal rings (in characteristic zero) but also for pairs (cf. \[KMM\]), and these โsingularities of pairsโ play a very important role in birational algebraic geometry. (To wit, see Kollรกrโs lecture note \[Ko\].) Therefore we are tempted to define corresponding โF-singularities of pairs.โ
In this paper, we define a few variants of โF-singularities of pairs,โ namely, strong F-regularity, divisorial F-regularity and F-purity of pairs, which are expected to correspond singularities of pairs called Kawamata log terminal (KLT), purely log terminal (PLT) and log canonical (LC).
The significance of our main theorem (Theorem 3.3) is enhanced by considering F-singularities of pairs. We also prove some results on F-singularities of pairs, which are analoguous to the results proved for singularities of pairs in characteristic zero (cf. \[Ko\]). If we keep in mind the (expected) correspondence of F-singularities and zero characteristic singularities of pairs, we find that Theorem 4.8 is parallel to the behavior of singularities of pairs under finite covering. Also, Theorem 4.9 is an analog of the so-called โinversion of adjunction,โ and Proposition 2.10 corresponds to some fundamental properties of โlog canonical thresholds.โ
Our main theorem for the no boundary case ($`\mathrm{\Delta }=0`$) was first proved by the second-named author several years ago \[W3\], and this result remained unpublished. Then the first-named author proposed a generalization to F-singularities of pairs, and the present paper was worked out.
In this paper, all rings will be commutative and Noetherian. Also, since we are interested only in normal local rings, we assume all rings are integral domains.
## 1. Preliminaries
Let $`A`$ be an integral domain of characteristic $`p>0`$ and let $`F:AA`$ be the Frobenius endomorphism given by $`F(a)=a^p`$. We always use the letter $`q`$ for a power $`q=p^e`$ of $`p`$. Since $`A`$ is assumed to be reduced, we can identify the following three maps:
$$F^e:AA,A^qA,AA^{1/q}.$$
The ring $`A`$ is called F-finite, if $`F:AA`$ (or $`AA^{1/p}`$) is a finite map. For example, $`A`$ is F-finite if it is essentially of finite type over a perfect field or it is a complete local ring with perfect residue field. We always assume $`A`$ is F-finite throughout this paper.
###### Demonstration Definition 1.1
(1) (Hochster and Roberts \[HR\]) $`A`$ is F-pure if the map $`AA^{1/p}`$ (hence $`AA^{1/q}`$ for every $`q=p^e`$, $`e1`$) splits as an $`A`$-module homomorphism.
(2) (Hochster and Huneke \[HH2\]) $`A`$ is strongly F-regular if for every $`cA`$ which is not in any minimal prime of $`A`$ (or equivalently, $`c0`$ in our case), there exists $`q=p^e`$ such that the $`A`$-module homomorphism $`AA^{1/q}`$ sending 1 to $`c^{1/q}`$ splits.
F-pure and (strongly) F-regular rings have many nice properties. But here, we will only list the following. Interested readers could refer \[HH1\], \[HH2\], \[HR\], \[FW\] or \[W2\] (in the last reference, โF-regularโ should read โstrongly F-regularโ).
###### Remark Remark 1.2
(1) Let $`(A,๐ช)`$ be a local ring and $`E=E_A(A/๐ช)`$ be the injective envelope of the $`A`$-module $`A/๐ช`$. Then the splitting of the $`A`$-homomorphism $`AA^{1/q}`$ sending 1 to $`c^{1/q}`$ is equivalent to the injectivity of the map $`EE_AA^{1/q}`$ sending $`\xi E`$ to $`\xi c^{1/q}E_AA^{1/q}`$ (cf. Proposition 2.4).
(2) There are notions of F-regular and weakly F-regular rings also defined by Hochster and Huneke \[HH1\], which we do not define here. On the other hand, a local ring $`(A,๐ช)`$ of dimension $`d`$ is called F-rational \[FW\], if $`A`$ is CohenโMacaulay and if for every $`c0A`$, there exists $`q=p^e`$ such that the map $`cF^e:H_๐ช^d(A)H_๐ช^d(A)`$, or equivalently, the mapping
$$H_๐ช^d(A)=H_๐ช^d(A)_AAH_๐ช^d(A)_AA^{1/q}H_๐ช^d(A^{1/q}),$$
sending $`\xi H_๐ช^d(A)`$ to $`\xi c^{1/q}=(cF^e(\xi ))^{1/q}`$, is injective.
Then the following implications hold:
$$\text{strongly F-regular }\text{F-regular }\text{weakly F-regular }\text{F-rational }\text{normal.}$$
Also, for a Gorenstein local ring, strongly F-regular is equivalent to F-rational.
(3) (\[Mc\], \[Wi\]) For $``$-Gorenstein local rings, strongly F-regular and weakly F-regular are equivalent. (A normal local ring is $``$-Gorenstein if its canonical module has a finite order in the divisor class group.) In fact, this remains true if the ring has an isolated non-$``$-Gorenstein point.
(4) (\[HR\], \[HH1\], \[HH2\]) If a ring is (strongly or weakly) F-regular (resp. F-pure), so are its pure subrings. (A ring homomorphism $`AB`$ is pure if for every $`A`$-module $`M`$, the map $`MM_AB`$ is injective).
(5) (\[MS1\], \[W2\]) Conversely, if $`(A,๐ช)(B,๐ซ)`$ is a finite local homomorphism of normal local rings which is รฉtale in codimension 1 and if $`A`$ is strongly F-regular (resp. F-pure), so is $`B`$.
###### Demonstration Notation 1.3
Let $`A`$ be a normal domain with quotient field $`L`$. A $``$-Weil divisor on $`Y=\mathrm{Spec}A`$ is a linear combination $`D=_{i=1}^r\alpha _iD_i`$ of irreducible reduced subschemes $`D_iY`$ of codimension 1 with coefficients $`\alpha _i`$. The round-down and round-up of $`D`$ is defined by $`D=_{i=1}^r\alpha _iD_i`$ and $`D=_{i=1}^r\alpha _iD_i`$, respectively. We also denote
$$A(D)=\{fL|\text{div}_Y(f)+D0\}.$$
Clearly $`A(D)=A(D)`$, and this is a divisorial (i.e., finitely generated reflexive) submodule of $`L`$. Conversely, any divisorial submodule $`I`$ of $`L`$ is written as $`I=A(D)`$ for some unique integer coefficient Weil divisor $`D`$. We denote by $`I^{(m)}`$ the reflexive hull of $`I^m`$. If $`I=A(D)=A(D)`$, then $`I^{(m)}=A(mD)`$.
For any ideal $`IA`$ and $`q=p^e`$, we denote by $`I^{[q]}`$ the ideal of $`A`$ generated by the $`q`$th powers of elements of $`I`$. Also, the notation $`()^{1/q}`$ will show that the module under consideration is an $`A^{1/q}`$-submodule of $`L^{1/q}`$. For example, $`(I^{[q]})^{1/q}=IA^{1/q}`$.
###### Demonstration 1.4. Frobenius map and its splitting
Let $`A`$ be a normal domain of characteristic $`p>0`$ and let $`D`$ be an effective Weil divisor on $`Y=\mathrm{Spec}A`$. Then for any $`q=p^e`$, we have a natural inclusion map $`ฤฑ:AA(D)^{1/q}`$, which is identified with the $`e`$-times Frobenius $`F^e:AA`$ followed by the inclusion map $`AA(D)`$. If $`DqD_0`$ for some effective divisor $`D_0`$, then $`ฤฑ`$ is factorized as
$$AA(D_0)=A(D_0)_AAA(D_0)_AA^{1/q}A(qD_0)^{1/q}A(D)^{1/q}.$$
Hence, if $`ฤฑ`$ splits as an $`A`$-module homomorphism, then $`AA(D_0)`$ also splits, and this implies $`D_0=0`$. This kind of argument is used frequently in the sequel. Note also that a splitting is preserved under localization, in particular, a localization of a strongly F-regular ring is again strongly F-regular.
## 2. F-regularity and F-purity of Pairs
###### Demonstration Definition 2.1
Let $`A`$ be an F-finite normal domain of characteristic $`p>0`$ and $`\mathrm{\Delta }`$ be an effective $``$-Weil divisor on $`Y=\mathrm{Spec}A`$.
(1) We say that the pair $`(A,\mathrm{\Delta })`$ is F-pure if the inclusion map $`AA((q1)\mathrm{\Delta })^{1/q}`$ splits as an $`A`$-module homomorphism for every $`q=p^e`$.
(2) $`(A,\mathrm{\Delta })`$ is strongly F-regular if for every nonzero element $`cA`$, there exists $`q=p^e`$ such that $`c^{1/q}AA((q1)\mathrm{\Delta })^{1/q}`$ splits as an $`A`$-module homomorphism.
(3) $`(A,\mathrm{\Delta })`$ is divisorially F-regular if for every nonzero element $`cA`$ which is not in any minimal prime ideal of $`A(\mathrm{\Delta })A`$, there exists $`q=p^e`$ such that $`c^{1/q}AA((q1)\mathrm{\Delta })^{1/q}`$ splits as an $`A`$-module homomorphism.
Definition 2.1 includes Definition 1.1 as the special case $`\mathrm{\Delta }=0`$. Namely, $`A`$ is F-pure if and only if $`(A,0)`$ is F-pure, and $`A`$ is strongly F-regular if and only if $`(A,0)`$ is strongly F-regular, or equivalently, $`(A,0)`$ is divisorially F-regular.
We collect some basic properties in the following.
###### Proposition 2.2
Let $`A`$ be an F-finite normal domain of characteristic $`p>0`$ and $`\mathrm{\Delta }`$ be an effective $``$-Weil divisor on $`Y=\mathrm{Spec}A`$.
###### Demonstration Proof
(1) It is obvious that strongly F-regular implies divisorially F-regular. Also, divisorially F-regular is equivalent to F-pure if $`dimA=1`$. If $`dimA2`$, we can choose a non-unit $`c0`$ of $`A`$ which is not in any minimal prime of $`A(\mathrm{\Delta })`$. Then the map $`(c^q)^{1/q^{}}AA((q^{}1)\mathrm{\Delta })^{1/q^{}}`$ does not split for every $`q^{}q`$, because it factors through $`cAA`$ (cf. 1.4). Hence, if $`(A,\mathrm{\Delta })`$ is divisorially F-regular and if $`q`$ is any power of $`p`$, then there is a power $`q^{}>q`$ such that $`(c^q)^{1/q^{}}AA((q^{}1)\mathrm{\Delta })^{1/q^{}}`$ splits. This implies that $`AA((q^{}1)\mathrm{\Delta })^{1/q^{}}`$ splits, and so does $`AA((q1)\mathrm{\Delta })^{1/q}`$. Consequently, $`(A,\mathrm{\Delta })`$ is F-pure.
(2) The sufficiency is clear. To show the necessity, let $`cA`$ be any nonzero element, and choose $`d0A(\mathrm{\Delta })`$ so that $`dA(q\mathrm{\Delta })A((q1)\mathrm{\Delta })`$ for every $`q=p^e`$. If $`(A,\mathrm{\Delta })`$ is strongly F-regular, then there exists a power $`q^{}`$ of $`p`$ such that $`A\mathrm{@}>(cd)^{1/q^{}}>>A((q^{}1)\mathrm{\Delta })^{1/q^{}}`$ splits. Since this map is factorized into $`A\stackrel{c^{1/q^{}}}{}A(q^{}\mathrm{\Delta })^{1/q^{}}\stackrel{d^{1/q^{}}}{}A((q^{}1)\mathrm{\Delta })^{1/q^{}}`$, the map $`A\stackrel{c^{1/q^{}}}{}A(q^{}\mathrm{\Delta })^{1/q^{}}`$ also splits. On the other hand, since $`A`$ is F-pure, the map $`A^{1/q^{}}A^{1/qq^{}}`$ splits for all $`q=p^e`$, and so does $`A(q^{}\mathrm{\Delta })^{1/q^{}}A(qq^{}\mathrm{\Delta })^{1/qq^{}}`$, too. Hence the map $`A\stackrel{c^{1/q^{}}}{}A(qq^{}\mathrm{\Delta })^{1/qq^{}}`$ splits for all $`q=p^e`$. Since this map is factorized into $`A\mathrm{@}>c^{1/qq^{}}>>A(qq^{}\mathrm{\Delta })^{1/qq^{}}\mathrm{@}>c^{(q1)/qq^{}}>>A(qq^{}\mathrm{\Delta })^{1/qq^{}}`$, the map $`A\mathrm{@}>c^{1/q}>>A(q\mathrm{\Delta })^{1/q}`$ splits for all $`q=p^eq^{}`$.
To prove (3) (resp. (4)), assume to the contrary that $`\mathrm{\Delta }`$ has a component $`\mathrm{\Delta }_0`$ with coefficient $`1`$ (resp. $`>1`$). Then there is a $`q=p^e`$ such that the coefficient of $`q\mathrm{\Delta }`$ (resp. $`(q1)\mathrm{\Delta }`$) in $`\mathrm{\Delta }_0`$ is at least $`q`$. Then the map $`AA(q\mathrm{\Delta })^{1/q}`$ (resp. $`AA((q1)\mathrm{\Delta })^{1/q}`$) factors through $`AA(\mathrm{\Delta }_0)A(q\mathrm{\Delta }_0)^{1/q}`$, which does not split. This implies that $`(A,\mathrm{\Delta })`$ cannot be strongly F-regular (resp. F-pure).
(5) The only being non-trivial is the assertion for divisorial F-regularity in the case $`\text{Supp}(\mathrm{\Delta })\text{Supp}(\mathrm{\Delta }^{})\mathrm{}`$. To prove this, we may assume without loss of generality that there is a unique irreducible component $`\mathrm{\Delta }_0`$ of $`\mathrm{\Delta }`$ such that $`\mathrm{\Delta }_0\text{Supp}(\mathrm{\Delta }^{})`$. Let $`c0A`$ be any element which is in $`๐ญ=A(\mathrm{\Delta }_0)`$ but not in any minimal prime of $`A(\mathrm{\Delta }^{})`$, and let $`\nu =v_๐ญ(c)`$, the value of $`c`$ at $`๐ญ`$. Then by prime avoidance, we can choose $`d0cA(\nu \mathrm{\Delta }_0)`$ which is not in any minimal prime of $`A(\mathrm{\Delta })`$. This implies that $`dA((q1)\mathrm{\Delta }^{})cA((q1)\mathrm{\Delta })`$ for $`q0`$. Now, if $`(A,\mathrm{\Delta })`$ is divisorially F-regular, then there exists $`q=p^e`$ such that $`A\stackrel{d^{1/q}}{}A((q1)\mathrm{\Delta })^{1/q}`$ splits, and we can show as in (1) that this is true for infinitely many $`q`$. Since this map factors into $`A\stackrel{c^{1/q}}{}A((q1)\mathrm{\Delta }^{})^{1/q}\stackrel{(d/c)^{1/q}}{}A((q1)\mathrm{\Delta })^{1/q}`$ if $`q0`$, the map $`A\stackrel{c^{1/q}}{}A((q1)\mathrm{\Delta }^{})^{1/q}`$ splits for some $`q`$. $`\mathrm{}`$
###### Remark Remark 2.3
In the definition of โF-pure,โ it seems apparently natural to refer the map $`AA(q\mathrm{\Delta })^{1/q}`$ instead of $`AA((q1)\mathrm{\Delta })^{1/q}`$. (We have seen in Proposition 2.2 (2) that this makes no difference for strong F-regularity.) But it does make crucial difference for F-purity. Let us say $`(A,\mathrm{\Delta })`$ is strongly F-pure if $`AA(q\mathrm{\Delta })^{1/q}`$ splits for every $`q=p^e`$. Then the proof of Proposition 2.2 (3) shows that if $`(A,\mathrm{\Delta })`$ is strongly F-pure, then $`\mathrm{\Delta }=0`$. But this is too much stronger than what we want for โF-purity.โ Note also that we have the implication โstrongly F-regular $``$ strongly F-pure $``$ F-pure,โ and that there is no implication between strong F-purity and divisorial F-regularity.
###### Proposition 2.4
Let $`(A,๐ช)`$ be a $`d`$-dimensional F-finite normal local ring of characteristic $`p>0`$ and let $`\mathrm{\Delta }`$ be an effective $``$-Weil divisor on $`Y=\mathrm{Spec}A`$. Then:
###### Demonstration Proof
The map $`c^{1/q}AA((q1)\mathrm{\Delta })^{1/q}`$ splits as an $`A`$-module homomorphism if and only if the map
$$\mathrm{Hom}_A(A((q1)\mathrm{\Delta })^{1/q},A)\stackrel{c^{1/q}}{}\mathrm{Hom}_A(A((q1)\mathrm{\Delta })^{1/q},A)\mathrm{Hom}_A(A,A)=A$$
is surjective. By the local duality, the Matlis dual of $`\mathrm{Hom}_A(A((q1)\mathrm{\Delta })^{1/q},A)\mathrm{Hom}_A(A(qK_A+(q1)\mathrm{\Delta })^{1/q},K_A)`$ is $`H_๐ช^d(A(qK_A+(q1)\mathrm{\Delta })^{1/q})`$, so that the surjectivity of the above map is equivalent to the injectivity of
$$H_๐ช^d(K_A)H_๐ช^d(A(qK_A+(q1)\mathrm{\Delta })^{1/q})\stackrel{c^{1/q}}{}H_๐ช^d(A(qK_A+(q1)\mathrm{\Delta })^{1/q}).$$
Since this map is identified with $`cF^e:H_๐ช^d(K_A)H_๐ช^d(A(qK_A+(q1)\mathrm{\Delta }))`$, we obtain the assertions (1), (2) and (3) at once. $`\mathrm{}`$
###### Corollary 2.5
If $`A`$ is a normal toric ring over a perfect field $`k`$ of characteristic $`p>0`$ and if $`\mathrm{\Delta }`$ is a reduced toric divisor, then the pair $`(A,\mathrm{\Delta })`$ is F-pure.
###### Demonstration Proof
Let $`M=^d`$, $`N=\mathrm{Hom}_{}(M,)`$ and denote the duality pairing of $`M_{}=M_{}`$ with $`N_{}=N_{}`$ by $`,:M_{}\times N_{}`$. Let $`A`$ be the toric ring defined by a rational polyhedral cone $`\sigma N_{}`$ and let $`D_1,\mathrm{},D_s`$ be the toric divisors of $`\mathrm{Spec}A`$ corresponding to the primitive generators $`n_1,\mathrm{},n_sN`$ of $`\sigma `$, respectively. To show the F-purity of the pair $`(A,\mathrm{\Delta })`$, we may assume that $`\mathrm{\Delta }=_{i=1}^sD_i`$, by Proposition 2.2 (5). Then one has $`K_A=\mathrm{\Delta }`$, so that $`A(qK_A+(q1)\mathrm{\Delta })=K_A`$. Therefore
$$H_๐ช^d(K_A)=H_๐ช^d(A(qK_A+(q1)\mathrm{\Delta }))=\underset{m\stackrel{ห}{\sigma }M}{}kx^m,$$
where $`\stackrel{ห}{\sigma }=\{mM_{}|m,n_i0\text{ for }i=1,\mathrm{},s\}`$, and the $`e`$-times Frobenius map $`F^e:H_๐ช^d(K_A)H_๐ช^d(A(qK_A+(q1)\mathrm{\Delta }))`$ is given by $`F^e(x^m)=x^{qm}`$ $`(m\stackrel{ห}{\sigma }M)`$. Since this map is injective, $`(A,\mathrm{\Delta })`$ is F-pure by Proposition 2.4 (1). $`\mathrm{}`$
###### Proposition 2.6
(Fedder-type criteria \[F\]) Let $`(R,๐ช)`$ be an F-finite regular local ring of characteristic $`p>0`$ and $`IR`$ be an ideal such that $`A=R/I`$ is normal. Let $`gRI`$ be an element whose image $`\overline{g}A`$ defines a reduced divisor $`\mathrm{div}_Y(\overline{g})`$ on $`Y=\mathrm{Spec}A`$, and consider a $``$-divisor $`\mathrm{\Delta }=t\mathrm{div}_Y(\overline{g})`$ for nonnegative $`\alpha `$. We put $`r_e=t(q1)`$ for each $`q=p^e`$. Then:
###### Demonstration Proof
The proof is essentially the same as those in \[F\], \[Gl\]. First, we may assume without loss of generality that $`(R,๐ช)`$ is a complete regular local ring. We will describe the map $`cF^e:H_๐ช^d(K_A)H_๐ช^d(A(qK_A+(q1)\mathrm{\Delta }))`$ in Proposition 2.4. Since $`A((q1)\mathrm{\Delta })=\overline{g}^{r_e}A`$ and since $`H_๐ช^d(K_A^{(q)})E_A(A/๐ชA)_AA^{1/q}`$ by \[W2, Theorem 2.5\], This map is viewed as
$$c\overline{g}^{r_e}F^e:E_A(A/๐ชA)E_A(A/๐ชA)_AA^{1/q}.$$
$`\mathrm{2.6.1}`$
Let $`E_R=E_R(R/๐ช)`$, $`E_A=E_A(A/๐ชA)`$ and let $`n=dimR`$. Then $`E_RH_๐ช^n(R)`$ since $`R`$ is Gorenstein, and we can identify $`E_R`$ with $`E_R_RR^{1/q}H_๐ช^n(R^{1/q})`$ via the identification of $`R`$ with $`R^{1/q}`$. Also, $`E_A\mathrm{Hom}_R(R/I,E_R)(0:I)_{E_R}E_R`$, and $`E_A`$ and $`E_R`$ have the 1-dimensional socle in common. Let $`z`$ be a generator of the socle of $`E_R`$ and let $`z^{}`$ be the corresponding socle generator in $`E_A`$.
Now, $`RR^{1/q}`$ is flat since $`R`$ is regular, so that via the identification $`R^{1/q}R`$, we have
$`E_A_RR^{1/q}`$ $`\mathrm{Hom}_{R^{1/q}}(R/I_RR^{1/q},E_R_RR^{1/q})`$ $`\mathrm{2.6.2}`$
$`\mathrm{Hom}_R(R/I^{[q]},E_R)(0:I^{[q]})_{E_R}`$
in $`E_R_RR^{1/q}E_R`$. Accordingly we have
$$E_A_AA^{1/q}E_A_RR^{1/q}_{R^{1/q}}A^{1/q}(0:I^{[q]})_{E_R}_RA\frac{(0:I^{[q]})_{E_R}}{I(0:I^{[q]})_{E_R}},$$
and the image of $`z^{}E_A`$ by the map (2.6.1) is $`cg^{r_e}z^q`$ mod $`I(0:I^{[q]})_{E_R}`$, where $`z^q`$ denotes the image of $`z`$ by the $`e`$-times Frobenius on $`E_RH_๐ช^n(R)`$. We also know from (2.6.2) (by putting $`I=๐ช`$) that $`z^qE_R`$ generates $`(0:๐ช^{[q]})_{E_R}`$.
Now the map (2.6.1) is injective if and only if $`cg^{r_e}z^qI(0:I^{[q]})_{E_R}`$ if and only if $`cg^{r_e}(0:๐ช^{[q]})_{E_R}I(0:I^{[q]})_{E_R}`$. Since $`I(0:I^{[q]})_{E_R}=(0:(I^{[q]}:I)_R)_{E_R}`$ by the Matlis duality, this is equivalent to saying that $`cg^{r_e}(I^{[q]}:I)๐ช^{[q]}`$. $`\mathrm{}`$
###### Corollary 2.7
(cf. \[Ha1\]) Let $`(R,๐ช)`$ be an F-finite regular local ring of characteristic $`p>0`$, and let $`f_1,\mathrm{},f_s,g๐ช`$ be an $`R`$-regular sequence. Let $`A=R/(f_1,\mathrm{},f_s)`$ and consider an effective $``$-divisor $`\mathrm{\Delta }=t\mathrm{div}_Y(\overline{g})`$ on $`Y=\mathrm{Spec}A`$. Put $`r_e=t(p^e1)`$. Then:
###### Demonstration Proof
The proof is easy, but we remark one point which might be overlooked. Put $`D=\mathrm{div}_Y(\overline{g})`$. Then the condition $`(f_1\mathrm{}f_s)^{q1}g^{r_e}๐ช^{[q]}`$ holds if and only if $`AA(t(q1)D)^{1/q}`$ splits as an $`A`$-module homomorphism by the preceding argument. So, to prove (1), we have to show the equivalence of the following two conditions.
The implication (i) $``$ (ii) is clear since $`(q1)tD(q1)tD`$. But (ii) $``$ (i) is not apparently clear if $`D`$ is non-reduced. To show this implication, let $`q`$ be any power of $`p`$, and choose $`q^{}`$ such that $`q^{}(q1)tDt(qq^{}1)D`$. Then by (ii), the map $`AA(t(qq^{}1)D)^{1/qq^{}}`$ splits, and this map is factorized into $`AA((q1)tD)^{1/q}A(q^{}(q1)tD)^{1/qq^{}}A(t(qq^{}1)D)^{1/qq^{}}`$. Hence $`AA((q1)tD)^{1/q}`$ splits. $`\mathrm{}`$
###### Remark Remark 2.8
We have Fedder-type criteria also for strong F-purity (see Remark 2.3 for a definition), by setting $`r_e=tp^e`$ in (1) of Proposition 2.6 and Corollary 2.7. The criteria for strong F-regularity also works if we put $`r_e=tp^e`$.
###### Demonstration Example 2.9
(1) Let $`A`$ be a regular local ring and let $`\mathrm{\Delta }=t\mathrm{div}(x_1\mathrm{}x_i)`$ on $`\mathrm{Spec}A`$, where $`x_1,\mathrm{},x_i`$ are part of regular parameters of $`A`$. If $`t1`$ (resp. $`t<1`$), then $`(A,\mathrm{\Delta })`$ is F-pure (resp. strongly F-regular).
(2) Let $`A=k[[X,Y,Z]]/(XYZ^2)`$ and denote the images of $`Z`$ in $`A`$ by $`z`$. Let $`\mathrm{\Delta }=t\mathrm{div}(z)`$. Corollary 2.7 tells us that if $`t1`$ (resp. $`t<1`$), then $`(A,\mathrm{\Delta })`$ is F-pure (resp. strongly F-regular).
(3) Let $`A=k[[x,y]]`$ and let $`\mathrm{\Delta }=\frac{5}{6}\mathrm{div}(x^2y^3)`$. Then in any characteristic $`p>0`$, $`(A,\mathrm{\Delta })`$ is F-pure but not strongly F-regular. But $`(A,\mathrm{\Delta })`$ is strongly F-pure if and only if $`p1`$ mod $`3`$.
The following is a variant of Fedderโs result \[F\], see also \[Ko, Lemma 8.10\].
###### Proposition 2.10
Let $`A=k[[x_1,\mathrm{},x_d]]`$ be a $`d`$-dimensional complete regular local ring over a perfect field $`k`$ of characteristic $`p>0`$ and let $`fA`$ be a nonzero element of multiplicity $`n`$, i.e., $`f๐ช^n๐ช^{n+1}`$, where $`๐ช`$ is the maximal ideal of $`A`$. For nonnegative $`t`$ we have:
###### Demonstration Remark
Note that if $`n=\mathrm{deg}f_n`$ is not divisible by $`p`$, then $`f_nJ`$, so that $`J`$ is equal to the Jacobian ideal $`(f_n,J)`$ of $`V=(f_n=0)^{d1}`$. So the smoothness of $`V`$ implies that $`(x_1^{m_1},\mathrm{},x_d^{m_d})J`$ for some $`m_1,\mathrm{},m_d`$.
###### Demonstration Proof
The assertions (1) and (2) immediately follows from Corollary 2.7. To prove (3) we use the Fedder-type critrion for strong F-purity (2.8):
$$(A,t\mathrm{div}(f))\text{ is strongly F-pure }f^{r_e}๐ช^{[p^e]}\text{ for all }e,$$
where $`r_e=tp^e`$. Since $`f_n^{r_e}๐ช^{[p^e]}`$ implies $`f^{r_e}๐ช^{[p^e]}`$, we may assume that $`f`$ is a homogeneous polynomial of degree $`n`$ in $`x_1,\mathrm{},x_d`$.
Assume $`p\mu :=(m_1+\mathrm{}+m_d)/(dnt)`$ and let $`j_e`$ be the integer such that $`f^{j_e+1}๐ช^{[p^e]}`$ but $`f^{j_e}๐ช^{[p^e]}`$. Then it is sufficient to prove that $`r_ej_e`$ for every $`e0`$. Assume to the contrary that there exists an $`e`$ such that $`r_e>j_e`$, and choose the smallest one among all such $`e`$. Then $`e>0`$, since $`r_0=0`$ by the assumption $`t<1`$. Also, by the minimality of $`e`$, we have $`f^{r_{e1}}๐ช^{[p^{e1}]}`$, and this implies that $`f^{pr_{e1}}๐ช^{[p^e]}`$. Indeed, if $`f^{pr_{e1}}๐ช^{[p^e]}`$, then $`f^{r_{e1}}=(f^{pr_{e1}})^{1/p}๐ช^{[p^{e1}]}A^{1/p}A=๐ช^{[p^{e1}]}`$, because $`AA^{1/p}`$ is pure. Hence $`pr_{e1}j_e<r_e`$, and $`j_e+1`$ is not divisible by $`p`$ since $`r_epr_{e1}p1`$. Therefore, by differentiating $`f^{j_e+1}๐ช^{[p^e]}`$ by $`x_i`$, we have $`f^{j_e}f/x_i๐ช^{[p^e]}`$ for $`i=1,\mathrm{},d`$, so that
$$f^{j_e}๐ช^{[p^e]}:J๐ช^{[p^e]}:(x_1^{m_1},\mathrm{},x_d^{m_d})๐ช^{[p^e]}+(x_1^{p^em_1}\mathrm{}x_d^{p^em_d})A.$$
Then $`f^{j_e}`$ must have a nonzero term in $`(x_1^{p^em_1}\mathrm{}x_d^{p^em_d})A`$, since $`f^{j_e}๐ช^{[p^e]}`$. Hence $`p^ed_{i=1}^dm_i=\mathrm{deg}f^{j_e}=nj_e<nr_entp^e`$. This inequality, together with the assumption $`t<d/n`$, implies $`p^e<\mu `$, which contradicts to $`p\mu `$. $`\mathrm{}`$
## 3. Main Theorem
First, we recall the definition of log terminal and log canonical singularities.
Let $`f:XY`$ be a proper birational morphism of normal varieties over a field $`k`$ and let $`E=_{i=1}^sE_i`$ be the exceptional divisor of $`f`$ with irreducible components $`E_1,\mathrm{},E_s`$. For a $``$-Weil divisor $`D`$ on $`Y`$ (resp. on $`X`$), we denote by $`f_{}^1(D)`$ (resp. $`f_{}D`$) the strict transform of $`D`$ in $`X`$ (resp. in $`Y`$). A $``$-Weil divisor $`D`$ on $`Y`$ is said to be $``$-Cartier if $`rD`$ is a Cartier divisor for some integer $`r>0`$. Then the pull-back $`f^{}(rD)`$ of $`rD`$ is also a Cartier divisor on $`X`$, and we can define the pull-back $`f^{}D`$ of $`D`$ as a $``$-Cartier divisor by $`f^{}D=\frac{1}{r}f^{}(rD)`$.
We denote the dualizing sheaves of $`X`$ and $`Y`$ by $`\omega _X`$ and $`\omega _Y`$, respectively, and fix $`\omega _X`$ and $`\omega _Y`$ as divisorial subsheaves of the rational function field $`k(X)=k(Y)`$ so that they coincide with each other outside the exceptional locus of $`f`$. Then the canonical divisor $`K_X`$ of $`X`$ (resp. $`K_Y`$ of $`Y`$) is also fixed by $`\omega _X=๐ช_X(K_X)`$ (resp. $`\omega _Y=๐ช_Y(K_Y)`$), and the strict transform of $`K_X`$ in $`Y`$ is $`f_{}K_X=K_Y`$.
Now let $`\mathrm{\Delta }`$ be an effective $``$-Weil divisor on $`Y`$ such that $`K_Y+\mathrm{\Delta }`$ is $``$-Cartier and denote by $`\stackrel{~}{\mathrm{\Delta }}`$ the strict transform $`f_{}^1(\mathrm{\Delta })`$ of $`\mathrm{\Delta }`$ in $`X`$. Let $`r>0`$ be an integer such that $`r(K_Y+\mathrm{\Delta })`$ is Cartier. Then $`r(K_X+\stackrel{~}{\mathrm{\Delta }})`$ and $`rf^{}(K_Y+\mathrm{\Delta })=f^{}(r(K_Y+\mathrm{\Delta }))`$ have integer coefficients, and coincide with each other outside the exceptional locus of $`f`$. Hence $`r(K_X+\stackrel{~}{\mathrm{\Delta }})=rf^{}(K_Y+\mathrm{\Delta })+_{j=1}^sb_jE_j`$ for some $`b_1,\mathrm{},b_s`$, and we have
$$K_X+\stackrel{~}{\mathrm{\Delta }}=f^{}(K_Y+\mathrm{\Delta })+\underset{j=1}{\overset{s}{}}a_jE_j,\text{where}a_j=\frac{b_j}{r}(j=1,\mathrm{},s).$$
$`\mathrm{3.1.1}`$
We call $`a_j`$ the discrepancy of $`E_j`$ with respect to $`(Y,\mathrm{\Delta })`$. In the following definition, we keep this notation for any desingularization $`f:XY`$ under consideration.
###### Demonstration Definition 3.1
Let $`Y`$ be a normal variety of characteristic zero, $`\mathrm{\Delta }`$ be an effective $``$-Weil divisor, and assume that $`K_Y+\mathrm{\Delta }`$ is $``$-Cartier.
(1) $`(Y,\mathrm{\Delta })`$ is said to be Kawamata log terminal (or KLT for short), if $`\mathrm{\Delta }=0`$ and if for every resolution of singularities $`f:XY`$ and every $`f`$-exceptional divisor $`E_j`$, the discrepancy $`a_j`$ defined in (3.1.1) satisfies $`a_j>1`$.
(2) $`(Y,\mathrm{\Delta })`$ is purely log terminal (or PLT for short), if for every resolution of singularities $`f:XY`$ and every $`f`$-exceptional divisor $`E_j`$, one has $`a_j>1`$.
(3) $`(Y,\mathrm{\Delta })`$ is log canoninal (or LC for short), if for every resolution of singularities $`f:XY`$ and every $`f`$-exceptional divisor $`E_j`$, one has $`a_j1`$.
When the pair $`(Y,0)`$ is KLT, or equivalently, PLT in this case (resp. LC), we say that $`Y`$ has log terminal (resp. log canonical) singularities.
###### Remark Remark 3.2
(1) Clearly, we have the implications โKLT $``$ PLT $``$ LC,โ and if $`(Y,\mathrm{\Delta })`$ is LC, then $`\mathrm{\Delta }`$ is reduced, i.e., every coefficients of $`\mathrm{\Delta }`$ is less than or equal to 1. Also, conditions (1) and (3) of Definition 3.1 is checked by referring some log resolution of $`(Y,\mathrm{\Delta })`$, that is, a resolution of singularities $`f:XY`$ such that the union of the exceptional set and $`\text{Supp}(f^1\mathrm{\Delta })`$ is a simple normal crossing divisor. Namely, $`(Y,\mathrm{\Delta })`$ is KLT (resp. LC) if and only if $`\mathrm{\Delta }=0`$ (resp. $`\mathrm{\Delta }`$ is reduced) and there exists a log resolution $`f:XY`$ such that $`a_j>1`$ (resp. $`a_j1`$) for every $`f`$-exceptional divisor $`E_j`$. Also, $`(Y,\mathrm{\Delta })`$ is PLT if and only if there exists a log resolution $`f:XY`$ such that $`f_{}^1\mathrm{\Delta }`$ is smooth and that $`a_j>1`$ for every $`f`$-exceptional divisor $`E_j`$.
(2) Log terminal (resp. log canonical) singularities behave similarly as strongly F-regular (resp. F-pure) rings under finite covers. Namely, if $`(A,๐ช)(B,๐ซ)`$ is a finite local homomorphism of normal local rings which is รฉtale in codimension 1, then $`A`$ is log terminal (resp. log canonical) if and only if so is $`B`$.
(3) If $`A`$ has log terminal singularities, then $`A`$ has rational singularities. Conversely, if $`A`$ has Gorenstein rational singularities, then $`A`$ has log terminal singularities.
###### Theorem 3.3
Let $`(A,๐ช)`$ be an F-finite normal local ring of characteristic $`p>0`$ and $`\mathrm{\Delta }`$ be an effective $``$-Weil divisor on $`Y=\mathrm{Spec}A`$ such that $`K_Y+\mathrm{\Delta }`$ is $``$-Cartier. Let $`f:XY=\mathrm{Spec}A`$ be a proper birational morphism with $`X`$ normal. Let $`E=_{j=1}^sE_j`$ be the exceptional divisor of $`f`$ and let
$$K_X+\stackrel{~}{\mathrm{\Delta }}=f^{}(K_Y+\mathrm{\Delta })+\underset{j=1}{\overset{s}{}}a_jE_j$$
as in $`(\mathrm{3.1.1})`$. If the pair $`(A,\mathrm{\Delta })`$ is F-pure $`(`$resp. divisorially F-regular$`)`$, then $`a_j1`$ $`(`$resp. $`a_j>1)`$ for every $`j=1,\mathrm{},s`$.
We will give two different proofs to this theorem, in which we relax the properness assumption of $`f:XY=\mathrm{Spec}A`$ as follows: $`f`$ factorizes into $`f:X\overline{X}\stackrel{\overline{f}}{}Y`$, where $`\overline{f}:\overline{X}Y`$ is a proper birational morphism with $`\overline{X}`$ normal, and $`X`$ is an open subset of $`\overline{X}`$ with $`\text{codim}(\overline{X}X,\overline{X})2`$. This assumption is preserved if we replace $`X`$ by the nonsingular locus $`X_{\text{reg}}`$ of $`X`$, and we do not lose any information about the discrepancies $`a_j`$ by this replacement, since $`X_{\text{reg}}`$ intersects every $`E_j`$. Hence we may assume without loss of generality that $`X`$ is Gorenstein in what follows.
###### Demonstration First proof
We will divide the proof into six steps. First (since $`\overline{f}`$ is proper and birational), $`f`$ is an isomorphism on an open set $`U`$ of $`Y`$ with $`\text{ codim}(YU,Y)2`$. Since $`K_A^{(i)}`$ is reflexive, we have
$$H^0(X,\omega _X^i)H^0(U,\omega _U^i)=K_A^{(i)}$$
for every $`i`$. Also, we can choose a nonzero element $`bA`$ such that
$$bK_A^{(i)}H^0(X,\omega _X^{(i)})\text{for each }i=0,1,\mathrm{},r1.$$
$`\mathrm{3.3.1}`$
(1) Since $`X`$ is Gorenstein, we have
$`om_{๐ช_X}(๐ช_X^{1/q},๐ช_X)`$ $`om_{๐ช_X}(๐ช_X^{1/q},\omega _X)_{๐ช_X}\omega _X^1`$ $`\mathrm{3.3.2}`$
$`(\omega _X)^{1/q}\omega _X^1๐ช_X((1q)K_X)^{1/q}.`$
by the adjunction formula for the finite map $`๐ช_X๐ช_X^{1/q}`$. By a similar argument on $`U`$, we also have an isomorphism
$$\alpha :\mathrm{Hom}_A(A^{1/q},A)\stackrel{}{}(K_A^{(1q)})^{1/q},$$
which induces $`\mathrm{Hom}_A(A((q1)\mathrm{\Delta })^{1/q},A)A((1q)K_A(q1)\mathrm{\Delta })^{1/q}`$ and is compatible with (3.3.2).
(2) Let us fix an embedding $`๐ช_X(iK_X)L`$, where $`L=k(X)`$ is the rational function field of $`X`$, and put $`A(r(K_A+\mathrm{\Delta }))=wA`$. Then $`A((mr+i)(K_A+\mathrm{\Delta }))=w^mA(i(K_A+\mathrm{\Delta }))`$, and $`\omega _X^r=๐ช_X(rK_X)=๐ช_X(r(f^{}(K_Y+\mathrm{\Delta })+a_jE_j\stackrel{~}{\mathrm{\Delta }}))=f^{}w๐ช_X(ra_jE_j\stackrel{~}{\mathrm{\Delta }})`$.
(3) Now, assume that $`(A,\mathrm{\Delta })`$ is F-pure and let $`\varphi :A((q1)\mathrm{\Delta })^{1/q}A`$ be a splitting of $`AA((q1)\mathrm{\Delta })^{1/q}`$. Then $`\varphi `$ induces a splitting $`\stackrel{~}{\varphi }:L^{1/q}L`$ of $`LL^{1/q}`$. Let $`X^{}=XZ`$, where $`Z=\text{Supp}((\stackrel{~}{\varphi }(๐ช_X^{1/q})+๐ช_X)/๐ช_X)`$. Then $`X^{}`$ is an open subset of $`X`$ since $`๐ช_X^{1/q}`$ is a coherent $`๐ช_X`$-module (we always assume that $`A`$ is F-finite). Also $`\stackrel{~}{\varphi }`$ induces an $`๐ช_X^{}`$-linear map $`๐ช_X^{}^{1/q}๐ช_X^{}`$, which we denote by the same letter $`\varphi `$. Then $`\varphi :๐ช_X^{}^{1/q}๐ช_X^{}`$ gives an F-splitting of $`X^{}`$ since the composition map $`๐ช_X^{}๐ช_X^{}^{1/q}\stackrel{\mathit{\varphi }}{}๐ช_X^{}`$ is the identity (cf. Lemma 2 of \[MS\]).
(4) Let $`\varphi `$ be as in (3) and fix an irreducible component $`E_j`$ of $`E`$. We want to examine if $`\varphi `$ is defined at the generic point of $`E_j`$ or not. Let $`\xi `$ be the generic point of $`E_j`$ and $`\eta `$ be a regular parameter of $`๐ช_{X,\xi }`$. We write $`q1=mr+i`$ with $`0i<r`$. Then by (3.3.1) we have
$$\alpha (\varphi )w^mA(iK_Ai\mathrm{\Delta })b^1w^mH^0(X,\omega _X^{(i)}).$$
(5) Now, we will show that $`a_j1`$ if $`(A,\mathrm{\Delta })`$ is F-pure. Assume, on the contrary, that $`a_j<1`$. Then, since $`ra_j`$, we have $`a_j11/r`$, and for a fixed integer $`s`$, we can take $`q0`$ so that $`mra_jq+s`$. This means that if $`\kappa `$ is a local generator of the $`๐ช_X^{1/q}`$-module $`om_{๐ช_X}(๐ช_X^{1/q},๐ช_X)`$ at $`\xi `$, then the stalk $`\varphi _\xi `$ of $`\varphi `$ at $`\xi `$ lies in $`[\eta ^q๐ช_{X,\xi }]^{1/q}\kappa `$. To see this let $`s=v_{E_j}(b)`$, the value of $`bA`$ at $`\xi `$, and let $`mra_jq+s`$. Then $`๐ช_{X.\xi }^{1/q}\alpha _\xi (\kappa )=[w^m\eta ^{mra_j}\omega _{X,\xi }^{(i)}]^{1/q}[w^m\eta ^{qs}\omega _{X,\xi }^{(i)}]^{1/q}`$ by (2), and $`\alpha (\varphi )_\xi [\eta ^sw^m\omega _{X,\xi }^{(i)}]^{1/q}`$ by (4). It follows that $`\varphi _\xi [\eta ^q๐ช_{X,\xi }]^{1/q}\kappa `$, whence $`\varphi _\xi (๐ช_{X,\xi }^{1/q})\eta ๐ช_{X,\xi }`$. This implies that $`\xi X^{}`$ but $`\varphi `$ is not an F-splitting at $`\xi `$. This contradiction concludes that $`a_i1`$.
(6) If $`(A,\mathrm{\Delta })`$ is divisorially F-regular, then for any $`c0A`$ which is not in any minimal prime of $`A(\mathrm{\Delta })`$, there exists an $`A`$-homomorphism $`\psi :A((q1)\mathrm{\Delta })^{1/q}A`$ sending $`c^{1/q}`$ to $`1`$. Then $`\varphi =c^{1/q}\psi `$ gives a splitting of $`AA((q1)\mathrm{\Delta })^{1/q}`$, and we may think $`\varphi c^{1/q}\mathrm{Hom}_A(A((q1)\mathrm{\Delta })^{1/q},A)[cA((1q)K_A(q1)\mathrm{\Delta })]^{1/q}`$. Let us choose $`c`$ so that the value $`t=v_{E_j}(c)`$ of $`c`$ at $`\xi `$ satisfies $`tr+s=r+v_{E_j}(b)`$. Then, arguing as in the F-pure case, we see that if $`a_j=1`$, then $`\varphi _\xi [\eta ^q๐ช_{X,\xi }]^{1/q}\kappa `$ and again $`\varphi `$ cannot be an F-splitting at $`\xi `$.
This completes the proof of the theorem. $`\mathrm{}`$
###### Demonstration Second proof
Next we shall give an alternative proof of Theorem 3.3 with different flavor. As we have seen in the first proof, the adjunction formula gives
$$om_{๐ช_X}(๐ช_X^{1/q},๐ช_X)(\omega _X^{(1q)})^{1/q}๐ช_X((1q)K_X)^{1/q},$$
$$\mathrm{Hom}_A(A((q1)\mathrm{\Delta })^{1/q},A)A((1q)K_A(q1)\mathrm{\Delta })^{1/q}.$$
(1) Now assume that $`(A,\mathrm{\Delta })`$ is F-pure and let $`\varphi :A((q1)\mathrm{\Delta })^{1/q}A`$ be a splitting of $`AA((q1)\mathrm{\Delta })^{1/q}`$. We regard $`\varphi \mathrm{Hom}_A(A((q1)\mathrm{\Delta })^{1/q},A)A((1q)K_A(q1)\mathrm{\Delta })^{1/q}`$ as a rational section of the sheaf $`\omega _X^{(1q)}`$ and consider the corresponding divisor on $`X`$,
$$D=D_\varphi =(\varphi )_0(\varphi )_{\mathrm{}},$$
where $`(\varphi )_0`$ (resp. $`(\varphi )_{\mathrm{}}`$) is the divisor of zeros (resp. poles) of $`\varphi `$ as a rational section of $`\omega _X^{(1q)}`$. Clearly, $`D`$ is linearly equivalent to $`(1q)K_X`$ and $`(\varphi )_{\mathrm{}}`$ is an $`f`$-exceptional divisor. Hence $`f_{}D`$ is linearly equivalent to $`(1q)K_Y`$ and $`f_{}D(q1)\mathrm{\Delta }`$. We denote $`X^{}=X\text{Supp }(\varphi )_{\mathrm{}}`$. Then $`\varphi `$ lies in
$$\mathrm{Hom}_{๐ช_X^{}}(๐ช_X^{}^{1/q},๐ช_X^{})H^0(X^{},\omega _X^{}^{(1q)}),$$
and gives an F-splitting of $`X^{}`$.
(2) We show that the coefficient of $`D`$ in each irreducible component is $`q1`$. Assume to the contrary that there exists an irreducible component of $`D`$, say $`D_0`$, whose coefficient is $`q`$. Then $`\varphi `$ lies in
$$\mathrm{Hom}_{๐ช_X^{}}(๐ช_X^{}(qD_0)^{1/q},๐ช_X^{})H^0(X^{},\omega _X^{}^{(1q)}(qD_0)),$$
and gives a splitting of the map $`๐ช_X^{}๐ช_X^{}(qD_0)^{1/q}`$. But this map factors through $`๐ช_X(D_0)`$, and $`๐ช_X^{}๐ช_X^{}(D_0)`$ never splits as an $`๐ช_X^{}`$-module homomorphism, since $`D_0`$ intersects $`X^{}`$. Consequently, every coefficient of $`D`$ must be $`q1`$.
(3) Let
$$B=\frac{1}{q1}D\stackrel{~}{\mathrm{\Delta }}.$$
Then $`B`$ is $``$-linearly equivalent to $`(K_X+\stackrel{~}{\mathrm{\Delta }})`$, so that $`f_{}B`$ is $``$-linearly equivalent to $`f_{}(K_X+\stackrel{~}{\mathrm{\Delta }})=(K_Y+\mathrm{\Delta })`$. Hence $`f_{}B`$ is $``$-Cartier, and we can define the pull-back $`f^{}f_{}B`$. Since $`B+_{i=1}^sa_iE_i`$ is $``$-linearly equivalent to $`f^{}(K_Y+\mathrm{\Delta })`$, $`(Bf^{}f_{}B)+_{i=1}^sa_iE_i`$ is an $`f`$-exceptional divisor which is $``$-linearly trivial relative to $`f`$. Hence
$$(Bf^{}f_{}B)+\underset{i=1}{\overset{s}{}}a_iE_i=0.$$
$`\mathrm{3.3.3}`$
(4) Now $`f_{}D(q1)\mathrm{\Delta }(q1)\mathrm{\Delta }(q1)\mathrm{\Delta }\mathrm{\Delta }^{}`$ for some (effective $``$-) Cartier divisor $`\mathrm{\Delta }^{}`$ on $`Y`$ which is independent of $`q`$. This implies $`f_{}B{\displaystyle \frac{1}{q1}}\mathrm{\Delta }^{}`$, whence $`f^{}f_{}B{\displaystyle \frac{1}{q1}}f^{}\mathrm{\Delta }^{}`$. Therefore, if we take $`q=p^e`$ sufficiently large, then the coefficient of $`f^{}f_{}B`$ in $`E_i`$ is greater than $`1/r`$. On the other hand, we have seen in (2) that the coefficient of $`B`$ in $`E_i`$ is at most 1. Since $`ra_i`$, it follows from (3.3.3) that $`a_i1`$.
(5) Assume now that $`(A,\mathrm{\Delta })`$ is divisorially F-regular. For a $``$-Cartier divisor $`\mathrm{\Delta }^{}`$ as in (4), we choose $`c0A`$ which is not in any minimal prime of $`A(\mathrm{\Delta })`$ so that the value $`v_{E_i}(c)`$ is greater than the coefficient of $`f^{}\mathrm{\Delta }^{}`$ in $`E_i`$. Then there exists an $`A`$-linear map $`\psi :A((q1)\mathrm{\Delta })^{1/q}A`$ sending $`c^{1/q}`$ to $`1`$, and $`\varphi =c^{1/q}\psi `$ gives a splitting of $`AA((q1)\mathrm{\Delta })^{1/q}`$. Let $`D_\varphi `$ and $`D_\psi `$ be the divisors defined by $`\varphi `$ and $`\psi `$ as rational sections of $`\omega _X^{(1q)}`$, respectively. Then the coefficient of $`D_\varphi `$ in $`E_i`$ is $`q1`$ by (2), so that the coefficient of $`D_\psi +f^{}\mathrm{\Delta }^{}`$ in $`E_i`$ is $`<q1`$, since $`D_\varphi =D_\psi +\text{div}_X(c)`$. Let $`B=(1/(q1))D_\psi \stackrel{~}{\mathrm{\Delta }}`$ and argue as in the F-pure case. Then we have $`a_i>1`$, as required. $`\mathrm{}`$
###### Remark Remark 3.4
The both proofs show that if the map $`AA((q1)\mathrm{\Delta })^{1/q}`$ sending $`1`$ to $`c^{1/q}`$ has a splitting for a single element $`c`$ in a sufficiently high power of $`H^0(X,๐ช_X(E))`$, then we have $`a_j>1`$ for every $`j`$. In the case $`\mathrm{\Delta }=0`$, if the map $`A\stackrel{c^{1/q}}{}A((q1)\mathrm{\Delta })^{1/q}`$ splits for an element $`c`$ with $`v_{E_i}(c)>0`$, then $`a_i>1`$.
###### Remark Remark 3.5
It follows that if $`(A,\mathrm{\Delta })`$ is strongly F-regular and if $`K_Y+\mathrm{\Delta }`$ is $``$-Cartier, then $`\mathrm{\Delta }=0`$ and in (3.1.1), one has $`a_j>1`$ for every $`j`$. But the above proof says something about strong F-regularity even when $`K_Y+\mathrm{\Delta }`$ is not $``$-Cartier.
Let $`f:XY=\mathrm{Spec}A`$ be a resolution of singularities admitting an $`f`$-ample Cartier divisor $`H`$ supported on the exceptional locus of $`f`$, and assume that $`(A,\mathrm{\Delta })`$ is strongly F-regular. Then we can prove that there exists a $``$-divisor $`G`$ on $`X`$ satisfying the following conditions:
To see this, for an effective $``$-divisor $`\mathrm{\Delta }^{}`$ as in (4) above, choose $`c0A`$ such that $`\text{div}_X(c)f^{}\mathrm{\Delta }^{}`$ and that this is a strict inequality for the coefficients of $`E_j`$โs. Then there is an $`A`$-linear map $`\theta :A((q1)\mathrm{\Delta })^{1/q}A`$ sending $`c^{2/q}`$ to $`1`$, by strong F-regularity. Let $`\psi =c^{1/q}\theta `$ and define the divisor $`D_\psi `$ as in the step (5) above. Now we put $`G={\displaystyle \frac{1}{q1}}D_\psi +\epsilon H`$ for $`\epsilon `$ with $`0<\epsilon 1`$. Then $`G`$ satisfies conditions (ii) and (iii). Also, the coefficient of $`G`$ in each irreducible component $`G_j`$ is $`1`$ and is $`>1`$ if $`G_j`$ is $`f`$-exceptional (resp. $`0`$ if $`G_j`$ is not $`f`$-exceptional). Since $`\mathrm{\Delta }=0`$ and ampleness is an open condition, we can perturb coefficients of $`G`$ slightly so that $`G`$ satisfies condition (i) as well as (ii) and (iii).
###### Demonstration Example 3.6
Let us take a look at two typical examples in the case $`\mathrm{\Delta }=0`$.
(1) Let $`A`$ be the localization of $`k[X,Y,Z,W]/(X^4+Y^4+Z^4+W^4)`$ at the unique graded maximal ideal, where $`k`$ is a field of characteristic $`p`$. By a criterion by R. Fedder \[F\], $`A`$ is F-pure if (and only if) $`p1`$ (mod 4). On the other hand, if $`f:X\mathrm{Spec}A`$ is the blowing-up of the maximal ideal of $`A`$, then $`X`$ is regular with $`\omega _X๐ช_X(E_0)`$, where $`E_0`$ is the exceptional divisor of $`f`$. Then the discrepancy of $`E_0`$ is $`a_0=1`$.
(2) Let $`A`$ be the localization of the $`r`$th Veronese subring $`k[X_1,\mathrm{},X_n]^{(r)}`$ of $`k[X_1,\mathrm{},X_n]`$ at the unique graded maximal ideal, where $`k`$ is a field of characteristic $`p`$. Then $`A`$ is $``$-Gorenstein with index $`r/(r,n)`$ and is surely strongly F-regular being a pure subring of a regular ring. Again, let $`X`$ be the blowing-up of the maximal ideal of $`A`$. Then $`X`$ is regular and if $`E_0`$ is the exceptional divisor, we have $`a_0=1+n/r`$.
Now, let us discuss rings essentially of finite type over a field of characteristic zero. Our goal is the following
###### Theorem 3.7
Let $`A`$ be a normal local ring essentially of finite type over a field $`k`$ of characteristic zero and let $`\mathrm{\Delta }`$ be an effective $``$-Weil divisor on $`Y=\mathrm{Spec}A`$ such that $`K_Y+\mathrm{\Delta }`$ is $``$-Carier. If $`(A,\mathrm{\Delta })`$ is of strongly F-regular (resp. divisorially F-regular, F-pure) type, that means the reduction modulo $`p`$ of $`(A,\mathrm{\Delta })`$ is strongly F-regular (resp. divisorially F-regular, F-pure) for infinitely many prime $`p`$, then $`(Y,\mathrm{\Delta })`$ is KLT (resp. PLT, LC).
###### Demonstration Proof
Let $`Z`$ be a normal affine variety over $`k`$ such that $`A๐ช_{Z,z}`$ for some $`zZ`$. We may assume that $`\mathrm{\Delta }`$ comes from a $``$-divisor $`\mathrm{\Delta }_Z`$ such that $`K_Z+\mathrm{\Delta }_Z`$ is $``$-Cartier. Let $`g:\stackrel{~}{X}Z`$ be a log resolution of $`(Z,\mathrm{\Delta }_Z)`$ (such that $`g_{}^1\mathrm{\Delta }_Z`$ is smooth, in order to prove PLT). Then since $`Z,\mathrm{\Delta }_Z,g,\stackrel{~}{X}`$ are defined by finite elements of $`k`$, we can choose a finitely generated $``$-subalgebra $`R`$ of $`k`$, schemes $`Z_R,\stackrel{~}{X}_R`$ of finite type over $`R`$, a $``$-divisor $`\mathrm{\Delta }_R`$ and an $`R`$-morphism $`g_R:\stackrel{~}{X}_RZ_R`$, which give $`Z,\stackrel{~}{X},\mathrm{\Delta }_Z`$ and $`g`$ after tensoring $`k`$ over $`R`$. Then, taking suitable open subset of $`\mathrm{Spec}R`$, we may assume $`Z_R,\stackrel{~}{X}_R`$ and each irreducible component of $`\mathrm{\Delta }_R`$ are flat over $`R`$, $`\stackrel{~}{X}_R`$ is smooth over $`R`$ with each irreducible component of the exceptional divisor and $`Z_R`$ is normal over $`R`$ (cf. \[EGA, IV, 12.1.7\]). Also, we can choose $`R`$ to be regular with trivial dualizing module. Then we may assume that the numbers $`a_i`$ in the formula $`K_X+\stackrel{~}{\mathrm{\Delta }}=f^{}(K_Y+\mathrm{\Delta })+_{i=1}^ra_iE_i`$ are preserved in the fibers of $`g_R:\stackrel{~}{X}_RZ_R`$ over $`\mathrm{Spec}R`$ in an open neighborhood of the generic point of $`\mathrm{Spec}R`$. By our assumption, this open neighborhood contains a maximal ideal $`๐ญ`$ of $`R`$, such that the base change of $`(Z_R,\mathrm{\Delta }_R)`$ to $`k(๐ญ)`$ over $`R`$ is strongly F-regular (resp. divisorially F-regular, F-pure) at $`\overline{z}`$, a specialization of $`z`$ (to be F-regular or F-pure is preserved under base field extension or restriction and so depends only on characteristic). This implies $`a_i>1`$ (resp. $`a_i1`$) by Theorem 3.3 and our assertion is proved. $`\mathrm{}`$
## 4. Applications
###### Demonstration 4.1. The graded case
(cf. \[W1, W2\]) Let $`R=_{n0}R_n`$ be a normal graded ring over a perfect field $`R_0=k`$ of characteristic $`p>0`$. Given a homogeneous element $`T`$ of degree $`1`$ in the quotient field of $`R`$, there is an ample $``$-Cartier divisor $`D`$ on $`X=\mathrm{Proj}R`$ such that
$$R=R(X,D)=\underset{n0}{}H^0(X,๐ช_X(nD))T^n.$$
If $`D=_{i=1}^s(e_i/d_i)D_i`$ for distict prime divisors $`D_1,\mathrm{},D_s`$ and coprime integers $`d_i`$ and $`e_i`$ with $`d_i>0`$ ($`1is`$), we denote $`D^{}=_{i=1}^s((d_i1)/d_i)D_i`$ and call it the โfractional partโ of $`D`$ (\[W1\]). Then $`K_R^{(r)}`$ is free if and only if $`r(K_X+D^{})`$ is linearly equivalent to $`bD`$ for some integer $`b`$.
If $`Z=\mathrm{Spec}_X\left(_{n0}๐ช_X(nD)T^n\right)`$ is the โgraded blowing-upโ of $`\mathrm{Spec}R`$ and if $`E_0X`$ is the exceptional divisor of this blowing-up, its discrepancy is $`a_0=1b/r`$. Consequently, if $`R`$ is strongly F-regular (resp. F-pure) then $`b<0`$ (resp. $`b0`$). (In this case, although $`Z`$ is not Gorenstein in general, the number $`a_0`$ is preserved after we make more blowing-ups and get a Gorenstein (or regular) scheme.)
Now let $`\mathrm{\Delta }=_{j=1}^rt_j\mathrm{\Delta }_j`$ be an effective $``$-Weil divisor on $`\mathrm{Spec}R`$ which is stable under the $`k^{}`$-action, i.e., each irreducible component $`\mathrm{\Delta }_j`$ of $`\mathrm{\Delta }`$ is defined by a homogeneous prime ideal of $`R`$ of height $`1`$. Then there exist effective $``$-Weil divisors $`\mathrm{\Gamma }_1,\mathrm{},\mathrm{\Gamma }_r`$ on $`X=\mathrm{Spec}R`$ such that for every $`i`$ and $`j=1,\mathrm{},r`$,
$$R(i\mathrm{\Delta }_j)=\underset{n}{}H^0(X,๐ช_X(i\mathrm{\Gamma }_j+nD))T^n.$$
Let $`\mathrm{\Gamma }=_{j=1}^rt_j\mathrm{\Gamma }_j`$ and $`D^{}`$ be the โfractional partโ of $`D`$ as above. Then we can rephrase Proposition 2.4 using this language, cf. \[W2\].
###### Proposition 4.2
Let the notation be as in 4.1 and let $`dimR=d+12`$.
###### Demonstration Example 4.3
Let $`R=R(^1,D)`$ with $`D=\frac{1}{2}((0)+(\mathrm{}))`$ and let $`\mathrm{\Delta }`$ be the divisor defined by $`R(\mathrm{\Delta })=_{n0}H^0(^1,๐ช((1)+nD))T^n`$. Then $`(R,\mathrm{\Delta })`$ is F-pure if and only if $`p2`$. This is a special case of the type (c) in 4.4 below.
###### Demonstration 4.4. Two-dimensional F-regular and F-pure pairs with integer coefficient boundary
Combining Theorem 3.3 and the technique used in \[Ha2\], we can classify F-regular and F-pure pairs $`(A,\mathrm{\Delta })`$ such that $`A`$ is a two-dimensional normal local ring of characteristic $`p>0`$ and that $`\mathrm{\Delta }`$ is an integer coefficient effective Weil divisor on $`\mathrm{Spec}A`$. The case $`\mathrm{\Delta }=0`$ is treated in \[Ha2\], so we treat here only the case $`\mathrm{\Delta }0`$.
Two-dimensional LC pairs with nonzero reduced boundaries are classified in terms of the dual graph of the union of the exceptional divisor $`E`$ and the strict transform $`\stackrel{~}{\mathrm{\Delta }}`$ of $`\mathrm{\Delta }`$ for the minimal log resolution $`f:XY=\mathrm{Spec}A`$, and such a pair is one of the following three types (\[A\], \[Ka\]):
Here, a blank circle $``$ (resp. a solid circle $``$) denotes an irreducible component $`E_i^1`$ of $`E`$ (resp. an irreducible component of $`\stackrel{~}{\mathrm{\Delta }}=f_{}^1(\mathrm{\Delta })`$), and the numbers โ$`2`$โ outside the circles in (c) mean that the corresponding components have self-intersection number equal to $`2`$. Among the above three types, type (a) is PLT but types (b) and (c) are not PLT.
Theorem 3.3 tells us that a two-dimensional F-pure (resp. divisorially F-regular) pair with reduced boundary is of type (a), (b) or (c) (resp. type (a)) as above. We can show that the converse is true but one exception.
###### Theorem 4.5
Let $`(A,๐ช)`$ be a two-dimensional normal local ring essentially of finite type over an algebraically closed field $`k=A/๐ช`$ of characteristic $`p>0`$ and let $`\mathrm{\Delta }`$ be a nonzero reduced Weil divisor on $`Y=\mathrm{Spec}A`$. Let $`f:XY`$ be the minimal resolution with exceptional divisor $`E`$ and let $`\stackrel{~}{\mathrm{\Delta }}=f_{}^1(\mathrm{\Delta })`$.
###### Demonstration Proof
(1) Let the dual graph of $`E\stackrel{~}{\mathrm{\Delta }}`$ be of type (a) as in 4.4, and fix any nonzero element $`cA`$ which is not in $`A(\mathrm{\Delta })=H^0(X,๐ช_X(\stackrel{~}{\mathrm{\Delta }}))`$. By Proposition 2.4, our goal is to show that there exists $`q=p^e`$ such that the map
$$cF^e:H_๐ช^2(K_A)H_๐ช^2(A(qK_A+(q1)\mathrm{\Delta }))$$
$`\mathrm{4.5.1}`$
is injective. Blowing up at $`E\stackrel{~}{\mathrm{\Delta }}`$ finitely many times if necessary, we may assume that $`\stackrel{~}{\mathrm{\Delta }}`$ does not intersect the strict transform of $`\mathrm{div}_Y(c)`$ on $`X`$ and that the graph of $`E\stackrel{~}{\mathrm{\Delta }}`$ is still of type (a) as follows.
$$\stackrel{~}{\mathrm{\Delta }}\text{}E_0\text{}E_1\text{}\mathrm{}\text{}E_l$$
Note that $`f:XY`$ may not be the minimal resolution any longer, but $`E_i`$ is a $`(1)`$-curve only if $`i=0`$. By Lemma 3.9 of \[Ha2\], one has an effective $`f`$-exceptional $``$-divisor
$$D=E_0+\frac{d1}{d}E_1+(\text{terms of }E_2,\mathrm{},E_l)$$
with $`d`$ such that $`(K_X+D)E_0=11/d`$ and $`(K_X+D)E_j=0`$ for $`j=1,\mathrm{},l`$. Let $`D^{}=DE_0`$.
Let $`n`$ be the integer with $`cH^0(X,๐ช_X(nE_0))H^0(X,๐ช_X((n+1)E_0))`$. Since $`H_๐ช^2(K_A)H_E^2(\omega _X)`$ by $`H^i(X,\omega _X)=0`$ $`(i=1,2)`$, the map (4.5.1), followed by $`H_๐ช^2(A(qK_A+(q1)\mathrm{\Delta }))H_E^2(๐ช_X(q(K_X+D^{})+(q1)(\stackrel{~}{\mathrm{\Delta }}+E_0)nE_0))`$, is
$$cF^e:H_E^2(\omega _X)H_E^2(๐ช_X(q(K_X+D^{})+(q1)(\stackrel{~}{\mathrm{\Delta }}+E_0)nE_0)),$$
$`\mathrm{4.5.2}`$
and it is sufficient to show that this map is injective for some $`q=p^e`$.
Let $`P,QE_0`$ be the points of intersection of $`E_0`$ with $`E_1,\stackrel{~}{\mathrm{\Delta }}`$, respectively, and let $`๐ก^{}=D^{}|_{E_0}=\frac{d1}{d}P`$ and $`๐=E_0|_{E_0}`$ as ($``$-)divisors on $`E_0^1`$. Then $`cH^0(X,๐ช_X(nE_0))`$ restricts to a nonzero element $`\overline{c}H^0(E_0,๐ช_{E_0}(n๐))`$, and $`\overline{c}H^0(E_0,๐ช_{E_0}(n๐Q))`$ since $`\stackrel{~}{\mathrm{\Delta }}f_{}^1\mathrm{div}_Y(c)=\mathrm{}`$.
Now if $`q=p^e`$ is sufficiently large, then $`q(K_X+D)(q1)\stackrel{~}{\mathrm{\Delta }}+nE_0`$ is $`f`$-nef, so that $`R^1f_{}\omega _X(q(K_X+D)(q1)\stackrel{~}{\mathrm{\Delta }}+nE_0)=0`$ (cf. \[Ha2, Lemma 3.3\]), or dually, $`H_E^1(๐ช_X(q(K_X+D)+(q1)\stackrel{~}{\mathrm{\Delta }}nE_0))=0`$. Hence we have the following commutative diagram with the vertical arrows being injective for $`q=p^e0`$.
$$\begin{array}{ccc}H_E^2(\omega _X)& \stackrel{cF^e}{}& H_E^2(๐ช_X(q(K_X+D^{})+(q1)(\stackrel{~}{\mathrm{\Delta }}+E_0)nE_0))\\ & & & & \\ H^1(E_0,\omega _{E_0})& \stackrel{\overline{c}F^e}{}& H^1(E_0,๐ช_{E_0}(q(K_{E_0}+๐ก^{})+(q1)Q+n๐))\end{array}$$
Note also that the map $`H^1(E_0,\omega _{E_0})H_E^2(\omega _X)`$ on the left is identified with $`k=A/๐ชE_A(A/๐ช)`$, whence an essential extension. By computing ฤech cohomologies on $`E_0^1`$, we can verify that the map $`\overline{c}F^e`$ at the bottom is injective for $`q=p^e0`$, from which follows the injectivity of the map (4.5.2) at the top.
(2) If the graph is of type (b) or (c), we use the anti-discrepancy of $`f`$ in place of $`D`$ in the proof of (1) above, i.e., $`D=f^{}(K_Y+\mathrm{\Delta })(K_X+\stackrel{~}{\mathrm{\Delta }})`$. Let $`D_0=D`$, $`D^{}=DD_0`$ and $`๐ก^{}=D^{}|_{D_0}`$, and argue as in (1) for $`c=1`$. Then it follows that $`(A,\mathrm{\Delta })`$ is F-pure if and only if the map
$$F^e:H^1(D_0,\omega _{D_0})H^1(D_0,๐ช_{D_0}(q(K_{D_0}+๐ก^{})+(q1)\stackrel{~}{\mathrm{\Delta }}|_{D_0}))$$
is injective for all $`q=p^e`$. We can show as in \[Ha2, Claim 4.8.1\] that this map is injective if and only if the graph of $`E\stackrel{~}{\mathrm{\Delta }}`$ is of type (b), or of type (c) and $`p2`$. See also \[Ha2\] for details. $`\mathrm{}`$
###### Demonstration 4.6. Strongly F-regular rings vs. admissible singularities
The correspondence of $``$-Gorenstein (strongly) F-regular rings and log terminal singularities (in the case $`\mathrm{\Delta }=0`$) is now well established. Namely, we have seen that a ring of characteristic zero has log terminal singularities if it is of (strongly) F-regular type and $``$-Gorenstein. The converse of this implication is also proved to be true \[Ha3\]. Thus we are tempted to consider what non-$``$-Gorenstein strongly F-regular rings are. We have no established answer to this question, but there is a candidate which is expected to correspond strongly F-regular rings even in non-$``$-Gorenstein case.
In \[N\], Nakayama introduced the notion of admissible singularities as an analog of log terminal singularities in the absence of $``$-Gorensteinness. Let $`Y`$ be a normal variety of characteristic zero and let $`\mathrm{\Delta }`$ be an effective $``$-Weil divisor on $`Y`$. Then the pair $`(Y,\mathrm{\Delta })`$ is said to be strictly admissible if there exist a birational morphism $`f:XY`$ from a nonsingular $`X`$ and a $``$-divisor $`G`$ on $`X`$ satisfying the following conditions:
On the other hand, we observed in Remark 3.5 that if $`(Y,\mathrm{\Delta })`$ is a strongly F-regular pair in characteristic $`p>0`$ and $`f:XY`$ is a resolution, then we can construct a $``$-divisor $`G`$ satisfying conditions (i), (ii) and (iii). But we do not have condition (iv), and even worse, this construction depends on characteristic $`p`$. However, the following example may be a positive evidence to the correspondence of strong F-regularity and admissible singularity.
###### Demonstration Example 4.7
Let $`X`$ be a smooth Fano variety of characteristic zero (i.e., a smooth projective variety with ample anti-canonical divisor $`K_X`$), and $`D`$ be any ample Cartier divisor on $`X`$. Let $`R=R(X,D)`$ and let $`f:Z\mathrm{Spec}R`$ be the โgraded blowing-upโ as in (4.1). The exceptional set of $`f`$ is a smooth divisor $`EX`$, so if we put $`G=(\epsilon 1)E`$ with $`0<\epsilon 1`$, conditions (i), (ii), (iv) in (4.6) are satisfied. Also, if we choose $`\epsilon 1`$, then $`(GK_Z)|_E=((K_Z+E)+\epsilon E)|_E=K_E\epsilon D`$ is ample since $`EX`$ is Fano, so that condition (iii) is satisfied. Consequently, $`R`$ has an admissible singularity whereas it is not $``$-Gorenstein in general.
On the other hand, $`R(X,K_X)`$ has a Gorenstein log terminal singularity, so that it is of F-regular type \[Ha3\]. However, since the strong F-regularity of $`R(X,D)`$ depends only on the โfractional partโ $`D^{}`$ of $`D`$ \[W2\] and the fractional parts of $`D`$ and $`K_X`$ are both $`0`$ in this case, $`R=R(X,D)`$ also has strongly F-regular type.
We study more about similarity of F-regularity and F-purity of pairs and โsingularities of pairsโ in characteristic zero \[Ko\], \[Sh\]. The following theorem generalizes \[W2, Theorem 2.7\]. See Remark 1.2 (5).
###### Theorem 4.8
Let $`(A,๐ช)(B,๐ซ)`$ be a finite local homomorphism of F-finite normal local rings which is รฉtale in codimension 1. Let $`\mathrm{\Delta }_A`$ be an effective $``$-Weil divisor on $`\mathrm{Spec}A`$ and let $`\mathrm{\Delta }_B=\pi ^{}\mathrm{\Delta }_A`$ be the pull-back of $`\mathrm{\Delta }_A`$ by the induced morphism $`\pi :\mathrm{Spec}B\mathrm{Spec}A`$. If $`(A,\mathrm{\Delta }_A)`$ is F-pure $`(`$resp. divisorially or strongly F-regular$`)`$, then so is $`(B,\mathrm{\Delta }_B)`$, too.
###### Demonstration Proof
Let $`d=dimA=dimB`$. Since $`AB`$ is รฉtale in codimension 1, the natural maps $`A^{1/q}_ABB^{1/q}`$ and $`A(qK_A+(q1)\mathrm{\Delta }_A)_ABB(qK_B+(q1)\mathrm{\Delta }_B)`$ are isomorphic in codimension 1. Hence, by \[W2, Lemma 2.2\], we have
$`H_๐ช^d(A(qK_A+(q1)\mathrm{\Delta }_A)^{1/q})`$ $`_AB`$
$`H_๐ช^d(A(qK_A+(q1)\mathrm{\Delta }_A)^{1/q}_AB)`$
$`H_๐ช^d(A(qK_A+(q1)\mathrm{\Delta }_A)^{1/q}_{A^{1/q}}A^{1/q}_AB)`$
$`H_๐ช^d(A(qK_A+(q1)\mathrm{\Delta }_A)^{1/q}_{A^{1/q}}B^{1/q})`$
$`H_๐ช^d((A(qK_A+(q1)\mathrm{\Delta }_A)_AB)^{1/q})`$
$`H_๐ซ^d(B(qK_B+(q1)\mathrm{\Delta }_B)^{1/q}).`$
Therefore, tensoring the map
$$F^e:H_๐ช^d(K_A)H_๐ช^d(A(qK_A+(q1)\mathrm{\Delta }_A))$$
$`\mathrm{4.8.1}`$
with $`B`$ over $`A`$ yields
$$F^e:H_๐ซ^d(K_B)H_๐ซ^d(B(qK_B+(q1)\mathrm{\Delta }_B)).$$
$`\mathrm{4.8.2}`$
Now, if $`(A,\mathrm{\Delta }_A)`$ is F-pure, then for every $`q=p^e`$, the map (4.8.1) is injective, and even a splitting injective map, since $`H_๐ช^d(K_A)E_A(A/๐ช)`$ is an injective $`A`$-module. Hence the map (4.8.2) is also injective for every $`q=p^e`$, so that $`(B,\mathrm{\Delta }_B)`$ is F-pure.
Similarly, if $`(A,\mathrm{\Delta }_A)`$ is strongly F-regular, then for every $`c0A`$, there exists $`q=p^e`$ such that the map
$$cF^e:H_๐ซ^d(K_B)H_๐ซ^d(B(qK_B+(q1)\mathrm{\Delta }_B))$$
is injective. This is true for every $`c0B`$, since $`AB`$ is a finite extension of normal domains, so that $`cBA0`$. Hence $`(B,\mathrm{\Delta }_B)`$ is strongly F-regular.
To prove the assertion for divisorial F-regularity, note that if $`cB`$ is not in any minimal prime ideal of $`B(\mathrm{\Delta }_B)`$, then there is an element of $`cBA`$ which is not in any minimal prime ideal of $`A(\mathrm{\Delta }_A)`$. Then the argument for strong F-regularity works also for divisorial F-regularity. $`\mathrm{}`$
###### Theorem 4.9
(cf. \[Ko, Theorem 7.5\]) Let $`(A,๐ช)`$ be an F-finite normal local ring of characteristic $`p>0`$ with $`Y=\mathrm{Spec}A`$ and let $`x๐ช`$ be a nonzero element.
###### Demonstration Proof
Let $`B=A/xA`$ and $`d=dimA`$. We will look at the Frobenius actions on $`E_A=E_A(A/๐ช)H_๐ช^d(K_A)`$ and $`E_B=E_B(B/๐ชB)H_๐ช^{d1}(K_B)`$ in slightly different ways in proving (1) and (2), respectively.
(1) Assume that $`(A,\mathrm{div}_Y(x))`$ is divisorially F-regular. Then for every element $`cA`$ which is not in any minimal prime ideal of $`xA`$, there exists $`q=p^e`$ such that the map $`cF^e:H_๐ช^d(K_A)H_๐ช^d(A(qK_Y+(q1)\mathrm{div}_Y(x)))`$ is injective, or equivalently, the map
$$\alpha :E_A=E_A_AAE_A_AA^{1/q}$$
sending $`\xi E_A`$ to $`\xi (cx^{q1})^{1/q}E_A_AA^{1/q}`$ is injective (see the proof of Proposition 2.6). On the other hand, we have an inclusion $`ฤฑ:E_B(0:x)_{E_A}E_A`$, and this map gives rise to a well-defined $`A^{1/q}`$-homomorphism
$$ศท:E_B_BB^{1/q}E_A_AA^{1/q}$$
sending $`\xi \overline{a}^{1/q}E_B_BB^{1/q}`$ to $`ฤฑ(\xi )(ax^{q1})^{1/q}E_A_AA^{1/q}`$, where $`\overline{a}`$ denotes the image of $`aA`$ in $`B=A/xA`$. Then we have the following commutative diagram.
$$\begin{array}{ccc}E_B& \stackrel{ฤฑ}{}& E_A\\ 1\overline{c}^{1/q}& & \alpha & & \\ E_B_BB^{1/q}& \stackrel{ศท}{}& E_A_AA^{1/q}\end{array}$$
Hence the injectivity of the map $`\alpha `$ implies that $`1_{E_B}\overline{c}^{1/q}:E_BE_B_BB^{1/q}`$ is injective. Consequently, we have that $`B`$ is strongly F-regular by Remark 1.2 (1).
(2) Assume that $`B`$ is strongly F-regular. Then $`B`$ is a CohenโMacaulay normal domain, so that $`A`$ and $`K_A`$ are also CohenโMacaulay. For any $`q=p^e`$ and $`cAxA`$, we consider the following commutative diagram with exact rows,
$$\begin{array}{ccccccccc}0& & K_A& \stackrel{x}{}& K_A& & K_A/xK_A& & 0\\ & & cx^{q1}F^e& & cF^e& & cF^e& & & \\ 0& & K_A^{(q)}& \stackrel{x}{}& K_A^{(q)}& & K_A^{(q)}/xK_A^{(q)}& & 0,\end{array}$$
where $`K_A/xK_AK_B`$. Since $`H_๐ช^{d1}(K_A^{(q)}/xK_A^{(q)})H_๐ช^{d1}(K_B^{(q)})`$ by the normality of $`B`$, this diagram induces the following commutative diagram of local cohomologies.
$$\begin{array}{ccc}H_๐ช^{d1}(K_B)& & H_๐ช^d(K_A)\\ cF^e& & cx^{q1}F^e& & \\ H_๐ช^{d1}(K_B^{(q)})& & H_๐ช^d(K_A^{(q)})\end{array}$$
Here, the map $`H_๐ช^{d1}(K_B)H_๐ช^d(K_A)`$ upstairs is an essential extension, and if $`q1`$ (mod $`r`$), then the map $`H_๐ช^{d1}(K_B^{(q)})H_๐ช^d(K_A^{(q)})`$ at the bottom is also injective since $`K_AK_A^{(q)}`$ is CohenโMacaulay.
Now, since $`B`$ is strongly F-regular, the map $`cF^e:H_๐ช^{d1}(K_B)H_๐ช^{d1}(K_B^{(q)})`$ in the above diagram is injective for all $`q=p^e0`$, by (2) of Propositions 2.2 and 2.4. Since $`r`$ is not divisible by $`p`$, we can choose $`q=p^e`$ so that $`q1`$ (mod $`r`$) and that $`cF^e:H_๐ช^{d1}(K_B)H_๐ช^{d1}(K_B^{(q)})`$ is injective. Then the diagram implies that the map $`cx^{q1}F^e:H_๐ช^d(K_A)H_๐ช^d(K_A^{(q)})`$ is injective for such a $`q=p^e`$, or equivalently,
$$cF^e:H_๐ช^d(K_A)H_๐ช^d(A(qK_Y+(q1)\mathrm{div}_Y(x)))$$
is injective. Hence $`(A,\mathrm{div}_Y(x))`$ is divisorially F-regular. $`\mathrm{}`$
###### Remark Remark 4.10
(1) In the situation of Theorem 4.9 (2), assume in addition that $`B=A/xA`$ is normal and CohenโMacaulay. Then we can also prove that $`(A,\mathrm{div}_Y(x))`$ is F-pure if and only if $`B=A/xA`$ is F-pure.
(2) The proof of Theorem 4.9 suggests a more general statement as follows: Let $`(A,๐ช)`$ and $`x๐ช`$ be as in Theorem 4.9. Let $`S=\mathrm{div}_Y(x)\mathrm{Spec}A/xA`$ and $`\mathrm{\Delta }`$ be an effective $``$-Weil divisor on $`Y=\mathrm{Spec}A`$ such that $`r(K_Y+\mathrm{\Delta })`$ is Cartier for some positive integer $`r`$ which is not divisible by $`p`$. Then $`(A,S+\mathrm{\Delta })`$ is divisorially F-regular if and only if $`(A/xA,\mathrm{\Delta }|_S)`$ is strongly F-regular.
If we replace โdivisorially F-regularโ and โstrongly F-regularโ in this assertion by โPLTโ and โKLTโ respectively, we find the so-called โinversion of adjunctionโ in characteristic zero \[Ko, Theorem 7.5\], \[Sh\].
###### Corollary 4.11
(cf. \[AKM\]) Let $`(A,๐ช)`$ be an F-finite local ring of characteristic $`p>0`$ and let $`x๐ช`$ be a nonzero element. Assume that $`A`$ is $``$-Gorenstein and that the order $`r`$ of the canonical class in the divisor class group of $`Y=\mathrm{Spec}A`$ is not divisible by $`p`$. If $`A/xA`$ is strongly F-regular, then $`A`$ is also strongly F-regular.
###### Demonstration Proof
If $`A/xA`$ is strongly F-regular, then $`A/xA`$ is F-rational. This implies that $`A`$ is F-rational, whence normal \[FW\], \[HH3, 4.2\]. Thus we can apply Theorem 4.9, which implies that $`(A,\mathrm{div}_Y(x))`$ is divisorially F-regular. It follows that $`(A,0)`$ is divisorially F-regular, meaning that $`A`$ is strongly F-regular. $`\mathrm{}`$
In fact, we do not have to assume that the index $`r`$ of $`A`$ is not divisible by $`p`$ in Corollary 4.11 (see Aberbach, Katzman and MacCrimmon \[AKM\]). On the other hand, an example by Singh \[Si\] shows that Corollary 4.11 fails in the absense of $``$-Gorensteinness.
## 5. Open Problems
So far, we have seen many positive evidences to the correspondence of F-purity (resp. divisorial, strong F-regularity) and LC (resp. PLT, KLT) property, which enable us to ask about the converse of Theorem 3.7.
###### Demonstration 5.1
Let $`A`$ be a normal ring essentially of finite type over a field of characteristic zero and let $`\mathrm{\Delta }`$ be a $``$-Weil divisor on $`Y=\mathrm{Spec}A`$.
###### Demonstration Conjecture 5.1.1
If $`(Y,\mathrm{\Delta })`$ is KLT (resp. PLT), then it is of open strongly F-regular type (resp. open divisorially F-regular type), i.e., the modulo $`p`$ reduction of $`(A,\mathrm{\Delta })`$ is strongly F-regular (resp. divisorially F-regular) for all $`p0`$.
Conjecture 5.1.1 is true when $`\mathrm{\Delta }=0`$ \[Ha3\]. Also, Theorem 4.8 and \[Ko, Theorem 7.5\] shows that a pair $`(A,\mathrm{div}_Y(x))`$ is PLT if and only if it is of divisorially F-regular type and $`A`$ is $``$-Gorenstein, since $`A/xA`$ is log terminal if and only if it is of F-regular type and $``$-Gorenstein.
###### Demonstration Problem 5.1.2
If $`(Y,\mathrm{\Delta })`$ is LC, then is it always of (dense) F-pure type, i.e., the modulo $`p`$ reduction of $`(A,\mathrm{\Delta })`$ is F-pure for infinitely many $`p`$?
###### Demonstration 5.2. Log canonical thresholds
(\[Ko, ยง8\], \[Sh\]) An affirmative answer to (5.1) suggests a Frobenius characterization of an important invariant called the log canonical threshold. Let $`Y`$ be a variety in characteristic zero with only log terminal singularity at a point $`yY`$ and $`\mathrm{\Delta }`$ be an effective $``$-Cartier divisor on $`Y`$. The log canonical threshold of $`\mathrm{\Delta }`$ at $`yY`$ is defined by
$$LCTh_y(Y,\mathrm{\Delta })=\text{sup}\{t|(Y,t\mathrm{\Delta })\text{ is LC at }yY\}.$$
Theorem 3.3 tells us that this invariant is greater than or equal to the supremum of $`t`$ such that reduction modulo $`p`$ of $`(๐ช_{Y,y},t\mathrm{\Delta })`$ is F-pure for infinitely many $`p`$.
###### Demonstration Conjecture 5.2.1
In the notation as above, the following are equal to each other:
We expect that $`LCTh_y(Y,\mathrm{\Delta })`$ is computable in terms of characteristic $`p`$ method. For example, by a Fedder-type criterion we can show that $`(k[[x,y]],t\mathrm{div}(x^2y^3))`$ is F-pure if and only if $`t5/6`$ (Example 2.9 (3)), and $`5/6`$ is nothing but the log canonical threshold of the hypersurface $`x^2y^3=0`$ in $`๐ธ_k^2`$ at the origin. Also, compare Proposition 2.10 with \[Ko, Lemma 8.10\].
###### Demonstration 5.3
Finally, we propose other open problems.
###### Demonstration Problem 5.3.1
If $`(A,\mathrm{\Delta })`$ is of strongly F-regular type ($`K_Y+\mathrm{\Delta }`$ is not necessarily $``$-Cartier), then does the pair $`(Y,\mathrm{\Delta })`$ have admissible singularities? How about the converse implication?
###### Demonstration Problem 5.3.2
Let $`(A,๐ช)`$ be a $`d`$-dimensional normal local ring of characteristic $`p>0`$ and let $`\mathrm{\Delta }`$ be an effective $``$-Weil divisor on $`Y=\mathrm{Spec}A`$ such that $`\mathrm{\Delta }=0`$. Define the tight closure of zero submodule in $`E=E_A(A/๐ช)H_๐ช^d(K_A)`$ with respect to the pair $`(A,\mathrm{\Delta })`$, $`0_E^\mathrm{\Delta }E=H_๐ช^d(K_A)`$, by
$`z0_E^\mathrm{\Delta }`$ $``$
$`c0A\text{ such that }cz^q=0\text{ in }H_๐ช^d(A(qK_A+(q1)\mathrm{\Delta }))\text{ for }q=p^e0,`$
where $`z^q`$ denotes the image of $`zE`$ by $`F^e:H_๐ช^d(K_A)H_๐ช^d(A(qK_A+(q1)\mathrm{\Delta }))`$. We ask for a geometric interpretation of the ideal $`\tau (A,\mathrm{\Delta })=\text{Ann}_A(0_E^\mathrm{\Delta })`$. Especially, if $`K_Y+\mathrm{\Delta }`$ is $``$-Cartier, is $`\tau (A,\mathrm{\Delta })`$ equal to the multiplier ideal of the pair $`(Y,\mathrm{\Delta })`$ in characteristic $`p0`$? The answer to this question is affirmative if $`\mathrm{\Delta }=0`$ (\[Ha4\], \[S2\]). See also \[HH1, 3\] for tight closure, and \[E\] for multiplier ideals.
Department of Mathematical Sciences, Waseda University, Okubo, Shinjuku-ku, Tokyo 169โ8555, Japan
E-mail: nharamn.waseda.ac.jp
Department of Mathematics, College of Humanities and Sciences, Nihon University, Sakura-josui, Setagaya-ku, Tokyo 156โ0045, Japan
E-mail: watanabemath.chs.nihon-u.ac.jp |
warning/0002/nucl-th0002064.html | ar5iv | text | # The role of electron-screening deformations in solar nuclear fusion reactions and the solar neutrino puzzle
## I Introduction
As the quest for the solution to the solar neutrino puzzle continues, all the parameters associated with the production and detection of solar neutrinos are exhaustively investigated. Naturally, stellar thermonuclear reactions have attracted a lot of attention in the nuclear physics community since not only do they govern the evolutionary stages of a star but also they can possibly hide the solution to the notorious discrepancy between theoretically and experimentally calculated solar neutrino fluxes.
The effects of the electron-ion screening on stellar fusion cross sections have drawn a lot of attention , especially in relation to the solar neutrino puzzle. Various screening prescriptions have been suggested, each of which has some inherent inadequacies. Salpeterโs weak screening formula can be safely used in the study of $`pp`$ solar fusion reactions where the weak screening regime is obeyed, though a recent study showed that, for most practical purposes, it can reasonably be applied to other solar reactions as well. On the other hand, the formula given by Mitler is considered the most reliable as it describes fairly well all the screening regimes, weak screening ($`WES`$) , intermediate ($`IS`$) and strong screening ($`SS`$), though the assumption of a constant electron density around the fusing nuclei is rather arbitrary.
As regards the solar neutrino puzzle, the prevalent belief is that, under the present standard solar model (SSM), the screening effect is under control. However, various non-standard solar models have been proposed, in an attempt to explain the discrepancy between the theoretical and experimental neutrino fluxes. Some of them are plausible, while others are exaggerated. Among those attempts, the existence of a primordial magnetic field in the solar interior has received relatively little attention, while its relation to the screened fusion reactions themselves has never been studied.
This lack of interest stems from the disheartening fact that if a large magnetic field (of the order of $`10^9`$ Gauss) is added to the equation of solar structure, the predicted event rate of the Cl<sup>37</sup> experiment is increased by a factor of two. As the magnetic field has the opposite effect from the one that is desired, no further investigation has been made by the same author. Admittedly, some other studies appeared which showed that a combination of differential rotation and strong magnetic fields in the sun can actually reduce the $`Cl^{37\text{ }}`$ signal but none of them has been adapted as a component of the SSM. Nevertheless, as it has been recently shown , a superstrong magnetic field can accelerate hydrogen fusion reactions in stars and in the laboratory, an effect which could influence the solar neutrino fluxes.
This is due to the fact that strong magnetic fields ($`B>10^9G)`$ compress hydrogen atoms both perpendicular and parallel to the field direction. The magnetic field ties the electrons to the field lines so that their response to a Coulomb attraction is essentially restricted to a one-dimensional motion parallel to the field. It is therefore plausible to assume that the presence of such a field would modify the screening effect of the electron cloud just as it happens on the surface of neutron stars. Such screening deformations can be studied qualitatively by means of the Debye-Hรผckel model. Actually, in the presence of a large magnetic field the DebyeโHรผckel radius which represents a spherical distribution of the electron cloud around the nuclei will become orientation dependent. This effects causes, inevitably, reaction rates and neutrino fluxes to be orientation dependent themselves.
The purpose of this paper is to investigate the dependence of thermonuclear fusion reaction rates on the deformations of the ionic-electron cloud that screens the reacting nuclei and its possible consequences on solar neutrino fluxes. The layout of the paper is as follows: In ยงII the advantages and disadvantages of the Debye-Hรผckel potential in the study of solar nuclear reactions are briefly investigated. In ยงIII the formalism for an axially deformed screening cloud is established, underlining its effects on the $`pp`$ reaction rates and the associated $`pp`$ neutrino fluxes. In ยงIV there is given a measure of the uncertainty of the solar neutrino fluxes due to the presence of such screening deformations. Finally ยงV investigates some potential sources of deformation, while the main results of the present paper are summarized in ยงVI.
## II Advantages and disadvantages of the Debye-Hรผckel potential
In the stellar plasma gravitational compression and quantum mechanical tunneling combine in order to achieve the classically impossible fusion between light nuclei. The electron gas that surrounds the nuclei acts as a catalyst in the reaction, by lowering the repulsive Coulomb barrier which prevents atomic nuclei from approaching each other.
In the framework of the Debye-Hรผckel model each nuclei is assumed to polarize its neighborhood creating a spherically symmetric but inhomogeneously charged ionic cloud around it. In this model, the potential $`V\left(r\right)`$ of a given nucleus of charge $`Z_1`$ can be found by the equation of Poisson :
$$^2V\left(r\right)=r_D^2V\left(r\right)$$
(1)
with
$$r_D^2=\frac{4\pi e^2}{kT}\left(\underset{i}{}Z_i^2n_i+n_e\theta _e\right)$$
(2)
where $`r_D`$ is the Debye-Hรผckel radius, $`\theta _e`$ is the electron degeneracy factor, and $`n`$ the number densities of ions $`\left(n_i\right)`$ with atomic number $`Z_i`$ and electrons $`\left(n_e\right),`$ respectively. The solution of Eq. $`\left(\text{1}\right)`$ is a Yukawa potential which acts in the vicinity of the ion $`Z_1e`$ and has the form:
$$V_D\left(r\right)=\frac{Z_1e}{r}\mathrm{exp}\left(\frac{r}{r_D}\right)$$
(3)
A measure of the stellar plasma response to Coulomb interactions is the plasma coupling parameter, which is defined as:
$$\mathrm{\Gamma }_{ij}=\frac{Z_iZ_je^2}{akT}$$
(4)
where $`a`$ is the mean interionic distance, $`Z_{ij\text{ }}`$the atomic numbers of the reactants and $`kT`$ the thermal kinetic energy. Although the domains of the weak screening $`\left(WES\right)`$, the intermediate screening $`\left(IMS\right)`$ and strong screening $`\left(SS\right)`$ are not precisely defined , in an ion fluid one can define them as follows:
$$WES:\mathrm{\Gamma }_{ij}<<1,IMS:\mathrm{\Gamma }_{ij}1,SS:\mathrm{\Gamma }_{ij}>>1$$
(5)
It is now easy to show that the assumption $`\mathrm{\Gamma }_{ij}<<1`$ of the $`\left(WES\right)`$ regime for the proton-atom reaction channel between the test nuclei $`Z_2e`$ and the generators of the above potential, i.e. nuclei $`Z_1e`$ plus ionic cloud, can be written:
$$\frac{Z_1Z_2e^2}{r_DkT}<<1$$
(6)
The above condition has generated a lot of controversy in the study of solar fusion cross sections as it is violated for solar nuclear reactions with $`Z_1Z_2`$ larger than unity.
Another interesting fact is that according to the above model for small $`r`$ the ion density now vanishes while the electron density diverges. By noting this, Mitler assumed the above model to be valid only beyond some radius $`r_1`$ while for shorter distances he assumed the that the electron density around the nucleus is constant and equal to the mean electron density in the plasma, $`n_e\left(\mathrm{}\right)`$.
However this is also an oversimplification, especially for the heavier nuclei of astrophysical interest where it underestimates the electron charge density around the ion. For instance near a $`Be^7`$ nucleus in central solar conditions we have a density , roughly $`3.8en_e\left(\mathrm{}\right)`$.
Moreover, the model in question is forced to undergo another compromise when calculating the thermalized cross section $`\sigma v^{sc}`$which appears in the screened reaction rate between nuclei $`i`$ and $`j`$ :
$$r_{ij}^{sc}=\left(1+\delta _{ij}\right)^1n_in_j\sigma v^{sc}$$
(7)
where $`n_i,n_j`$ are the number densities of nuclei $`i`$ and $`j`$respectively, and $`\delta _{ij}`$ is the Kronecker delta. As the WKB integral involved cannot be found analytically, we inevitably resort to a linear approximation of potential $`\left(\text{3}\right)`$ . However even if we donโt resort to that approximation and take into account non-linearities of the Debye-Hรผckel potential this would add a negligible contribution to Salpeterโs $`\left(WES\right)`$ prescription. Note that whenever the $`\left(WES\right)`$ condition is challenged one should always investigate non-linear corrections as will be shown in the study that follows.
Finally, the spherical symmetry around the two reactants is taken for granted by the standard Debye-Huckel model. The fact that no deformations are assumed for the ionic-electron cloud around the point like nucleus can obviously be the source of uncertainties in the calculation of stellar reaction rates. It is the principle objective of this paper to investigate the effects of such deformations on reaction rates and the associated solar neutrino fluxes, without focusing in detail on the sources of deformations themselves.
## III Electron screening deformations
In the adiabatic model the target and the projectile nuclei are assumed to be surrounded by a static, spherical electron cloud, whose electron charge density falls off exponentially with respect to the distance from the center of the cloud which is the nucleus itself. In central solar conditions the mean ion velocity $`u_i=\left(8kT/\pi \mu \right)^{1/2}`$ is roughly fifty times smaller than the mean electron velocity $`p_e^2/2m_e\left(3/2\right)kT`$ thus justifying the fact that as the nucleus moves the electron cloud has enough time to re-arrange itself so that it practically screens the nucleus at all times. However oscillations of the ionic cloud are inevitable due to their speed which is much lower than that of the electrons. For the main sequence stars Mitler showed that in the framework of the standard solar model (SSM) the distortion of the common charge clouds has only a small effect on the screening calculations $`\left(2\%\right).`$
Moreover, the possible presence of a strong magnetic field is bound to cause substantial deformations of the electron cloud by compressing it both parallel and perpendicular to the field direction. On the other hand the existence of other, as yet unspecified, sources of deformations cannot be ruled out.
Therefore treating the screening cloud as a rigid (albeit inhomogeneous) sphere is an assumption which must by further investigated, especially when primordial magnetic fields are considered.
Studies of heavy nuclei fusion reactions have shown that theoretical predictions of cross section can be greatly improved by assuming rotations and deformations of the fusing nuclei. It is therefore plausible to consider similar effects in the study of screened thermonuclear reactions where the electron cloud is assumed to be deformed. In fact this deformation can be parametrized in the framework of the liquid-drop model so that the Debye-Hรผckel radius is considered a measure of the electron cloud. The deformed DH radius is now:
$$r_D(\theta ,\varphi )=r_D^{\left(0\right)}\left[1+\underset{m}{}\beta _mY_2^m(\theta ,\varphi )\right]$$
(8)
where $`r_D^{\left(0\right)},`$ takes care of volume conservation and $`Y_2^m`$ is the usual spherical harmonic function. For simplicity and reasons that will soon become clear only quadrupole deformations will be considered. Moreover, we assume that the deformation is axially symmetric and take the $`z`$ axis along the axis of symmetry. Disregarding the rotational degree of freedom we obtain the surface shape
$$r_D\left(\theta \right)=r_D^{\left(0\right)}\left[1+\beta Y_2^0\left(\mathrm{cos}\theta \right)\right]$$
(9)
where the angle $`\theta `$ is measured from the axis of symmetry i.e. the $`z`$ axis. Note that $`r_D^{\left(0\right)}\left(\beta \right)r_D^{\left(0\right)}\left(\beta \right).`$
For $`\beta >0`$ the single axis is larger than the double axis and the cloud is a prolate spheroid, that is cigar-shaped. For $`\beta <0`$ the single axis is smaller than the double axis and the cloud is an oblate spheroid i.e. disk-shaped.
Note that the weak screening approximation restricts the possible values of $`\beta .`$ A reasonable assumption which stems from Eq.$`\left(\text{6}\right)`$ is that in the $`WES`$ regime the following relation must be fulfilled:
$$\frac{Z_1Z_2e^2}{r_D\left(\theta \right)kT}0.1$$
(10)
In solar conditions for the $`pp`$ reaction where the use of the $`WES`$ formalism is incontrovertible we have
$$\frac{e^2}{r_DkT}0.05$$
(11)
Therefore , for all orientations, the following inequality must hold :
$$0.8\beta 0.8$$
(12)
where we have disregarded the contribution of volume conservation which is always less than $`5\%`$. The deformation parameter can now take all the above values without violating the $`WES`$ condition $`\left(\text{6}\right)`$ , thus rendering the use of the deformed screening formalism legitimate.
If we take into account the nuclear potential $`V_N\left(r\right)`$ then the total potential of the reaction is
$$V(r;\theta ;\beta )=V_N\left(r\right)+V_D(r;\theta ;\beta )+\frac{\mathrm{}^2}{2\mu r^2}l\left(l+1\right)$$
(13)
where the centrifugal term is assumed to be independent of the orientation. This assumption is immaterial here as we will only consider very low-energy reactions where $`s`$-interactions dominate. In that case the orientation dependent cross section of the nuclear fusion reaction is given by:
$$\sigma (E;\theta ;\beta )=\frac{S\left(E\right)}{E}P(E;\theta ;\beta )$$
(14)
where
$$P(E;\theta ;\beta )=\mathrm{exp}\left[\frac{2\sqrt{2\mu }}{\mathrm{}}_R^{r_c(\theta ;\beta )}\sqrt{V_D(r;\theta ;\beta )E}๐r\right]$$
(15)
and $`r_c(\theta ;\beta )`$ is the classical turning point given by
$$V_D(r_c;\theta ;\beta )=E$$
(16)
The thermonuclear reaction which can be studied safely by means of the above formalism is the one that dominates the solar neutrino production namely: $`H^1(p,e^+\nu _e)H^2.`$ For that reaction, in the undeformed weak screening case, it turns out that in the region of the maximum energy production $`R=0.09R_\text{ }`$the Gamow peak is $`E_0=5.599keV`$ , the classical turning point is roughly $`r_c=0.01r_{D\text{ , }}`$while the $`WES`$ enhancement factor is to good approximation: $`f_0^{WES}=1.049.`$
Actually, non-linear screening corrections can be safely neglected. Along the whole profile of orientations and $`\beta `$ parameters the contribution of higher order terms to the shifts of the screening corrections and the Gamow peak has been found negligible. Hence, screening deformation effects can be simply represented by Salpeterโs $`WES`$ formula
$$f_D(\theta ;\beta )=\mathrm{exp}\left(\frac{e^2}{kTr_D\left(\theta \right)}\right)$$
(17)
where the DH radius $`r_D\left(\theta \right)`$ is now orientation dependent. Finally we obtain:
$$f_D(\theta ;\beta )=\left(f_0^{wes}\right)^{g^1(\theta ;\beta ,)}$$
(18)
where $`g(\theta ;\beta )`$ is the ratio $`r_D(\theta ;\beta )/r_D`$:
$$g(\theta ;\beta )=\left\{\frac{1}{2}_1^1\left[1+\beta \sqrt{\frac{5}{16\pi }}\left(3u^21\right)\right]^3๐u\right\}^{1/3}\left[1+\beta \sqrt{\frac{5}{16\pi }}\left(3\mathrm{cos}^2\theta 1\right)\right]$$
(19)
In Fig.1 the orientation dependent DH radius $`r_D(\theta ;\beta )`$, measured in units of $`r_D,`$ is plotted in polar coordinates with respect to the azimuthal angle $`\theta `$ and the deformation parameter $`\beta `$ . It is obvious that along the axis of symmetry of the cloud, $`z`$-axis $`\left(\theta =0\right),`$ a positive $`\beta `$parameter โstretches outโ the ionic cloud (a prolate spheroid shape) while a negative $`\beta `$ parameter โsucks inโ the cloud (an oblate spheroid shape). As one would expect, for both positive and negative parameters the larger the absolute value of $`\beta `$ the more pronounced the deformation.
Note that Fig. 1 actually represents the deformation factor $`g(\theta ;\beta )`$ while the corresponding classical turning point is $`r_c(\theta ;\beta )=261g(\theta ;\beta )fm.`$ For example for a $`\beta =0.8`$ deformation a proton cruising along the z-axis in the plasma with an energy equal to the Gamow peak $`E_0=5.56`$ $`keV`$ will come up against the potential wall at a distance roughly $`1.43`$ times further than in the undeformed case $`\left(r_c=261fm\right)`$, that is $`r_c\left(\theta =0,\beta =0.8\right)=373fm`$. On the other hand for a $`\beta =0.8`$ the classical turning point is reduced to $`r_c\left(\theta =0,\beta =0.8\right)=123fm.`$
The screened Coulomb potential $`V(r;\theta ;\beta )`$ can be visualized by means of Fig. 2 where we have plotted in polar coordinates the deformed shapes of the potential at a distance $`r=r_D`$ from the origin. At that distance the potential contours $`V_D(r_D;\theta ;\beta )`$ of Fig. 2 are only a function of the orientation. Hence, as a proton enters the ionic cloud of the Hydrogen atom on its way to fusion, according to the angle at which it enters the cloud it will experience a different (thicker or thinner) potential wall . The thicker the wall, the most improbable the reaction and of course the smaller the reaction rate.
The orientation dependent acceleration of the reaction is displayed in Fig. 3 where the screening factor is plotted with respect to the azimuthal angle and the deformation parameter for the $`pp`$ reaction at $`R=0.09R_{}`$. The reaction rate can be $`1.1`$ times faster if the proton enters a disk-shaped cloud $`\left(\beta =0.8\right)`$ at zero angle. On the other hand a much slower reaction is obtained for a cigar-shaped ionic cloud (almost no enhancement at all for $`\beta =0.8)`$.
The impact of screening deformation on the neutrino fluxes will be discussed in a more quantitative way in the section that follows.
## IV Solar neutrino fluxes.
In the solar region of maximum energy production $`\left(R/R_{}=0.09\right)`$ the thermal kinetic energy is $`kT=1.161keV`$ while for the $`pp`$ reaction the standard (undeformed) weak-screening factor is $`f_0^{WES}=1.049.`$ On the other hand for a $`\beta =0.4`$ deformation and an angle of impact of $`\theta =0`$ the screening factor is $`f_D=1.067.`$ This corresponds to an acceleration of the (undeformed) weakly screened $`pp`$ reaction by roughly $`1.7\%`$which in turn reflects on the cross-section factor given by $`S_D=f_DS.`$ As the principal source of energy in the Sun is the $`pp`$ reaction this acceleration would influence both the solar structure and the neutrino fluxes by reducing the central temperature and density in order to conserve luminosity. (An account of what happens in the sun if the cross-section factor $`S_{pp}`$ increases can be found in Ref. .).
In most solar evolution codes the $`pp`$ screening factor is evaluated by means of Salpeterโs formula which has been proved to be valid and accurate in standard conditions. In the deformed case the quantity $`f_D`$ should be used, instead. We can obtain an approximation of the uncertainties introduced due to the presence of such deformations by using the proportionality formulae which relate neutrino fluxes to screening factors. In order to isolate the $`pp`$ uncertainty, we will assume that except for the $`pp`$ reaction all the other neutrino-producing reactions remain unaffected by the deformations, thus obtaining a minimum of the total associated uncertainty.
For various solar fusion reactions the ratios of the deformed neutrino fluxes $`\mathrm{\Phi }^D`$ to the ones obtained in the $`WES`$ regime $`\mathrm{\Phi }^{WES},`$ are as follows:
$`H^1(p,e^+\nu _e)H^2`$
$$\left(\frac{\mathrm{\Phi }_{pp}^D}{\mathrm{\Phi }_{pp}^{WES}}\right)_{pp}=\left(\frac{f_D}{f_0^{WES}}\right)^{0.14}$$
(20)
$`H^1(pe^{},\nu _e)H^2`$
$$\left(\frac{\mathrm{\Phi }_{hep}^D}{\mathrm{\Phi }_{hep}^{WES}}\right)_{pp}=\left(\frac{f_D}{f_0^{WES}}\right)^{0.08}$$
(21)
$`Be^7(e^{},\nu _e)Li^7:`$
$$\left(\frac{\mathrm{\Phi }_{Be^7}^D}{\mathrm{\Phi }_{Be^7}^{WES}}\right)_{pp}=\left(\frac{f_D}{f_0^{WES}}\right)^{0.97}$$
(22)
$`Be^7(p,\gamma )B^8(e^+,\nu _e)B^8:`$
$$\frac{\mathrm{\Phi }_B^D}{\mathrm{\Phi }_B^{WES}}=\left(\frac{f_D}{f_0}\right)^{2.6}$$
(23)
$`N^{13}\left(e^+\nu _e\right)C^{13}`$ and $`O^{15}(e^+,\nu _e)N^{15}:`$
$$\frac{\mathrm{\Phi }_{N,O}^D}{\mathrm{\Phi }_{N,O}^{WES}}=\left(\frac{f_D}{f_0^{WES}}\right)^{22/8}$$
(24)
According to the above formulae, the presence of a screening deformation of $`\beta =0.4`$ and an angle of impact of $`\theta =0`$ can perturb the solar neutrino fluxes calculated in the SSM by at least $`4.5\%`$ for the $`N,O`$ and $`B^8`$ neutrinos and by less than $`1\%`$ for the less sensitive $`pp,hep,`$and $`B^7`$neutrinos. Such a source of deformation would also perturb the electron cloud around the heavier nuclei involved in neutrino production, thus modifying their screening factors which we arbitrarily considered constant here.
It is very tempting to investigate what happens if the deformations are stronger. For a collision along the z-axis with $`\beta =0.8`$ we obtain $`f_D=1.16`$ and the corresponding uncertainties are at least $`23\%`$ for the $`N,O`$ and $`B^8`$ neutrinos, $`9.2\%`$ for the $`Be^7`$ neutrinos and less than $`2\%`$ for the $`pp,hep`$ ones. Adding the effects of the screening factors of the other reactions, which have been disregarded so far, the uncertainties can be dramatic. It seems therefore that the presence of screening deformations in the sun can tune the predicted neutrino fluxes in order to reduce the observed deficit.
It is crucial to study what deviations from the SSM parameters such deformations can induce. As was previously noted conservation of luminosity implies a reduction in the central temperature $`T_c`$ as a result of an increase in $`S_{pp},`$ since $`L_{}S_{pp}T_c^8`$. This implies:
$`{\displaystyle \frac{T_c^D}{T_c^{WES}}}=\left({\displaystyle \frac{f_D}{f_0^{WES}}}\right)^{\frac{1}{8}}`$
Therefore, for the weaker (stronger) of the above considered deformations the $`\left(WES\right)`$ central temperature of the sun would have to decrease by $`0.2\%`$ $`\left(1.2\%\right)`$. On the other hand, as the ratio $`\rho /T^3`$ is approximately constant along the solar profile, the new central density would now be given by:
$`{\displaystyle \frac{\rho _c^D}{\rho _c^{WES}}}=\left({\displaystyle \frac{f_D}{f_0^{WES}}}\right)^{\frac{3}{8}}`$
that is roughly decreased by $`0.6\%\left(3.7\%\right)`$ with respect to the $`\left(WES\right)`$ assumption. Both values represent reasonable deviations from the $`\left(WES\right)`$ SSM considering that the $`\left(WES\right)`$ assumption itself causes a $`0.6\%`$ deviation from the $`\left(NOS\right)`$ central temperature and another $`1.8\%`$ deviation from the $`\left(NOS\right)`$ central density.
## V Screening deformation sources
Having established the general formalism for the investigation of screening deformations, a discussion of potential sources of such deformations is imperative. In a study of the effects of superstrong magnetic fields in the stellar plasma there was shown clearly that even at zero energies the screening electron cloud is deformed , in the sense that it becomes compressed perpendicular and parallel to the field direction. This deformation can dramatically accelerate hydrogen fusion reactions not only in neutrons stars but also in the laboratory as was shown recently,. It is very plausible therefore to assume that the presence of a superstrong magnetic field in the solar interior would cause similar effects.
From the present work it is obvious that if a substantial correction to the solar neutrino fluxes is to be made by means of screening deformations those have to be larger than $`\beta =4`$. To gain an idea of what kind of solar magnetic field would cause such deformation in the center of the sun we can use the results of Ref. where it shown clearly that any magnetic field weaker than the Intense Magnetic Field Regime:
$`B_{IMF}=4.7\times 10^9G`$
would have no effect on hydrogen fusion reactions. Hence, in order to decrease the solar neutrino fluxes by means of magnetically catalyzed thermonuclear fusion, the magnetic field required must be stronger than $`10^{10}G.`$
Admittedly such a superstrong field cannot be easily justified. After the disheartening result that the presence of a strong magnetic $`\left(10^9G\right)`$ in the solar interior increases the predicted neutrino fluxes (doubles the $`Cl^{37}`$ signal by increasing the pressure gradient in the sun) few investigators have looked into the matter. This is also due to some additional arguments which indicate the implausibility of a solar magnetic field larger than $`10^9G.`$ Such arguments include the limiting strength set by Chandrasekhar and Fermi, stability reasons and magnetic buoyancy, though there is a very interesting work which argues that a combination of a differential rotation and magnetic field can reduce the $`Cl^{37\text{ }}`$signal opposing the results of Ref. . Note that Helioseismology is another concern when considering such a strong magnetic field in the sun. Nevertheless, such a field, despite stability and buoyancy counter-arguments, can have been formed by the interstellar magnetic field which was frozen into the matter out of which the sun was formed, or there may be an unspecified mechanism of continuous generation.
It is now obvious that, regarding the solar neutrino puzzle, even if we accept the presence of a superstrong magnetic field the corrections induced are much smaller than the ones required to reconcile theory and experiment. Therefore, taking also into account other counter-arguments, it seems that magnetically induced screening deformations (magnetic catalysis) cannot possibly be the answer to the neutrino puzzle.
Another plausible source of screening deformations is the fact that in the solar center the average interionic spacing is $`a2.8\times 10^9cm,`$which is similar to the Debye radius. Therefore the spherical electron cloud assumption is not well justified. The cloud is more likely to assume an ellipsoidal distribution around the two reactants, like a fissioning nucleus . However, the argument that the incoming fusing nucleus will be carrying its own cloud, suggested in Ref. , and Ref. , thus increasing the deformations, doesnโt seem to have a significant effect as in that case the total of the nucleus-nucleus, nucleus-cloud, and cloud-cloud interactions would be:
$$V\left(r\right)=\frac{Z_1Z_2e^2}{r}\frac{3}{2}\frac{Z_1Z_2e^2}{r_D}$$
(25)
where the cloud-cloud interaction, in analogy to the recent results for the laboratory$`\text{[15]},`$ would be much smaller than the screening shift itself. Therefore the screening shift would have been at least $`50\%`$ larger than the $`\left(WES\right)`$ one used in standard DH theory thus accelerating for example the hydrogen fusion reaction by roughly $`2.5\%`$ with respect to the WES regime. This would result in an increase of the $`\left(NOS\right)`$ $`pp`$ neutrino fluxes of the order of $`1\%`$ and a similar decrease in the $`\left(NOS\right)`$ central solar temperature. It is therefore obvious that any screening deformations due to SSM cloud-cloud interactions in the solar plasma are negligible.
Finally a non-Maxwellian distribution of velocities, such as the flat Maxwellian, where the relative particle motion is frozen to $`0K`$ along a specific direction , will establish a preferential direction of motion for ions thus inducing electron cloud deformations. In fact any kind of deviation from statistical equilibrium is a source of screening deformations. For example it has been shown that a progressive depletion of the Maxwellian tail can yield a neutrino counting rate below 1 SNU but a physical cause for that distribution has not been found. Taking into account other counter-arguments to the existence of a non-Maxwellian distribution in the sun this non-standard source of deformation has to be deferred to a forthcoming article where the issue will be studied in detail.
## VI Conclusions
In this work we investigate the response of the proton-proton reaction to electron-ion screening deformations in the solar plasma. Those deformations are studied in the framework of the Debye-Hรผckel model and the results show that they can induce an orientation-dependent thermalized cross section which causes the solar neutrino fluxes to be orientation-dependent themselves. Therefore, in principle, screening deformations can influence the solar neutrino fluxes with reasonable deviations from the macroscopic values of the SSM.
Various potential deformation sources are discussed but none of them is found capable of inducing deformations strong enough to have a significant impact on the SSM neutrino fluxes. However, the existence of other, as yet unspecified, deformation sources cannot be ruled out. It seems therefore necessary to further investigate the possible presence of such sources which could cause a substantial degree of uncertainty to the solar neutrino fluxes.
Regarding the novelties of the present paper they can be summarized as follows:
The effects of deformations have received a minimal attention by only two authors (Ref. and Ref. ). They both concluded that the effect is small but none of them studied deformations in a non-standard solar model as they assumed SSM conditions from the very beginning of their calculations. The present paper studies screening deformations in a way completely independent of the source and solar model conditions. From now on, each time a non-standard source of deformation appears, it can possibly be connected to the deformation parameter $`\beta `$ introduced here, just as it happens in heavy ion fusion reactions. However, that task is admittedly not a trivial one. Moreover, the effects of magnetically catalyzed fusion on the SSM have been discussed here for the first time and eventually overruled as a solution to the solar neutrino puzzle. Finally, here for the first time there is shown that, if sufficiently strong, screening deformation can actually perturb considerably the SSM neutrino fluxes with reasonable changes in the central solar temperature and density. That result warrants further research into the origin and the effects of such screening deformation sources in the solar interior.
## VII ACKNOWLEDGMENTS
This work was financially supported by the Hellenic State Grants Foundation (IKY) under contract #135/2000. It was initiated at the Hellenic War College while its revised version was written at ECT during a nuclear physics fellowship. The author would like to thank the director of ECT Prof. R. Malfliet for his kind hospitality and support.
FIGURE CAPTIONS
Figure 1.
The variation of the orientation dependent DH radius $`r_D(\theta ;\beta )`$, measured in units of $`r_D,`$ with respect to the azimuthal angle $`\theta `$ and the deformation parameter $`\beta `$ in polar coordinates.
Figure 2.
The orientation dependent DH potential $`V_D(r;\theta ;\beta ),`$ measured in units of the screening term of Eq.$`\left(\text{11}\right)`$, at a distance $`r=r_D`$ with respect to the azimuthal angle $`\theta `$ and the deformation parameter $`\beta `$ in polar coordinates. The effect is calculated for the $`pp`$ reaction in the region of the maximum energy production $`R=0.09R_{}`$ where $`r_D=25719fm`$ and $`\frac{e^2}{r_DkT}0.056`$ $`keV.`$
Figure 3.
The variation of the orientation dependent screening factor $`f_D(\theta ;\beta )`$ with respect to the azimuthal angle $`\theta `$ and the deformation parameter $`\beta `$ in polar coordinates. The effect is calculated for the $`pp`$ reaction in the region of the maximum energy production $`R=0.09R_{}`$ where $`r_D=25719fm.`$ |
warning/0002/hep-th0002202.html | ar5iv | text | # 1 Introduction
## 1 Introduction
To explore the duality between large $`N`$ gauge theories and supergravity it is important to study cases with less supersymmetry and theories which are non-conformal . In this letter we study an example of such a theory. This theory is obtained on D1-branes at a conifold singularity. The conifold preserves $`1/4`$ of the supersymmetries of the full type IIB string theory. The theory on the D1-brane is a supersymmetric gauge theory in $`1+1`$ dimensions with $`4`$ supercharges.
We construct the supergravity solution of this configuration. We investigate the decoupling limit and find the domains of validity of supergravity description and the super-Yang-Mills description. We see that the infra-red limit of the super-Yang-Mills corresponds to matrix string theory in the background of the conifold.
Thus the infra-red limit of the super-Yang-Mills on a single D1-brane at the conifold should correspond to world sheet of a fundamental string propagating in the background of the conifold. We study the infra-red dynamics of the D1-brane gauge theory using the methods developed for $`(4,4)`$ theories by following and. The theory on the D1-brane at the conifold has $`(2,2)`$ supersymmetry. There is a 1-1 map from the moduli space of the Higgs branch of the D1-brane gauge theory to the conifold. The throat region of the Higgs branch corresponds to the singularity at the origin of the conifold. Though our theory has only $`(2,2)`$ supersymmetry most of the methods developed by to study $`(4,4)`$ theories go through. Their method involves using the Coulomb variables to give an effective description of the throat region of the Higgs branch. In theories with $`8`$ supercharges the Coulomb branch moduli space metric can receive correction only up to 1-loop. We do not have such facility for the case of $`(2,2)`$ theories. Nevertheless scale invariance constraints the metric to a form which enables us to extract the effective degrees of freedom. The matching of the R-symmetries in the ultra-violet and the infra-red works out just as in the case of $`(4,4)`$ theories.
Using these methods we are able to show that the throat region of the Higgs branch in the infra-red is captured by a $`๐ฉ=(2,2)`$ superconformal field theory consisting of a linear dilaton with background charge $`Q=\sqrt{2}`$ and a compact scalar. This agrees with the world sheet descriptions of strings at the conifold.
Then we consider $`N`$ D1-branes at the conifold. The infra-red limit of this theory corresponds to matrix string theory in the background of the conifold. The conifold provides a background in which the ultra-violet $`U(1)_R`$ symmetry is realized in the infra-red world sheet symmetry of the matrix string theory. We construct the leading interaction in the form of the twist operator. We note that the leading interaction is marginal.
The organization of this letter is as follows. In section 2 we study the decoupling limit of the N D1-branes at the conifold and investigate the domains of validity of the supergravity and the gauge theory. Section 3 analyses the infra-red dynamics of the gauge theory of a single D1-brane at the conifold. Section 4 formulates matrix string theory in the background of the conifold. We conclude in section 5 The appendix contains details of the supergravity solution.
## 2 Supergravity and the large $`N`$ limit of the D1-brane theory at the conifold
In this section we study the supergravity solution of $`N`$ D1-branes at a conifold singularity in the decoupling limit. We consider the configuration in which the D1-branes are aligned along the $`x^1`$ co-ordinate. The supergravity solution is given by (The verification of this solution is given in the Appendix.)
$`ds^2`$ $`=`$ $`f^{1/2}(dx_0^2+dx_1^2)+f^{1/2}[dr_1^2+r_1^2d\chi ^2+dr_2^2`$
$`+`$ $`{\displaystyle \frac{r_2^2}{9}}(d\psi +\mathrm{cos}\theta _1d\varphi _1+\mathrm{cos}\theta _2d\varphi _2)^2+{\displaystyle \frac{r_2^2}{6}}{\displaystyle \underset{i=1}{\overset{2}{}}}(d\theta _i^2+\mathrm{sin}^2\theta _id\varphi _i^2)]`$
$`e^{(\mathrm{\Phi }\mathrm{\Phi }_{\mathrm{}})}`$ $`=`$ $`f^{1/2}`$
$`B_{01}`$ $`=`$ $`{\displaystyle \frac{1}{2}}(f^11)`$
$`f`$ $`=`$ $`1+{\displaystyle \frac{NCg_s\alpha ^3}{(r_1^2+r_2^2)^3}}`$
The co-ordinates transverse to the D1-brane are $`r_1,\chi ,r_2,\psi ,\theta _i,\varphi _i`$. $`i`$ runs from 1 to 2. $`r_1`$, $`\chi `$ are polar co-ordinates of $`R^2`$. The remaining coordinates parameterize the conifold. The angular part of the conifold is parametrized by $`\psi ,\theta _i,\varphi _i`$, the radial part is parametrized by $`r_2`$. $`(\theta _1,\varphi _1)`$ and $`(\theta _2,\varphi _2)`$ parameterizes $`S^2\times S^2`$ as polar co-ordinates. $`\psi [0,4\pi ]`$ parameterizes the $`U(1)`$ fiber over $`S^2\times S^2`$. $`C`$ is fixed by charge quantization. It is given by $`C=864\pi ^2/(16+15\pi )`$. This configuration preserves $`4`$ supersymmetries out of the $`32`$ supersymmetries of Type IIB supergravity.
We study the decoupling limit of this system as in . It takes the form
$$U_1=\frac{r_1}{\alpha ^{}}=\text{fixed},U_2=\frac{r_2}{\alpha ^{}}=\text{fixed},g_{YM}^2=\frac{1}{2\pi }\frac{g_s}{\alpha ^{}}=\text{fixed},\alpha ^{}0$$
(2)
The metric and the dilaton of the supergravity solution in this limit is given by
$`{\displaystyle \frac{ds^2}{\alpha ^{}}}`$ $`=`$ $`{\displaystyle \frac{U^3}{g_{YM}\sqrt{2\pi NC}}}(dx_0^2+dx_1^2)+{\displaystyle \frac{g_{YM}\sqrt{2\pi NC}}{U^3}}[dU_1^2+U_1^2d\chi ^2+dU_2^2+`$
$`+`$ $`{\displaystyle \frac{U_2^2}{9}}(d\psi +\mathrm{cos}\theta _1d\varphi _1+\mathrm{cos}\theta _2d\varphi _2)^2+{\displaystyle \frac{U_2^2}{6}}{\displaystyle \underset{i=1}{\overset{2}{}}}(d\theta _i^2+\mathrm{sin}^2\theta _id\varphi _i^2)]`$
$`e^\mathrm{\Phi }`$ $`=`$ $`\sqrt{{\displaystyle \frac{8\pi ^3NCg_{YM}^6}{U^6}}}`$
where $`U^2=U_1^2+U_2^2`$.
Let us now discuss the domains of validity of various descriptions of this system. We will use the co-ordinate $`U`$ to set the energy scale at which we wish to look at the system. We will start with the high energies. For $`Ug_{YM}\sqrt{N}`$ the super-Yang-Mills perturbation theory can be trusted.
To find out when the supergravity solution given in (2) can be trusted let us estimate the curvature of the solution. Using the equation of motion an estimate of the curvature in string units is given by
$`\alpha ^{}R`$ $``$ $`g^{U_1U_1}_{U_1}\mathrm{\Phi }_{U_1}\mathrm{\Phi }+g^{U_2U_2}_{U_2}\mathrm{\Phi }_{U_2}\mathrm{\Phi }`$
$``$ $`{\displaystyle \frac{U}{g_{YM}\sqrt{N}}}`$
To trust supergravity the curvature should be small. Furthermore we need to ensure that the expansion in string coupling is valid. This is requires $`e^\mathrm{\Phi }`$ in (2) to be small. Thus the supergravity solution is valid for $`g_{YM}N^{1/6}Ug_{YM}\sqrt{N}`$. In addition to this condition, we must have $`UU_2^{2/3}(Ng_{YM}^2)^{1/6}`$. The latter condition arises from the fact that there is a curvature singularity at $`U_2=0`$ <sup>1</sup><sup>1</sup>1This point was raised by G. Horowitz and N. Itzhaki.
In the region $`Ug_{YM}N^{1/6}`$ we can use S-duality to study the solution. Performing S-duality on the near horizon solution in (2) we obtain
$`{\displaystyle \frac{ds^2}{\stackrel{~}{\alpha }^{}}}`$ $`=`$ $`{\displaystyle \frac{U^6}{g_{YM}^44\pi ^2NC}}(dx_0^2+dx_1^2)+{\displaystyle \frac{1}{2\pi g_{YM}^2}}[dU_1^2+U_1^2d\chi ^2+dU_2^2+`$
$`+`$ $`{\displaystyle \frac{U_2^2}{9}}(d\psi +\mathrm{cos}\theta _1d\varphi _1+\mathrm{cos}\theta _2d\varphi _2)^2+{\displaystyle \frac{U_2^2}{6}}{\displaystyle \underset{i=1}{\overset{2}{}}}(d\theta _i^2+\mathrm{sin}^2\theta _id\varphi _i^2)]`$
$`e^\mathrm{\Phi }`$ $`=`$ $`\sqrt{{\displaystyle \frac{8\pi ^3NCg_{YM}^6}{U^6}}}`$
where $`\stackrel{~}{\alpha }^{}=g_s\alpha ^{}`$. The supergravity solution in (2) is the near horizon geometry of $`N`$ fundamental strings at the conifold. An estimate of the curvature in string units can be performed as before, this gives the following
$$\stackrel{~}{\alpha }^{}R\frac{g_{YM}^2}{U^2}$$
(6)
This shows that this S-dual supergravity description is valid for $`Ug_{YM}`$. As before we also have an additional condition $`\alpha ^{}U_21`$ so as to avoid the curvature singularity at $`U_2=0`$. For small $`U`$ we have no supergravity description regardless of $`N`$.
If one is able to capture the infra-red behaviour of the super-Yang-Mills on the D1-branes then it is clear that one obtains a non-perturbative description of propagation of fundamental strings in the background of the conifold. This is analogous to case of D1-branes in flat space. The infra-red behaviour of super-Yang-Mills with gauge group $`U(N)`$ and $`16`$ supercharges provides a nonperturbative description of string theory . It is interesting to compare this also with the infra-red behaviour of the D1/D5 system. There the infra-red dynamics of the D1-branes captures the DLCQ of the little string theories in the Higgs branch. The Coulomb branch conformal field theory gives a non-perturbative description of srings in the background of Neveu-Schwarz 5-branes.
## 3 Infra-red dynamics of a single D1-brane
In this section we show that the gauge theory of a single D1-brane at the conifold in the infra-red flows to a superconformal field theory of a string in the background of the conifold.
### 3.1 The gauge theory
The gauge theory on the D1-brane at the conifold consists of a $`U(1)\times U(1)`$ gauge theory with $`(2,2)`$ supersymmetry in $`1+1`$ dimensions . The matter content of this theory consists of two sets of chiral multiplet, $`A_i`$, and $`B_i`$ with $`i=1,2`$. The $`A`$โs and $`B`$โs are charged as $`(1,1)`$ and $`(1,1)`$ respectively. The diagonal $`U(1)`$ decouples. It corresponds to the free $`U(1)`$ gauge multiplet on the single D1-brane. The $`2`$ scalars of this gauge multiplet represent motion of the D1-brane along $`r_1`$ and $`\chi `$. Under the relative $`U(1)`$ the $`A`$โs have charge $`+1`$ and the $`B`$โs have charge $`1`$. The D-term of the $`U(1)`$ vector multiplet is given by
$$D=|A_1|^2+|A_2|^2|B_1|^2|B_2|^2$$
(7)
We consider the case in which both the Fayet-Iliopoulos term and the theta term in the Lagrangian are set to zero. The conifold is realized as the moduli space of vacuum of the Higgs branch of this theory. Setting the D-term to zero and dividing by the gauge group $`U(1)`$ realizes the conifold. The complex coordinates of the conifold are given by
$$z_1=A_1B_1,z_2=A_2B_2,z_3=A_1B_2,z_4=A_2B_1$$
(8)
with
$$z_1z_2z_3z_4=0$$
(9)
Therefore the infrared theory is a superconformal field theory with the conifold as its target space. The central charge in the Higgs branch is given by counting the gauge invariant degrees of freedom. This is seen to be $`9`$. To isolate the description of the conifold at the singularity we describe the conifold as follows.
If point $`(a_1,a_2,a_3,a_4)`$ satisfies (9), and $`a_i0`$ then one can obtain another solution which is given by $`\sigma ^{1/2}(a_1,a_2,a_3,a_4)`$. Here $`\sigma `$ is a complex number. This particular scaling is chosen so that the complex polynomial describing the conifold in (9) is homogeneous of degree $`1`$. Therefore the conifold can be described by the space
$$\sigma \times (z_1z_2z_3z_4=0)/\sigma $$
(10)
where the space $`(z_1z_2z_3z_4)/\sigma `$ is $`2`$ complex dimensional hypersurface $`z_1z_2z_3z_4=0`$ in the $`3`$ complex dimensional weighted projective space $`WCP_{\frac{1}{2},\frac{1}{2},\frac{1}{2},\frac{1}{2}}^3`$.
The central charge of the $`๐ฉ=(2,2)`$ supeconformal field theory on the hypersurface in the weighted projective space is zero. Thus it does not contain any degrees of freedom. The entire central charge $`9`$ of the superconformal field theory of the conifold thus resides on the superconformal field theory on the one dimensional complex space parametrized by $`\sigma `$. It is clear this space is endowed with a nontrivial metric. To obtain this metric we examine the theory at the origin of the Higgs branch. At the origin of the Higgs branch there is a โsingularityโ as the Fayet-Iliopoulos term and the theta term are set to zero. This corresponds to the $`z_i=0`$ point on the conifold. In the infra-red we can describe the throat region of the Higgs branch using the Coulomb variables. This is because in the infra-red, the vector multiplet is an auxiliary field. The kinetic terms of the vector multiplet decouples as they are irrelevant. Then the Higgs fields can be written in terms of the vector multiplet using equations of motion. This method of describing the Higgs branch using the Coulomb variables was done in of $`N=(4,4)`$ gauge theories. We obtain the metric on the complex line parametrized by $`\sigma `$ by appealing to the description of the Higgs branch in terms of the Coulomb branch.
To proceed with the analysis we write down the Lagrangian for the relative $`U(1)`$. We follow the convention of but work with Euclidean world sheet metric. We set $`y^0=iy^2`$ in the formulae of .
$`L`$ $`=`$ $`L_{\mathrm{matter}}+L_{\mathrm{gauge}}`$ (11)
$`L_{\mathrm{matter}}`$ $`=`$ $`{\displaystyle d^2yD_\mu \overline{A}_iD^\mu A_i}+D_\mu \overline{B}_iD^\mu B_i+2\sigma \overline{\sigma }\overline{A}_iA_i+2\sigma \overline{\sigma }\overline{B}_iB_i`$
$``$ $`i\overline{\psi }_i^A(D_1+iD_2)\psi _i^Ai\overline{\psi }_i^B(D_1+iD_2)\psi _i^B`$
$``$ $`i\overline{\psi }_{+i}^A(D_1iD_2)\psi _{+i}^Ai\overline{\psi }_{+i}^B(D_1iD_2)\psi _{+i}^B`$
$`+`$ $`\sqrt{2}(\overline{\sigma }\overline{\psi }_{+i}^A\psi _i^A+\sigma \overline{\psi }_i^A\psi _{+i}^A)\sqrt{2}(\overline{\sigma }\overline{\psi }_{+i}^B\psi _i^B+\sigma \overline{\psi }_i^B\psi _{+i}^B)`$
$`+`$ $`i\sqrt{2}\overline{A}_i(\psi _i^A\lambda _+\psi _{+i}^A\lambda _{})+i\sqrt{2}A_i(\overline{\lambda }_{}\overline{\psi }_{+i}^A\overline{\lambda }_+\overline{\psi }_i^A)`$
$`+`$ $`i\sqrt{2}\overline{B}_i(\psi _i^B\lambda _+\psi _{+i}^B\lambda _{})+i\sqrt{2}B_i(\overline{\lambda }_{}\overline{\psi }_{+i}^B\overline{\lambda }_+\overline{\psi }_i^B)`$
$``$ $`D\overline{A}_iA_i+D\overline{B}_iB_i`$
$`L_{\mathrm{gauge}}`$ $`=`$ $`{\displaystyle \frac{1}{g_{YM}^2}}{\displaystyle }d^2y({\displaystyle \frac{1}{2}}F_{01}^2{\displaystyle \frac{1}{2}}D^2`$
$``$ $`i\overline{\lambda }_+(_1i_2)\lambda _+i\overline{\lambda }_{}(_1+i_2)\lambda _{}+_\mu \overline{\sigma }^\mu \sigma )`$
The superpotential on a single D1-brane is zero, therefore we have set the $`F`$-terms to zero. The vector multiplet consists of fields $`F_{01},\sigma ,\lambda _+,\lambda _{}`$. $`\sigma `$ is a complex scalar corresponding to the components of a $`4`$ dimensional gauge field along $`2`$ compact directions. The gauginos $`\lambda _+,\lambda _{}`$ are complex Weyl fermions. The super partners of the chiral multiplet $`A_i,B_i`$ are $`\psi _{i+,}^A,\psi _{i+,}^A`$ respectively. They are also complex Weyl fermions in $`2`$ dimensions.
The scaling dimension of the Yang-Mills coupling is $`1`$. Therefore in the infra-red the coupling tends to $`\mathrm{}`$. This ensures that the kinetic term for the vector multiplet $`L_{\mathrm{gauge}}`$ decouples in the infra-red. In the Higgs branch the scaling dimension of the scalars $`A_i,B_i`$ is zero, and its superpartners have dimension $`1/2`$. The scalars $`\sigma `$ and the gauge boson have scaling dimension $`1`$. Its superpartners have scaling dimension $`3/2`$. This is more evidence that the operators in $`L_{\mathrm{gauge}}`$ are irrelevant. Therefore in the infra-red, the Lagrangian is restricted to only $`L_{\mathrm{matter}}`$.
We now can integrate over the auxiliary vector multiplet in $`L_{\mathrm{matter}}`$. This forces the D-term to be set to zero and one obtains the Higgs branch as a $`๐ฉ=(2,2)`$ superconformal field theory over the conifold. However in order to describe the theory near the singularity we will follow the method of . Here the vector multiplets are regarded as composite operators on the Higgs branch. They are roughly given by
$$\sigma =\frac{1}{\sqrt{2}}\frac{\overline{\psi }_i^B\psi _{+i}^B\overline{\psi }_i^A\psi _{+i}^A}{\overline{A}_iA_i+\overline{B}_iB}$$
(12)
This amounts to integrating out the chiral multiplets in $`L_{\mathrm{matter}}`$. This was argued in to be valid at large values of $`\sigma `$. From (12) we see that is valid roughly for small values of the chiral multiplets $`A_i,B_i`$. Thus by large values of $`\sigma `$ we are probing the singularities of the Higgs branch. We will discuss the systematics of the expansion as we perform the 1-loop computation.
### 3.2 The one-loop calculation
The terms in the Lagrangian which are relevant for the 1-loop calculation are
$`L`$ $`=`$ $`{\displaystyle d^2y_\mu \overline{A}_i^\mu A_i}+_\mu \overline{B}_i^\mu B_i+2\sigma \overline{\sigma }\overline{A}_iA_i+2\sigma \overline{\sigma }\overline{B}_iB_i`$
$``$ $`i\overline{\psi }_i^A(_1+i_2)\psi _i^Ai\overline{\psi }_i^B(_1+i_2)\psi _i^B`$
$``$ $`i\overline{\psi }_{+i}^A(_1i_2)\psi _{+i}^Ai\overline{\psi }_{+i}^B(_1i_2)\psi _{+i}^B`$
$`+`$ $`\sqrt{2}(\overline{\sigma }\overline{\psi }_{+i}^A\psi _i^A+\sigma \overline{\psi }_i^A\psi _{+i}^A)\sqrt{2}(\overline{\sigma }\overline{\psi }_{+i}^B\psi _i^B+\sigma \overline{\psi }_i^B\psi _{+i}^B)`$
Integrating out the chiral multiplets to 1-loop gives the following terms in the action
$`S_{1\mathrm{loop}}`$ $`=`$ $`4\mathrm{ln}\left[\text{det}(\mathrm{}+2\overline{\sigma }\sigma )\right]+2\mathrm{ln}\left[\text{det}\left(\begin{array}{cc}i(_1+i_2)& +\sqrt{2}\sigma \\ +\sqrt{2}\overline{\sigma }& i(_1i_2)\end{array}\right)\right]`$ (16)
$`+`$ $`2\mathrm{ln}\left[\text{det}\left(\begin{array}{cc}i(_1+i_2)& \sqrt{2}\sigma \\ \sqrt{2}\overline{\sigma }& i(_1i_2)\end{array}\right)\right]`$ (19)
On simplification one gets
$`S_{1\mathrm{loop}}`$ $`=`$ $`\text{Tr}\mathrm{ln}\left[1+\left(\begin{array}{cc}0& \frac{1}{\mathrm{}+2\sigma \overline{\sigma }}[\sqrt{2}(_1+i_2)]\sigma \\ \frac{1}{\mathrm{}+2\sigma \overline{\sigma }}[\sqrt{2}(_1i_2)]\overline{\sigma }& 0\end{array}\right)\right]`$ (22)
$`+`$ $`\text{Tr}\mathrm{ln}\left[1+\left(\begin{array}{cc}0& \frac{1}{\mathrm{}+2\sigma \overline{\sigma }}[\sqrt{2}(_1+i_2)]\sigma \\ \frac{1}{\mathrm{}+2\sigma \overline{\sigma }}[i\sqrt{2}(_1i_2)]\overline{\sigma }& 0\end{array}\right)\right]`$ (25)
From (22) it is clear that the expansion parameter is $`\frac{d\sigma }{\sigma ^2}`$. We expect this expansion to be valid for $`d\sigma \sigma ^2`$. Thus, we obtain a good description of the Higgs branch in terms of the Coulomb variables for large $`\sigma `$, which according to (12) corresponds to regions near the singularity. After further simplification the leading order in the velocity expansion is given by
$`S_{1\mathrm{loop}}`$ $`=`$ $`4{\displaystyle d^2y(_1+i_2)\sigma (_1i_2)\overline{\sigma }\text{Tr}\left[\frac{1}{(\mathrm{}+2\sigma \overline{\sigma })^2}\right]}`$
$`=`$ $`4{\displaystyle d^2y(_1+i_2)\sigma (_1i_2)\overline{\sigma }\frac{d^2k}{4\pi ^2}\frac{1}{(k^2+2\sigma \overline{\sigma })^2}}`$
$`\sigma `$ is a complex scalar which represents the 2 coordinates of the coulomb branch. Writing these in polar coordinates and performing the integration we obtain the following metric on the moduli space.
$$ds^2=\frac{dr^2}{2\pi r^2}+\frac{d\theta ^2}{2\pi }$$
(27)
There is also a torsion given by
$$B_{r\theta }=\frac{1}{2\pi r}$$
(28)
The torsion is a pure gauge term. The space is topologically $`R\times S^1`$. This does not have any nontrivial closed 2-cycles. Thus there is no obstacle for gauging away the torsion. The gauge transformation is given by
$$B_{r\theta }=_r\mathrm{\Lambda }_\theta _\theta \mathrm{\Lambda }_r$$
(29)
Setting $`\mathrm{\Lambda }_r=0`$ gives $`\mathrm{\Lambda }_\theta =\mathrm{ln}r/2\pi `$. We have performed only a 1-loop calculation for the moduli space metric. The moduli space metric in (27) can be argued to be of the form given by the 1-loop result by scale invariance. $`\sigma `$ is a scalar with dimension $`1`$ in the ultra-violet gauge theory. The only scale invariant term one can write for the moduli space metric is $`d\sigma d\overline{\sigma }/\sigma \overline{\sigma }`$. We have determined the coefficient by the 1-loop calculation. A similar 1-loop calculation was done for the $`(4,4)`$ relevant for the D1/D5 system in . In $`(4,4)`$ gauge theories a non-renormalization theorem determines the moduli space metric by the 1-loop result . For the $`(2,2)`$ case there is no such theorem, but conformal invariance constraints the metric up to a numerical coefficient. To determine the infra-red degrees of freedom, we do not need this coefficient.
### 3.3 The infra-red degrees of freedom
It is clear from the moduli space metric in (27) that in the infra-red the dimension of the field $`r`$ is not determined. We define a scalar field $`\varphi `$ as $`r=e^{\varphi /2}`$. Then the kinetic term for $`\varphi `$ is just that of a free field. The field $`\varphi `$ can behave like a linear dilaton. The dimension of $`r`$ is specified only if one knows the background charge of the linear dilaton $`\varphi `$. We determine the background charge by requiring that the central charge of the infra-red $`๐ฉ=(2,2)`$ superconformal field theory to be $`9`$ which as we saw before was the central charge of the Higgs branch. We will justify this using R-symmetries in the section 3.4.
The bosonic fields capturing the infra-red dynamics are the linear dilaton $`\varphi `$, the compact scalar $`\theta `$. The radius of the compact scalar is 2. This is the value determined by the 1-loop calculation. The bosonic part of the action is given by
$$L=\frac{1}{8\pi }d^2y\sqrt{g}(g^{\mu \nu }_\mu \varphi _\nu \varphi QR\varphi +g^{\mu \nu }_\mu \theta _\nu \theta )$$
(30)
where $`R`$ is the world sheet curvature and $`Q`$ is the background charge of the linear dilaton. We have redefined $`\theta `$ so that we can use the $`\alpha ^{}=2`$ convention. The superpartners of $`\varphi `$ and $`\theta `$ are free fermions as the curvature of the moduli space $`R_{r\theta r\theta }`$ is zero.
Now we can evaluate the central charge. The total central charge of the infra-red superconformal field theory is
$$c=3/2+3Q^2+3/2$$
(31)
Demanding that the central charge be $`9`$ gives $`Q^2=2`$. The sign of the background charge is determined by requiring that the singularity is at the strong coupling region.This fixes the background charge to be $`Q=\sqrt{2}`$. At this point we mention that in the $`(4,4)`$ case the sign was determined by the fact that the conformal field theory of the Higgs branch and the conformal field theory of the Coulomb branch could be considered as two different subalgebras of the large $`๐ฉ=(4,4)`$ superconformal algebra.
We note that demanding that the central charge be $`9`$ does not determine the infra-red conformal field theory completely. The central charge of a conformal field theory is unchanged if it is deformed by a marginal operator. One obvious marginal deformation is the radius of the compact scalar. Thus we cannot determine the radius of the compact scalar in the infra-red limit. We determined the radius only using a 1-loop calculation. This certainly can change in the infra-red limit.
To be explicit we write down the holomorphic generators of the infra-red $`๐ฉ=(2,2)`$ superconformal field theory.
$`\overline{G}`$ $`=`$ $`\psi _z\overline{X}+_z\psi `$ (32)
$`G`$ $`=`$ $`\overline{\psi }_zX+_z\overline{\psi }`$
$`J_R`$ $`=`$ $`\psi \overline{\psi }+i\sqrt{2}_zX^2`$
$`T`$ $`=`$ $`_z\overline{X}_zX{\displaystyle \frac{1}{\sqrt{2}}}_z^2X^1+{\displaystyle \frac{1}{2}}(_z\psi \overline{\psi }+\psi _z\overline{\psi })`$
where
$$X=\frac{X^1+iX^2}{\sqrt{2}},\overline{X}=\frac{X^1iX^2}{\sqrt{2}},\psi =\frac{\psi ^1+i\psi ^2}{\sqrt{2}}\overline{\psi }=\frac{\psi ^1i\psi ^2}{\sqrt{2}}.$$
(33)
The field $`X^1`$ corresponds to the linear dilaton $`\varphi `$ and $`X^2`$ corresponds to the compact scalar $`\theta `$. The fermions $`\psi ^1,\psi ^2`$ are the superpartners of $`X^1`$ and $`X^2`$ respectively. There is a similar set of anti-holomorphic generators.
### 3.4 Comparison of the R-symmetries
We now compare the R-symmetries of the ultra-violet and the infra-red theory and show that our identification of the infra-red degrees of freedom is justified.
In the D1-brane gauge theory there is a $`U(1)_L\times U(1)_R`$ R-symmetry. The fields and their charges under $`U(1)_L`$ are as follows , $`(\psi _{+i}^A,\psi _{+i}^B,e^{i\theta },\lambda _{})`$ have charges $`(1,1,1,1)`$. The rest of the fields are uncharged under $`U(1)_L`$. The fields $`(\psi _i^A,\psi _i^B,e^{i\theta },\lambda _+)`$ have charges $`(1,1,1,1)`$ under $`U(1)_R`$. The rest of the fields are uncharged under $`U(1)_R`$. The infra-red behaviour of state localized far along the Higgs branch is approximately free. In a $`๐ฉ=(2,2)`$ free super conformal field theory with a flat metric the R-symmetry does not act on the bosons. The $`U(1)_L`$ and $`U(1)_R`$ does not act on the bosons of the chiral multiplet $`A_i`$ and $`B_i`$. Thus it is natural to identify the R-symmetry of the D1-brane gauge theory with that of the $`๐ฉ=(2,2)`$ superconformal field theory of the Higgs branch. As the $`๐ฉ=(2,2)`$ superconformal algebra relates the R-symmetry to the central charge we are justified in requiring that the central charge of the infra-red superconformal field theory of the Higgs branch be $`9`$.
The Coulomb branch is parametrized by the bosons $`|\sigma |,e^{i\theta }`$ As $`e^{i\theta }`$ is charged under $`U(1)_L\times U(1)_R`$ this cannot be the R-symmetry of the conformal field theory of the Coulomb branch. It must be as the theory flows to the infra-red on the Coulomb branch an R-symmetry is developed. This is similar to the case of $`(4,4)`$ gauge theories . In these theories the $`SU(2)_R`$ symmetry of the gauge theory in the ultra-violet is a candidate for the R-symmetry of the conformal field theory of the Coulomb branch. This symmetry is enhanced to $`SU(2)\times SU(2)`$ as the theory flows to the infra-red in the Coulomb branch.
Let us now examine from the infra-red degrees of freedom of the Higgs branch in the throat region whether the R-symmetry acts similar to $`U(1)_L\times U(1)_R`$. Let us focus on the holomorphic part, the anti-holomorphic part follows similarly. The field $`(e^{iX^2/\sqrt{2}},\psi )`$ are charged as $`(1,1)`$. This is what is expected under the identification of $`U(1)_L\times U(1)_R`$ as the R-symmetry of the conformal field theory of the Higgs branch.
### 3.5 Fundamental strings at the conifold
We have seen in section 2 that using arguments of that the infra-red theory of the D1-brane at the conifold should correspond to that of the world sheet of fundamental strings in the background of the conifold. We wish to compare the infra-red theory obtained with what is known about string propagation at the conifold.
String propagation at singularities have been studied recently in a series of works . It is seen from these works that string propagation at the conifold is described by a linear dilaton theory with back ground charge $`Q=\sqrt{2}`$ and a compact scalar. The effective degrees of freedom of the infra-red D1-brane gauge theory precisely matches with this.
String propagation at the resolved conifold $`z_1^2+z_2^2+z_3^2+z_4^2=\mu `$, has been discussed in . In these series of works it was argued that the world sheet theory was given by a $`๐ฉ=(2,2)`$ $`SL(2,R)/U(1)`$ Kazama-Suzuki model at level 3. The $`SL(2,R)/U(1)`$ model away from the origin consists of a linear dilaton with $`Q=\sqrt{2}`$ and a compact scalar at the self dual radius . This also agrees with the infra-red degrees of freedom of the D1-brane gauge theory.
## 4 Matrix String theory in a Conifold background
In this section we use the Lagrangian of $`N`$ D1-branes in a conifold background to formulate matrix string theory in this background. To formulate matrix string theory we need that the spatial coordinate of the two dimensional Yang-Mills to be compact . To the knowledge of the author, matrix string theory has not been formulated in a background with 8 supersymmetries.
The Lagrangian of $`N`$ D1-branes at a conifold is that constructed in dimensionally reduced to two dimensions. It consists of a $`U(N)\times U(N)^{}`$ gauge theory with $`(2,2)`$ supersymmetry. We will use $`d=4`$ $`N=1`$ supersymmetry nomenclature to classify our fields. There are two gauge multiplets corresponding to the two gauge groups. The bosonic fields of the gauge multiplet consist of two bosons transforming in the adjoint of the corresponding gauge group. The fermions of the gauge multiplet are complex Weyl fermions in two dimensions. They are spinors of $`SO(2)`$ the symmetry in the transverse directions parametrized by $`r_1`$ and $`\chi `$. This can be seen from the fact that they arise from dimensional reduction of complex Weyl fermions of $`4`$ dimension. They transform in the adjoint representation of the gauge group. We list the fields of the gauge multiplets below.
$`\text{Bosons}A_\mu ,A_\mu ^{},X_a,X_a^{}`$ (34)
$`\text{Fermions}\lambda _+,\lambda _{},\lambda _+^{}\lambda _{}^{}`$
where $`a=1,2`$. The primes over the field variables indicate that they transform under the gauge group $`U(N)^{}`$. There are 4 chiral multiplets arranged in two sets $`A_i`$ and $`B_i`$, $`i=1,2`$. The $`A_i`$ transform as $`U(N)\times \overline{U(N)^{}}`$, while the $`B_i`$ transform as $`\overline{U(N)}\times U(N)^{}`$. We use the the capital $`A`$โs and $`B`$โs to indicate the superfields as well as the bosonic component. The fields of the chiral multiplets are
$`\text{Superfield}A_iA_i,\psi _i^A\psi _{+i}^A`$ (35)
$`\text{Superfield}B_iB_i,\psi _i^B\psi _{+i}^B`$
The Lagrangian has a superpotential given by
$$W=\frac{1}{2}ฯต^{ij}ฯต^{kl}\text{Tr}A_iB_kA_jB_l$$
(36)
We choose the gauge coupling of the two gauge groups to be identical. The bosonic potential is given by
$`U`$ $`=`$ $`g_{YM}^2{\displaystyle \underset{i}{}}\left|{\displaystyle \frac{W}{A_i}}\right|^2+g_{YM}^2{\displaystyle \underset{i}{}}\left|{\displaystyle \frac{W}{B_i}}\right|^2+{\displaystyle \frac{1}{2g_{YM}^2}}\text{Tr}D^2+{\displaystyle \frac{1}{2g_{YM}^2}}\text{Tr}D^2`$
$`+`$ $`A_i^{}(X_1^2+X_2^2)A_i++A_i^{}(X_1^2+X_2^2)A_i+B_i^{}(X_1^2+X_2^2)B_i+B_i^{}(X_1^2+X_2^2)B_i`$
$`+`$ $`{\displaystyle \frac{1}{2\stackrel{~}{g}_{YM}^2}}[X_1,X_2]^2+{\displaystyle \frac{1}{2g_{YM}^2}}[X_1^{},X_2^{}]^2`$
Perturbative Type IIA string theory is realized out of the matrix description in the $`g_{YM}\mathrm{}`$. From the superpotential we see that this limit selects out a vacuum. One such vacuum is in which all the $`A`$โs and $`B`$โs are diagonal. We analyze the theory around this vacuum. The gauge group is broken down to $`U(1)^N`$ in this vacuum. Each of the $`U(1)`$ corresponds to the center of mass $`U(1)`$ for the single D1-brane considered in section 3 The Weyl group of $`U(N)\times U(N)^{}`$ acts on the vacuum as
$`ASAS^{}`$ (38)
$`BS^{}BS^{}`$
We see that the Weyl group of $`U(N)\times U(N)^{}`$ transform $`A`$ and $`B`$ to values which are diagonal with the entries permuted only if $`S`$ and $`S^{}`$ are the same element of the Weyl group. Such transformations takes one vacuum to another. The gauge invariant vacuum is given by identifying these. The $`D`$ terms for the relative $`U(1)`$ for each of the D1-brane reduce to
$`|A_1^m|^2+|A_2^m|^2|B_1^m|^2|B_2^m|^2=0`$ (39)
where $`m=1,\mathrm{}N`$. $`N`$ copies of the conifold is realized as the moduli space of vacuum. The complex coordinates of the $`N`$ copies of the conifold are given by
$$z_1^m=A_1^mB_1^m,z_2^m=A_2^mB_2^m,z_3^m=A_1^mB_2^m,z_4=A_2^mB_1^m$$
(40)
Thus the conformal field theory that describes the infra-red limit of this gauge theory is a sigma model on the orbifold target space
$$\frac{(R^2\times C)^N}{S(N)}$$
(41)
where $`S(N)`$ is the symmetric group and $`C`$ stands for the conifold. The $`R^2`$ refers to the transverse directions parametrized by $`r_1`$ and $`\chi `$. These arise from the values of scalars $`X_a^m`$ corresponding to the diagonal $`U(1)`$โs. The question of finding backgrounds for matrix string theory where an $`U(1)_R`$ is present in the ultraviolet gauge theory which appears as the world sheet $`U(1)_R`$ was raised recently in . In the ultraviolet of this matrix theory, there is a superpotential. Therefore the chiral multiplets are charged under the $`U(1)_R`$ in the ultrviolet. Thus this R-symmetry cannot be the $`U(1)_R`$ of the world sheet theory in the infrared <sup>2</sup><sup>2</sup>2The author thanks E. Silverstein and Y. S. Song for pointing out that for this matrix theory too, the question in is unresolved, correcting the erroneous conclusion in the earlier draft..
Using the results of section 3, the conformal field theory on the conifold near the singularity that captures the infra-red limit of the gauge theory is given by the orbifold
$$\frac{(R^2\times R_\varphi \times S^1)^N}{S(N)}$$
(42)
where $`R_\varphi `$ stands for the linear dilaton with background charge $`Q=\sqrt{2}`$ and $`S^1`$ refers to the compact scalar. At this point let us examine the domain of validity of the super conformal field theory on the orbifold (42). The super conformal field theory on the orbifold (42) is valid when $`g_{YM}\mathrm{}`$ and $`z0`$. Here $`z`$ stands for all the co-ordinates in (40). Now, it is important that there exists a domain in these limits that the mass of the off-diagonal chiral multiplets can be neglected. The mass of the off-diagonal chiral multiplets roughly goes as $`g_{YM}^2z^2`$. Thus, the super conformal field theory on the orbifold is valid in the limits $`g_{YM}\mathrm{}`$, $`z0`$ and $`g_{YM}^2z^2\mathrm{}`$.
We would like to construct the leading interaction vertex represented by the twist operator corresponding to the $`Z_2`$ conjugacy class of the permutation group. For this we focus on the orbifold
$$\frac{(R^2\times R_\varphi \times S^1)^2}{Z_2}$$
(43)
Going over to center of mass and relative coordinate, conformal field theory on the following target space is realized
$$(R^2\times R_\varphi \times S^1)\times \frac{(R^3\times S^1)}{Z_2}$$
(44)
Here the linear dilaton $`R_\varphi `$ has background charge $`Q=\sqrt{2}\times \sqrt{2}`$. From the fact that the orbifold is a $`(R^3\times S^1)/Z_2`$, the interaction vertex is represented by the twist operator that is marginal. This is unlike the case of D1-branes in flat space where the leading interaction was irrelevant . The fact that there is a marginal operator in the infrared in this case is puzzling. This results perhaps from the fact that the theory is strongly coupled at the singularity. The issue of whether this marginal operator is turned on or not in the infrared theory is important. If the operator is turned on there is no weak coupling limit and persumably the infrared behaviour does not look like a perturbative matrix string theory. It is important to resolve this issue further.
## 5 Conclusions
We have used methods developed for the analysis of infra-red dynamics of $`(4,4)`$ gauge theories to study the infra-red dynamics of the $`(2,2)`$ gauge theory on a D1-brane at the conifold. We showed that the infra-red dynamics is captured by a $`๐ฉ=(2,2)`$ superconformal field theory consisting of a linear dilaton with background charge $`Q=\sqrt{2}`$ and a compact scalar. This agreed with the expectation that the infra-red theory should correspond to that of a fundamental string at the conifold. We mention that these methods can be used to analyze infra-red dynamics of $`(2,2)`$ theories with one dimensional Coulomb branch.
The Lagrangian of $`N`$ D1-branes at the conifold was used to formulate matrix string theory on this background. We note that the leading interaction represented by the twist operator in this case is marginal unlike the case of D1-branes in flat space.
Acknowledgements
The author acknowledges useful discussions with Nissan Itzhaki, Kiril Krasnov, Juan Maldacena and Ashoke Sen. He thanks Rajesh Gopakumar, Gary Horowitz and Joe Polchinksi for a careful reading of the manuscript, useful comments and discussions. He also acknowledges Avinash Dhar, Gautam Mandal and Spenta Wadia for a careful reading of the manuscript, comments and advice. The work of the author is supported by NSF grant PHY97-22022.
## Appendix A The supergravity solution
We follow the following strategy to verify the supergravity solution in (2). We convert the solution given in (2) into the Einstein metric and into the conventions of . In these conventions we make the following ansatz for the supergravity solution.
$`ds^2`$ $`=`$ $`f^{3/4}(dx_0^2+dx_1^2)+f^{1/4}[dr_1^2+r_1^2d\chi ^2+dr_2^2`$
$`+`$ $`{\displaystyle \frac{r_2^2}{9}}(d\psi +\mathrm{cos}\theta _1d\varphi _1+\mathrm{cos}\theta _2d\varphi _2)^2+{\displaystyle \frac{r_2^2}{6}}{\displaystyle \underset{i=1}{\overset{2}{}}}(d\theta _i^2+\mathrm{sin}^2\theta _id\varphi _i^2)]`$
$`e^\mathrm{\Phi }`$ $`=`$ $`f^{1/2}`$
$`B_{01}`$ $`=`$ $`f^1`$
where $`f`$ is an unknown function. We then show that the equations of motion of the various field just reduce to the Laplacian for the function $`f`$ in the coordinates transverse to the D1-brane.
We first substitute this ansatz in the dilaton equation. The dilaton eqation in the convention of is given by
$`_{MN}(\sqrt{g}g^{MN}_N\mathrm{\Phi }){\displaystyle \frac{1}{24}}\sqrt{g}e^\mathrm{\Phi }H^2`$ (46)
$`={\displaystyle \frac{\kappa ^2T_2}{2}}{\displaystyle d^2\xi \sqrt{\gamma }\gamma ^{ij}_iX^M_jX^Ng_{MN}e^{\mathrm{\Phi }/2}\delta ^{10}(xX)}`$
Substituting the values of the field given in (A) in the static gauge we find that the dilaton equation (46) reduces to
$$\frac{1}{2f}_{r_1}(K_{r_1}f)+\frac{1}{2f}_{r_2}(K_{r_2}f)=T_2\kappa ^2K\frac{108}{(4\pi )^3f}\frac{\delta (r_1)}{r_1}\frac{\delta (r_2)}{r_2^5}$$
(47)
where
$$K=\frac{r_1r_2^5}{108}\mathrm{sin}\theta _1\mathrm{sin}\theta _2$$
(48)
(47) is the Laplacian for the transverse space. It is clear that the solution of this is given by
$$f=A+\frac{B}{(r_1^2+r_2^2)^3}$$
(49)
where $`A`$ and $`B`$ are constants.
We now verify that the antisymmetric tensor equation of motion and the Einstein equation also reduces to the Laplacian in transverse space (47). The antisymmetric tensor equation is
$$_M(\sqrt{g}e^\mathrm{\Phi }H^{MNO})=2\kappa ^2T_2d^2\xi ฯต^{i_1i_2}_{i_1}X^N_{i_2}X^O\delta ^{10}(xX)$$
(50)
Substituting the ansatz in (A) the equation for the anisymmetric tensor reduces to
$$_{r_1}(K_{r_1}f)+_{r_2}(K_{r_2}f)=2T_2\kappa ^2K\frac{108}{(4\pi )^3}\frac{\delta (r_1)}{r_1}\frac{\delta (r_2)}{r_2^5}$$
(51)
(47) is identical to the above equation. Therefore the same solution (49) satisfies it.
The Einstein equation is given in (3.15) of . After some tedious but straight forward calculations the components of Einstein equation along the D1-brane reduce to
$$\frac{1}{2f^{7/4}}_{r_1}(K_{r_1}f)+\frac{1}{2f^{7/4}}_{r_2}(K_{r_2}f)=T_2\kappa ^2K\frac{108}{(4\pi )^3f^{7/4}}\frac{\delta (r_1)}{r_1}\frac{\delta (r_2)}{r_2^5}$$
(52)
This is the same as (47). The remaining components of the Einstein equation reduce to identities for the ansatz in (A). The brane field equations are also automatically satisfied.
To fix the constants $`A`$ and $`B`$ we go back to the string metric and into the conventions of . $`A=1`$ as we require that at infinity the metric reduce to that of the conifold. The constant $`B`$ by demanding that the net charge of the D1-branes is quantized. This gives that
$$B=Ng_s\alpha ^3\frac{864\pi ^2}{16+15\pi }$$
(53)
where $`N`$ is an integer. |
warning/0002/cond-mat0002450.html | ar5iv | text | # Ferromagnetism in III-V and II-VI semiconductor structures
## I Introduction
Because of complementary properties of semiconductor and ferromagnetic material systems, a growing effort is directed toward studies of semiconductor-magnetic nanostructures. Applications in sensors and memories as well as for computing using electron spins can be envisaged . The hybrid nanostructures, in which both electric and magnetic field are spatially modulated, are usually fabricated by patterning of a ferromagnetic metal on the top of a semiconductor or by incorporation of ferromagnetic clusters into a semiconductor matrice . In such devices, the stray fields can control charge and spin dynamics in the semiconductor. At the same time, spin-polarized electrons from the metal could be injected into the semiconductor. The efficiency of such a process appears, however, to be prohibitively low \[5-7\].
Already the early studies of Cr spinels and rock-salt Eu and Mn-based chalcogenides led to the observation of a number of outstanding phenomena associated with the interplay between ferromagnetic cooperative phenomena and semiconducting properties. The discovery of ferromagnetism in zinc-blende III-V and II-VI Mn-based compounds allows one to explore physics of previously not available combinations of quantum structures and magnetism in semiconductors. For instance, a possibility of changing the magnetic phase by light in (In,Mn)As/(Al,Ga)Sb and (Cd,Mn)Te/(Cd,Zn,Mg)Te heterostructures was put into the evidence. The injection of spin-polarized carriers from (Ga,Mn)As to a (In,Ga)As quantum well in the absence of an external magnetic field was demonstrated, too . It is then important to understand the ferromagnetism in these semiconductors, and to ask whether the Curie temperatures $`T_C`$ can be raised to above 300 K from the present 110 K observed for Ga<sub>0.947</sub>Mn<sub>0.053</sub>As . In this paper, we outline briefly the main ingredients of a model put recently forward to describe quantitatively the hole-mediated ferromagnetism in tetrahedrally coordinated semiconductors . We also list a number of counterintuitive experimental findings, which are explained by the model. Finally, we present the relevant chemical trends and discuss various suggestions concerning the design of novel ferromagnetic semiconductor systems. The recent comprehensive reviews present many other aspects of III-V , II-VI as well as of IV-VI magnetic semiconductors, which are not discussed here.
## II Origin of ferromagnetism
Since we aim at quantitative description of experimental findings, the proposed theoretical approach makes use of empirical facts and parameters wherever possible. In this section, we discuss those effects, which are regarded as crucial in determining the magnitude of ferromagnetic couplings in p-type magnetic semiconductors .
### II.1 Charge and spin state of Mn ions
We consider tetrahedrally coordinated semiconductors, in which the magnetic ion Mn occupies the cation sublattice, as found by extended x-ray absorption fine structure (EXAFS) studies in the case of Cd<sub>1-x</sub>Mn<sub>x</sub>Te and Ga<sub>1-x</sub>Mn<sub>x</sub>As . The Mn provides a localized spin and, in the case of III-V semiconductors, acts as an acceptor. These Mn acceptors compensate the deep antisite donors commonly present in GaAs grown by low-temperature molecular beam epitaxy, and produce a p-type conduction with metallic resistance for the Mn concentration $`x`$ in the range $`0.04x0.06`$ \[18,24-26\]. According to optical studies, Mn in GaAs forms an acceptor center characterized by a moderate binding energy $`E_a=110`$ meV, and a small magnitude of the energy difference between the triplet and singlet state of the bound hole $`\mathrm{\Delta }ฯต=8\pm 3`$ meV. This small value demonstrates that the hole introduced by the divalent Mn in GaAs does not reside on the d shell or forms a Zhang-Rice-like singlet , but occupies an effective mass Bohr orbit . Thus, due to a large intra-site correlation energy $`U`$, (Ga,Mn)As can be classified as a charge-transfer insulator, a conclusion consistent with photoemission spectroscopy . At the same time, the p-d hybridization results in a spin-dependent coupling between the holes and the Mn ions, $`H_{pd}=\beta N_o๐๐บ`$. Here $`\beta `$ is the p-d exchange integral and $`N_o`$ is the concentration of the cation sites. The analysis of both photoemission data and magnitude of $`\mathrm{\Delta }ฯต`$ leads to the exchange energy $`\beta N_o1`$ eV. Similar values of $`\beta N_o`$ are observed in II-VI diluted magnetic semiconductors with comparable lattice constants . This confirms Harrisonโs suggestion that the hybridization matrix elements depend primarily on the intra-atomic distance . According to the model in question, the magnetic electrons remain localized at the magnetic ion, so that they do not contribute to charge transport. This precludes Zenerโs double exchange as the mechanism leading to ferromagnetic correlation between the distant Mn spins. At the same time, for some combinations of transition metals and hosts, the โchemicalโ attractive potential introduced by the magnetic ion can be strong enough to bind the hole on the local orbit . In an intermediate regime, the probability of finding the hole around the magnetic ion is enhanced, which result in the apparent increase of $`|\beta N_o|`$ with decreasing $`x`$ .
In addition to the carrier-spin interaction, the p-d hybridization leads to the superexchange, a short-range antiferromagnetic coupling between the Mn spins. In order to take the influence of this interaction into account, it is convenient to parameterize the dependence of magnetization on the magnetic field in the absence of the carriers, $`M_o(H)`$, by the Brillouin function, in which two empirical parameters, the effective spin concentration $`x_{eff}N_o<xN_o`$ and temperature $`T_{eff}>T`$, take the presence of the superexchange interactions into account . The dependencies $`x_{eff}(x)`$ and $`T_{AF}(x)T_{eff}(x)T`$ are known for (Zn,Mn)Te, while $`xx_{eff}`$ and $`T_{AF}0`$ for (Ga,Mn)As, as explained below.
### II.2 Electronic states near metal-insulator transition
Because of ionized impurity and magnetic scattering, the effective mass holes introduced by Mn in III-V compounds or by acceptors such as N or P in the case of II-VI DMS, are at the localization boundary. It is, therefore, important to discuss the effect of localization on the onset of ferromagnetism. The two-fluid model constitutes the established description of electronic states in the vicinity of the Anderson-Mott metal-insulator transition (MIT) in doped semiconductors. According to that model, the conversion of itinerant electrons into singly occupied impurity states with increasing disorder occurs gradually, and begins already on the metal side of the MIT. This leads to a disorder-driven static phase separation into two types of regions: one populated by electrons in extended states, and another containing singly occupied impurity-like states. The latter controls the magnetic response of doped non-magnetic semiconductors and gives rise to the presence of BMPs on both sides of the MIT in magnetic semiconductors . Actually, the formation of BMPs shifts the MIT towards the higher carrier concentrations . On crossing the MIT, the extended states become localized. However, according to the scaling theory of the MIT, their localization radius $`\xi `$ decreases rather gradually from infinity at the MIT towards the Bohr radius deep in the insulator phase, so that on a length scale smaller than $`\xi `$ the wave function retains an extended character. Such weakly localized states are thought to determine the static longitudinal and Hall conductivities of doped semiconductors.
The central suggestion of the recent model is that the holes in the extended or weakly localized states mediate the long-range interactions between the localized spins on both sides of the MIT in the III-V and II-VI magnetic semiconductors. As will be discussed below, $`T_C`$ is proportional to the spin susceptibility of the carrier gas $`\chi _s`$ which, in turn, is proportional to the thermodynamic spin density-of-states $`\rho _s`$. Like other thermodynamic quantities, $`\rho _s`$ does not exhibit any critical behavior at the MIT. The quantitative renormalization of $`\rho _s`$ by disorder will depend on its actual form, for instance, on the degree of compensation. The enhancement of $`\rho _s`$ by the carrier-carrier interactions can be described by the Fermi-liquid parameter $`A_F`$, $`\rho _sA_F\rho _s`$. The value of $`A_F=1.2`$, as evaluated by the local-spin-density approximation for the relevant hole concentrations, has been adopted for the computations.
The participation of the same set of holes in both charge transport and the ferromagnetic interactions is shown, in (Ga,Mn)As and in (Zn,Mn)Te , by the agreement between the temperature and field dependencies of the magnetization deduced from the extraordinary Hall effect, $`M_H`$, and from direct magnetization measurements, $`M_D`$, particularly in the vicinity of $`T_C`$. However, below $`T_C`$ and in the magnetic fields greater than the coercive force, while $`M_H`$ saturates (as in standard ferromagnets), $`M_D`$ continues to rise with the magnetic field . This increase is assigned to the BMPs, which interact weakly with the ferromagnetic liquid. To gain the Coulomb energy, the BMPs are preferentially formed around close pairs of ionized acceptors. In the case of p-(Ga,Mn)As this leads to a local ferromagnetic alignment of neighbor Mn d<sup>5</sup> negative ions , so that $`xx_{eff}`$ and $`T_{AF}0`$. By contrast, BMPs in p-(Zn,Mn)Te are not preferentially formed around Mn pairs and encompass more Mn spins for a given $`x`$, as the small binding energy $`E_a=54`$ meV corresponds to a relatively large localization radius. The presence of a competition between the ferromagnetic and antiferromagnetic interactions in II-VI compounds, and its absence in III-V materials, constitutes the important difference between those two families of magnetic semiconductors.
### II.3 The structure of the valence band
The meaningful model of the hole-mediated ferromagnetism has to take into account the complex nature of the valence band in semiconductors resulting from $`kp`$ and spin-orbit interactions. Therefore, the hole dispersion and wave functions are computed by diagonalizing the 6x6 Kohn-Luttinger matrix together with the recalculated p-d exchange contribution . The model is developed for zinc-blende and wurzite semiconductors, allows for warping, quantizing magnetic fields, strain, and arbitrary orientations of $`๐ด`$. Because of the spin-orbit interaction, the exchange splitting of the valence band depends on the relative orientation of the magnetization and the hole wave vector. This mixing of orbital and spin degrees of freedom accounts for substantial differences between the effects of the exchange interaction upon properties of the electron and hole liquids. The tabulated values of the effective-mass constants \[49-52\] are taken as input parameters. According to interband magnetoptics and photoemission studies $`\beta N_o=1.1\pm 0.1`$ and $`1.2\pm 0.2`$ eV for (Zn,Mn)Te and (Ga,Mn)As, respectively.
### II.4 Zener model of carrier-mediated ferromagnetism
Zener first proposed the model of ferromagnetism driven by the exchange interaction between carriers and localized spins. However, this model was later abandoned, as neither the itinerant character of the magnetic electrons nor the quantum (Friedel) oscillations of the electron spin-polarization around the localized spins were taken into account, both of these are now established to be critical ingredients for the theory of magnetic metals. In particular, a resulting competition between ferromagnetic and antiferromagnetic interactions leads rather to a spin-glass than to a ferromagnetic ground state. In the case of semiconductors, however, the mean distance between the carriers is usually much greater than that between the spins. Under such conditions, the exchange interaction mediated by the carriers is ferromagnetic for most of the spin pairs, which reduces the tendency towards spin-glass freezing. Actually, for a random distribution of the localized spins, the mean-field value of the Curie temperature $`T_C`$ deduced from the Zener model is equal to that obtained from the Ruderman, Kittel, Kasuya, and Yosida (RKKY) approach, in which the presence of the Friedel oscillations is explicitly taken into account .
The starting point of the model is the determination how the Ginzburg-Landau free-energy functional $`F`$ depends on the magnetization $`M`$ of the localized spins. As mentioned above, the hole contribution to $`F`$, $`F_c[M]`$ is computed by diagonalizing the 6x6 Kohn-Luttinger matrix together with the p-d exchange contribution, and by the subsequent computation of the partition function $`Z`$. This model takes the effects of the spin-orbit interaction into account, a task difficult technically within the RKKY approach, as the spin-orbit coupling leads to non-scalar terms in the spin-spin Hamiltonian. Moreover, the indirect exchange associated with the virtual spin excitations between the valence bands, the Bloembergen -Rowland mechanisms , is automatically included.
The remaining part of the free energy functional, that of the localized spins, is given by
$$F_S[M]=_0^M๐M_oH(M_o),(1)$$
(1)
where $`H(M_o)`$ is the inverse function of the experimental dependence of the magnetization on the magnetic field H in the absence of the carriers. By minimizing $`F=F_c[M]+F_S[M]`$ with respect to $`M`$ at given $`T`$, $`H`$, and hole concentration $`p`$, one obtains $`M(T,H)`$ within the mean-field approximation, an approach that is quantitatively valid for long-range exchange interactions.
### II.5 Curie temperature
As described above, $`T_C`$ can be computed by minimizing the free energy, and without referring to the explicit form of the carrier periodic wave functions $`u_{i๐}`$. Since near $`T_C`$ the relevant magnetization $`M`$ becomes small, the carrier free energy, and thus $`T_C`$, can also be determined from the linear response theory. The resulting $`T_C`$ assumes the general form
$$T_C=x_{eff}N_oS(S+1)\beta ^2A_F\rho _s(0,T_C)/12k_BT_{AF}.$$
(2)
In the absence of disorder
$$\rho _s(๐,T)=8\underset{ij๐}{}\frac{|u_{i๐}|s_z|u_{j๐+๐}|^2f_i(๐)[1f_j(๐+๐)]}{[E_j(๐+๐)E_i(๐)]},$$
(3)
where $`z`$ is the direction of magnetization and $`f_i(๐)`$ is the Fermi-Dirac distribution function for the $`i`$-th valence band subband. It is seen that for spatially uniform magnetization (the ferromagnetic case), only the terms corresponding to $`q=0`$ (i.e., the diagonal terms) determine $`T_C`$. For the strongly degenerate carrier gas as well as neglecting the spin-orbit interaction, $`\rho _s`$ is equal to the total density-of-states $`\rho `$ for intra-band charge excitations, where $`\rho =m_{DOS}^{}k_F/\pi ^2\mathrm{}^2`$.
Equation (2) for $`T_C`$ with $`\rho _s=\rho `$ has already been derived by a number of equivalent methods . It is straightforward to generalize the model for the case of the carriers confined to the $`d`$-dimensional space . The tendency towards the formation of spin-density waves in low-dimensional systems as well as possible spatial correlation in the distribution of the magnetic ions can also be taken into account. The mean-field value of the critical temperature $`T_๐`$, at which the system undergoes the transition to a spatially modulated state characterized by the wave vector $`๐`$, is given by the solution of the equation,
$$\beta ^2A_F(๐,T_๐)\rho _s(๐,T_๐)๐๐ป\chi _o(๐,T_๐,๐ป)|\varphi _o(๐ป)|^4=4g^2\mu _B^2.$$
(4)
Here $`๐`$ spans the $`d`$-dimensional space, $`\varphi _o(๐ป)`$ is the envelope function of the carriers confined by a $`(3d)`$-dimensional potential well $`V(๐ป)`$; $`g`$ and $`\chi _o`$ denote the Landรฉ factor and the $`๐`$-dependent magnetic susceptibility of the magnetic ions in the absence of the carriers, respectively. Within the MFA, such magnetization shape and direction will occur in the ordered phase, for which the corresponding $`T_๐`$ attains the highest value.
## III Comparison of the model to selected experimental results
### III.1 Magnetic circular dichroism in (Ga,Mn)As
In the case of II-VI DMS, detail information on the exchange-induced spin-splitting of the bands, and thus on the coupling between the effective mass electrons and the localized spins has been obtained from magnetooptical studies . A similar work on (Ga,Mn)As led to a number of surprises. The most striking was the opposite order of the absorption edges corresponding to the two circular photon polarizations in (Ga,Mn)As comparing to II-VI materials. This behavior of circular magnetic dichroism (MCD) suggested the opposite order of the exchange-split spin subbands, and thus a different origin of the sp-d interaction in these two families of DMS. A new light on the issue was shed by studies of photoluminescence (PL) and its excitation spectra (PLE) in p-type (Cd,Mn)Te quantum wells . As shown schematically in Fig. 1, the reversal of the order of PLE edges corresponding to the two circular polarizations results from the Moss-Burstein effect, that is from the shifts of the absorption edges associated with the empty portion of the valence subbands in the p-type material. This model was subsequently applied to interpret qualitatively the magnetoabsorption data for metallic (Ga,Mn)As . Surprisingly, however, the anomalous sign of the MCD was present also in non-metallic (Ga,Mn)As, in which EPR signal from occupied Mn acceptors was seen . It has, therefore, been suggested that the exchange interaction between photo- and bound-holes is responsible for the anomalous sign of the MCD in those cases . The presence of such a strong exchange mechanism is rather puzzling, and it should be seen in non-magnetic p-type semiconductors. At the same time, according to our two-fluid model, the co-existence of strongly and weakly localized holes is actually expected on the both sides of the MIT. Since the Moss-Burstein effect operates for interband optical transitions involving weakly localized states, it leads to the sign reversal of the MCD, also on the insulating side of the MIT.
Another striking property of the MCD is a different temperature dependence of the normalized MCD at low and high photon energies in ferromagnetic (Ga,Mn)As . This observation was taken as an evidence for the presence of two spectrally distinct contributions to optical absorption . A quantitative computation of MCD spectra has recently been undertaken . The theoretical results demonstrate that because of the Moss-Burstein effect, the magnetization-induced splitting of the bands leads to a large energy difference between the positions of the absorption edges corresponding to the two opposite circular polarizations. This causes an unusual dependence of the low-energy onset of MCD on magnetization, and thus on temperature. These considerations lead to a quantitative agreement with the experimental findings, provided that the actual hole dispersion and wave functions are taken for the computation of MCD.
### III.2 Curie temperature in (Ga,Mn)As, (Zn,Mn)Te and (Cd,Mn)Te quantum well
The most interesting property of Ga<sub>1-x</sub>Mn<sub>x</sub>As epilayers is the large magnitude of $`T_C`$, of the order of 100 K for the Mn concentration $`x`$ as low as 5% . Because of this high $`T_C`$, the spin-dependent extraordinary contribution to the Hall resistance $`R_H`$ persists up to 300 K, making an accurate determination of the hole density difficult \[18, 24-26\]. However, the recent measurement of $`R_H`$ up to 27 T and at 50 mK yielded an unambiguous value of $`p=3.5\times 10^{20}`$ cm<sup>-3</sup> for a metallic Ga<sub>0.947</sub>Mn<sub>0.053</sub>As sample, in which $`T_C`$ = 110 K is observed . As shown in Fig. 2, the present model explains, with no adjustable parameters, such a high value of $`T_C`$.
The studied epilayers of (Zn,Mn)Te:N were on the insulating side of the metal-insulator transition (MIT), as the bound magnetic polaron (BMP) formation enhances localization. Nevertheless, if the concentration of acceptors was sufficiently high, the ferromagnetic Curie-Weiss temperatures $`T_{CW}`$ were observed as well as magnetic hysteresis were detected below $`T_CT_{CW}`$ . At the same time, the values of $`T_C`$ were much lower than those characterizing (Ga,Mn)As. However, a comparison of the experimental and calculated values of $`T_C`$ for (Zn,Mn)Te as a function of $`x`$ and $`p`$ (Fig. 3) demonstrates that the present model is capable to explain the magnitude of $`T_C`$ except for the samples with the smallest $`x`$. In the latter, $`p/xN_o`$ is as large as 0.6, so that precursor effects of Friedel oscillations and Kondo correlation are expected at low temperatures .
Two effects appear to account for the greater $`T_C`$ values in p-(Ga,Mn)As than in p-(Zn,Mn)Te at given $`p`$ and $`x`$. First is the smaller magnitude of the spin-orbit splitting between the $`\mathrm{\Gamma }_8`$ and $`\mathrm{\Gamma }_7`$ bands in arsenides, $`\mathrm{\Delta }_o=0.34`$ eV, in comparison to that of tellurides, $`\mathrm{\Delta }_o=0.91`$ eV. Once the Fermi energy $`E_F`$ approaches the $`\mathrm{\Gamma }_7`$ band, the density-of-states effective mass increases, and the reduction of the carrier spin susceptibility by the spin-orbit interaction is diminished. The computed value of $`T_C`$ for $`p=3\times 10^{20}`$ cm<sup>-3</sup> is greater by a factor of four in (Ga,Mn)As than that evaluated in the limit $`\mathrm{\Delta }_o>>E_F`$. The other difference between the two materials is the destructive effect of antiferromagnetic interactions, which operate in II-VI compounds but are of minor importance in III-V materials, as explained in Sec. 2.
The model discussed above describes also the magnitude of $`T_C`$ and its dependence on $`x`$ in modulation-doped quantum wells of p-(Cd,Mn)Te, if $`A_F=2`$ is assumed , an expected value for the relevant densities of the two-dimensional hole liquid. A good description of $`T_C(p)`$ is also obtained, provided that disorder broadening of the density-of-states at low $`p`$ is taken into account . Whether the ground state corresponds to uniform magnetization or rather to a spin-density wave is under study now.
### III.3 Effects of strain
Already early studies of a ferromagnetic phase in (Ga,Mn)As epilayers demonstrated the existence of substantial magnetic anisotropy . Magnetic anisotropy is usually associated with the interaction between spin and orbital degrees of freedom of the magnetic electrons. According to the model in question, these electrons are in the d<sup>5</sup> configuration. For such a case the orbital momentum $`L=0`$, so that no effects stemming from the spin-orbit coupling could be expected. To reconcile the model and the experimental observations, we note that the interaction between the localized spins is mediated by the holes, characterized by a non-zero orbital momentum. An important aspect of the present model is that it does take into account the anisotropy of the carrier-mediated exchange interaction associated with the spin-orbit coupling in the host material, an effect difficult to include within the standard approach to the RKKY interaction.
The computed effect of the cubic anisotropy on $`T_C`$ has been found to be small; differences between $`T_C`$ values calculated for various orientations of magnetization in respect to crystallographic axes are below 0.1 K in (Ga,Mn)As . The corresponding differences are, however, greater in the presence of epitaxial strain; of the order of 1 K for 1% biaxial strain in the (001) plane. Thus, such a strain can control the orientation of the easy axis. According to the computation for the relevant hole concentrations, the easy axis is in the plane for the case of unstrained or compressively strained films but under tensile strain the easy axis takes the direction. These expectations are corroborated by the experimental study, in which appropriate substrates allowed to control the direction of strain . It worth noting that similarly to strain, also confinement of the holes affects the magnetic anisotropy in accord with the theoretical model, the easy axis is oriented along the growth direction in the ferromagnetic p-(Cd,Mn)Te quantum wells .
## IV Other perspective materials
In view of the general agreement between experiment and theory for $`T_C`$ and the magnetic anisotropy, it is tempting to extend the model for material systems that might be suitable for fabrication of novel ferromagnetic semiconductors. For instance, the model suggests immediately that $`T_C`$ values above 300 K could be achieved in Ga<sub>0.9</sub>Mn<sub>0.1</sub>As, if such a large value of $`x`$ would be accompanied by a corresponding increase of the hole concentration. Figure 4 presents the values of $`T_C`$ computed for various tetrahedrally coordinated semiconductors containing 5% of Mn and $`3.5\times 10^{20}`$ holes per cm<sup>3</sup> . In addition to adopting the tabulated values of $`\gamma _i`$ and $`\mathrm{\Delta }_o`$ \[49-52\] the same value of $`\beta =\beta `$\[(Ga,Mn)As\] for all group IV and III-V compounds was assumed, which results in an increase of $`|\beta N_o|a_o^3`$, where $`a_o`$ is the lattice constant, a trend known to be obeyed within the II-VI family of magnetic semiconductors . By extending the model for wurzite semiconductors, $`T_C`$ values for parameters of ZnO (Fig. 3) and for wurzite GaN (not shown) were evaluated. For the employed parameters the magnitude of $`T_C`$ for the cubic GaN (Fig. 4) is by 6% greater than that computed for the wurzite structure.
The data (Fig. 4) demonstrate that there is much room for a further increase of $`T_C`$ in p-type magnetic semiconductors. In particular, a general tendency for greater $`T_C`$ values in the case of lighter elements stems from the corresponding increase in p-d hybridization and reduction of spin-orbit coupling. It can be expected that this tendency is not altered by the uncertainties in the values of the relevant parameters. Important issues of solubility limits and self-compensation as well as of the transition to a strong-coupling case with decreasing $`a_o`$ need to be addressed experimentally. We note in this context that since, in general, III-V compounds can easier be doped by impurities that are electrically active, whereas II-VI materials by transition metals, a suggestion has been put forward to grow magnetic III-V/II-VI short period superlattice .
Finally, we address the important question about the feasibility of synthesizing n-type or intrinsic ferromagnetic semiconductors. A work on (Zn,Mn)O:Al is relevant in this context. One should not forget, however, about the existence of, e.g., europium chalcogenides and chromium spinels, whose ferromagnetism is not driven by free carriers. Actually, a theoretical suggestion has been made that superexchange in Cr-based II-VI compounds can lead to a ferromagnetic order. Desired material properties, such as divergent magnetic susceptibility and spontaneous magnetization, can also be achieved in the case of a strong antiferromagnetic super-exchange interaction. The idea here is to synthesize a ferrimagnetic system that would consist of antiferromagnetically coupled alternating layers containing different magnetic cations, e.g., Mn and Co.
The above list of possibilities is by no means exhausting. With no doubt we will witness many unforeseen developments in the field of ferromagnetic semiconductors in the near future.
## Acknowledgments
The work at Tohoku University was supported by the Japan Society for the Promotion of Science and by the Ministry of Education, Japan; the work in Poland by State Committee for Scientific Research, Grant No. 2-P03B-02417, and by Foundation for Polish Science. We are grateful to F. Matsukura and Y. Ohno for collaboration on III-V magnetic semiconductors. One of us (T D.) thanks J. Jaroszyลski and P. Kossacki in Warsaw and J. Cibert and D. Ferrand in Grenoble for collaboration on II-VI magnetic semiconductors. |
warning/0002/astro-ph0002445.html | ar5iv | text | # NICMOS Imaging of the Damped Ly-๐ผ Absorber at ๐ง=1.89 toward LBQS 1210+1731 : Constraints on Size and Star Formation Rate
## 1 INTRODUCTION
Damped Lyman-$`\alpha `$ absorption systems detected in spectra of high-redshift quasars are believed to be the progenitors of present-day galaxies, because they show high H I column densities (log $`N_{HI}20.0`$) and display absorption lines of several heavy elements. However, there are various competing ideas regarding the nature of the galaxies underlying the DLAs. Wolfe et al. (1986) suggested that the DLAs are rotating proto-disks. This suggestion has also been made by Prochaska & Wolfe (1997, 1998), based on asymmetric line profiles of the heavy-element absorption lines in DLAs. On the other hand, gas-rich dwarf galaxies have also been suggested as candidate objects for the DLAs (York et al. 1986; Matteucci, Molaro, & Vladilo 1997). Recently, Jimenez, Bowen, & Matteucci (1999) have suggested that high-redshift DLAs may arise in low-surface brightness galaxies. The lack of substantial chemical evolution found in studies of element abundances in DLAs (e.g., Pettini et al. 1999; Kulkarni, Bechtold, & Ge 2000a) also shows that the currently known population of DLAs seems to be dominated by metal-poor objects, so DLAs may consist of dwarf or low-surface brightness galaxies with modest star formation rates. Unfortunately, it is hard to determine what type of galaxies underlie the DLAs, since most previous efforts to directly image the high-redshift DLAs have failed. A few detections have been made at low redshifts, which showed those DLAs to arise in low surface-brightness galaxies (see, e.g., Steidel et al. 1995a, 1995b; LeBrun et al. 1997). But high-redshift DLAs with $`z_{abs}<z_{em}`$ have proven hard to detect, and the question of the nature of galaxies giving rise to these DLAs is still open.
Many of the previous attempts to detect the emission from DLAs concentrated on the Ly-$`\alpha `$ emission, which is an expected signature from a star-forming region (e.g. Smith et al. 1989; Hunstead, Pettini, & Fletcher 1990; Lowenthal et al. 1995). There have been only a few Ly-$`\alpha `$ detections of DLAs so far. Mรธller & Warren (1998) and Mรธller, Warren, & Fynbo (1998) detected Ly-$`\alpha `$ emission in the fields of two DLAs at $`z=2.81`$ and $`z=1.93`$. However, both of these DLAs have $`z_{abs}z_{em}`$ and may be different from the general population of intervening DLAs. Djorgovski et al. (1996) and Djorgovski (1997) reported Ly-$`\alpha `$ emitting objects with $`R25`$ (and inferred SFRs of a few M yr<sup>-1</sup>) in fields of a few DLAs, located at 2-3 $`\mathrm{}`$ from the quasar. However, the Ly-$`\alpha `$ technique cannot definitively measure the star formation rates of the DLAs because of the generally unquantifiable effects of dust extinction in the systems. The lack of detections in the other Ly-$`\alpha `$ studies of interevening DLAs could indicate either that DLAs have low star formation rates (SFR) or that the emission is extinguished by dust. As pointed out by Charlot & Fall (1991), even small quantities of dust are sufficient to extinguish the Ly-$`\alpha `$ emission, since resonant scattering greatly increases the path length of Ly-$`\alpha `$ photons attempting to escape from an H I cloud. Indeed, observations of reddening of background quasars and evidence for depletion of Cr, Fe, Ni etc. relative to Zn suggest the presence of a small amount of dust in DLAs (see, e.g., Pei, Fall, & Bechtold 1991; Pettini et al. 1997; Kulkarni, Fall, & Truran 1997). Thus, it is hard to constrain the SFRs in DLA galaxies using the non-detections or weak detections of Ly-$`\alpha `$ emission.
The issues of dust and SFR in high redshift DLAs are also important in view of recent claims based on mid-IR and far-IR observations that a large fraction of the star formation at high redshifts is hidden from us by dust obscuration (e.g., Elbaz et al. 1998; Clements et al. 1999). One way to discern whether the previous non-detections of Ly-$`\alpha `$ were due to low SFR or presence of dust is to search for longer wavelength emission lines less affected by dust extinction and not subject to resonant scattering. The ground-based near-IR spectroscopic survey of Bunker et al. (1999), which searched for redshifted H-$`\alpha `$ emission in $`11\mathrm{}\times 2.5\mathrm{}`$ regions around 6 quasars with DLAs at $`z>2`$ and reached 3 $`\sigma `$ detection levels of 6-18 M yr<sup>-1</sup>, failed to detect any redshifted H-$`\alpha `$ emission from the DLAs in their sample. Some of the ground-based narrow-band photometric surveys for H-$`\alpha `$ emission from DLAs have also failed to detect any emission line objects in the DLA fields (e.g., Teplitz, Malkan, & McLean 1998, who however found H-$`\alpha `$ emitters in the fields of some weaker non-DLA metal line systems). Some other narrow-band searches for H-$`\alpha `$ emission have revealed multiple objects in the DLA fields separated by several arcseconds from the quasar (2-12$`\mathrm{}`$ for Bechtold et al. 1998, 9-120 $`\mathrm{}`$ for Mannucci et al. 1998). These surveys, which had 3 $`\sigma `$ detection limits of $`5`$ M yr<sup>-1</sup> (Bechtold et al. 1998) or $`10`$ M yr<sup>-1</sup> (Mannucci et al. 1998), found the H-$`\alpha `$ emitting objects to have a wide range of inferred SFRs (10-20 M yr<sup>-1</sup> for Bechtold et al. 1998, 6-90 M yr<sup>-1</sup> for Mannucci et al. 1998). The relatively large separations of these emission line objects from the quasars indicates that they are not the DLA absorbers themselves, but star-forming companions in the same larger structure (e.g. sheet or filament) as the DLA. None of these ground-based surveys has been able to probe the regions very close to the quasar sightline (angular separations $`<2`$ $`\mathrm{}`$), because of the limitations imposed by seeing in these studies. While these studies offer interesting information about the environments of the DLAs, high sensitivity diffraction-limited imaging is necessary for the detection of the DLA absorbers themselves (to probe small angular separations), and thus for determining the morphology and SFRs of the DLAs. The HST WFPC2 study of Le Brun et al. (1997) has detected candidates with angular separations $`<2`$ $`\mathrm{}`$ in broad band images for six DLAs at $`z<1`$ and one DLA at $`z=1.78`$. However, the information obtained from this study about the nature of high-redshift DLAs is limited since no narrow-band images were obtained and since the sample contained only 1 DLA at $`z>1`$. As mentioned earlier, the HST WFPC2 study of Mรธller & Warren (1998 and references therein) detected Ly-$`\alpha `$ emission in a $`z_{abs}>z_{em}`$ DLA, but this DLA may differ from intervening DLAs.
To summarize, many previous attempts to detect emission from high-redshift intervening DLAs have failed. The few detections so far consist mainly of either weak Ly-$`\alpha `$ detections (which cannot constrain the SFR completely) or detections of H-$`\alpha `$ companions at fairly large angular separations from the quasars. There are only four objects detected so far in fields of high-$`z`$ interevening DLAs at small angular separations. These objects have impact parameters between 4.3 and 11.5 h$`{}_{}{}^{1}{}_{70}{}^{}`$ kpc (where $`H_0=70`$ h<sub>70</sub> km s<sup>-1</sup> Mpc<sup>-1</sup>), and are promising candidates for the DLAs in those sightlines (see Mรธller & Warren 1998 and references therein; the other DLA impact parameter data listed in Mรธller & Warren 1998 are biased toward $`z_{abs}z_{em}`$ DLAs.). To further increase the number of promising candidates for high-redshift intervening DLAs, it is necessary to carry out more deep high spatial resolution near-infrared searches for DLAs.
We have obtained deep diffraction-limited images of three DLAs at $`z2`$ with the Near Infrared Camera and Multi-Object Spectrometer (NICMOS) onboard the Hubble Space Telescope (HST). Here we describe our NICMOS observations of the quasar LBQS 1210+1731 ($`z_{em}=2.543\pm 0.005`$; Hewett, Foltz, & Chaffee 1995), which has a spectroscopically known damped Ly-$`\alpha `$ absorber ($`z_{abs}=1.8920`$ and log $`N_{HI}=20.6`$; Wolfe et al. 1995). Our observations have the unique benefit of combining high near-IR sensitivity and high spatial resolution with a more stable and quantifiable PSF than is currently possible with ground-based observations. A further feature of some of our observations is the use of the NICMOS coronagraph, which greatly decreases the scattered light background outside of the coronagraphic hole and therefore allows a study of the environment of the DLA. Our analysis indicates two objects at 0.25 $`\mathrm{}`$ and 0.7 $`\mathrm{}`$ from the quasar that we cannot explain as any known artifacts of the PSF. We believe that these objects are likely to be real and may be associated with the DLA and its companions, at impact parameters of 1.5 and 3.8 h$`{}_{}{}^{1}{}_{70}{}^{}`$ kpc. We have thus probed regions far closer to the quasar sight-line than in most previous studies of high-redshift intervening DLAs, and the two objects we report mark the closest detected high-redshift intervening DLA candidates yet to any quasar sight line. Sections 2, 3, and 4 describe the observations, data reduction, and the subtraction of the quasar point spread functions. Our results are described in section 5. Section 6 describes various tests of our data analysis procedures, carried out to investigate whether the features seen after PSF subtraction are real. A summary of the results of the various data analysis tests is given in subsection 6.12. (Readers interested mainly in the scientific discussion can go directly from section 5 to subsection 6.12.) Finally, sections 7 and 8 discuss the implications of our observations for sizes, environment, and star-formation rates of DLA galaxies.
## 2 OBSERVATIONS
The field of LBQS 1210+1731 was first observed on 1998, July 22 from 07:14 UT to 16:40 UT, using NICMOS camera 2 (pixel scale $`0.076\mathrm{}`$, field of view $`19.45\mathrm{}\times 19.27\mathrm{}`$). A sequence of spatially offset broad-band images was obtained in multiaccum mode with the F160W (H) filter (central wavelength 1.5940 $`\mu `$m, FWHM 0.4030 $`\mu `$m). Field offsetting was accomplished with a 5-point spiral dither pattern in steps of $`7.5`$ pixels, using the Field Offset Mirror (FOM) internal to NICMOS. The exposures at each dwell point were 512 s long, giving a total integration time of 2560 s. See MacKenty et al. (1997) for a detailed description of NICMOS imaging modes and options. The multiaccum observations consisted of non-destructive readouts in the โstep32โ readout timing sequence, i.e. โmultiaccumโ readouts separated logarithmically up to 32 s and linearly in steps of 32 s beyond that. In addition, narrow-band images were obtained in the filter F190N (central wavelength 1.9005 $`\mu `$m, FWHM 0.0174 $`\mu `$m), in which the redshifted H-$`\alpha `$ emission from the DLA, if present, would lie. Four-point spiral dither patterns in steps of 7.5 pixels, with a 704 s โstep64โ multiaccum exposure at each dwell point, were repeated in five successive orbits, resulting in a total integration time of 14,080 s. The spatial resolution of the F160W and F190N images is 0.14 $`\mathrm{}`$ (1.8 pixels) and 0.17 $`\mathrm{}`$ (2.1 pixels) FWHM, respectively. Thus, camera 2 is almost critically sampled at the wavelengths used for our observations.
Finally, broad-band images in the F160W filter were also obtained using the camera 2 coronagraph on 1998, July 29 from 10:14 to 13:07 UT. These consisted of an initial pair of 92 s long target-acquisition images, which were followed by placement of the target in the coronagraphic hole (0.3 $`\mathrm{}`$ or 4 pixels in geometrical radius) and then integration of the object for a total of 4960 s (5 exposures of 480 s each in the first orbit and 5 exposures of 512 s each in the second orbit, all using the step32 multiaccum timing sequence). No dithering was used, of course, for the coronagraphic observations. The NICMOS coronagraph is comprised of two optical elements, a 165 $`\mu `$m physical diameter hole in the camera 2 field divider mirror at the reimaged HST f/24 optical telescope assembly (OTA) focus and a Lyot stop at a cold pupil in the cryostat. The coronagraphic system significantly reduces both scattered and diffracted energy from the occulted targetโs point spread function core by factors of 4-6, compared to direct imaging (Schneider et al. 1998; Lowrance et al. 1998). Thus our coronagraphic images have higher sensitivity than the non-coronagraphic images for detecting those foreground damped Ly-$`\alpha `$ absorber or associated companions that are much fainter than the quasar and lie outside the coronagraphic hole.
To circumvent image artifacts known as โbarsโ in all our camera 2 images, cameras 1 and 3 were run in parallel, as discussed by Storrs (1997).
## 3 REDUCTION OF IMAGES
The images were reduced using the IRAF package Nicred 1.8, developed specifically for the reduction of multiaccum NICMOS data (McLeod 1997). The dark image used was that made from on-orbit dark exposures taken during the NICMOS calibration program. For the non-coronagraphic images, the flat-field image used was made from on-orbit exposures taken with the internal calibration lamps during the NICMOS Cycle-7 calibration program. For the coronagraphic images, the flat-field image was made with target-acquisition data taken just before the coronagraphic exposures. This ensures that the coronagraphic hole is in the same position on the detector for the flat as for the quasar data, which is critical for studying faint objects close to the edge of the coronagraphic hole. (The standard calibration flats are not adequate for this purpose because the position of the coronagraphic hole on the detector changed with time, and a flat exposure taken at another time had the hole in a different place.)
First, the exposures at each individual dither position were reduced using Nicred 1.8. Briefly, the steps followed by Nicred 1.8 are as described below:
1. subtraction of the zeroth read from successive reads, both for the quasar data and the dark data,
2. dark subtraction, read by read,
3. linearity correction, cosmic ray rejection, and fitting of slope to the successive reads in the multiaccum data, to get count rates in ADU s<sup>-1</sup>,
4. correction of non-uniform bias level across the array (โthe pedestal effectโ),
5. repeating step 3 on the bias-corrected image to get more accurate count rates,
6. flatfielding using the appropriate flats,
7. subtracting the median of each row from that row and then likewise for columns, to remove bands caused by bias jumps during simultaneous use of amplifiers of other cameras in parallel, and thus to improve the flatness of the background,
8. fixing bad pixels using bi-cubic spline interpolation across the neighboring pixels.
The images for the different dither positions were registered by cross-correlating with the IRAF task xregister. The quasar was used as the reference object since no other point sources were available in our images. Finally, the registered images were averaged together using a bad-pixel mask that took out any remaining bad pixels, and rejecting pixels deviating by more than 3 $`\sigma `$ from the average of the five F160W images, using averaged sigma-clipping.
For the F190N images, where there were five exposures (one in each orbit) at each of the four dither positions, we first median-combined the five exposures at each position separately, and then registered and median-combined the four positions together to make the final image. For the coronagraphic F160W images, where there were five exposures at the same position in each of the two orbits, we averaged the exposures in each orbit separately and then took a weighted average (weighting by exposure times) of the combined exposures from the two orbits.
In an attempt to gain the higher spatial resolution made possible by the half-integral dithers (in steps of 7.5 pixels), we also experimented with magnification (repixelization) of the images at the individual dither positions before combining them. The images for each individual dither position processed as per steps 1-8 above were magnified (i.e. numerically resampled) by factors of two each in x and y directions. A cubic spline interpolation was used to divide the pixels into subpixels, with the flux kept conserved. As discussed further in section 6.7, our results do not depend much on whether or not the magnification is done.
Figs. 1, 2, and 3 show the final reduced images for the non-coronagraphic F160W, non-coronagraphic F190N, and coronagraphic F160W observations. The orientations of Figs. 1 and 2 agree exactly while they differ from that of Fig. 3 by only 2.026 degrees. All three images have an essentially zero background. The F190N image shows a weak residual flat field and nonuniformities in the corners caused by amplifier glow. This effect is much less noticeable in the reduced F160W image. We believe that the F190N image is limited by the quality of the F190N flat field available to us. The F190N flat field, made from six 192 s long in-flight exposures to calibration lamps, has a count rate of 37.72 ADU s<sup>-1</sup>, while the F160W flat, made from nine 24 s exposures, has a count rate of 1113 ADU s<sup>-1</sup>. The rms deviation in the count rate per pixel is about $`2\%`$ for each of the six frames combined to make the F190N flat, while it is about $`0.02\%`$ for each of the nine frames combined to make the F160W flat. The lack of a better F190N flat is unfortunate. However, this should not be a serious problem for the quasar and DLA images, since they lie in the central part of the array. Figures 1 and 2 show the quasar point source along with the diffraction pattern. The coronagraphic image in Fig. 3 shows the quasar light to be reduced greatly, although not completely. To study whether there is any additional underlying faint emission from the DLA in any of these images, we need to subtract the respective PSFs.
## 4 SUBTRACTION OF THE QUASAR POINT SPREAD FUNCTION
### 4.1 SELECTION OF THE PSF STAR
Reference point spread functions for subtraction were obtained by using observations of stars in the same filter / aperture combinations as those employed for the quasar imaging. PSF star observations were not included in our own observations since we wanted to maximize the use of the available HST observing time for imaging of the quasar fields. We therefore used PSF star observations from other programs (in particular the stellar images from the photometric monitoring program carried out during Cycle 7 NICMOS calibration) for constructing the reference PSFs for subtraction. Such directly observed PSFs, when exposed to high S/N, are expected to provide better match to the quasar data than the theoretical Tiny Tim PSFs (since the observed PSFs incorporate any real optical effects not simulated in Tiny Tim). We have also actually experimented with the use of calculated Tiny Tim PSFs and find that they do indeed provide poorer match to the quasar than the observed stellar PSFs.
For the non-coronagraphic images, the PSF observations were chosen such that the telescope focus โbreathingโ (Bely 1993) values matched as closely as possible the values for the DLA observations. This is important because changes in the HST focus translate into corresponding changes in the fine structure of the PSF. To estimate the OTA focus positions for the epoch of the quasar and PSF star observations, we used the HST focus ephemerides provided by STScI (Hershey 1998; Hershey & Mitchell 1998). For the non-coronagraphic F160W and F190N images, we used the PSF star P330E, observed on July 8, 1998 and May 29, 1998 respectively. We also studied the effect of PSF variations on our results by using PSF observations of P330E with a range of different โbreathingโ focus positions obtained on different dates, and also by using observations of other PSF stars. (See section 6.2 and 6.3 for a detailed description.) The F160W non-coronagraphic image of P330E, made by combining four exposures of 3 s each, had a count rate of 108.10 ADU s<sup>-1</sup> at the maximum of the first Airy ring. The corresponding quasar image, made by combining five exposures of 512 s each, had a count rate of 1.91 ADU s<sup>-1</sup> at the maximum of the first Airy ring. For the F190N filter, the P330E image, made by combining three exposures of 64 s each, had 0.055 ADU s<sup>-1</sup> at the maximum of the first Airy ring. The corresponding count rate was 0.0013 ADU s<sup>-1</sup> for the F190N quasar image, made by combining 20 exposures of 704 s each.
For the coronagraphic observations, the choice of the PSF star was guided by the requirement that the position of the star in the coronagraphic hole be as close as possible to that of the quasar in our observations. This is critical, because even when the target-acquisition flight software succeeds in acquiring the target and putting it in the coronagraphic hole, there are usually some small residual differences between the actual position where the target is placed and the desired position of the target in the hole, i.e. the โlow scatter pointโ of the coronagraph (see Schneider 1998 for details). The PSF wings and โglintsโ from the edge of the coronagraphic hole depend sensitively on the precise position of the point source within the hole. We therefore used the observations of star GL83.1 for which we had coronagraphic observations (from another NICMOS GTO program), with the star placed at a position within 0.04 pixels of the position of the quasar LBQS 1210+1731 in our data. The observations of GL83.1 were taken on August 1, 1998 at a breathing value close to that for our quasar coronagraphic observations.
For all the โprimaryโ PSF star choices, the proximity of the observation dates with those of our quasar observations also ensures that the plate scale of the camera is the same for the PSF and the quasar observations.
### 4.2 SUBTRACTION OF THE PSF STAR
All of the observations of the PSF star P330E were obtained in 4-point spiral dither patterns in steps of 4.0 $`\mathrm{}`$ ($`52.6`$ pixels). The dithers for the PSF star were obtained by using actual spacecraft movements, while the dithers for LBQS 1210+1731 were obtained by moving the field offset mirror (FOM) internal to NICMOS. But the use of the FOM should not cause any differences between the combined quasar PSF and the reference star PSF, since we registered all of the quasar exposures individually to a common reference before combining. The PSF star observations were analyzed in exactly the same manner as the quasar observations, following the procedure outlined in section 3. The same interpolation scheme was used for resampling of the PSF star and quasar images. The difference in the dithering steps for the quasar and the PSF star may give rise to difference in actual sampling of the quasar and PSF star images. But, as described in section 6.7, we have verified that the difference images are reproduced well when both the quasar and PSF star images are numerically resampled by a factor of two. The final reduced PSF star images were subtracted from the corresponding quasar images after suitable scaling and registration, using the IDL program โIDP-3โ (Lytle et al. 1999). The scale factors were chosen using the relative intensities of the PSF wings in the quasar image and the PSF star image. For the coronagraphic image, the relative intensities of the PSF โglintsโ near the edge of the hole were also used in determining the PSF scaling factor. All the parameters (i.e., relative x and y alignment of the PSF star image with respect to the quasar image and the intensity scaling factor for the PSF star image) were fine-tuned iteratively to obtain the minimum variance in roughly 3 $`\mathrm{}`$ x 3 $`\mathrm{}`$ subregions (around the quasar) in the PSF-subtracted image. Radial plots of the quasar image, the aligned and scaled PSF image, and the difference of the two were also examined to check the alignment and scaling of the PSF. Figs. 4a, 5a, 6a show zoomed $`3`$$`\times 3`$โณ subregions around the quasar, from the non-coronagraphic F160W, non-coronagraphic F190N, and coronagraphic F160W images shown in Figs. 1, 2, 3, respectively. Figs. 4b, 5b, 6b show the PSF-subtracted versions of Figs. 4a, 5a, and 6a, respectively, using the closest matching PSFs available.
## 5 RESULTS
### 5.1 NON-CORONAGRAPHIC F160W IMAGES
Fig. 4b shows the F160W image after subtraction of the PSF image of star P330E obtained on July 8, 1998. The fidelity of the PSF subtraction is seen from the fact that the diffraction pattern disappears completely and most of the residual image contains a random mixture of positive and negative values. The radially symmetric residuals may be explained, for the most part, by a mismatch between the SED of the quasar and that of the PSF star (spectral type G2V). These color terms lead to small differences in the structure and size-scale of the PSFs. These differences are non-negligible under the $`25\%`$ bandpass of the F160W filter, but are negligible under the 1$`\%`$ bandpass of the F190N filter. See section 6.3 for further discussion. The main asymmetric residual is the emission feature to the โlower rightโ of the center, about 3 pixels (0.26 $`\mathrm{}`$) away from the center. (This feature is seen more clearly if the data are sub-sampled by a factor of 2, as discussed further in section 6.7 and Fig. 15.) There is no correspondingly strong and symmetrically located negative feature in the image, and the bright knot can not be made to disappear after reregistration of the PSF and quasar images or rescaling of the PSF image without causing large negative residuals elsewhere (see section 6.9 and Fig. 18). We cannot completely rule out that this โknotโ is an artifact in the PSF. However, given the significant excess over a number of pixels, it is likely that it is a real feature. This feature (which we name object โO1โ) is about 0.40 $`\mathrm{}`$ long. If this emission knot is associated with the damped Ly-$`\alpha `$ absorber at $`z=1.892`$, then it is $``$ 2.4 h$`{}_{}{}^{1}{}_{70}{}^{}`$ kpc long for $`q_0=0.5`$, or 3.2 h$`{}_{}{}^{1}{}_{70}{}^{}`$ kpc long for $`q_0=0.1`$. In section 6, we examine in detail the question of whether or not O1 is real. In this section we discuss the properties of O1 assuming that it is real and is associated with the DLA absorber.
The faintness and diffuse nature of object O1 make its photometry rather difficult. We estimated the flux from this object in the PSF-subtracted image, using three different procedures, and then took an average of the three values.
Since accurate aperture photometry is difficult, we first estimated the flux by subtracting the PSF star from object O1, now multiplying the star by a factor large enough to make object O1 look indistinguishable from noise. This PSF multiplying factor can then be used directly to estimate the flux of object O1, since the PSF star P330E is a well-calibrated NICMOS photometric standard. From this method, we deduce that object O1 is $`2.55\times 10^4`$ as bright as the PSF star P330E. This implies a flux of 3.22 ADU s<sup>-1</sup> or 7.1 $`\mu `$Jy in the F160W filter. To convert the count rate to flux, we used the NICMOS photometric calibration factor of $`2.190\times 10^6`$ Jy/(ADU s<sup>-1</sup>) for the F160W filter. This factor was derived using the solar-type photometric standard star P330E (that we also used for PSF subtraction).
As a rough check of the above flux value, we also did aperture photometry on a circular aperture 4 pixels in diameter centered on O1 using the IRAF task apphot. A constant background value was subtracted as the sky value. (This constant, estimated as the average of the mean values per pixel of about 20 10x10 subregions in different parts of the image, had a very low mean value and therefore made negligible change to the final flux values.) This yields 2.55 ADU s<sup>-1</sup>, i.e., 5.6 $`\mu `$Jy before correcting for aperture effects. For reference, the 1 $`\sigma `$ noise level in the PSF-subtracted image is about 0.13 ADU s<sup>-1</sup> per pixel (0.28 $`\mu `$Jy per pixel) in a circular annulus 0.2 $`\mathrm{}`$ wide centered at 0.3 $`\mathrm{}`$ from the quasar center. The corresponding noise levels at 0.5 $`\mathrm{}`$, 0.7 $`\mathrm{}`$, 0.9 $`\mathrm{}`$, and 1.1 $`\mathrm{}`$ from the quasar center are 0.023, 0.014, 0.012, and 0.012 ADU s<sup>-1</sup> per pixel (i.e., 0.051, 0.031, 0.027, and 0.026 $`\mu `$Jy per pixel respectively).
The 4-pixel diameter circular aperture covers most of the region of emission in O1 and avoids the residuals near the center of the quasar and very narrow features that we think arise from residual PSF differences. This region however does not include the pixels at the extreme ends of the โmajor axisโ of O1, which is about 6 pixels long. We therefore also estimated the flux by doing a pixel-by-pixel addition over the region actually occupied by O1, which gives 2.63 ADU s<sup>-1</sup>, i.e., 5.8 $`\mu `$Jy. The flux values estimated by both the aperture photometry methods need to be corrected for the fact that a significant fraction of the energy of even a point source lies outside the radius of 2 pixels. Using aperture photometry on the standard star P330E, we estimate that the aperture correction factor between radii $`r=2`$ and $`r=7.5`$ pixels is 1.625. A further factor of 1.152 has been estimated for camera 2 filter F160W for the aperture correction from a 7.5-pixel radius aperture to the total flux, based on standard NICMOS photometric calibrations made with the standard star P330E. Thus, the total aperture correction factor is 1.872 for the second method. We note, however, that there is a roughly 10 $`\%`$ uncertainty in the aperture correction factor. Schneider et al. have estimated the above correction factor to be 2.08. Taking an average, we adopt an aperture correction factor of 1.98 $`\pm 0.1`$. Since the region used in the third method is approximately also 4 pixels in diameter (although slightly bigger near the ends of the โmajor axisโ of O1), the aperture correction is (at least) 1.98 in this case. We therefore use this factor for the third method of flux estimation also, although it is hard to be sure of the exact aperture correction in this case. After the aperture corrections, we derive flux values of 11.0 $`\mu `$Jy with the second method and 11.4 $`\mu `$Jy with the third method.
On averaging the three flux values derived above, we estimate a flux of $`9.8\pm 2.4`$ $`\mu `$Jy for the flux from object O1. This corresponds to $`m_{F160W}=20.11_{0.24}^{+0.30}`$ (taking the zero magnitude to correspond to 1083 Jy in the Johnson system). Here we have used equal weights for the three values while averaging, although we note that the value obtained by subtracting a point source is likely to be more accurate than the other two values. Note that the error estimate indicates the standard deviation among the three flux estimates obtained by the three methods, and thus reflects the uncertainties in the size and shape of object O1. For comparison, the 1 $`\sigma `$ uncertainty in the background near O1 is 0.051 $`\mu `$Jy per pixel, or $`0.2`$ $`\mu `$Jy over the region of $`12`$ pixels occupied by O1.
The observed F160W flux corresponds to a luminosity (at mean rest frame wavelength of 0.55 $`\mu `$m) of about $`1.5\times 10^{10}`$ h$`{}_{}{}^{2}{}_{70}{}^{}`$ L for $`q_0=0.5`$ and about $`2.8\times 10^{10}`$ h$`{}_{}{}^{2}{}_{70}{}^{}`$ L for $`q_0=0.1`$. Thus, object O1 is fainter than an L galaxy at $`z=1.89`$ by 0.2-0.9 magnitudes. If O1 is not the DLA, the DLA is even fainter. Our results here are consistent with those of Djorgovski (1997), who reported a possible counterpart to the $`z=4.10`$ DLA toward DMS 2247-0209. That DLA candidate is located $`3.3\mathrm{}`$ from the quasar (i.e. 22 h$`{}_{}{}^{1}{}_{70}{}^{}`$ kpc for q$`{}_{0}{}^{}=0.1`$), with an inferred continuum luminosity of 0.5 L.
### 5.2 NON-CORONAGRAPHIC F190N IMAGES
At a redshift of $`z_{DLA}=1.892`$, any H-$`\alpha `$ emission would be expected to lie at $`\lambda _{obs}=1.898`$ $`\mu `$m, which is very close to the center and mean $`\lambda `$ of 1.900 $`\mu `$m for the filter F190N. Thus, the narrow-band images in filter F190N are expected to reveal any redshifted H-$`\alpha `$ emission from the DLA. Fig. 5b shows the PSF-subtracted F190N image using the PSF image of the star P330E observed on May 29, 1998. The residual image shows an emission feature in roughly the same place ($`0.28`$ $`\mathrm{}`$ away from the quasar center to the lower right) and with roughly the same size (0.42 $`\mathrm{}`$ long) as the feature seen in the non-coronagraphic F160W image. This feature is more clearly evident if the images are sub-sampled by a factor of 2 (as discussed in sec. 6.7 and Fig. 16). As in the F160W image, this feature also does not disappear after realigning or rescaling the PSF. This suggests that feature O1 may be a real object. If O1 is associated with the DLA absorber at $`z=1.892`$, the absorber is $``$ 2.5 h$`{}_{}{}^{1}{}_{70}{}^{}`$ kpc long for $`q_0=0.5`$, or $``$ 3.4 h$`{}_{}{}^{1}{}_{70}{}^{}`$ kpc long for $`q_0=0.1`$.
As in the case of the broad-band images, the photometry of O1 is rather difficult. We do it in three different ways and take an average. In an attempt to get a flux estimate free of the uncertain aperture correction factor, we first subtracted the standard star P330E from O1, scaling the star such that the feature O1 just disappears. This method gives a flux in O1 of $`4.11\times 10^4`$ times that of P330E. This corresponds to a flux of 0.204 ADU s<sup>-1</sup>, i.e. 9.1 $`\mu `$Jy. Here we have used the NICMOS photometric calibration factor of $`4.455\times 10^5`$ Jy/(ADU s<sup>-1</sup>) for the F190N filter. Aperture photometry yields a flux of 0.114 ADU s<sup>-1</sup> in a 4-pixel diameter circular aperture centered on the center of O1. A pixel-by-pixel addition over the region occupied by O1 gives a flux of 0.132 ADU s<sup>-1</sup>. Based on the photometry of the standard star P330E, we estimate that the aperture correction factor between $`r=2`$ and $`r=7.5`$ is 1.691. The aperture correction factor between $`r=7.5`$ and the total flux is expected to be 1.159 for the F190N filter. This implies a total aperture correction factor of 1.960 for the $`r=2`$ values. But a $`10\%`$ uncertainty exists in the aperture correction factor, similar to that discussed above for the F160W images. We therefore adopt an average aperture correction factor of 2.06. Applying this aperture correction, we get flux values of 0.234 ADU s<sup>-1</sup> (10.4 $`\mu `$Jy) and 0.272 ADU s<sup>-1</sup> (12.2 $`\mu `$Jy), respectively for the second and third methods. Averaging the three flux values obtained by the three methods, we get $`10.6\pm 1.5`$ $`\mu `$Jy. The error bar of 1.5 $`\mu `$Jy denotes only the standard deviation among the three values and thus reflects the uncertainties arising from the lack of knowledge about the size and shape of O1. For comparison, the 1$`\sigma `$ noise levels in the F190N image (after PSF subtraction) at $`r=0.5`$ $`\mathrm{}`$, 0.7 $`\mathrm{}`$, 0.9 $`\mathrm{}`$, and 1.1 $`\mathrm{}`$ from the quasar are 0.0027, 0.0019, 0.0018, and 0.0016 ADU s<sup>-1</sup> per pixel, i.e., 0.119, 0.084, 0.081, and 0.073 $`\mu `$Jy per pixel, respectively. Thus the 1 $`\sigma `$ sky noise uncertainty in the total summed flux over the $`12`$ pixel region occupied by O1 is $``$ 0.4 $`\mu `$Jy (using the noise estimates just outside O1 at $`r=0.5`$ โณ).
The expected F190N continuum must be subtracted from the observed flux in order to determine if a statistically significant redshifted H-$`\alpha `$ excess exists. We estimate the continuum under the F190N filter by scaling the F160W image using the relative photometric calibration of the two filters. We find that, in fact, this expected continuum flux agrees almost completely with the observed F190N flux. After subtraction of the expected F190N continuum image (scaled from the F160W image) from the observed F190N image, we find a very marginal excess of 0.0074 ADU s<sup>-1</sup>. With the aperture correction, this corresponds to 0.68 $`\mu `$Jy. The 1 $`\sigma `$ noise level in the F190N-F160W image is 0.0026 ADU s<sup>-1</sup> per pixel just outside the location of O1. This noise level corresponds to a 1 $`\sigma `$ uncertainty of 0.4 $`\mu `$Jy in the total flux summed over the region occupied by O1. The slight excess at the location of O1 in the F190N-F160W image is thus not statistically significant. We therefore conclude that the contribution to the F190N flux from redshifted H-$`\alpha `$ emission is negligible. It is not likely that we could have missed the H-$`\alpha `$ emission from O1. The H-$`\alpha `$ emission from the DLA could lie outside the F190N bandpass only if the DLA galaxy is lower in velocity by more than 980 km s<sup>-1</sup> or higher in velocity by more than 1770 km s<sup>-1</sup> from the absorption redshift. Such offsets are higher than the observed internal velocity dispersion in any typical single galaxy.
Integrating over the FWHM of the F190N filter, assuming no dust extinction, and using the prescription of Kennicutt (1983) for conversion of H$`\alpha `$ luminosity to SFR, the nominal marginal excess of 0.68 $`\mu `$Jy in the F190N-F160W image corresponds to an SFR of 1.1 $`h_{70}^2`$ M yr<sup>-1</sup> for $`q_0=0.5`$, or 2.0 $`h_{70}^2`$ M yr<sup>-1</sup> for $`q_0=0.1`$. To derive a better estimate of the uncertainty in the H-$`\alpha `$ flux, we experimented with subtractions of the PSF-subtracted F190N and F160W images. In a 4-pixel region (roughly the size of our resolution element), an H-$`\alpha `$ emission strength of about 0.016 ADU s<sup>-1</sup> (0.71 $`\mu `$Jy) would yield S/N = 3. With an aperture correction factor of 3.41, this corresponds to a total 3 $`\sigma `$ flux limit of 0.054 ADU s<sup>-1</sup> or 2.4 $`\mu `$Jy. This translates into a 3 $`\sigma `$ upper limit on the SFR of 4.0 $`h_{70}^2`$ M yr<sup>-1</sup> for $`q_0=0.5`$ or 7.4 $`h_{70}^2`$ M yr<sup>-1</sup> for $`q_0=0.1`$. (We consider the possibility of dust extinction in section 7.3 below.)
### 5.3 CORONAGRAPHIC F160W IMAGES
An F160W coronagraphic image of the quasar is shown in Fig. 6a, in which the coronagraphic hole is masked out. Almost all the flux seen in this reduced coronagraphic image is due to residual scattered light from the quasar, and โglintsโ from the edge of the hole. After subtraction of a reference PSF image using observations of the star GL83.1, these artifacts disappear almost entirely (Fig. 6b). The bright emission feature about 0.25 $`\mathrm{}`$ to the lower right of the quasar center, seen in Figs. 4b and 5b, is just inside the coronagraphic hole and is therefore not seen in Fig. 6b. However, the coronagraph is very effective in reducing the quasar light outside of the coronagraphic hole, and can therefore be used to look at other objects in the field.
A weak feature (which we name object โO2โ) remains after PSF subtraction (to the โtop leftโ of the hole, about $`0.7`$ $`\mathrm{}`$ away from the quasar center). This feature is dominated by four knots of continuum emission. No artifacts resembling this feature have been seen in the coronagraphic images of PSF stars from other NICMOS GTO programs. By contrast, the knots seen to the โlower leftโ are known artifacts in the coronagraphic PSF. โO2โ is detected in the same place if the data for each of the two orbits are analyzed separately, which suggests that it may be real. It is not likely to be a trail of a cosmic ray event, since it is present in the images over a period of two orbits (much longer than typical time-scales for the decay of cosmic ray persistence in the NICMOS detectors). The knots in feature O2 are much weaker than the peak in feature O1, but are about 2-3 times the rms noise in the background. O2 has a total linear size of about 9-10 pixels, i.e. about 0.7-0.8 $`\mathrm{}`$ . The 1 $`\sigma `$ noise levels per pixel in the PSF subtracted image at 0.3 $`\mathrm{}`$, 0.5 $`\mathrm{}`$, 0.7 $`\mathrm{}`$, 0.9 $`\mathrm{}`$, and 1.1 $`\mathrm{}`$ from the quasar center are 0.032 ADU s<sup>-1</sup>, 0.012 ADU s<sup>-1</sup>, 0.011 ADU s<sup>-1</sup>, 0.0088 ADU s<sup>-1</sup>, and 0.0094 ADU s<sup>-1</sup>, i.e., 0.069, 0.027, 0.024, 0.019, and 0.021 $`\mu `$Jy per pixel, respectively. Compared to the non-coronagraphic F160W image, these noise levels indicate factors of 4.06, 1.88, 1.30, 1.41, and 1.26 improvements, respectively, in the 1 $`\sigma `$ sensitivities at 0.3 $`\mathrm{}`$, 0.5 $`\mathrm{}`$, 0.7 $`\mathrm{}`$, 0.9 $`\mathrm{}`$, and 1.1 $`\mathrm{}`$ from the quasar center. These factors are much smaller than those typically reported for NICMOS coronagraphic performance, because the low signal from our faint quasar makes our observations read noise dominated.
The results of coronagraphic imaging (e.g. appearance of object O2) are not expected to be very sensitive to the data reduction procedures. No dithering was used between the coronagraphic exposures, to ensure that the quasar always remained in the coronagraphic hole. Therefore the individual coronagraphic exposures were not registered before they were combined. The quasar was acquired with on-board target acquisition and placed in the coronagraphic hole at the beginning of the first orbit. The quasar was placed in the same position in the second orbit. Guide star acquisition was done at the beginning of each of the two orbits using the same guide stars. Therefore we believe that there are no significant offsets between the quasarโs positions in the hole in the various coronagraphic exposures. Indeed, as mentioned above, the coronagraphic images obtained in each orbit separately show the object O2, which appears similar in both the orbits. The fact that features in the coronagraphic PSF other than object O2 disappear after the PSF subtraction also suggests that object O2 is not the result of misregistration of the individual exposures. We therefore believe that object O2 is likely to be real.
This feature O2 is detected marginally in the non-coronagraphic F160W image (Fig. 4b) due to the higher scattered light from the quasar in that image. We note that, while of lower sensitivity, the faint compact emission features to the top left of the quasar in this image are at positions similar to those of the O2 knots in the coronagraphic image. Object O2 is not seen in the narrow-band image in Fig. 5b. However, considering that it is much fainter than object O1, it is not entirely surprising that any emission from O2 is not detected in the narrow-band images (which are about 10 times less sensitive, at the separation of O2 from the quasar). Considering this, and its faintness and larger angular distance from the quasar compared to O1, it is not completely clear whether the feature O2 has any connection with the DLA. But it may be associated with the DLA or its companions. It is also possible that objects O1 and O2 are associated with the host galaxy of the quasar rather than the DLA. We discuss this possibility further in section 7.4. If O2 is indeed associated with the DLA at $`z=1.892`$, then it has a size of 4-5 h$`{}_{}{}^{1}{}_{70}{}^{}`$ kpc for $`q_0=0.5`$.
## 6 IS OBJECT O1 REAL?
In view of the low S/N of object O1 and its small angular separation from the quasar in our non-coronagraphic F160W and F190N images, we carried out a number of tests on the images to investigate whether O1 is real or merely an artifact of the data reduction or PSF subtraction procedures. Here we describe these tests, listing the potential sources of error that we investigated in each case, and the corresponding results.
### 6.1 Is minimum variance the right criterion in PSF subtraction?
We have registered and normalized the PSF star by varying the position and multiplicative scaling factor of the PSF star so as to minimize the variance in the region of interest near the quasar. This seems to be the most objective way of judging the goodness-of-fit of the PSF subtraction. To determine whether any bias could be caused by the use of the minimum-variance criterion, we also verified that the results from this method are closely consistent with K. McLeodโs method of forcing the intensity at the first Airy minimum to zero (see McLeod, Rieke, & Storrie-Lombardi 1999). The PSF star position given by the two methods for the optimum PSF subtractions in each case agree to within 0.002 pixels. The PSF normalization factors from the two methods agree to within about 2 $`\%`$. In either case, our broad conclusions about the nature of the residuals after PSF subtraction (including feature O1) are the same for both the methods. Therefore we believe that our strategy of minimizing the variance is sound.
### 6.2 Telescope breathing effects?
To examine how sensitive the detection of the main emission knot O1 is to the fine structure of the PSF subtracted, we created a suite of difference images using a variety of reference PSFs for the non-coronagraphic broad-band and narrow-band images. We particularly sought to investigate the effects of HST breathing focus changes on our results. The changes in HST focus consist of two components. First, there is a long-term slow change caused by shrinkage in the Optical Telescope Assembly (OTA) of HST due to moisture desorption, which is periodically corrected by secondary mirror moves. In addition, short-term focus variations on the time scale of the HST orbit, arising from thermally driven displacements of the OTA secondary, can be even larger in magnitude than desorption correction compensations.
Fig. 7 shows the effect of using various observations of the PSF star P330E on the non-coronagraphic F160W images. Fig. 7a is the same as Fig. 4b, while Figs. 7b, 7c, and 7d show the results obtained by subtracting images of P330E taken on different dates and with different breathing values. The breathing values denote the position of the HST secondary mirror in units of $`\mu `$m with respect to a common reference, i.e., with respect to the best focus of WFPC2 planetary camera (see Hershey & Mitchell 1998). All four panels of Fig. 7 show the feature O1 in roughly the same place with other variations being much smaller in amplitude than O1. Fig 8 shows the difference of the PSFs used in making Fig. 7. Figs. 8a, 8b, 8c show, respectively, (PSF for fig. 7a - PSF for fig. 7b), (PSF for fig. 7a - PSF for fig. 7c), and (PSF for fig. 7a - PSF for fig. 7d). The differences among the residuals in the different panels of Fig. 8 arise partly from breathing variations. But we note that the different dates for the reference PSF observations imply the use of different guide stars, and hence the PSF star would have landed on different pixels in these different images. Therefore the intra-pixel response function, in addition to focus changes, could also account for some of the differences between the different PSF images. <sup>1</sup><sup>1</sup>1We further note that the breathing models of Hershey & Mitchell (1998) have some uncertainty. This could give rise to some residuals in our difference images arising from differential inaccuracies in the breathing values for the quasar and the PSF star images predicted by the models. In any case, the symmetric nature of the residuals in Fig. 8 and the absence of the knot O1 in these images suggests that the latter feature is present in the quasar images, and not an artifact in any individual PSF image.
Fig. 9 illustrates the effect of telescope breathing focus variations on the F190N images, with three different observations of PSF star P330E. Fig. 9a is the same as Fig. 5b. The corresponding PSF star differences are shown in Fig. 10. Figs. 10a and 10b show, respectively, (PSF for Fig. 9a - PSF for Fig. 9b), and (PSF for Fig. 9a - PSF for Fig. 9c). The feature O1 is detected in the same place in all the panels of Fig. 9, and most of it is not seen in the PSF star differences (Fig. 10).
### 6.3 Using different PSF stars: Color mismatch between quasar and PSF star?
Color terms in the PSFs are potentially important sources of error in the difference images. For our primary PSF subtractions we have used the solar analogue P330E ($`m_{F110w}m_{F160W}=0.44`$, $`m_{F160W}m_{F222M}=0.08`$) as described in the previous section. However, we also experimented with a red PSF star BRI0021 ($`m_{F110w}m_{F160W}=1.17`$, $`m_{F160W}m_{F222M}=0.80`$). Fig. 11 shows the effect of using four different PSF stars (P330E, BRI 0021, Q1718PSF, and GSC4) in the top left, top right, bottom left, and bottom right panels, respectively. Note that the breathing values for the four observations are quite different, which could explain the differences in the appearance of O1. In any case, all the images show an excess residual at the location of O1, while such excesses are not seen in the differences of the PSFs themselves (shown in Fig. 12). Thus object O1 is probably not an artifact caused by color mismatch between the quasar and the PSF star.
### 6.4 Calibration defects for column 127 influencing the centroids?
Two of our five dither positions for the non-coronagraphic F160W images had the quasar image near column 127 (in camera 2 detector coordinates). This column is well-known to be โphotometrically challengedโ. A โbad stripeโ in this column results if the dark frame used for the calibration is not a perfect match to the dark current in the actual observations. We corrected for the โbad stripeโ in this column by including it in the bad-pixel mask used while imcombining the five dither positions. However, potentially this column may influence the centroids of the images at the two dither positions and hence the centroid of the imcombined image. To explore this possibility, we looked at each of the three remaining dither positions not affected by column 127. Fig. 13 shows central regions of the PSF-subtracted images for these three individual dither positions in top left, top right, and bottom left panels. The bottom right panel shows the result of PSF subtraction for the image obtained by combining only these three dither positions. For obtaining the minimum-variance solutions for these PSF subtractions, we have excluded the PSF cores while determining the variance. The first dither position as well as the combined image (bottom right panel) show an asymmetric excess emission near the position of object O1. This suggests that the feature O1 is not caused by errors arising from column 127, since none of the dither positions considered in Fig. 13 include this column.
### 6.5 Persistence effects from the quasar image at previous dither positions?
The experiments with PSF subtraction on the individual dithers described in test (6.4) above also help to show that image persistence effects are not important in causing feature O1. This is because even the very first dither position (which should not suffer from quasar persistence effects) shows the presence of an asymmetric feature at the location of O1 (see the top left panel of Fig. 13). Also, the quasar LBQS 1210+1731 is faint, so it is not likely to cause persistence effect. The fact that O1 has roughly equal intensity in all dither positions and does not go away even in the final dither position implies that O1 is also not a left over persistence image from a bright object or cosmic ray detected before the start of our observations.
### 6.6 Difference between โcamera 1-2 focusโ vs. โcamera 2 focusโ?
Our quasar images were obtained with the NICMOS internal focusing mechanism optimized for parallel camera 1 and 2 operations, whereas all of our reference PSFs were taken at the camera 2 exclusive focus. A very slight deviation from confocality in the two cameras results in a wavefront error of 0.049 $`\mu `$m mm<sup>-1</sup> of focal dispersion. The โfocus errorโ in camera 2 at the critically sampled $`\lambda `$ (1.75 $`\mu `$m) is $`\lambda /33`$ with the focus at the common position. While small, this โfocus errorโ can affect the fine structure of the PSF to a very small degree, as higher order aberrations also change (with an aggregate power of about half the focus error).
To investigate whether this effect can give rise to feature O1, we took two approaches. In the first approach, we used images of a star actually observed at the compromise focus โcamera 1-2โ. There were no systematic PSF star measurements made at this focus position during the NICMOS calibration program. However, we found a star in one of our images of the galaxy cluster CL0939+47 taken for another NICMOS GTO program. Fig. 14a shows the PSF subtraction results obtained for our quasar image using this observed PSF at the compromise focus. The fact that some of the features in O1 are still seen while some disappear suggests that some of the O1 features (e.g. the blob to the right of the core) could be real.
In the second approach, we constructed simulated NICMOS camera 2 PSFs using the Tiny Tim program (version 4.4, Krist & Hook 1997). We constructed simulated PSFs for the two PAM (pupil alignment mirror) positions corresponding to the two focii at the time of our quasar observations on July 22, 1998 and the P330E observations on July 8, 1998. These two simulated Tiny Tim PSFs differ only in this focus position and both used the same values of rms jitter (0.007 $`\mathrm{}`$), same x and y pixel positions for placement of PSF star center, same pixel size, etc. After making these two simulated PSFs, we corrected them for the slight relative difference in the actual x and y plate scales (interpolated in time for the dates of the observations for LBQS 1210+1731 and the observations for P330E). This slight repixelization corrects for the fact that Tiny Tim creates images with equal x and y pixel scales, whereas the actual x and y pixel scales differ by $`0.9\%`$. The difference of the two simulated Tiny Tim PSFs corrected for the unequal x and y pixel scales is shown in Fig. 14b. The difference does show some residuals along the diagonal directions. However, these are symmetric in shape on both sides of the center, and are accompanied by much larger residuals in the core of the image. Fig. 14c shows, on the same stretch as Fig. 14a, the difference of the two simulated Tiny Tim PSFs after normalizing each to match the quasar. The residuals in Fig. 14c are much weaker than object O1. A simple relative translation between the images for the two focii cannot give rise to a feature as significant as O1 without causing a much larger residual in the core. We therefore conclude that while the difference in the PAM positions for the quasar and the PSF star could cause some of the residuals in our PSF subtractions, they cannot be the major source of these residuals.
To pursue this analysis further, we took the ratio of the two simulated Tiny Tim PSFs after correction for plate-scales <sup>2</sup><sup>2</sup>2Here, by the ratio of the simulated Tiny Tim PSFs, we mean the ratio of the PSF with the PAM position for LBQS 1210+1731 to the PSF with the PAM position for the star P330E, the same two PSFs whose difference is shown Figs. 14b and 14c., and multiplied this ratio by the actual observed P330E PSF to make our โbest-guessโ PSF. The resultant PSF has the advantages of combining the correct focus (PAM) position (because of the Tiny Tim simulation), the best estimate of breathing (because of use of the actual observation of P330E which matches closely in breathing with the LBQS 1210+1731 data), and any other actual optical effects that Tiny Tim does not simulate adequately. In the bottom right panel of Fig. 14, we show the resultant image obtained after subtracting this โbest-guessโ synthetic P330E PSF from the LBQS 1210+1731 data. Once again, excess emission is seen at the position of O1. This suggests that O1 is not caused by artifacts of relative focus difference (โcamera 1-2โ focus vs. โcamera 2โ focus) between the quasar and the PSF star.
### 6.7 Errors in imcombining or interpolating the dithers?
We used Nicred 1.8 to interpolate the registered dithers onto a grid of single camera 2 pixels, or onto a grid of half integer camera 2 pixels. This magnification (repixelization) or lack thereof made little difference in the result. This is clear from Fig. 15, which shows the F160W images made on using grids of single camera 2 pixels (Fig. 15a), and half integer camera 2 pixels (Fig. 15b). Figs. 15c and 15d show the same figures with pixel replication instead of cubic convolution interpolation in the IDP3 display. The similarity between the left and right panels is reassuring and results from camera 2 being nearly critically sampled at 1.6 $`\mu `$m. The same was also found to be true for the F190N images. Fig. 16 shows the F190N images made on using grids of single camera 2 pixels and half integer camera 2 pixels (Figs. 16 a and 16b shown with cubic convolution and Figs. 16c and 16d shown with pixel replication). Both Figs. 15 and 16 suggest that object O1 is likely to be real and does not arise from interpolation errors.
### 6.8 Effects of asymmetries or saturation in the core?
Based on our experience with other NICMOS GTO data, the core of the PSF often shows paired positive and negative residuals after PSF subtraction. In case the core of the quasar and PSF star images have some asymmetries which might mimic an O1-like feature after PSF subtraction, we have also done the PSF subtraction without including the core for the variance calculation. The right panel in Fig. 17 shows the non-coronagraphic F160W image obtained after minimizing the variance in the PSF subtraction without including the core of the image for variance calculation. The region thus excluded, shown with a circular mask, covers the region up to the first Airy minimum. The left panel in Fig. 17 shows the non-coronagraphic F160W image obtained when the core is included (same as the image shown in Fig. 4b, except that the central core is masked in the display only for easy comparison with the right panel of Fig. 17). The similarity between the two panels of Fig. 17 (presence of a feature at the position of O1) suggests that asymmetries in the core are not the source of O1.
This same experiment also shows that saturation in the core of the quasar image cannot be a major problem (since the results obtained by including or excluding the core agree very closely). The individual 512-s exposures at each of the 5 dither positions in the F160W image of our quasar have peaks of about 15000 ADU or 83000 e<sup>-</sup> in the quasar central pixel. Thus we do expect that they should not be saturated, given the $`98\%`$ linearity saturation limit of 173,000 e<sup>-</sup> for camera 2.
### 6.9 Errors in alignment of PSF star with respect to the quasar?
Fig. 18 shows the effects of shifting the PSF star by 0.1 pixel in various directions relative to the quasar on the difference F160W images. Fig. 18e is the optimum minimum-variance solution (same as Fig. 4b). Figs. 18b, 18h, 18d, and 18f show the PSF subtractions obtained after shifting the PSF star by 0.1 pixel in the top, bottom, left, and right directions, respectively, with respect to the quasar. Figs. 18a, 18c, 18g and 18i show the corresponding results on shifting the PSF star by 0.1 pixel in the upper left, upper right, lower left and lower right directions, respectively, with respect to the quasar. The large residuals in the core caused by even the slight shifts illustrate that the PSF star is very well aligned with respect to the quasar in the optimum PSF subtraction (Fig. 18e). We have shown shifts of 0.1 pixel in Fig. 18 to make the changes easier to view. But, judging by the minimum in the variance, we believe that our relative alignment of the quasar and PSF star images is good to at least 0.01 pixel. Thus, errors in alignment of the PSF star with respect to the quasar should be insignificant.
### 6.10 Errors in stacking the individual dither positions?
The different dither positions are registered in Nicred 1.8 using cross-correlation. To check the accuracy of the image registration, we compared the centroids of the images at the various dither positions after registration. The 1 $`\sigma `$ variation among the centroid values of the different dither images was found to be about 0.04-0.06 pixels, for both the quasar and the PSF star. The centroid values for any individual dither position calculated from different methods (tasks imexam, center, and starfind in IRAF) were also found to agree within about 0.06 pixels. Thus, there is a small uncertainty in the centroid values, but it does not seem large enough to cause a feature such as O1. The fact that the individual dither positions show some excess at the position of O1 (Fig. 13) also suggests that O1 is not a spurious feature resulting from stacking errors.
### 6.11 Comparison with other data
(a) Comparison of PSF stars with each other: PSF stars seem to subtract very well from each other, with the same caveat about color and breathing. There is no hint of a $`1\%`$ residual at the position corresponding to the feature O1 (See fig. 12.)
(b) Same reduction on other quasar data: We reduced the data for quasar Q1718+4807 at z=1.084 (a quasar without a DLA absorber) from another NICMOS GTO program, using the same reduction and PSF subtraction procedures as we have used for LBQS 1210+1731. We do not see the object O1 there.
We have also reduced the data for the other quasars with DLAs from our sample, which will be described in separate papers (Kulkarni et al. 2000b, 2000c). Comparing the results for LBQS 1210+1731 with the results for those quasars, we find that some of the residuals in the PSF subtractions appear similar, while some of the features are different. This suggests that part of the emission at the position of O1 is likely to be real, although part of it could be some artifact that we have not yet understood despite the large number of data analysis experiments described above.
### 6.12 Summary of Results from Various Data Analysis Tests
Overall, we conclude that the best-fitting PSF and several others with reasonably close breathing values suggest a possible detection of an object (object โO1โ) located at about 0.25 $`\mathrm{}`$ from the quasar center, in both the F160W and F190N images. The appearance and properties of this object are more sensitive to the important step of PSF subtraction than to other data reduction steps such as flat fielding. However, our extensive tests suggest that this object is not an artifact of color or focus mismatch or spatial misalignment between the quasar and PSF star images. It is also not caused by image persistence or saturation or by the procedures used for interpolation or stacking of the individual images. We therefore believe that object O1 is likely to be real.
The most relevant broad and narrow band summary images showing object O1 are the zoomed, magnified, PSF subtracted images in Figs. 15 (b) and 16 (b). The small angular separation of O1 from the quasar suggests that it is likely to be associated with the DLA absorber. The corresponding impact parameter is 1.5 h$`{}_{}{}^{1}{}_{70}{}^{}`$ kpc for q<sub>0</sub> = 0.5 or 2.0 h$`{}_{}{}^{1}{}_{70}{}^{}`$ kpc for q<sub>0</sub> = 0.1. We have thus probed regions far closer to the quasar sight-line than most previous studies of high-redshift intervening DLAs. Object O1 marks the closest detected high-redshift DLA candidate yet to any quasar sight line. Object O1 is 0.4$`\mathrm{}`$ long. If O1 is the DLA at $`z=1.89`$, this translates into 2.4 h$`{}_{}{}^{1}{}_{70}{}^{}`$ kpc for q<sub>0</sub> = 0.5 or 3.2 h$`{}_{}{}^{1}{}_{70}{}^{}`$ kpc for q<sub>0</sub> = 0.1. It has a luminosity (at mean rest frame wavelength of 0.55 $`\mu `$m) of about $`1.5\times 10^{10}`$ h$`{}_{}{}^{2}{}_{70}{}^{}`$ L for $`q_0=0.5`$ and about $`2.8\times 10^{10}`$ h$`{}_{}{}^{2}{}_{70}{}^{}`$ L for $`q_0=0.1`$. Obejct O1 is thus fainter than an L galaxy at $`z=1.89`$ by 0.2-0.9 magnitudes. The comparison of the broad and narrow band fluxes implies a nominal statistically insignificant SFR of 1.1 h$`{}_{}{}^{2}{}_{70}{}^{}`$ $`M_{}`$ yr<sup>-1</sup>, with a 3 $`\sigma `$ upper limit of 4.0 h$`{}_{}{}^{2}{}_{70}{}^{}`$ $`M_{}`$ yr<sup>-1</sup>, for q<sub>0</sub> = 0.5.
Another fainter object O2 which consists of 4 knots of continuum emission is also seen in our images. (See Fig. 6b.) This object, at angular separation of 0.65 $`\mathrm{}`$ from the quasar (well outside the first Airy ring of the quasar PSF) is also not a known artifact of the PSF. It is thus also likely to be real and may be a companion to the DLA. The spatial extent of O2 is 4-5 h$`{}_{}{}^{1}{}_{70}{}^{}`$ kpc and its projected impact parameter is 3.8 h$`{}_{}{}^{1}{}_{70}{}^{}`$ kpc. Object O2, like object O1, is also closer to the quasar sightline than most other high-redshift DLA candidates detected before.
We note, however, that because of the faintness and proximity of O1 to the quasar, we cannot completely rule out the possibility that this feature could partly be some as yet unknown artifact of the PSF (that is not simulated by Tiny Tim either). If any such errors are the actual cause of O1, then the DLA absorber and the quasar host galaxy are even fainter than O1. In that case, we can use our images to put very sensitive upper limits on the size and brightness of both the DLA absorber and the quasar host. We discuss the implications of our observations in the following section.
## 7 DISCUSSION
The most important result from our observations is that there are no large bright galaxies close to the quasar in the field of the DLA absorber toward LBQS 1210+1731. Feature O1 is the most likely candidate for any object associated with the DLA. In subsections 7.1, 7.2, and 7.3, we assume that object O1 is associated with the DLA to derive constraints on various properties of DLAs. But we also consider alternative possibilities in subsection 7.4, mainly the possibility that O1 may be associated with the host galaxy of the quasar.
### 7.1 CONSTRAINTS ON SIZES AND MORPHOLOGY OF DLAs
Our observations show no evidence for a big, well-formed galaxy as expected in some scenarios for the DLAs \[e.g., the proto-spiral model suggested by Wolfe et al. (1986), Prochaska & Wolfe (1997, 1998), Jedamzik & Prochaska (1998)\]. Feature O1 has an estimated size of 2-3 $`h_{70}^1`$ kpc, while feature O2, if real, consists of small knots spread over about 4-5 $`h_{70}^1`$ kpc. Thus, these data suggest that the absorber is compact and clumpy, as expected in the hierarchical picture of galaxy formation. However, it is hard to be completely certain of the morphology, partly because of the sensitivity of the detailed image structure to the various factors discussed in sec. 6. Furthermore, it is possible that O1 and O2 are the brightest regions within a bigger galaxy, the rest of which we cannot see. Thus, we cannot completely rule out the large disk scenario, although the compact sizes and low SFRs suggest that the hierarchical picture may be favored. Analysis of the other DLAs from our sample and further deeper observations will help to more definitively distinguish between the large disk vs. hierarchical models.
### 7.2 CONSTRAINTS ON ENVIRONMENT OF DLA ABSORBERS
Apart from features O1 and O2 very close to the quasar, our images show two prominent galaxies in the non-coronagraphic F160W image (one in the upper left corner or west of the quasar and the other at the middle of the left edge of the image or roughly north of the quasarโ see Fig. 1). There is also a third very weak feature to the left (roughly north) of the quasar, a little less than half the way along the line joining the quasar and the galaxy at the middle left edge. The galaxy west of the quasar is just barely apparent in the non-coronagraphic F190N image, while the other two objects are not seen in the non-coronagraphic F190N image. The two prominent galaxies in the non-coronagraphic F160W image are off the field of the coronagraphic image, while the prominent galaxy seen at the bottom edge (northeast) of the coronagraphic F160W image is off the field of the non-coronagraphic images. It is possible that the faint feature to the left (north) of the quasar is spurious. But on running maximum-entropy and Lucy deconvolutions of the images, all the three objects (including the faint feature) in the F160W image were found to remain significant. These objects are likely to be galaxies in the same group as the DLA, although we do not have redshift information on them. In any case, they have fairly large impact parameters (4.52 $`\mathrm{}`$, 11.00 $`\mathrm{}`$, and 10.96 $`\mathrm{}`$ for the faint feature, the galaxy to the west of the quasar, and the galaxy to the north of the quasar, respectively). At the redshift of the DLA absorber, these impact parameters would correspond to 26.8, 65.1 and 64.9 h$`{}_{}{}^{1}{}_{70}{}^{}`$ kpc respectively, for q$`{}_{0}{}^{}=0.5`$. For q$`{}_{0}{}^{}=0.1`$, the corresponding values are 36.6, 89.0, and 88.6 h$`{}_{}{}^{1}{}_{70}{}^{}`$ kpc. These large values make it unlikely for any of these features to be the DLA absorber itself.
If real, the continuum emission knots in object O2 may be highlighting the brightest regions in a companion to the DLA galaxy. The roughly filamentary morphology may indicate an edge-on disk galaxy or a part of a spiral arm. Alternatively, it may suggest individual star-forming sub-galactic clumps formed in a filamentary over-dense region, similar to the filamentary arrangements of galaxies and sub-galactic units found in numerical simulations of structure formation. It is interesting to note that the WFPC2 observations of a $`z=2.811`$ DLA by Mรธller & Warren (1998) also indicate a filamentary arrangement of 3 bright objects, although on a much larger scale (separation of 21 $`\mathrm{}`$ ). The angular separation of their closest object from the quasar was 1.17 $`\mathrm{}`$, whereas for our features O1 and O2, the angular separations are $`0.26`$โณ and $`0.65`$โณ, respectively. \[We note, however, an important difference between our DLA and the DLA studied by Mรธller & Warren. The latter has a redshift very close to that of the quasar ($`z_{em,CIV}=2.77`$, $`z_{em,[OIII]}=2.788`$, and $`z_{em,H\alpha }=2.806`$). Therefore it is likely to be associated with the quasar and may not be representative of DLA galaxies in general.\]
The 1 $`\sigma `$ noise levels far away from the quasar are 0.011 ADU s<sup>-1</sup> per pixel for our PSF-subtracted non-coronagraphic F160W image and 0.0088 ADU s<sup>-1</sup> per pixel for the PSF-subtracted coronagraphic F160W image. These levels translate into 0.024 $`\mu `$Jy per pixel and 0.019 $`\mu `$ Jy per pixel respectively. The corresponding 1 $`\sigma `$ noise equivalent magnitudes for the non-coronagraphic and coronagraphic F160W images are 26.6 magnitudes per pixel (21.0 magnitudes per square arcsecond) and 26.9 magnitudes per pixel (21.3 magnitudes per square arcsecond), respectively. For comparison, the Hubble Deep Field F160W images had a 1 $`\sigma `$ noise level of $`1.22\times 10^9`$ Jy per Camera 3 pixel (Thompson et al. 1999). Thus, for the field galaxies far from the quasar in the PSF-subtracted F160W observations, our images are about 5.1-5.4 magnitudes less deep than the Hubble Deep Field images. <sup>3</sup><sup>3</sup>3Before doing the PSF subtraction, the 1 $`\sigma `$ noise levels far away from the quasar are $`0.0046`$ ADU s<sup>-1</sup> per pixel (0.010 $`\mu `$Jy per pixel or 27.6 mag per pixel) for the non-coronagraphic F160W image and $`0.0084`$ ADU s<sup>-1</sup> per pixel (0.018 $`\mu `$Jy per pixel or 26.9 mag per pixel) for the coronagraphic F160W image. The process of PSF subtraction decreases the 1 $`\sigma `$ deviations by a large factor near the quasar, but increases the noise far away from the quasar. This is because of the use of actually observed PSF star images (with high but finite S/N) for PSF subtraction, which contribute to the noise level. But for reasons mentioned earlier, it is better to use observed PSF star images rather than Tiny Tim models to get good matches to the quasar PSF. The higher noise level in the coronagraphic F160W image before PSF subtraction compared to the non-coronagraphic F160W image seems to arise from the use of the target acquisition flat rather than the higher-S/N standard flat used for the non-coronagraphic image. In any case, our images both before and after PSF subtraction do not show any field galaxies other than those mentioned above.
Our images do not show any objects other than object O2 in the close vicinity of the DLA. From the galaxy number count-magnitude relation based on deep NICMOS images (Yan et al. 1998), about 1 galaxy is expected for $`H<21`$ in the camera 2 field. Thus our observations are consistent with these predictions within the uncertainties. There is no sign of strong clustering of galaxies around the DLA.
### 7.3 CONSTRAINTS ON STAR-FORMATION RATE AND DUST IN DLAs
It is quite surprising, given the high sensitivity of our observations and the reasonably high rest-frame V-band luminosity of object O1, that O1 shows almost no detectable H-$`\alpha `$ emission. The lack of significant H-$`\alpha `$ emission in our images puts fairly tight constraints on the star-formation rate in the DLA toward LBQS 1210+1731, i.e. a $`3\sigma `$ upper limit of $`4.0`$ h$`{}_{}{}^{2}{}_{70}{}^{}`$ M yr<sup>-1</sup> for $`q_0=0.5`$, if no dust is assumed. This is by far the most severe existing constraint on the SFR in high-$`z`$ DLAs. For comparison, the near-IR spectroscopic survey of Bunker et al. (1999), aimed at detecting H$`\alpha `$ from DLAs, gave typical upper limits of $``$ 15 M yr<sup>-1</sup>, for q$`{}_{0}{}^{}=0.5`$ and $`H_0=70`$ km s<sup>-1</sup> Mpc<sup>-1</sup>. In Fig. 19, we compare the result from our data (shown as a filled triangle) with the $`3\sigma `$ upper limits from Bunker et al. (1999) (shown as unfilled triangles). Our limit on the SFR marks an improvement by a factor of 3 over the tightest constraints of Bunker et al. (1999) on the SFR in DLA galaxies. The curve in Fig. 19 shows the predicted average SFR(z) in a DLA expected if DLAs are proto-disks, as derived by Bunker et al. (1999) using the closed-box model of Pei & Fall (1995) for the global star formation rate. It is clear that our upper limit on the SFR is much lower than the predicted value at $`z=1.89`$. We note that the low SFR estimated here is consistent with the result of Djorgovski (1997) who reported SFR of $`0.7`$ M yr<sup>-1</sup> in the $`z=4.1`$ DLA toward DMS 2247-0209, on the basis of a weak Ly-$`\alpha `$ emission line (assuming no dust extinction). (We note, however, that our limit is less sensitive to dust extinction uncertainties owing to the use of H-$`\alpha `$ rather than Ly-$`\alpha `$ emission.)
In principle, the lack of detectable H-$`\alpha `$ emission from the DLA could be because of dust extinction, in which case the actual SFR could be higher. In order to reconcile our upper limit of 4.0 M yr<sup>-1</sup> for q$`{}_{0}{}^{}=0.5`$, $`H_0=70`$ km s<sup>-1</sup> Mpc<sup>-1</sup>, with the expectation of the closed-box proto-disk model of 38.6 M yr<sup>-1</sup>, an optical depth $`\tau _{0.66\mu m}2.3`$ would be required at the rest frame H-$`\alpha `$ line, if a simple screen of dust in front of the DLA is assumed to extinguish the H-$`\alpha `$ emission. For extinction curves similar to those in the Milky Way, the Small Magellanic Cloud, or the Large Magellanic Clouds, this would imply $`\tau _B3.4`$ at $`\lambda _B=4400`$ ร
. To have such high extinction, the DLA would be required to have a mean dust-to-gas ratio $`k\tau _B(10^{21}/N_{HI})8.7`$. Even if the HI column density is assumed to be higher by a factor of $`3`$ at the position of O1 compared to the $`N_{HI}`$ detected in the DLA line (since the projected separation of O1 from the quasar would indicate that the DLA absorbing region may be a scale length away from the peak of emission from O1), one still requires a mean dust-to-gas ratio $`k3`$. This is much higher than the mean dust-to-gas ratio of 0.8 for the Milky Way, or the typical value of $``$ 0.03-0.1 for the DLA galaxies, suggested by observations of background quasar reddening and heavy element depletions (see, e.g., Pei et al. 1991; Pettini et al. 1997 and references therein).
It is, however, possible that the dust may be intermingled with the gas very close to the stars in the DLA. However, given the low dust-to-gas ratios seen in DLA absorbers, it is hard to imagine that most of the Ly-continuum photons could be absorbed even before H-$`\alpha `$ photons can be produced. Thus, there is a good chance that the lack of H-$`\alpha `$ emission could be indeed because of low SFR.
In the absence of dust obscuration, it follows from Fig. 19 that our results, as well as those of Bunker et al. (1999), indicate SFRs much lower than the expectations of the proto-disk model. This together with the compact sizes seen in our images again suggests that the observations do not agree with the proto-disk models. It is possible that O1 is a dwarf galaxy. Star formation in dwarf galaxies is inferred to proceed in bursts separated by quiescent periods lasting up to several Gyr (e.g., Grebel 1998). It may be that we are observing object O1 during such a relatively quiescent stage. It is also possible that O1 is a low-surface brightness galaxy, since such galaxies show lower SFR.
### 7.4 ALTERNATIVE POSSIBILITIES
Finally, it is possible that object O1 is not the DLA absorber, but that it arises mostly in the quasar host galaxy. We cannot test this possibility further because we do not have narrow-band images in filters tuned to $`z_{em}=2.543`$. However, we cannot rule out this possibility either. If O1 is in fact the host galaxy of the quasar, then it would have a luminosity (at rest frame 0.45 $`\mu `$m) of $`2.9\times 10^{10}`$ $`h_{70}^2`$ L for $`q_0=0.5`$ or $`6.4\times 10^{10}`$ $`h_{70}^2`$ L for $`q_0=0.1`$. The images would then suggest that the quasar host is not a large galaxy with or without interactions, but rather shows a compact morphology. The strongest feature in the quasar host would then be off-center with respect to the quasar nucleus, which has been observed in other quasars. If O1 is in fact the quasar host, then the limits on the luminosity and SFR in the DLA are even more severe than our estimates in sections 5.1 and 5.2. Conversely, if O1 is the DLA galaxy, then the constraints on the quasar host are more severe than those given above.
It is also possible that O1 is an interloper galaxy at an even lower redshift than the DLA. However, there is no spectroscopic evidence available for this based on the available spectra. Ultraviolet archival spectra with HST or IUE (which would contain any potential DLA line at a lower redshift) are not available, while the ground-based optical spectra are only medium-resolution. We therefore do not consider this possibility further.
## 8 CONCLUSIONS AND FUTURE WORK
With deep diffraction-limited NICMOS images of LBQS 1210+1731, we have probed regions far closer to the quasar sight-line than in most previous studies of high-redshift intervening DLAs. The two objects we report mark the closest detected high-redshift DLA candidates yet to any quasar sight line. Our continuum and H$`\alpha `$ images of the $`z=1.89`$ DLA toward LBQS 1210+1731 suggest that this DLA is not a big galaxy with high SFR, but may be compact (2-3 $`h_{70}^1`$ kpc in size), probably consisting of multiple sub-units. Assuming no dust extinction of H-$`\alpha `$ emission, we place a 3 $`\sigma `$ upper limit of 4.0 $`h_{70}^2`$ M yr<sup>-1</sup> on the star formation rate, for $`q_0=0.5`$ . Our continuum and H$`\alpha `$ observations are consistent with the hierarchical models, in which DLAs arise in several sub-galactic clumps or dwarf galaxies, which eventually come together to form the present-day galaxies (see, e.g., York et al. 1986; Matteucci et al. 1997). Indeed, theoretical simulations of merging proto-galactic fragments in cold dark matter cosmologies (e.g., Haehnelt et al. 1998), low surface brightness galaxies (e.g., Jimenez et al. 1999), and collapsing halos with merging clouds (e.g., McDonald & Miralda-Escudโe 1999) have also been found to reproduce the observed properties of DLAs (asymmetric line profiles of metal absorption lines, metallicities, H I content etc.) The small sizes of high-$`z`$ DLAs suggested by our observations are also consistent with the small sizes of galaxies seen in other independent high-redshift observations, e.g., in the NICMOS Hubble Deep Field observations (Thompson et al. 1999). Together, these observations may be indications that while star formation had begun long before $`z=2`$ resulting in some chemical enrichment, most of the dynamical assembly of galaxies as we know them today occurred more recently, and at $`z2`$, the various constituent units were still coming together. However, it cannot be ruled out that the DLA toward LBQS1210+1731 is a large low surface brightness galaxy with a low SFR, which is below our detection limit even in the F160W image.
We point out that our conclusions are, nevertheless, based on detailed observations of only one high-$`z`$ DLA. It is quite possible that different DLAs have different rates of evolution because of different physical conditions. Indeed, this is suggested by the large scatter in the metallicity-redshift relation of DLAs (see, e.g., Pettini et al. 1999 and references therein). The NICMOS observations of other DLAs from our sample are currently being analyzed and will help to explore the generality of our conclusions. To improve the statistics of the DLA imaging studies, it is necessary to obtain high spatial resolution near-IR images of more high-redshift DLAs. It would be very valuable to carry out a deeper near-IR imaging survey of more DLAs with HST, if the NICMOS cryocooler or the near-IR channel of WFC3 becomes available in the near future. A major advantage of such HST observations will be a relatively stable PSF compared to that currently achieved with any ground-based telescope, which is crucial for the detection of DLAs. It will also be of great interest to complement the HST observations with observations from adaptive optics systems on large ground-based telescopes. Although these systems will not initially have the relatively stable PSF offered by HST, they will be able to achieve even higher spatial resolution and higher imaging sensitivity. Such future space and ground-based observations will provide further insight into the structure and nature of DLA galaxies, and thereby help to constrain theoretical models of the formation and evolution of galaxies.
###### Acknowledgements.
This project was supported by NASA grant NAG 5-3042 to the NICMOS Instrument Definition Team. It is a pleasure to thank Nicholas Bernstein and Keith Noll for their assistance in the scheduling of our observations. We thank Elizabeth Stobie, Dyer Lytle, Earl OโNeil, Irene Barg, and Anthony Ferro for software and computer support. We also thank Andrew Bunker for making his model star formation rate versus redshift curves available to us ahead of publication.
FIGURE CAPTIONS
FIG. 1โ NICMOS camera 2 non-coronagraphic 1.6 $`\mu `$m broad-band image of the field of LBQS 1210+1731. The color scheme is indicated with the bar on the bottom of the image. The flux scale in ADU s<sup>-1</sup> is indicated on the color bar. Image Y axis is -121.961 degrees east of north.
FIG. 2โ NICMOS camera 2 non-coronagraphic 1.9 $`\mu `$m narrow-band image of the field of LBQS 1210+1731. Image Y axis is -121.961 degrees east of north.
FIG. 3โ NICMOS camera 2 coronagraphic 1.6 $`\mu `$m broad-band image of the field of LBQS 1210+1731. The quasar has been placed in the coronagraphic hole. Image Y axis is -123.987 degrees east of north.
FIG. 4โ Zoomed-in $`2.74`$$`\times 2.71`$โณ region of the NICMOS camera 2 non-coronagraphic 1.6 $`\mu `$m broad-band image of the field of LBQS 1210+1731, (a) before PSF subtraction (top), (b) after PSF subtraction (bottom). The residual feature is labeled as O1 in the bottom panel.
FIG. 5โ Zoomed-in $`2.69`$$`\times 2.66`$โณ region of the NICMOS camera 2 non-coronagraphic 1.9 $`\mu `$m narrow-band image of the field of LBQS 1210+1731, (a) before PSF subtraction (top), (b) after PSF subtraction (bottom). The residual feature is labeled as O1 in the bottom panel.
FIG. 6โ Zoomed-in $`2.66`$$`\times 2.71`$โณ region of the NICMOS camera 2 coronagraphic 1.6 $`\mu `$m broad-band image of the field of LBQS 1210+1731, (a) before PSF subtraction (top), (b) after PSF subtraction (bottom). The four residual features are labeled as O2 in the bottom panel.
FIG. 7โ Effect of HST โbreathingโ focus variations on the PSF subtracted F160W non-coronagraphic image. Zoomed-in $`2.74`$$`\times 2.71`$โณ region of the NICMOS camera 2 non-coronagraphic 1.6 $`\mu `$m broad-band image of the field of LBQS 1210+1731, on using images from 4 different observations of the PSF star P330E. The PSF star observation dates and breathing values are (July 8, 1998; 1.0), (August 9, 1998; 1.2), (September 7, 1998; 0.7), and (March 7, 1998; -1.7) respectively for the (a) top left, (b) top right, (c) bottom left, and (d) bottom right panels. The quasar observations were obtained on July 22, 1998 at breathing value of 2.2.
FIG. 8โ Differences of PSFs used in Fig. 7. (a) PSF for Fig. 7a - PSF for Fig. 7b (top left panel), (b) PSF for Fig. 7a - PSF for Fig. 7c (top right), (c) PSF for Fig. 7a - PSF for Fig. 7d (bottom left).
FIG. 9โ Effect of HST breathing focus variations on the PSF subtracted F190N non-coronagraphic image. Zoomed-in $`3.04`$$`\times 3.01`$โณ region of the NICMOS camera 2 non-coronagraphic 1.9 $`\mu `$m narrow-band image of the field of LBQS 1210+1731, on using images from 3 different observations of the PSF star P330E. The PSF star observation dates and breathing values are (May 29, 1998; 1.5), (March 7, 1998; 1.3), and (July 8, 1998; 0.7), respectively for the (a) top left, (b) top right, and (c) bottom left panels. The quasar observations were obtained on July 22, 1998 at mean breathing value of 2.3.
FIG. 10โ Differences of PSFs used in Fig. 9. (a) PSF for Fig. 9a - PSF for Fig. 9b (top left panel), (b) PSF for Fig. 9a - PSF for Fig. 9c (top right).
FIG. 11โ Effect of using different PSF stars on the PSF subtracted F160W non-coronagraphic image. Zoomed-in $`2.74`$$`\times 2.71`$โณ region of the NICMOS camera 2 non-coronagraphic 1.6 $`\mu `$m broad-band image of the field of LBQS 1210+1731, on using PSF stars P330E, BRI 0021, Q1718PSF, and GSC4. The PSF star observation dates and average breathing values are (July 8, 1998; 1.0), (December 20, 1997; -1.2), (July 21, 1998; 1.8) and (Nov. 11, 1997; -3.2), respectively for the (a) top left, (b) top right, (c) bottom left, and (d) bottom right panels. The quasar observations were obtained on July 22, 1998 at mean breathing value of 2.3. The differences between the different panels may be largely because of breathing differences. (The match with the quasar breathing value is poor for panels b and d, while it is best for panel c.) The negative feature near the right edge of panel (c) is because of a second star near the main PSF star used for subtraction.
FIG. 12โ Differences of PSFs used in Fig. 11. (a) PSF for Fig. 11a - PSF for Fig. 11b (top left panel), (b) PSF for Fig. 11a - PSF for Fig. 11c (top right), and (c) PSF for Fig. 11a - PSF for Fig. 11d (bottom left). The negative feature near the right edge of panel (b) is because of a second star near the main PSF star used for subtraction in panel (c) of Fig. 11.
FIG. 13โ Effect of using different individual dither positions on the PSF subtracted F160W non-coronagraphic image. Zoomed-in $`3.12`$$`\times 3.09`$โณ region of the field of LBQS 1210+1731, on using position 1, position 4, position 5 in the spiral-dither pattern (top left, top right and bottom left panels respectively). The bottom right panel shows the result of combining the three positions.
FIG. 14โ Effect of different focus positions on the PSF subtracted F160W non-coronagraphic image. Zoomed-in $`3.12`$$`\times 3.09`$โณ region of the NICMOS camera 2 non-coronagraphic 1.6 $`\mu `$m broad-band image of the field of LBQS 1210+1731. (a) The top left panel shows the PSF subtraction obtained on using PSF star image from the field of galaxy cluster CL0939+47 which has the same โcamera 1-2โ focus as our quasar data. (b) The top right panel shows the difference of two simulated Tiny Tim PSFs corresponding to the โcamera 2 focusโ and โcamera 1-2โ focus. (c) The bottom left panel shows, on the same intensity stretch as the top left panel, the difference of the two simulated Tiny Tim PSFs after normalizing each to match the quasar. (d) The bottom right panel shows the quasar image after subtracting a synthetic PSF made by multiplying the July 8, 1998 image of P330E by the ratio of the Tiny Tim models for the two focii. See the text for details.
FIG. 15โ Effect of different magnification schemes in data reduction and different interpolation schemes in image display on the PSF subtracted F160W non-coronagraphic image. Zoomed-in $`3.12`$$`\times 3.09`$โณ region of the F160W image for (a) no magnification in image analysis, bicubic interpolation in image display (top left), (b) magnification by a factor of 2 and bicubic interpolation in image display (top right), (c) no magnification and pixel replication in image display (bottom left), (d) magnification by a factor of 2 and pixel replication in image display (bottom right). Note the similarities between the magnified and unmagnified images, resulting from camera 2 being almost critically sampled at 1.6 $`\mu `$m.
FIG. 16โ Effect of different magnification schemes in data reduction and different interpolation schemes in image display on the PSF subtracted F190N non-coronagraphic image. Zoomed-in $`3.12`$$`\times 3.09`$โณ region of the F190N image for (a) no magnification in image analysis, bicubic interpolation in image display (top left), (b) magnification by a factor of 2 and bicubic interpolation in image display (top right), (c) no magnification and pixel replication in image display (bottom left), (d) magnification by a factor of 2 and pixel replication in image display (bottom right).
FIG. 17โ Effect of including or excluding the image core on the PSF subtracted F160W non-coronagraphic image. Zoomed-in $`2.74`$$`\times 2.71`$โณ region of the F160W image of the field of LBQS 1210+1731 obtained on (a) including the region indicated by the circular mask (left panel) and (b) excluding the region indicated by the circular mask (right panel). Note the similarity between the two panels.
FIG. 18โ Effect of shifting the PSF star by 0.1 pixel in various directions relative to the quasar on the PSF subtracted non-coronagraphic F160W image. The central panel (e) is the optimum minimum-variance solution (same as Fig. 4b). Top central (b) and bottom central (h) panels correspond to PSF star shifts of 0.1 pixel in +y and -y directions. Left central (d) and right central (f) panels correspond to PSF star shifts of 0.1 pixel in -x and +x directions. Top left (a), top right (c), bottom left (g), and bottom right (i) panels correspond to PSF star shifts of 0.1 pixel in the top left, top right, bottom left, and bottom right directions, respectively. The large residuals caused by the slight shifts illustrate how well centered the PSF star is with respect to the quasar in the optimum PSF subtraction. The regions shown are 2.74 $`\mathrm{}\times 2.71\mathrm{}`$ regions around the quasar.
FIG. 19โ Mean star formation rate in DLAs in M yr<sup>-1</sup> as a function of redshift, for $`q_0=0.5`$, $`H_0=70`$ km s<sup>-1</sup> Mpc<sup>-1</sup>. The filled triangle shows the upper limit from this work, while the unfilled triangles show the limits from Bunker et al. (1999). The curve shows the prediction from a closed-box model applied to proto-disk galaxies, as calculated by Bunker et al. (1999). Note that our SFR limit is a factor of 3 improvement over the tightest limits of Bunker et al. (1999), and that most data points are inconsistent with the proto-disk model. |
warning/0002/gr-qc0002001.html | ar5iv | text | # Nonminimal Global Monopoles and Bound Orbits
Abstract
We perform a numerical analysis of the gravitational field of a global monopole coupled nonminimally to gravity, and find that, for some given nonminimal couplings (in contrast with the minimal coupling case), there is an attractive region where bound orbits exist. We exhibit the behavior of the frequency shifts that would be associated with โrotation curvesโ of stars in circular orbits in the spacetimes of such global monopoles.
PACS number(s): 11.27.+d, 04.40.-b, 98.62.Gq
An interesting property of global monopoles is that they are configurations with energy density decreasing with the distance as $`r^2`$ . This is very suggestive in view of the fact that, naively at least, this is precisely what seems to be required to provide a natural explanation for the flatness of the rotation curves in spiral galaxies. Moreover by assuming the existence of a global monopole in a typical galaxy the total Newtonian mass contribution of the portion of the global monopole contained within $`r_{gal}`$ is found to be $`M\alpha r_{gal}/2`$, where $`\alpha =8\pi G\eta ^2`$ is the deficit angle of the global monopole spacetime. If we take the radial extent of a galaxy to be $`r_{gal}15`$ Kpc and consider a typical grand unified theory with a symmetry breaking scale $`\eta 10^{16}`$ GeV, this mass turns out to be $`10^{69}`$ GeV which is ten times the total mass due to the $`10^{11}`$ solar-mass stars in the typical galaxy: $`M_{stars}10^{68}`$ GeV. This again seems to be precisely what is called for in the observations . Finally, and what seems to be an added bonus, if we assume that the field of the monopole extends on average a distance of 10 galactic radii from the galaxy, where the configuration presumably coincides with that of the monopole centered in the neighboring galaxy, the total mass associated with each monopole (and thus with each galaxy) turns out to be 100 times that of the galaxy. Thus the total contribution of the monopoles to the average density in the universe is of the right order of magnitude to account for the inflationary prediction of a universe with critical density. This is in fact the basic argument that allows one to place bounds on the number density of global monopoles present in the universe , the exact nature of which depends of course on the value taken for $`\eta `$. Indeed global monopoles have been shown to be able to trigger inflation on their own for suitably high $`\eta `$ , , .
Unfortunately this picture doesnโt hold up to scrutiny at the next level of analysis because of two problems: 1) upon substitution of Newtonian gravity by general relativity, it turns out , that the linearly divergent mass has, at large distances, an effect analogous to that of a deficit solid angle plus that of a tiny mass at the origin. The mass of this โmonopole coreโ is about $`0.8\alpha `$. In fact Harari and Loustรณ showed that this small core mass is negative and produces a repulsive potential. They studied the motion of test particles, in the large distance region, in the spacetime of a global monopole, concluding that there are no bound orbits; 2) the fact that the monopole configuration is rather unique, in the sense that it is basically independent of the ordinary matter content in the corresponding galaxy, which conflicts with the fact that there is a rather large range of galactic masses for which the dark matter component is about 10 times more massive than the ordinary matter component .
In spite of these remarks the initially mentioned results are in our view still very suggestive, specially if we look at the remarkable degree of universality that seems to emerge from the systematic study of galactic dynamics : 1) the ordinary matter content is such a good indicator of the mass of the dark matter component; 2) the universality of the form of the radial distribution of the dark matter content; 3) other seemingly coincidental facts as scaling properties between dark and luminous galactic structure parameters. We, therefore, believe that the general idea of monopoles as seeds of galaxy formation, and models for the galactic (and perhaps cosmological) dark matter deserve to be further explored.
The purpose of this Letter is to illustrate the possibilities in these models by considering the simplest nontrivial modification of the global monopole picture: nonminimally coupled global monopoles. In this work we will show that the situation regarding the problems alluded to above, can change dramatically with the introduction of a nonminimal coupling.
Letโs consider the simplest such possibility, namely a theory of a triplet of scalar fields $`\varphi ^a`$, $`a=1,2,3`$, nonminimally coupled (NMC) to gravity with global O(3) symmetry which is broken spontaneously to U(1). The Lagrangian for the model is
$$=\sqrt{g}\left(\frac{1}{16\pi }R+F(R,\varphi ^a\varphi ^a)\right)\sqrt{g}\left[\frac{1}{2}(\varphi ^a)^2+V(\varphi ^a\varphi ^a)\right].$$
(1)
where $`V(\varphi ^a\varphi ^a)`$ is a potential that results in the breaking of the symmetry and is usually taken to be the Mexican hat potential $`V(\varphi ^a\varphi ^a)=\frac{\lambda }{4}(\varphi ^a\varphi ^a\eta ^2)^2`$, and $`F`$ is an arbitrary function of two variables. Note that in these theories the resulting potential for the scalar fields (the terms in the Lagrangian that are independent of the scalar field gradients) is
$`V(\varphi ^a\varphi ^a)_{res}=V(\varphi ^a\varphi ^a)F(R,\varphi ^a\varphi ^a).`$ (2)
Therefore, the ordinary matter would affect the location of the minimum of the effective potential through its effect on $`R`$, thus avoiding the scenario where the monopole configuration is universal. This opens the possibility to recover the correlation between the masses in the dark and ordinary matter components of the galaxy that were alluded before.
Next we focus on the existence or lack of existence of bound orbits in the corresponding spacetimes. To do this we will look at a specific example corresponding to the choice $`F=(\xi \varphi ^a\varphi ^a)R`$ (where $`\xi `$ is a constant) and proceed to construct the resulting spacetime and analyze the motion of test particles. We will specialize to the case of a spherical and static solution to the coupled field equations. The configuration describing a monopole is taken as usual
$`\varphi ^a`$ $`=`$ $`\eta f(r){\displaystyle \frac{x^a}{r}},`$ (3)
with $`x^ax^a=r^2`$, so that we will actually have a monopole solution if $`f1`$ at spatial infinity. We adopt the following metric:
$$ds^2=N^2(r)dt^2+A^2(r)dr^2+r^2d\theta ^2+r^2\mathrm{sin}^2\theta d\phi ^2,$$
(4)
and analyze solutions of the gravitational and scalar fields equations describing global monopole configurations and the resulting space-time. Owing to the complexity of the resulting equations, we will perform a numerical analysis. To do so we adopt the following variables:
$`\nu (\stackrel{~}{r})`$ $`=`$ $`\mathrm{ln}[N(\stackrel{~}{r})],`$ (5)
$`A(\stackrel{~}{r})`$ $`=`$ $`\left(1\alpha {\displaystyle \frac{2m(\stackrel{~}{r})}{\stackrel{~}{r}}}\right)^{1/2},`$ (6)
where
$$\alpha =\frac{\mathrm{\Delta }}{1+2\xi \mathrm{\Delta }},\mathrm{\Delta }=8\pi \eta ^2,$$
(7)
where we have introduced the dimensionless quantity $`\stackrel{~}{r}r/r_c\eta \lambda ^{1/2}r`$.
The solutions are obtained by the standard shooting method with boundary conditions ensuring a regular origin, i.e., $`f(0)=0`$, $`m(0)=0`$, $`m(0)_{,\stackrel{~}{r}}=\alpha /2`$, and that the spacetime is asymptotically flat but for a deficit angle , i.e., $`f(\mathrm{})=1`$, $`\nu (\mathrm{})=\frac{1}{2}\mathrm{ln}\left[1\alpha \right]`$. The โshooting parameterโ is taken to be $`f(0)_{,\stackrel{~}{r}}`$.
Then we calculate the ADM (gravitational) mass of these configurations, (see for a rigorous definition of the ADM mass for spacetimes that are asymptotically flat but for a deficit angle) which can be easily evaluated from the integral
$$M_{\mathrm{ADM}\alpha }=\mathrm{lim}_{\stackrel{~}{r}\mathrm{}}m(\stackrel{~}{r})=_0^{\mathrm{}}\left(4\pi \stackrel{~}{r}^2T_t^t(\stackrel{~}{r})\frac{\alpha }{2}\right)๐\stackrel{~}{r}.$$
(8)
For all values of $`\xi `$ that we analyzed the resulting ADM mass was negative as in the minimal coupling case (here we must remind the reader that the ADM mass for quasi-asymptotically flat spacetimes does not share the positiveness property of the ADM mass of asymptotically flat spacetimes). For $`\xi =2`$, for example (corresponding to the configuration we will analyze in more detail) the resulting ADM mass is $`2.5\alpha `$.
We now study the motion of test particles around a global monopole. To do this we consider the โequatorial geodesicsโ for the metric (4) (i.e., geodesics lying on the plane $`\theta =\pi /2`$). We denote the 4-velocity of a test particle as $`u^\mu =(\dot{t},\dot{r},0,\dot{\phi })`$, where the dot stands for derivation with respect to an affine parameter which in the case of a massive particle can be taken as the proper time. The resulting equation is
$$\frac{\mathrm{exp}[2\nu (\stackrel{~}{r})]}{\left[1\alpha \frac{2\stackrel{~}{m}(\stackrel{~}{r})}{\stackrel{~}{r}}\right]}(\dot{\stackrel{~}{r}})^2+\mathrm{exp}[2\nu (\stackrel{~}{r})]\left[\frac{L^2}{(\stackrel{~}{r})^2}+1\right]=E^2,$$
(9)
where $`E`$ and $`L`$ are the energy and the angular momentum per unit of mass of the test particle respectively. This equation shows that the radial motion is the same as that of a particle with a โposition dependent massโ with energy $`E^2/2`$ in ordinary one-dimensional, nonrelativistic mechanics moving in the effective potential
$$V_{eff}=\mathrm{exp}[2\nu (\stackrel{~}{r})]\left[\frac{L^2}{(\stackrel{~}{r})^2}+1\right].$$
(10)
In Figure 1 we show the effective potential $`V_{eff}`$ (see eq. (10)) for the deficit angles $`\alpha =0.43`$ and $`\alpha =0.125`$ respectively and for three values of the angular momentum; we note the presence of a potential well and therefore the existence of stable circular orbits.
The rotation curves (RC) of spiral galaxies are deduced from the red and blue shifts of the emitted radiation by stars moving in โcircular orbitsโ on both sides of the central region . We therefore evaluate the frequency shift for a light signal emitted from a โtest starโ in circular orbit in the spacetime of a global monopole and received by a static observer at infinity. The two maximum and minimum frequency shifts (associated with a receding or an approaching star respectively) turn out to be:
$$z_\pm =1\left[\frac{(1\alpha )}{N(\stackrel{~}{r})(N(\stackrel{~}{r})\stackrel{~}{r}N(\stackrel{~}{r})_{,\stackrel{~}{r}})}\right]^{1/2}\left[1\left(\frac{\stackrel{~}{r}N(\stackrel{~}{r})_{,\stackrel{~}{r}}}{N}\right)^{1/2}\right],$$
(11)
with $`z`$ defined as usual,
$$z\frac{\mathrm{\Delta }\omega }{\omega }=\frac{\omega _R\omega _E}{\omega _R},$$
(12)
where $`\omega _E`$ ($`\omega _R`$) is the emitted (detected) frequency. We can define $`z_D=\frac{1}{2}(z_+z_{})`$ which would be the quantity that is operationally identified with the velocity $`|v(\stackrel{~}{r})|`$ of the RC of spiral galaxies
$$z_D=\left[\frac{(1\alpha )}{N(\stackrel{~}{r})(N(\stackrel{~}{r})\stackrel{~}{r}N(\stackrel{~}{r})_{,\stackrel{~}{r}})}\right]^{1/2}\left[\frac{\stackrel{~}{r}N(\stackrel{~}{r})_{,\stackrel{~}{r}}}{N}\right]^{1/2}.$$
(13)
In Figure 2 we show $`z_+`$ (dashed line), $`z_{}`$ (solid line) and $`z_D`$ (dash-dotted line) as functions of $`\stackrel{~}{r}`$ for the cases $`\alpha =0.43`$ (left fig.) and $`\alpha =0.125`$ (right fig.) respectively. We note that even for this very simple model the figures that would correspond to the rotation curves contain a relatively โflat regionโ within the values of $`r`$ corresponding to the stable orbits (i.e. the behavior of $`z_D`$ near its maximum).
We are of course not suggesting that this simple model is in any way a realistic candidate for a natural explanation of the phenomena associated with the RC of spiral galaxies. However, the fact that the two main objections against the study of global monopoles in this context can be removed by such a simple modification as the introduction of a nonminimal coupling and the surprising coincidences mentioned at the beginning of this letter, provide very suggestive motivations for pursuing in the analysis of such possibilities. We must, nevertheless, emphasize that monopole-antimonopole pairs annihilate rather efficiently and therefore the viability of this scenario, as applied to the galactic rotation curves, will also depend strongly on the degree of asymmetry in the number of monopoles and antimonopoles in the early universe.
In conclusion we have shown that the qualitative difficulties of the original monopole scenario for explaining the rotation curves of galaxies can be solved by the introduction of nonminimal coupling. The challenge for the future is to look for a specific realization of this more complicated scenario that combines the qualitative success of the latter model with the quantitative coincidences of the former.
###### Acknowledgements.
We wish to thank R. M. Wald for helpful discussions and also to J. Guven for his comments. M.S. and D.S. would like to acknowledge partial support from DGAPA-UNAM project IN121298 and to thank the supercomputing department of DGSCA-UNAM. |
warning/0002/math0002085.html | ar5iv | text | # Why would multiplicities be log-concave ?
## 0 Introduction
The aim of these notes is to discuss some heuristic arguments, conjectures, and rigorous results related to the following phenomenon. Physical analogy, explained in Section 1, suggests that under certain circumstances the logarithms of multiplicities of irreducible representations can be expected to be concave as a function of the highest weight. In Section 2 we discuss some cases when this is known or expected to be the case and explore various implications of this concavity. In Section 3 we discuss the classical limit, in which much more general results can be established.
This text is not a survey. It is based on several talks I gave on the subject on various occasions and represents only my personal point of view. I hope that the somewhat informal style of these notes will make the basic ideas easier to explain. For missing details, the reader is referred to the original papers . For surveys on log-concavity in general, see .
I very much benefited from the discussions with a number of people, first of all, with my colleagues from the Institute of Problems of Information Transmission, especially R. Dobrushin, G. Olshanski, and S. Pirogov, and also with V. Ginzburg, W. Graham, A. Khovanskii, and A. Kirillov. In particular, the results of lead V. Ginzburg to conjecture that the push-forward of the Liouville measure on an arbitrary symplectic manifold under the moment map for a compact group action should be log-concave. Same conjecture, independently of , was proposed by A. Knutson (later, a counterexample to this conjecture was found in ; for positive results see ).
I would like to thank A. Buch for providing me with a program for computation of tensor product multiplicities.
## 1 Physical motivation: entropy and its concavity
### 1.1
Consider a quantum mechanical system, that is, a selfadjoint operator $`H`$ in a Hilbert space $`V`$. For simplicity, we assume that $`V`$ is spanned by the eigenvectors of $`H`$.
The multiplicity $`\mathrm{\Omega }(E)`$ of an eigenvalue $`E\mathrm{spec}(H)`$ measures how many states of our system have the energy $`E`$. In other words, fixing an energy level $`E`$ this does not determine the state of the system uniquely: there remain $`\mathrm{\Omega }(E)`$ possibilities. The size of this indeterminacy equals $`\mathrm{log}_2\mathrm{\Omega }(E)`$ bits of information.
In statistical physics, there is the basic relation<sup>1</sup><sup>1</sup>1This relation, in the form $`S=k\mathrm{log}W`$, is written on Boltzmannโs tombstone.
$$\overline{)S=k\mathrm{log}\mathrm{\Omega }}$$
where $`\mathrm{\Omega }`$ is the number of states with given values of macroscopic parameters such as energy, $`k`$ is the Boltzmann constant, and $`S`$ is the *entropy*, which measures the degree of disorder in the system or, in other words, the lack of information about the precise state of the system. We are thus led to think of
$$S(E)=\mathrm{log}\mathrm{\Omega }(E)$$
as of the entropy of the energy level $`E`$.
### 1.2
In statistical mechanics, the entropy is always a concave function of all additive macroscopic parameters such as energy $`E`$, volume $`V`$, or the number of particles $`N`$. There is a simple physical argument for this concavity and it goes as follows. Suppose we have two systems with parameters $`(E_1,V_1,\mathrm{})`$ and $`(E_2,V_2,\mathrm{})`$, respectively, contained in two reservoirs separated by an impervious wall:
Let us now bring them in contact by removing this wall. The energy and the volume of the new system will be $`E_1+E_2`$ and $`V_1+V_2`$, respectively, whereas the entropy will increase
$$S(E_1+E_2,\mathrm{})S(E_1,\mathrm{})+S(E_2,\mathrm{})$$
(1.1)
because of the additional disorder introduced by allowing the systems to mix. The net increase in entropy is called the *entropy of mixing* and its positivity reflects the irreversibility of mixing.
There is, however, one case when the mixing is clearly reversible and that is when the two systems were identical to begin with, that is, when
$$(E_1,V_1,\mathrm{})=(E_2,V_2,\mathrm{}),$$
in which case we can simply insert back the wall to recover the original situation. Thus, in this case the entropy of mixing vanishes. In other words,
$$S(2E,\mathrm{})=2S(E,\mathrm{}).$$
(1.2)
Combining (1.1) with (1.2) we get the concavity of the entropy<sup>2</sup><sup>2</sup>2 In thermodynamics, one has the relation $`\frac{S}{E}=\frac{1}{T}`$ where $`T`$ is the temperature. Therefore $`\frac{^2S}{E^2}<0`$ means that temperature rises when energy increases..
### 1.3
Of course, in order to apply statistical considerations one needs the system in question to have a very large or infinite number of degrees of freedom. Still, it is natural to ask whether in some interesting cases one can expect or, even better, prove the concavity of $`S(E)`$. Fortunately, interesting examples do exists.
An obvious limitation for the entropy concavity principle is that the concavity of $`S(E)`$ is clearly not preserved under direct sums. Hence our system has be in some sense *irreducible*. The concrete meaning of this irreducibility will be different in different context. In Section 2, the space $`V`$ will be an irreducible module of some ambient group. In Section 3, we will be dealing with group actions on irreducible algebraic varieties.
### 1.4
First, however, one has to modify the definition of concavity. Indeed, the support $`\mathrm{spec}(H)`$ of the function $`\mathrm{\Omega }(E)`$ is countable and hence $`S(E)`$ cannot be concave in the usual sense.
We suppose that $`\mathrm{spec}(H)`$ is contained in a lattice, which without loss of generality we can take to be $``$, and we define concavity to mean
$$S(\alpha E_1+(1\alpha )E_2)\alpha S(E_1)+(1\alpha )S(E_2),\alpha [0,1],$$
whenever the middle point $`\alpha E_1+(1\alpha )E_2`$ lies in the lattice $``$.
The abstract form of this convention is the following:
###### Definition 1.
Let $`F:๐ธ๐`$ be a function from a Abelian semigroup $`๐ธ`$ to an ordered Abelian semigroup $`๐`$. We say that this function is *concave* if
$$(p+q)F(C)pF(A)+qF(B)$$
for any $`A,B,C๐ธ`$ satisfying
$$(p+q)C=pA+qB,p,q_0.$$
In the case when
$$๐=(_0,\times )$$
is the multiplicative semigroup of nonnegative real numbers with the usual ordering, we also call the function $`F`$ logarithmically concave, or *log-concave* for short.
In our examples, $`๐ธ`$ will be usually isomorphic to $`^n`$ or $`^n`$, whereas the target semigroup $`๐`$ will occasionally be something more interesting.
### 1.5
Since the eigenvalues of $`H`$ are now integers, the time evolution $`e^{itH}`$ defines a representation of the standard circle $`๐^1`$ on $`V`$.
More generally, for any compact group<sup>3</sup><sup>3</sup>3 Here and in what follows we assume all compact groups to be *connected*. $`K`$, one can ask whether for some interesting representation $`V`$ of $`K`$ the multiplicities $`\mathrm{\Omega }(\lambda )`$ of irreducible representations $`V^\lambda `$
$$V=\underset{\lambda K^{}}{}\mathrm{\Omega }(\lambda )V^\lambda $$
form a log-concave function on the weight lattice of $`K^{}`$.
Examples of such representations will be discussed in Section 2. They are, in a sense, related to the โthermodynamicsโ of classical groups.
### 1.6
Now consider a Hamiltonian system of classical mechanics, that is, a manifold $`M^{2n}`$ with symplectic form $`\omega `$ and with an energy function
$$h:M^{2n}.$$
The form $`{\displaystyle \frac{\omega ^n}{n!}}`$ is a volume form on $`M^{2n}`$ which defines a measure (called the *Liouville measure*). Let $`\mathrm{\Omega }(E)`$ be the density of the push-forward of this measure under $`h`$
$$h_{}\left(\frac{\omega ^n}{n!}\right)=\mathrm{\Omega }(E)dE,E.$$
In other words, $`\mathrm{\Omega }(E)`$ tells us how many states of our system have the energy $`E`$. Again, we think of
$$S(E)=\mathrm{log}S(E)$$
as of the entropy<sup>4</sup><sup>4</sup>4 This entropy is not to be confused with the entropy of the dynamical system defined on $`h^1(E)`$ by the Hamiltonian flow $`\dot{x}=\{h,x\}`$. of the energy level $`E`$.
As in Section 1.5, this can be generalized to the situation of a Hamiltonian action of a compact group $`K`$ action on $`M^{2n}`$. Let
$$\varphi :M^{2n}\mathrm{Lie}(K)^{}$$
be the moment map for this action. For any $`\xi \mathrm{Lie}(K)^{}`$ the volume
$$\mathrm{\Omega }(\xi )=\mathrm{Vol}\varphi ^1(\xi )$$
measures how many points of $`X`$ have the energy $`\xi `$. This function is clearly invariant under the coadjoint action of $`K`$ on $`\mathrm{Lie}(K)^{}`$, so we can and will assume that $`\xi `$ lies in the positive Weyl chamber $`๐ฅ_+`$.
We can ask whether for some actions the function $`\mathrm{log}\mathrm{\Omega }(\xi )`$ is concave on $`๐ฅ_+`$. Observe that such a concavity implies, in particular, that the set
$$\mathrm{supp}\mathrm{\Omega }(\xi )=\varphi (M^{2n})๐ฅ_+,$$
is convex, which is a famous classical result .
It turns out that the supply of cases where $`\mathrm{log}\mathrm{\Omega }(\xi )`$ is concave is now much richer than in the quantum situation. As shown by W. Graham in , it includes all torus actions on compact Kรคhler manifolds. It also includes all actions on projective varieties, possibly singular. It was conjectured by V. Ginzburg and A. Knutson that it is true for any symplectic $`M^{2n}`$. This was shown to be not the case by Y. Karshon in .
We will discuss this classical situation in algebraic setting in Section 3.
## 2 Some results and conjectures on log-concavity of multiplicities
### 2.1
Again, we begin with a motivation, this time a historical one. Here is how the question of logarithmic concavity of multiplicities arose in the โthermodynamicsโ of classical groups.
Let $`U(\mathrm{})`$ denote the inductive limit of $`U(n)`$ with respect to standard embeddings $`U(n)U(n+1)`$ which can be visualized as follows:
The description of the characters<sup>5</sup><sup>5</sup>5 An abstract definition of characters is: indecomposable central continuous positive definite functions. More concretely, they are spherical functions of the Gelfand pair
$$U(\mathrm{})\times U(\mathrm{})\text{diag}U(\mathrm{})$$
or, equivalently, traces of factor representations of type $`\text{I}_n`$ or $`\text{II}_1`$. of $`U(\mathrm{})`$ is a fundamental result with nontrivial history. Voiculescu in proved that functions of the form
$$gdet\left[e^{\gamma ^+(g1)+\gamma ^{}(g1)}\frac{1+\beta _i^+(g1)}{1\alpha _i^+(g1)}\frac{1+\beta _i^{}(g^11)}{1\alpha _i^{}(g^11)}\right]$$
(2.1)
where
$$0\alpha _i^\pm ,0\beta _i^\pm 1,0\gamma ^\pm $$
are parameters, are characters of $`U(\mathrm{})`$ and conjectured that there are no other characters. It was observed by Boyer and, independently, by Vershik and Kerov that this conjecture is equivalent to the Schoenbergโs conjecture about the so-called totally positive sequences (see below) which was already established by Edrei in using some deep results about entire functions.
### 2.2
Vershik and Kerov also outlined a different and more direct proof which uses approximation of characters of $`U(\mathrm{})`$ by normalized characters of $`U(n)`$. It follows from a general principle due to Vershik (see and also ), that any character $`\chi `$ of $`U(\mathrm{})`$ is a limit of a sequence of normalized characters $`\chi _n`$ of $`U(n)`$ as $`n\mathrm{}`$ in the sense that
$$\chi _n|_{U(k)}\stackrel{\text{ uniformly }}{}\chi |_{U(k)}$$
(2.2)
for any fixed $`k=1,2,3,\mathrm{}`$. In , Vershik and Kerov gave necessary and sufficient conditions for the convergence of $`\{\chi _n\}`$ and identified the corresponding limits with functions (2.1).
This approximation principle is a materialization of certain general ergodic theory ideas and is closely akin to some standard constructions in statistical physics such as construction of Gibbsian measures in an infinite volume by a thermodynamic limit transition. In that case, one chooses a sequence of boxes which fill up the space (just as in the above visualization of $`U(\mathrm{})`$), for each box one picks some boundary condition which specifies a Gibbsian measure (in our case, $`\chi _n`$), and one requires convergence of the induced measures on all compact sets.
### 2.3
Note that the formula (2.1) is multiplicative in the eigenvalues of $`gU(\mathrm{})`$. This multiplicativity can be established a priori; as shown by Olshanski, see for example , such and more general multiplicativity are very characteristic for representations of infinite-dimensional classical groups.
Any character $`\chi `$ of $`U(\mathrm{})`$ is therefore uniquely determined by its restriction to $`U(1)`$
$$x(z)=\chi \left(\left[\begin{array}{c}z\\ & 1\\ & & \mathrm{}\end{array}\right]\right)=\underset{k}{}x_kz^k.$$
Conversely, any function $`gdetx(g)`$ is a character of $`U(\mathrm{})`$ provided that it is positive definite, which means that its restriction to any $`U(n)`$ is a nonnegative linear combination of the characters of $`U(n)`$, that is, of the rational Schur functions $`s_\lambda `$.
The identity
$$\underset{i=1}{\overset{n}{}}\left(\underset{k}{}x_kz_i^k\right)=\underset{\lambda =(\lambda _1\mathrm{}\lambda _n)^n}{}det\left[x_{\lambda _ii+j}\right]_{i,j=1\mathrm{}n}s_\lambda (z_1,\mathrm{},z_n),$$
shows that this positivity is equivalent to the positivity of some (in fact, all) minors of the infinite Toeplitz matrix $`\left[x_{ji}\right]_{i,j}`$, which is precisely Schoenbergโs definition of a *totally positive* sequence.
In particular, the positivity of $`2\times 2`$ minors means that
$$x_n^2x_{n1}x_{n+1}.$$
Thus, one knows a priori that the restriction of any character of $`U(\mathrm{})`$ to $`U(1)`$ has log-concave multiplicities.
### 2.4
The question whether the same is true before the limit, that is, whether the restriction of any irreducible representation of $`U(n)`$ to standard $`U(1)`$ has log-concave multiplicities, surfaced when we were working with G. Olshanski on a generalization of the VershikโKerov theorem . Originally, this log-concavity was needed to replace the uniform convergence in (2.2) by convergence of Taylor series, see Section 3 in . Eventually, in it was replaced by a more elementary argument, but nonetheless this log-concavity is a valid question with interesting answer.
As it turns out, for any representation $`V^\lambda `$ of $`U(n)`$ the multiplicity of the irreducible representation $`V^\mu `$ of the standard $`U(k)U(n)`$ is a log-concave function of the pair
$$(\lambda ,\mu )U(n)^{}U(k)^{}.$$
In fact, one can say more. Without loss of generality, let us assume that $`\lambda `$ is a partition and consider the space
$$V^{\lambda /\mu }=\mathrm{Hom}_{U(k)}(V^\mu V^\lambda )$$
whose dimension is the multiplicity in question. The space $`V^{\lambda /\mu }`$ is an $`U(nk)`$ module with character given by the skew Schur function $`s_{\lambda /\mu }`$
$$\mathrm{tr}_{V^{\lambda /\mu }}\left(\left[\begin{array}{c}z_1\\ & z_2\\ & & \mathrm{}\end{array}\right]\right)=s_{\lambda /\mu }(z_1,z_2,\mathrm{}).$$
One has the following
###### Theorem 1 ().
Suppose $`(\lambda _i,\mu _i)`$, $`i=1,2,3`$, are partitions such that
$$(\lambda _2,\mu _2)=\frac{1}{2}(\lambda _1,\mu _1)+\frac{1}{2}(\lambda _3,\mu _3).$$
Then the following polynomial has nonnegative coefficients:
$$s_{\lambda _2/\mu _2}^2s_{\lambda _1/\mu _1}s_{\lambda _3/\mu _3}_0[z_1,z_2,\mathrm{}].$$
(2.3)
The coefficients of the polynomial $`s_{\lambda /\mu }`$ correspond to standard tableaux of shape $`\lambda /\mu `$. In the proof of Theorem 1, one constructs a certain transformation on pairs of standard tableaux and proves that it is injective.
Similar results for orthogonal and symplectic groups are also established in .
### 2.5
It is likely that (2.3) is actually a nonnegative linear combination of Schur functions. One can propose a conjecture which would, among other things, imply this property.
Recall that the Littlewood-Richardson coefficients $`c_{\lambda \mu \nu }`$ are defined by
$$c_{\lambda \mu \nu }=dim\left(V^\lambda V^\mu V^\nu \right)^G$$
where the superscript $`G`$ stands for invariants of $`G=U(n)`$. If either of the arguments of $`c_{\lambda \mu \nu }`$ is not a dominant weight, we set $`c_{\lambda \mu \nu }=0`$ by definition. Often, one uses the numbers
$$c_{\mu \nu }^\lambda =c_{\lambda ^{}\mu \nu }$$
where $`\lambda ^{}`$ is the highest weight of the dual module $`\left(V^\lambda \right)^{}`$
$$(\lambda _1,\mathrm{},\lambda _n)^{}=(\lambda _n,\mathrm{},\lambda _1).$$
The numbers $`c_{\mu \nu }^\lambda `$ are coefficients in the expansions
$`V^{\lambda /\mu }`$ $`={\displaystyle \underset{\nu }{}}c_{\mu \nu }^\lambda V^\nu ,`$
$`V^\mu V^\nu `$ $`={\displaystyle \underset{\lambda }{}}c_{\mu \nu }^\lambda V^\lambda .`$
###### Conjecture 1.
The function
$$(\lambda ,\mu ,\nu )\mathrm{log}c_{\lambda \mu \nu }$$
is concave.
If true, this concavity would have some interesting applications. In particular, since
$$c_{\lambda \mu \nu }=c_{\lambda \nu \mu }$$
we conclude that
$$c_{\lambda ,\frac{\mu +\nu }{2},\frac{\mu +\nu }{2}}\stackrel{?}{}c_{\lambda \nu \mu },$$
provided $`\frac{\mu +\nu }{2}`$ is an integral weight. This is equivalent to the inclusion of representations
$$V^\nu V^\mu \stackrel{?}{}\left(V^{\frac{\mu +\nu }{2}}\right)^2,$$
(2.4)
which can be interpreted as saying that the representation valued function
$$V:\lambda V^\lambda $$
(2.5)
is concave with respect to the natural ordering and tensor multiplication of representations <sup>6</sup><sup>6</sup>6Remark that it follows from Weylโs dimension formula that
$$\lambda \mathrm{log}dimV^\lambda $$
is a concave function. That is, the function (2.5) considered as a function into just vector spaces without group action is concave with respect to the tensor product..
If (2.4) is true then we certainly have the following inclusion of $`U(n)`$-modules
$$\left(V^{\lambda _1}V^{\mu _1}\right)\left(V^{\lambda _3}V^{\mu _3}\right)\stackrel{?}{}\left(V^{\lambda _2}V^{\mu _2}\right)^2$$
(2.6)
for $`(\lambda _i,\mu _i)`$ as in Theorem 1. The last inclusion is equivalent to
$$V^{\lambda _1/\mu _1}V^{\lambda _3/\mu _3}\stackrel{?}{}\left(V^{\lambda _2/\mu _2}\right)^2.$$
(2.7)
Indeed, the equation (2.7) is equivalent to
$$\begin{array}{c}\left(\underset{\nu _1}{}c_{\lambda _1^{}\mu _1\nu _1}V^{\nu _1}\right)\left(\underset{\nu _3}{}c_{\lambda _3^{}\mu _3\nu _3}V^{\nu _3}\right)\stackrel{?}{}\hfill \\ \hfill \left(\underset{\nu _2}{}c_{\lambda _2^{}\mu _2\nu _2}V^{\nu _2}\right)\left(\underset{\nu _4}{}c_{\lambda _2^{}\mu _2\nu _4}V^{\nu _4}\right),\end{array}$$
(2.8)
whereas (2.6) says that
$$\begin{array}{c}\left(\underset{\nu _1}{}c_{\lambda _1\mu _1\nu _1^{}}V^{\nu _1}\right)\left(\underset{\nu _3}{}c_{\lambda _3\mu _3\nu _3^{}}V^{\nu _3}\right)\stackrel{?}{}\hfill \\ \hfill \left(\underset{\nu _2}{}c_{\lambda _2\mu _2\nu _2^{}}V^{\nu _2}\right)\left(\underset{\nu _4}{}c_{\lambda _2\mu _2\nu _4^{}}V^{\nu _4}\right).\end{array}$$
(2.9)
To get (2.8) from (2.9), take the dual space of everything, which will replace $`V^{\nu _i}`$ by $`V^{\nu _i^{}}`$, and then replace $`\lambda _i`$ by $`\lambda _i^{}`$ and $`\nu _i^{}`$ by $`\nu _i`$. The inclusion (2.7) is equivalent to Schur-positivity of (2.3).
Similarly, the conjecture and the symmetry
$$c_{\lambda \mu \nu }=c_{\nu \lambda \mu }$$
imply that
$$c_{\lambda ^{},\mu ^{},\nu ^{}}\stackrel{?}{}c_{\lambda \mu \nu }$$
provided the weight
$$\left(\begin{array}{c}\lambda ^{}\\ \mu ^{}\\ \nu ^{}\end{array}\right)=\left(\begin{array}{ccc}\alpha & 0& 1\alpha \\ 1\alpha & \alpha & 0\\ 0& 1\alpha & \alpha \end{array}\right)\left(\begin{array}{c}\lambda \\ \mu \\ \nu \end{array}\right),$$
is an integral weight and $`0\alpha 1`$.
### 2.6
Here is another implication of Conjecture 1 which is actually known to be true. Concavity of $`\mathrm{log}c_{\lambda \mu \nu }`$ implies that the support
$$\mathrm{supp}c_{\lambda \mu \nu }=\{(\lambda ,\mu ,\nu ),c_{\lambda \mu \nu }0\}$$
is convex. In particular, since it contains the origin $`(0,0,0)`$, it is *saturated*, meaning that
$$c_{k\lambda ,k\mu ,k\nu }0c_{\lambda ,\mu ,\nu }0,$$
(2.10)
for any $`k=2,3,\mathrm{}`$. In fact, since $`c_{0,0,0}=1`$, Conjecture 1 implies that
$$c_{k\lambda ,k\mu ,k\nu }\stackrel{?}{}(c_{\lambda ,\mu ,\nu })^k.$$
The saturation (2.10) turns out to be a very important property, see . It has been recently established by A. Knutson and T. Tao in , see also .
### 2.7
As already pointed out in Section 1.3, log-concavity is (in contrast to so many things in representation theory) not an additive property: it is totally destroyed by direct sums.
It seems however likely that log-concavity should be a *multiplicative* property, that is, it should behave nicely with respect to tensor products. For example, recall that it is well known and easy to prove that the convolution of two log-concave sequences is again log-concave. This is equivalent to saying that the set of $`U(1)`$-modules with log-concave multiplicities is closed under tensor products.
This multiplicativity principle fits together nicely with the above conjecture about tensor product multiplicities.
## 3 Log-concavity in the classical limit
### 3.1
Dealing with actual multiplicities may be a subtle business. Fortunately, many of these subtleties disappear in the classical limit and much more general results can be established.
Let us assume that the phase space of our classical system is an irreducible projective algebraic variety $`X^N`$ over $``$ which is stable under the action of a compact group $`KGL(N+1)`$. We write $`X`$ in place of $`M^{2n}`$ to stress the fact that we are now working with projective algebraic varieties which are allowed to be singular.
Even for singular $`X`$, the moment map
$$\varphi :X\mathrm{Lie}(K)^{}$$
is still well defined as the restriction of the moment map for the $`K`$-action on $`^N`$. It is well known (see e.g. Theorem 6.5 in ) that the function $`\mathrm{\Omega }(\xi )`$ from Section 1.6 describes the asymptotics of the multiplicities of $`K`$-modules in polynomials of very large degree on $`X`$.
More concretely, let
$$[X]=_{d=0}^{\mathrm{}}[X]_d$$
be the homogeneous coordinate ring of $`X`$. The space $`[X]_d`$ of degree $`d`$ polynomials on $`X`$ decomposes as a $`K`$-module
$$[X]_d=\underset{\lambda K^{}}{}\mathrm{\Omega }_d(\lambda )V^\lambda ,.$$
One can view $`\mathrm{\Omega }_k(\lambda )`$ as a measure on $`K^{}๐ฅ_+`$. After proper normalization, the measures $`\mathrm{\Omega }_k(k\lambda )`$ converge weakly to $`\mathrm{\Omega }(\xi )d\xi `$ as $`k\mathrm{}`$ where $`d\xi `$ is the Lebesgue measure on $`๐ฅ_+`$.
In other words, for any $`A๐ฅ_+`$, the integral $`_A\mathrm{\Omega }(\xi )๐\xi `$ describes the leading asymptotics of the sum $`_{\lambda kA}\mathrm{\Omega }_k(\lambda )`$ as $`k\mathrm{}`$. Hence, informally, $`\mathrm{\Omega }(\xi )`$ is the multiplicity $`\mathrm{\Omega }_k(k\lambda )`$ averaged over some infinitesimal neighborhood of $`\xi `$. Such an averaging over infinitesimally close energy levels is a very natural thing to do from the statistical physics perspective.
### 3.2
The function $`\mathrm{\Omega }(\xi )`$ depends not only on the $`K`$-action on $`X`$ as such but also on the embedding $`X^N`$, where $`^{N+1}`$ is a representation space of $`K`$ or, equivalently, of the complexification $`G`$ of $`K`$.
In intrinsic terms, such an embedding is a very ample invertible sheaf $`L`$ in the $`G`$-linearized Picard group $`\mathrm{Pic}^G(X)`$ of $`X`$. Write $`\mathrm{\Omega }(\xi ,L)`$ to stress the dependence on both $`\xi `$ and $`L`$. Because $`L`$ enters the definition of $`\mathrm{\Omega }(\xi ,L)`$ only via $`L^n`$, $`n\mathrm{}`$, the function $`\mathrm{\Omega }(\xi ,L)`$ is well-defined for any
$$L\mathrm{Pic}^G(X)_{}.$$
Since ample sheafs form a semigroup in $`\mathrm{Pic}_G(X)`$ it makes sense to ask whether $`\mathrm{log}\mathrm{\Omega }(\xi ,L)`$ is concave as a function of the pair $`(\xi ,L)`$.
In fact, we already saw an example of such a bivariate concavity in Section 2.4 where the multiplicities for restrictions from $`U(n)`$ to $`U(k)`$ turned out to be concave in the pair of highest weights.
### 3.3
In this setting, the log-concavity of $`\mathrm{\Omega }(\xi )`$ and $`\mathrm{\Omega }(\xi ,L)`$ was established in and , respectively, by using the classical Brunn-Minkowski inequality of convex analysis.
Here we want to use the same ideas to approach the problem from a slightly different angle. Instead of looking at the weak limit of measures $`\mathrm{\Omega }_k(k\lambda )`$, which involves averaging over infinitesimally close energy levels, we want to look at the asymptotics of the sequence $`\mathrm{\Omega }_k(k\lambda )`$ for some fixed $`\lambda `$. This is a natural thing to do from the representation theory point of view.
### 3.4
There is a standard trick which allows to dispose of the first variable $`\xi `$ in $`\mathrm{\Omega }(\xi ,L)`$ by enlarging the variety $`X`$. Indeed, by definition of the multiplicities $`\mathrm{\Omega }_k`$ we have
$`\mathrm{\Omega }_k(k\lambda ,L)`$ $`=dim\left(H^0(X,L^k)V^{k\lambda ^{}}\right)^G,`$
$`=dimH^0(X\times G/B,(LL_\lambda ^{})^k)^G,`$
where $`\lambda ^{}K^{}`$ is the highest weight of $`(V^\lambda )^{}`$, $`G/B=K/T`$ is the flag variety of $`K`$, the sheaf $`L_\lambda ^{}\mathrm{Pic}^G(G/B)`$ corresponds to the map of $`G/B`$ onto the orbit of the highest vector in $`P(V^\lambda )`$, and superscript $`G`$ denotes $`G`$-invariants. So, without loss of generality, we can assume that $`\xi =0`$.
### 3.5
Recall that, by definition,
$$[X//_LG]=\underset{k}{}H^0(X,L^k)^G$$
is the homogeneous coordinate ring of the Geometric Invariant Theory quotient $`X//_LG`$ corresponding to $`L\mathrm{Pic}^G(X)`$. We write $`X//_LG`$ in place of the standard $`X//G`$ to stress the dependence on $`L`$.
The sequence
$$\mathrm{\Omega }_k(0,L)=dimH^0(X,L^k)^G,k=0,1,2,\mathrm{}$$
(3.1)
may fail to have a $`k\mathrm{}`$ asymptotics for the following trivial reason. Consider the set
$$\{k,H^0(X,L^k)^G0\}_0.$$
Since $`X`$ is irreducible, it is a semigroup and it either contains all sufficiently large integers or lies in a proper subgroup if $``$. We want to avoid the latter case because in that case the sequence (3.1) does not have any asymptotics. So we will replace $`L`$ by a suitable power of $`L`$ in that case.
By replacing $`L`$ by its power one can also achieve that $`[X//_LG]`$ is generated by its degree $`1`$ graded component and so we can assume this as well. Thus, we have an embedding
$$X//_LG\left(\left(H^0(X,L)^G\right)^{}\right)$$
and we denote by $`\mathrm{deg}X//_LG`$ the degree of this embedding. It follows that in this case
$$\mathrm{\Omega }_k(0,L)\mathrm{deg}X//_LG\frac{k^d}{d!},d=dimX//_LG,$$
as $`k`$ goes to $`\mathrm{}`$. We now want to show that $`\mathrm{log}\mathrm{deg}X//_LG`$ is a concave function of $`L`$.
### 3.6
In fact, one can establish a more general fact. Let $`Y`$ be an irreducible algebraic variety of dimension $`d`$ and let $`๐=(Y)`$ be the field of rational functions on $`Y`$. Let $`S๐`$ be a $``$-linear subspace such that $`1S`$ and which generates $`๐`$ as a field. The embedding $`S๐`$ corresponds to a subvariety $`Y_SS^{}`$ which is birationally isomorphic to $`Y`$. Let $`\mathrm{deg}Y_S`$ denote its degree.
Given two such subspaces $`S_1`$ and $`S_2`$, denote by $`S_1S_2`$ the subspace generated by all products $`f_1f_2`$, where $`f_iS_i`$. We will show that
$$\sqrt[d]{\mathrm{deg}Y_{S_1S_2}}\sqrt[d]{\mathrm{deg}Y_{S_1}}+\sqrt[d]{\mathrm{deg}Y_{S_2}}$$
(3.2)
for any such pair $`S_1`$ and $`S_2`$. Since, clearly,
$$\mathrm{deg}Y_{S^2}=2^d\mathrm{deg}Y_S$$
the inequality (3.2) implies that $`\sqrt[d]{\mathrm{deg}Y_S}`$, and consequently, $`\mathrm{log}\mathrm{deg}Y_S`$ is a concave function of $`S`$.
### 3.7
In particular, (3.2) would imply the concavity of $`\mathrm{log}\mathrm{deg}X//_LG`$. Indeed, although the varieties $`X//_LG`$ may not be isomorphic for different $`L`$, they are always birationally isomorphic. Their common field of fractions is the field $`๐=(X)^G`$ of rational $`G`$-invariants.
Given some $`L_1`$ and $`L_2`$, pick some $`\varphi _iH^0(X,L_i)^G`$. Replacing the $`L_i`$โs if necessary by their multiples, we can assume that $`[X//_{L_i}G]`$ is generated by $`S_i=\varphi _i^1H^0(X,L_i)^G๐`$. Since the algebra $`[X//_{L_1L_2}G]`$ contains the algebra generated by $`S_1S_2`$, we get from (3.2) the desired lower bound on the asymptotics of the dimensions of the graded components of $`[X//_{L_1L_2}G]`$.
### 3.8
Now, in order to establish (3.2), we will construct convex sets $`\mathrm{\Delta }_S^d`$ of dimension $`d=dimY`$ such that
$$\mathrm{deg}Y_S=d!\mathrm{vol}\mathrm{\Delta }_S$$
and
$$\mathrm{\Delta }_{S_1S_2}\mathrm{\Delta }_{S_1}+\mathrm{\Delta }_{S_2}.$$
The inequality (3.2) will then follow immediately from the classical Brunn-Minkowski inequality, see e.g. . See also e.g. the appendix by A. Khovanskii in for a discussion of the relationship between classical inequalities of the convex analysis and algebraic geometry. For example, the Alexandrov-Fenchel inequality, which is stronger than the Brunn-Minkowski inequality, corresponds to the Hodge index theorem for surfaces.
### 3.9
The construction of $`\mathrm{\Delta }_S^d`$ is similar to the definition of a Newton polytope.
Choose a smooth point $`yY`$ which lies away from the singularities of the maps from $`Y`$ to $`Y_{S_i}`$ and their inverses. Choose a flag of subvarieties
$$YY^1\mathrm{}Y^d=y,\mathrm{codim}Y^k=k,$$
which are all smooth at $`y`$. Fix some local equation $`u_k`$ of $`Y^k`$ in $`Y^{k1}`$.
This data give rise to a map
$$๐0f๐ณ(f)=(๐ณ_1(f),\mathrm{},๐ณ_d(f))^d,$$
where
$`๐ณ_1(f)`$ $`=\mathrm{ord}_{Y^1}f,`$
$`๐ณ_2(f)`$ $`=\mathrm{ord}_{Y^2}\left(fu_1^{๐ณ_1(f)}\right)|_{Y^1},`$
$`๐ณ_3(f)`$ $`=\mathrm{ord}_{Y^3}\left(\left(fu_1^{๐ณ_1(f)}\right)|_{Y^1}u_2^{๐ณ_2(f)}\right)|_{Y^2}\mathrm{}`$
and so on. It is clear that $`๐ณ`$ is a valuation, that is,
$`๐ณ(fg)`$ $`=๐ณ(f)+๐ณ(g),`$ (3.3)
$`๐ณ(f+g)`$ $`\mathrm{min}\{๐ณ(f),๐ณ(g)\},`$
where the ordering on $`^d`$ is lexicographic.
It is also clear that the residue field of $`๐ณ`$ is isomorphic to $``$ and hence for any $``$-linear subspace of $`S๐`$ we have
$$dim_CS=|๐ณ(S0)|.$$
(3.4)
### 3.10
By definition, set
$$\mathrm{\Gamma }_S=\left\{(k,๐ณ(f)),fS^k0\right\}^{1+d}.$$
It follows from (3.3) that $`\mathrm{\Gamma }_S`$ is a semigroup.
Denote by $`\mathrm{\Lambda }_S^{1+d}`$ the lattice generated by $`\mathrm{\Gamma }_S`$. Let $`_S^{1+d}`$ be the closed convex cone generated by $`\mathrm{\Gamma }_S`$ and let $`\mathrm{\Delta }_S`$ be the intersection of $`_S`$ with the subspace $`(1,^d)^{1+d}`$. It is clear that
$$\mathrm{\Delta }_S=\overline{\left\{\frac{๐ณ(f)}{k},fS^k0\right\}},$$
where bar denotes closure.
Since the point $`y`$ corresponds to a smooth point of $`Y_S`$, there exist
$$f_0,f_1,\mathrm{},f_d=1[Y_S]$$
such that
$$๐ณ(f_k)=(\underset{k\text{ times}}{\underset{}{0,\mathrm{},0}},1,\mathrm{}).$$
Thus, $`\mathrm{\Delta }_S`$ contains a $`d`$-dimensional simplex and so
$$dim\mathrm{\Delta }_S=d$$
(3.5)
### 3.11
Let us now prove that
$$\mathrm{vol}\mathrm{\Delta }_S=d!\mathrm{deg}Y_S,$$
On the one hand, it follows that (3.4) and the definition of $`\mathrm{\Gamma }_S`$ that
$$\left|\mathrm{\Gamma }_S(k,^d)\right|=dim_CS^k\mathrm{deg}Y_S\frac{k^d}{d!},k\mathrm{}.$$
Since $`\mathrm{\Gamma }_S_S\mathrm{\Lambda }_S`$ we have
$$\frac{\mathrm{deg}Y_S}{d!}\mathrm{vol}\mathrm{\Delta }_S,$$
where the normalization of Lebesgue measure is given by the intersection of the lattice $`\mathrm{\Lambda }_S`$ with the subspace $`(1,^d)^{1+d}`$.
The inverse inequality will be deduced from the following results of Khovanskii . Choose a sequence of finitely generated subgroups
$$\mathrm{\Gamma }_1\mathrm{\Gamma }_2\mathrm{}\mathrm{\Gamma }_S$$
such that $`\mathrm{\Gamma }_S=\mathrm{\Gamma }_i`$ and each $`\mathrm{\Gamma }_i`$ generates the lattice $`\mathrm{\Lambda }_S`$. Let $`_i`$ denote the cone generated by $`\mathrm{\Gamma }_i`$ and let $`\mathrm{\Delta }_i`$ be the corresponding hyperplane section of $`_i`$. It is clear that
$$\mathrm{\Delta }_1\mathrm{\Delta }_2\mathrm{}\mathrm{\Delta }_S=\overline{\mathrm{\Delta }_i}.$$
(3.6)
It is a theorem of Khovanskii, see Proposition 3 in Section 3 of , that there exists vectors $`\gamma _i\mathrm{\Gamma }_i`$ such that
$$(_i+\gamma _i)\mathrm{\Lambda }_S\mathrm{\Gamma }_i.$$
It follows that
$$\frac{\left|\mathrm{\Gamma }_i(k,^d)\right|}{k^d}\mathrm{vol}\mathrm{\Delta }_i,k\mathrm{}.$$
Hence $`\mathrm{vol}\mathrm{\Delta }_i(d!)^1\mathrm{deg}Y_S`$ and (3.6) implies that
$$\mathrm{vol}\mathrm{\Delta }_S\frac{\mathrm{deg}Y_S}{d!},$$
as was to be shown.
### 3.12
The relation
$$\mathrm{\Delta }_{S_1S_2}\mathrm{\Delta }_{S_1}+\mathrm{\Delta }_{S_2}.$$
(3.7)
follows immediately from (3.3) and the definition of $`\mathrm{\Delta }_S`$. Now we are almost in position to finish the proof of (3.2) by applying the Brunn-Minkowski inequality. One remaining detail is that our normalization of the volume $`\mathrm{vol}\mathrm{\Delta }_S`$ depends on the lattice $`\mathrm{\Lambda }_S`$. It is, however, clear that
$$\mathrm{\Lambda }_{S_1S_2}\mathrm{\Lambda }_{S_1},\mathrm{\Lambda }_{S_2}.$$
Therefore, if we normalize the volume according to the lattice $`\mathrm{\Lambda }_{S_1S_2}`$ we have
$$\mathrm{vol}\mathrm{\Delta }_{S_1S_2}=\frac{\mathrm{deg}Y_{S_1S_2}}{d!},\mathrm{vol}\mathrm{\Delta }_{S_i}\frac{\mathrm{deg}Y_{S_i}}{d!}.$$
This and the Brunn-Minkowski inequality applied to (3.7)
$$\sqrt[d]{\mathrm{vol}\mathrm{\Delta }_{S_1S_2}}\sqrt[d]{\mathrm{vol}\mathrm{\Delta }_{S_1}}+\sqrt[d]{\mathrm{vol}\mathrm{\Delta }_{S_2}}$$
completes the proof of (3.2). |
warning/0002/hep-th0002107.html | ar5iv | text | # References
hep-th/0002107
AEI-2000-004
The Asymptotic Groundstate of SU(3) Matrix Theory
Jens Hoppe and Jan Plefka
Albert-Einstein-Institut
Max-Planck-Institut fรผr Gravitationsphysik
Am Mรผhlenberg 1, D-14476 Golm, Germany
hoppe,plefka@aei-potsdam.mpg.de
Abstract
The asymptotic form of a $`SU(3)`$ matrix theory groundstate is found by showing that a recent ansatz for a supersymmetric wavefunction is non-trivial (i.e. non-zero).
February 2000
Maximally supersymmetric $`SU(N)`$ gauge quantum mechanics in $`d=9`$ has in recent years received much attention due to its close relation<sup>1</sup><sup>1</sup>1based on to the eleven-dimensional supermembrane in the $`N\mathrm{}`$ limit, its description of the dynamics of $`N`$ D0 branes in superstring theory , as well as the M theory proposal of . In these physical interpretations the existence of a unique normalizable zero-energy groundstate is an important consistency requirement. An explicit construction of the vacuum state, though highly desirable, appears to be quite difficult. Another approach is to study the behavior of the wavefunction far out at infinity where the degrees of freedom in the Cartan-subalgebra become free and the remaining degrees of freedom form the zero energy vacuum state of supersymmetric harmonic oscillators . The full asymptotic groundstate was constructed for the $`SU(2)`$ model in , here we consider the $`SU(3)`$ case. Assuming that the Cartan-subalgebra degrees of freedom are asymptotically governed by a set of free effective supercharges $`Q_\alpha `$ a proposal was recently made as to which of the harmonic wavefunctions constructed in is annihilated by the $`Q_\alpha `$. In this letter we prove the non-triviality of this ansatz.
The asymptotic supersymmetry charge for the $`d=9`$ SU(3) model reads
$$Q_\alpha =i\mathrm{\Gamma }_{\alpha \beta }^a\left(\theta _\beta ^1\frac{}{x_1^a}+\theta _\beta ^2\frac{}{x_2^a}\right)$$
(1)
where $`x_1^a,x_2^a`$ $`(a=1,\mathrm{},9)`$ are the bosonic and $`\theta _\alpha ^1,\theta _\alpha ^2`$ $`(\alpha =1,\mathrm{},16)`$ are the fermionic degrees of freedom of the Cartan sector; we work with a real, symmetric representation of the Dirac matrices and our charge conjugation matrix equals unity. It is advantageous to go to the complex variables
$`\lambda =\frac{1}{\sqrt{2}}(\theta ^1+i\theta ^2)`$ $`z^a=\frac{1}{\sqrt{2}}(x_1^a+ix_2^a)`$
$`\lambda ^{}=\frac{1}{\sqrt{2}}(\theta ^1i\theta ^2)`$ $`\overline{z}^a=\frac{1}{\sqrt{2}}(x_1^aix_2^a).`$ (2)
Note that we have now divided the fermions into creation and annihilation operators, obeying the algebra
$$\{\lambda _\alpha ,\lambda _\beta ^{}\}=\delta _{\alpha \beta },$$
(3)
and we define the fermionic vacuum $`|`$ by $`\lambda _\alpha |=0`$. The completely filled state is denoted by $`|+=\frac{1}{16!}ฯต^{\alpha _1\mathrm{}\alpha _{16}}\lambda _{\alpha _1}^{}\mathrm{}\lambda _{\alpha _{16}}^{}|`$. Clearly $`|`$ and $`|+`$ are $`SO(9)`$ singlets. However, there is a third $`SO(9)`$ singlet state
$$|\mathbb{๐}=(\lambda ^{}\mathrm{\Gamma }^{ab}\lambda ^{})(\lambda ^{}\mathrm{\Gamma }^{bc}\lambda ^{})(\lambda ^{}\mathrm{\Gamma }^{cd}\lambda ^{})(\lambda ^{}\mathrm{\Gamma }^{da}\lambda ^{})|$$
(4)
in the half-filled sector. It can be shown that there are no further $`SO(9)`$ singlets. A further symmetry group acting on these states is the Weyl group, the discrete asymptotic remnant of the continous $`SU(3)`$ of the full system. The Weyl group for $`SU(3)`$ may be generated by two elements $`P`$ and $`C`$ , which act on the complex fermions $`\lambda `$ and $`\lambda ^{}`$ as
$`P:\lambda `$ $`\lambda ^{}`$ $`\lambda ^{}\lambda `$
$`C:\lambda `$ $`e^{{\scriptscriptstyle \frac{2\pi i}{3}}}\lambda `$ $`\lambda ^{}e^{{\scriptscriptstyle \frac{2\pi i}{3}}}\lambda ^{}.`$ (5)
As $`P`$ interchanges $`|+`$ and $`|`$, leaves the eight fermion sector invariant and the three $`SO(9)`$ singlets are known to form one two dimensional irreducible representation under the Weyl group and one singlet the state $`|\mathbb{๐}`$ has to be Weyl invariant. This is consistent with $`C`$ transforming $`|\pm `$ into $`\mathrm{exp}(\frac{2\pi i}{3})|\pm `$. So $`|\mathbb{๐}`$ is the unique $`SO(9)`$ and Weyl invariant state.
In the complex variables the supersymmetry charge (1) reads $`Q_\alpha =i(/\lambda )_\alpha i(\overline{}/\lambda ^{})_\alpha `$, where $`_a=d/dz^a`$ and $`\overline{}_a=d/d\overline{z}^a`$. We seek for an asymptotic groundstate $`|\mathrm{\Psi }`$ obeying $`Q_\alpha |\mathrm{\Psi }=0`$. Note that while $`Q_\alpha `$ squares to $`\overline{}`$, the condition $`\overline{}|\mathrm{\Psi }=0`$ does not imply $`Q_\alpha |\mathrm{\Psi }=0`$ due to the purely asymptotic considerations, i.e. $`|\mathrm{\Psi }`$ not being square integrable caused by its singularity at the origin.
Consider now the ansatz for $`|\mathrm{\Psi }`$
$$|\mathrm{\Psi }=ฯต^{\alpha _1\mathrm{}\alpha _{16}}Q_{\alpha _1}\mathrm{}Q_{\alpha _{16}}\frac{1}{(z\overline{z})^8}|\mathbb{๐}.$$
(6)
$`|\mathrm{\Psi }`$ is obviously annihilated by $`Q_\alpha `$, as $`Q_\alpha `$ squares to the Laplacian $`\overline{}`$ which in turn annihilates the harmonic function $`(z\overline{z})^8`$. Note that $`|\mathrm{\Psi }`$ is $`SO(9)\times `$Weyl invariant by construction. What remains to be shown, however, is that $`|\mathrm{\Psi }`$ is non-vanishing.
For this we consider the matrix element
$`|\mathrm{\Psi }`$ $`=`$ $`ฯต^{\alpha _1\mathrm{}\alpha _{16}}|[(/\lambda )_{\alpha _1}+(\overline{}/\lambda ^{})_{\alpha _1}]\mathrm{}[(/\lambda )_{\alpha _{16}}+(\overline{}/\lambda ^{})_{\alpha _{16}}]`$ (7)
$`(\lambda ^{}\mathrm{\Gamma }^{ab}\lambda ^{})(\lambda ^{}\mathrm{\Gamma }^{bc}\lambda ^{})(\lambda ^{}\mathrm{\Gamma }^{cd}\lambda ^{})(\lambda ^{}\mathrm{\Gamma }^{da}\lambda ^{})|{\displaystyle \frac{1}{(z\overline{z})^8}},`$
which we now need to normal order by making use of the anticommutator relation
$$\{(/\lambda )_\alpha ,(\overline{}/\lambda ^{})_\beta \}=\delta _{\alpha \beta }\overline{}+\mathrm{\Gamma }_{\alpha \beta }^{ab}_a\overline{}_b.$$
(8)
From the $`2^{16}`$ terms generated from expanding out the brackets in the first line of (7) only those containing 4 $`(\overline{}/\lambda ^{})`$ and 12 $`(/\lambda )`$ survive. Normal ordering of these terms then yields
$`|\mathrm{\Psi }`$ $``$ $`ฯต^{\alpha _1\mathrm{}\alpha _{16}}\mathrm{\Gamma }_{\alpha _1\alpha _2}^{a_1a_2}\overline{}_{a_1}_{a_2}\mathrm{\Gamma }_{\alpha _3\alpha _4}^{a_3a_4}\overline{}_{a_3}_{a_4}\mathrm{\Gamma }_{\alpha _5\alpha _6}^{a_5a_6}\overline{}_{a_5}_{a_6}\mathrm{\Gamma }_{\alpha _7\alpha _8}^{a_7a_8}\overline{}_{a_7}_{a_8}`$
$`|(/\lambda )_{\alpha _9}\mathrm{}(/\lambda )_{\alpha _{16}}(\lambda ^{}\mathrm{\Gamma }^{ab}\lambda ^{})(\lambda ^{}\mathrm{\Gamma }^{bc}\lambda ^{})(\lambda ^{}\mathrm{\Gamma }^{cd}\lambda ^{})(\lambda ^{}\mathrm{\Gamma }^{da}\lambda ^{})|{\displaystyle \frac{1}{(z\overline{z})^8}},`$
and the final contractions then result in
$`|\mathrm{\Psi }`$ $``$ $`ฯต^{\alpha _1\mathrm{}\alpha _{16}}\mathrm{\Gamma }_{\alpha _1\alpha _2}^{a_1a_2}\overline{}_{a_1}_{a_2}\mathrm{\Gamma }_{\alpha _3\alpha _4}^{a_3a_4}\overline{}_{a_3}_{a_4}\mathrm{\Gamma }_{\alpha _5\alpha _6}^{a_5a_6}\overline{}_{a_5}_{a_6}\mathrm{\Gamma }_{\alpha _7\alpha _8}^{a_7a_8}\overline{}_{a_7}_{a_8}`$ (9)
$`(/\mathrm{\Gamma }^{ab}/)_{\alpha _9\alpha _{10}}(/\mathrm{\Gamma }^{bc}/)_{\alpha _{11}\alpha _{12}}(/\mathrm{\Gamma }^{cd}/)_{\alpha _{13}\alpha _{14}}(/\mathrm{\Gamma }^{da}/)_{\alpha _{15}\alpha _{16}}{\displaystyle \frac{1}{(z\overline{z})^8}},`$
where the precise (non-zero) combinatorial coefficient in this relation is not of interest, as we only need to show the non-vanishing of $`|\mathrm{\Psi }`$. In order to proceed we note that
$$(/\mathrm{\Gamma }^{ab}/)_{[\alpha \beta ]}=\mathrm{\Gamma }_{\alpha \beta }^{ab}+4^{[a}\mathrm{\Gamma }^{b]c}{}_{\alpha \beta }{}^{}_{c}^{}.$$
(10)
Hence (9) may be reduced to a differential operator in $`_a`$ and $`\overline{}_a`$ of degree 16 acting on $`(z\overline{z})^8`$ provided we know the precise form of the 16 index tensor $`t_{(16)}^{a_1\mathrm{}a_{16}}`$
$$t_{(16)}^{a_1\mathrm{}a_{16}}=ฯต^{\alpha _1\mathrm{}\alpha _{16}}\mathrm{\Gamma }_{\alpha _1\alpha _2}^{a_1a_2}\mathrm{\Gamma }_{\alpha _3\alpha _4}^{a_3a_4}\mathrm{}\mathrm{\Gamma }_{\alpha _{15}\alpha _{16}}^{a_{15}a_{16}}.$$
(11)
Clearly $`t_{(16)}^{a_1\mathrm{}a_{16}}`$ must be expressable in form of a large string of space-indexed $`\delta `$-functions, the $`ฯต^{a_1\mathrm{}a_9}`$ symbol cannot appear. Interestingly enough this tensor also appears in the leading one-loop quantum correction to the M-theory effective action contracted with four Riemann tensors . Its precise form can be computed and is most conveniently written down in a form contracted with an antisymmetric auxiliary tensor $`X^{ab}`$
$`t_{(16)}^{a_1\mathrm{}a_{16}}X^{a_1a_2}\mathrm{}X^{a_{15}a_{16}}`$ $`=`$ $`1052^{19}[5(\text{tr}X^2)^4+384\text{tr}X^8256\text{tr}X^2\text{tr}X^6`$ (12)
$`+72(\text{tr}X^2)^2\text{tr}X^448(\text{tr}X^4)^2]`$
where the product of $`X`$ is to be understood in the matrix sense.
The knowledge of $`t_{(16)}^{a_1\mathrm{}a_{16}}`$ now enables us to finally evaluate (9) using (10), which is still rather involved and most effectively done with the help of the computer algebra system FORM . Our final result reads
$$|\mathrm{\Psi }(^2)^4[(\overline{})^2^2\overline{}^2]^2\frac{1}{(z\overline{z})^8}=(^2)^6(\overline{}^2)^2\frac{1}{(z\overline{z})^8},$$
(13)
which is non-vanishing and completes the proof of the non-triviality of (6).
Acknowledgement
J.H. would like to thank M. Bordemann and R. Suter for previous collaborations on the subject, J.P. thanks R. Helling and H. Nicolai for valuable discussions. |
warning/0002/cond-mat0002301.html | ar5iv | text | # Evidence for nodal quasiparticles in electron-doped cuprates from penetration depth measurements
## Abstract
The in-plane magnetic penetration depth, $`\lambda (T)`$, was measured down to 0.4 K in single crystals of electron-doped superconductors, Pr<sub>1.85</sub>Ce<sub>0.15</sub>CuO<sub>4-ฮด</sub> (PCCO) and Nd<sub>1.85</sub>Ce<sub>0.15</sub>CuO<sub>4-ฮด</sub> (NCCO). In PCCO, the superfluid density varies as $`T^2`$ from 0.025 up to roughly 0.3 $`T/T_c`$ suggestive of a d-wave state with impurities. In NCCO, $`\lambda (T)`$ shows a pronounced upturn for $`T<4`$ K due to the paramagnetic contribution of Nd<sup>3+</sup> ions. Fits to an s-wave order parameter over the standard BCS range ($`T/T_c`$ = 0.32) limit any gap to less than $`\mathrm{\Delta }_{min}(0)/T_c=0.57`$ in NCCO. For PCCO, the absence of paramagnetism permits a lower temperature fit and yields an upper limit of $`\mathrm{\Delta }_{min}(0)/T_c=0.2`$.
There is by now a consensus that the hole-doped high-$`T_c`$ cuprates exhibit $`d`$wave pairing symmetry . For electron-doped cuprates the issue remains unresolved. While most theories for the mechanism of high temperature superconductivity are insensitive to the sign of the carriers, some predict that n and p type materials will have different pairing symmetry, making its determination an important challenge . Early microwave measurements of the penetration depth in NCCO were interpreted within an $`s`$wave model . However, Cooper pointed out that the power law dependence for $`\lambda (T)`$ indicative of a nodal order parameter could be masked by a large paramagnetic contribution from $`Nd^{+3}`$ ions . Newer microwave measurements by Kokales et al., performed on the same sample used in this paper, have revealed an upturn and power-law temperature dependence and are consistent with our data . Measurements of $`\lambda (T)`$ using single grain boundary junctions have favored a gapped state. Some tunneling measurements favor an $`s`$wave order parameter, albeit with significant departures from an isotropic weak coupling BCS picture , while others report a zero-bias conductance peak . Half-integral flux indicative of d-wave pairing was recently reported in tricrystal experiments with both NCCO and PCCO films.
In this letter we report measurements of $`\lambda (T)`$ down to 0.4 K in single crystals of both NCCO and PCCO. Lower temperatures and higher resolution combine to permit a more precise determination of the temperature dependence of $`\lambda (T)`$ than any previously reported. In NCCO, a large paramagnetic contribution is observed below 4 K. In non-magnetic PCCO, we find an overall $`T^2`$ variation of the superfluid density up to $`T/T_c0.3`$, suggesting the presence of nodal quasiparticles in the presence of strong impurity scattering.
Single crystals of R<sub>1.85</sub>Ce<sub>0.15</sub>CuO<sub>4-ฮด</sub> (R=Nd or Pr) were grown using directional solidification technique and annealed in argon to achieve optimal superconducting properties . Penetration depth was measured using an 11 MHz tunnel-diode driven LC resonator . Samples were mounted on a movable sapphire stage with temperature controllable from 0.4 K to 100 K. The low noise level, $`\mathrm{\Delta }f_{min}/f_05\times 10^{10}`$, results in a sensitivity of $`\mathrm{\Delta }\lambda 0.5`$ ร
for our samples \[$`0.5\times 0.5\times 0.02`$ mm\]. The large anisotropy of these materials ($`\lambda _c/\lambda _{ab}`$ 30-80 ) forces one to apply the rf field perpendicular to the conducting planes. Otherwise, the frequency shift will be dominated by changes of the interplane penetration depth, for which there exists no straightforward connection to the pairing symmetry. A semi-analytical solution for the rf susceptibility of a platelet sample of square base $`2w\times 2w`$ and thickness $`2d`$ in this orientation was analyzed in detail in Ref. . At low temperatures the frequency shift, $`\mathrm{\Delta }ff(T)f(0)`$, is related to the change in penetration depth, $`\mathrm{\Delta }\lambda \lambda (T)\lambda (0)`$, via $`\mathrm{\Delta }f=G\mathrm{\Delta }\lambda `$, with the calibration constant $`G=V_sf_0/\left[2V_0(1N)R\right]`$, where $`N`$ is the effective demagnetization factor, $`V_s`$ is the sample volume, $`V_0`$ is the effective coil volume, $`f_0`$ is the resonance frequency in the absence of a sample and $`R0.2w`$ is an effective dimension . Although this result is similar to the known solution for an infinite slab in parallel field , the effective dimension $`R`$ differs significantly from $`R=w/2`$ obtained for an infinite bar. This difference is due to penetration of the magnetic field from the top and bottom surfaces. The sample and apparatus dependent constant $`\mathrm{\Delta }f_0V_sf_0/\left[2V_0(1N)\right]`$ is measured by moving the sample out of the coil in situ. The overall calibration was tested with samples of Nb, YBCO and BSCCO and gave $`d\lambda /dT`$ within 10 % of reported values . In order to determine the normalized superfluid density, $`\rho _s\left[\lambda (0)/\lambda (T)\right]^2`$, it is necessary to know the absolute magnitude of the penetration depth, $`\lambda (0)`$. Measured values of $`\lambda (0)`$ in NCCO vary from 1000 \- 2600 ร
. In PCCO the only reported value, $`\lambda (0)1000`$ ร
, was estimated from the measurements of the lower critical field $`H_{c1}`$ and is less reliable due to demagnetization and possible vortex effects . We recently developed a new technique for determining $`\lambda (0)`$ from the frequency shift obtained by warming a sample coated with a thin layer of Al above $`T_c(Al)`$. A detailed description of this procedure will be published elsewhere . This technique applied to PCCO gave $`\lambda (0)=2500\pm 100\AA `$ which will be used as an upper estimate in this paper.
In Fig.1 we plot the penetration depth for PCCO sample 1 along with several fits. The fitting range was 5 K to assure validity of the low temperature BCS expansion for an isotropic s-wave state, $`\mathrm{\Delta }\lambda =\lambda (0)\sqrt{\pi \mathrm{\Delta }(0)/2T}\mathrm{exp}(\mathrm{\Delta }(0)/T)`$. In each case, the small negative offset $`A=\lambda (0)\lambda (0.4K)`$ was determined as a fit parameter. The solid line is a fit with 188 data points to a power law, $`\mathrm{\Delta }\lambda `$ = $`A+BT^2`$ with $`B=3.70\pm 0.01\AA /K^2`$ and $`\chi ^28.2`$. The short dotted line shows the best fit to the BCS s-wave expression. With both $`\lambda (0)`$ and $`\mathrm{\Delta }(0)`$ as free parameters we obtained $`\mathrm{\Delta }(0)/T_c=0.538\pm 0.002`$ and $`\lambda (0)=426\pm 3`$ ร
. The s-wave fit is somewhat worse than the power law ($`\chi ^217.6`$) and gives an unrealistic value for $`\lambda (0)`$. The dash-dotted line shows the s-wave fit where $`\mathrm{\Delta }(0)/T_c`$ was fixed at the weak-coupling BCS value (1.76). In this case an unrealistically large $`\lambda (0)=13570\pm 50`$ ร
was obtained. For comparison, the dotted line is a plot of the BCS expression with $`\mathrm{\Delta }(0)/T_c=1.76`$ and set to a more realistic value of $`\lambda (0)=2500`$ ร
.
If the order parameter is an anisotropic s-wave, then the minimum gap value determines the low temperature asymptotic behavior. The BCS functional form for $`\mathrm{\Delta }\lambda `$ still holds, but with $`\mathrm{\Delta }_{min}`$ replacing the isotropic gap. The temperature range over which this asymptotic form is valid is now reduced accordingly. For an isotropic gap, $`\mathrm{\Delta }(0)/T_c=1.76`$, the range of validity in reduced temperature is $`(T/T_c)_{max}`$ = $`t_{max}0.32`$ . For an anisotropic gap, simple rescaling forces the range of validity down to $`t_{max}0.18\mathrm{\Delta }_{min}(0)/T_c`$. Without a priori knowledge of $`\mathrm{\Delta }_{min}(0)/T_c`$, we do not know $`t_{max}`$ and so it is necessary to successively reduce the range until the gap value obtained from the fit becomes independent of the range. Following this procedure for PCCO, with $`\lambda (0)=2500`$ ร
fixed, we find that $`\mathrm{\Delta }(min)/T_c`$ extrapolates to $`0.20\pm 0.05`$ as $`t_{max}0`$. The same procedure for Nb yields $`\mathrm{\Delta }(0)/T_c=1.74\pm 0.02`$, as expected. This stricter criterion would imply that any residual gap is less than $`11\%`$ of the isotropic BCS value.
The overall best fit for sample 1 was achieved for a $`T^{2.25\pm 0.01}`$ power law ($`\chi ^22.1`$). Although this may appear unphysical, it is important to recall that the integer power laws expected for a nodal order parameter strictly apply to the normalized superfluid density, $`\rho _s`$, and not the measured quantity, $`\mathrm{\Delta }\lambda `$. If $`\rho _s=1c_nT^n`$ then $`\mathrm{\Delta }\lambda `$ has corrections of order $`T^{2n}`$ and a fit to $`\mathrm{\Delta }\lambda `$ can result in an artificial intermediate power (e.g. 2.25). This distinction is clearly evident in high quality, untwinned YBCO above 10 K .
In Fig. 2 we plot, for sample 1, normalized superfluid density, $`\rho _s\left[1+\mathrm{\Delta }\lambda (T)/\lambda (0)\right]^2`$, vs. $`(T/T_c)^2`$ for three choices of $`\lambda (0)`$ spanning the range of reported values. $`\lambda (0)=`$1200 ร
yields $`1\rho _s(T/T_c)^2`$ up to 8.4 K, while larger choices for $`\lambda (0)`$ reduce the range of pure quadratic behavior. For comparison we show data taken for Nb in the same apparatus. Up to to $`T/T_c=0.5`$, the Nb data fits the low temperature BCS expansion perfectly with $`\mathrm{\Delta }(0)/Tc=1.74\pm 0.02`$, giving us confidence that the measurement technique is sound.
With an exponent of $`n=2.25`$, sample 1 is our weakest candidate for a nodal order parameter. Fig. 3 shows data for samples 1, 2 and 3 (offset for clarity) and the same Nb data with a fit to the BCS form. Samples 2 and 3 have power laws much closer to n = 2. In this plot we have chosen the largest value of $`\lambda (0)`$ in order to cast the power law model in the most unfavorable light (i.e., smallest range of pure quadratic behavior). We conclude that the superfluid density in PCCO is best described by a quadratic power law variation with temperature.
There are currently two theories for a quadratic power law in d-wave superconductors. Kosztin and Leggett showed that the divergence of the effective coherence length near the nodes of a $`d`$wave order parameter yields $`1\rho _sT^2`$ due to nonlocal electrodynamics . Nonlocality is predicted to arise for the orientation used here (field perpendicular to conducting planes) below $`T_{nonlocal}\xi (0)\mathrm{\Delta }(0)/\lambda (0)`$, where $`\xi (0)`$ is the coherence length. In electron-doped cuprates $`T_{nonlocal}`$ 0.5 K to 2.5 K within our current knowledge of superconducting parameters. Since we observe a quadratic temperature dependence up to 8-10 K in some samples, nonlocality is unlikely to be the source. A stringent test for nonlocality would require a comparison between this data and $`\lambda (T)`$ obtained from the $`H||ab`$ plane orientation. However, the $`H||ab`$ orientation involves the interplane penetration depth, as discussed earlier.
Impurity scattering in the unitary limit provides a more plausible explanation for the quadratic dependence of $`\rho _s(T)`$ . In the โdirty $`d`$waveโ scenario, $`\rho _s(T)`$ will cross over from a linear to quadratic temperature dependence below $`T^{}6\mathrm{ln}2\gamma /\pi `$, where $`\gamma 0.63\sqrt{\mathrm{\Gamma }\mathrm{\Delta }(0)}`$ and $`\mathrm{\Gamma }`$ is a scattering rate parameter, proportional to the impurity concentration . The slope $`d\lambda /dT^2|_{T0}\pi \lambda (0)/\left(6\gamma \mathrm{\Delta }(0)\right)`$. Casting this result in dimensionless form gives: $`\mathrm{\Gamma }/T_c0.28/\left[d\rho _s/d(T/T_c)^2\right]`$. In Fig. 4 we plot $`T_c`$ versus $`\mathrm{\Gamma }/T_c`$ for all five samples studied. Samples 1,2 and 3 are marked. The transitions in some samples were broad and two different criteria were used to estimate $`T_c`$ \- onset of the diamagnetic signal and the inflection point of the $`\mathrm{\Delta }\lambda (T)`$ curve. In general, the trend shows $`T_c`$ suppression with increased scattering rate as expected for a d-wave state with impurity scattering . $`\mathrm{\Gamma }/T_c`$ is at least 10 times larger than the scattering rate observed in clean YBCO .
We now discuss measurements of $`\lambda (T)`$ in NCCO. This compound has been studied much more thoroughly than PCCO and was cited as the first evidence for $`s`$wave pairing in e-doped materials . However, Nd<sup>3+</sup> ions introduce a large paramagnetic background and influence the measured penetration depth . With magnetic permeability $`\mu (T)=1+C/\left(\mathrm{\Theta }+T\right)`$, the measured penetration depth is given by $`\lambda (T)=\lambda _L(T)\sqrt{\mu (T)}`$, where $`\lambda _L`$ is the London penetration depth, $`C`$ is a Curie-Weiss constant, and $`\mathrm{\Theta }`$ is the characteristic temperature for antiferromagnetic interaction. To fit the data,we take $`\mathrm{\Theta }=1.2`$ K from neutron scattering and specific heat measurements. For $`C`$ we have chosen two representative values $`C=0.3`$ and $`0.05`$, calculated assuming the effective magnetic moment of Nd<sup>3+</sup> ions to be $`2.4\mu _B`$ and $`1\mu _B`$, respectively.
Figure 5 shows $`\mathrm{\Delta }\lambda (T)`$ measured in a single crystal of NCCO. The inset shows the low-temperature range. Below $`T4`$ K there is a pronounced upturn, which we attribute to the paramagnetic contribution of Nd<sup>3+</sup> ions. The upper solid line in the inset to Fig.5 shows the power law fit ($`\lambda _LT^n`$) which yields $`n=1.35\pm 0.03`$ and $`n=1.40\pm 0.03`$ for $`C=0.3`$ and $`0.05`$, respectively. (Fits for the two different values of $`C`$ are indistinguishable on this scale.) The lower line shows analogous fits to the low temperature $`s`$wave expression from which we obtain $`\mathrm{\Delta }(0)/T_c=0.569\pm 0.006`$ and $`\mathrm{\Delta }(0)/T_c=0.573\pm 0.006`$. Fits were obtained from data up to $`t=0.32`$. The higher temperature data is shown for completeness. Changing the value of $`\mathrm{\Theta }`$ from 1.2 to 2 also had a small effect on the fit parameters. The value of $`\mathrm{\Delta }(0)/T_c`$ is close to that obtained in PCCO. Again, for a strict test we should fit over a correspondingly reduced temperature range. However, the dominant paramagnetic contribution below $`t<0.16`$ renders this procedure meaningless. The n = 1.4 exponent obtained from the power law fit is closer to the clean d-wave limit and could imply that after correction for paramagnetism, unitary-limit scattering in NCCO is smaller than in PCCO.
In conclusion, we have measured the penetration depth $`\lambda (T)`$ in electron-doped PCCO and NCCO single crystals down to 0.4 K. In non-magnetic PCCO, $`\rho _s`$ decreases quadratically with temperature up to $`t`$ 0.3, consistent with a dirty d-wave (gapless) scenario. The correlation between $`T_c`$ and the rate of change of superfluid density is also consistent with this picture. In NCCO, a large paramagnetic contribution to the penetration depth was observed. $`\lambda _L(T)`$ was found to vary as $`T^{1.4}`$. For both materials, a fit over the same temperature range to an s-wave model sets an upper limit of $`\mathrm{\Delta }(0)/T_c=0.57`$ but requires unrealistically small values $`\lambda (0)`$. For PCCO, the test can be made more stringent and we reduce the upper limit to $`\mathrm{\Delta }_{min}(0)/T_c=0.2`$.
We thank J. R. Cooper, A. A. Abrikosov, S. M. Anlage, N. Goldenfeld, P. J. Hirschfeld, R. Klemm, L. Taillefer, A. Mourachkine, Y. Yeshurun and C. C. Tsuei for useful discussions and communications. This work was supported by the Science and Technology Center for Superconductivity Grant No. NSF-DMR 91-20000. Work in Maryland is supported by Grant No. NSF-DMR 97-32736. |
warning/0002/astro-ph0002314.html | ar5iv | text | # DETECTION OF MULTI-TeV GAMMA RAYS FROM MARKARIAN 501 DURING AN UNFORESEEN FLARING STATE IN 1997 WITH THE TIBET AIR SHOWER ARRAY
## 1 INTRODUCTION
Mrk 501 and Mrk 421 have been well detected as extra-galactic TeV $`\gamma `$-ray sources by Whipple and subsequent ground-based Cerenkov detectors (Ong (1998)). They are the so-called BL Lac objects, which are radio-loud AGNs (Active Galactic Nuclei) whose relativistic jets are aligned along our line of sight. Flux variability on various scales is a common feature of BL Lac objects as already seen in Mrk 501 and Mrk 421, and spectral variations of $`\gamma `$-rays coming from these sources are considered to be a very powerful tool for understanding the physics of BL Lac objects. When Mrk 501 was first detected by the Whipple Collaboration in 1995 (Quinn et al. (1996)), it showed rather low fluxes at a level significantly below the Crab flux. In March of 1997, however, this source went into a state of remarkably flaring activity and its high state lasted for almost half a year with highly variable and strong $`\gamma `$-ray emission. The maximum flux reached roughly 10 times that of the Crab. During this period, several groups (Protheroe et al. (1997)) observed strong $`\gamma `$-ray emission from this source with imaging atmospheric Cerenkov detectors. Independent measurements of the $`\gamma `$-ray spectrum seem to show a gradual softening towards higher energy, while the systematic uncertainties in the flux estimates remain too large to reach a common understanding. The energy spectrum and its shape are very important quantities for clarifying the mechanism of $`\gamma `$-ray production or particle acceleration at the source, and eventually to lead to the actual measurement of the intergalactic infrared or far-infrared background field (Stecker & de Jager (1993)). Hence, confirmation of the detection of $`\gamma `$-rays with a different technique will be strongly required.
The Tibet air shower array, operating since 1990, is located at Yangbajing (4300 m above sea level) in Tibet (Amenomori et al. (1992)). This array has a capability of detecting $`\gamma `$-rays in the TeV energy region with high efficiency and good angular resolution. Using this array, we have succeeded in detecting the Crab at the 5.5 $`\sigma `$ level (Amenomori et al. (1999)). In this paper we present the observation of multi-TeV $`\gamma `$-ray flares from Mrk 501 in 1997. The result obtained with well established air shower technique is important for comparing with those by imaging atmospheric Cerenkov telescopes.
## 2 EXPERIMENT
The Tibet air shower array consists of two overlapping arrays (Tibet-II and HD) as described elsewhere (Yuda (1996)). The Tibet-II array comprises 185 scintillation detectors (BICRON 408A) of 0.5 m<sup>2</sup> each placed on a 15 m square grid with an enclosed area of 36,900 m<sup>2</sup>, and the HD (high density) array is operating inside the Tibet-II array to detect cosmic ray showers with energies lower than 10 TeV. This HD array consists of 109 scintillation detectors (some of detectors are commonly used in both arrays), placed on a 7.5 m square grid covering an area of 5,175 m<sup>2</sup>. The detector arrangement of the Tibet air shower array is schematically shown in Fig. 1. Every detector, except those placed with a 30 m spacing on the outskirts of the inner detector matrix of the Tibet-II array, is equipped with a fast timing (FT) phototube (HPK H1161) and is thus referred to as an โFT-detectorโ. A lead plate of 5 mm thickness is placed on the top of each detector to improve the fast timing data by converting $`\gamma `$-rays in the showers to electron pairs. This lead converter typically increases the shower size by a factor of about 2 and improves the angular resolution by about 30 % (Amenomori et al. (1990)).
All the TDCs (time-to-digital converters) and ADCs (analog-to-digital converters) are regularly monitored by using a calibration module in the FASTBUS system at every 20 minutes. The length of each signal cable is also monitored by measuring a mismatched-reflection pulse from each detector. The data-taking system has been operating under any 4-fold coincidence in the FT-detectors, resulting in that the trigger rate of the events is about 200 Hz for the Tibet-II while being about 115 Hz for the HD array.
The observation presented here was made by using the data taken between 1997 February and 1997 August. The event selection was done by imposing the following three conditions to the recorded data : 1) Each of any four FT detectors should record a signal of more than 1.25 particles ; 2) among the four detectors recording the highest particles, two or more should be within each detector area of the Tibet II and HD arrays denoted by the dotted and solid lines, respectively, in Fig.1 ; and 3) the zenith angle of the incident direction should be less than 45. After data processing and quality cuts, the total number of events selected were $`5.5\times 10^8`$ for the HD array and $`1.0\times 10^9`$ for the Tibet-II array, respectively, with the effective running time of 155.3 days.
## 3 ARRAY PERFORMANCE
Since the background cosmic rays are isotropic and $`\gamma `$-rays from a source are apparently centered on the source direction, a bin size for collecting on-source data should be determined based on the arrayโs angular resolution so as to optimize the signal to noise ratio. In order to achieve a good resolution, a study of core-finding techniques and shower-front curvature corrections has been done (Amenomori et al. (1990)). The angular accuracy of the Tibet array can be checked thoroughly by observing the shadow that the Moon casts in the cosmic rays (Amenomori et al. (1993)). The Tibet II and HD arrays have a capability of observing the Moonโs shadow with good statistics. The mode energies of primary protons to be detected are about 3 TeV and about 8 TeV for the Tibet HD and II arrays, respectively. Hence the angular resolution of each array can be independently examined in respective energy regions. The statistical significance of the Moonโs shadow observed with both arrays becomes about 10$`\sigma `$ or more for half a year observation. From this observation, we estimated the angular resolution of both arrays to be better than 0.9 for all events. We have also found that the angular resolution scales with $`\rho `$, where $`\rho `$ stands for the sum of the number of shower particles per m<sup>2</sup> detected in each counter. The resolution increases with increasing $`\rho `$ as $`0.8^{}\times ((\rho )/20)^{0.3}`$ ( $`15<\rho <300`$ ).
The Moonโs shadow by the events with $`\rho `$ = 15-50 was found at the position shifted from the Moon center to the west by 0.32 ($`\pm 0.10^{}`$). The primary cosmic rays casting the Moonโs shadow are almost protons and the mean energy of protons capable of generating these events at Yangbajing is estimated to be about 4.7 TeV by the simulation. On the other hand, a proton of energy E impinging at normal angle on the Earth is deflected by the geomagnetic field and its deflection angle is calculated as $`\mathrm{\Delta }\theta E1.6^{}`$ TeV. So, the observed shift of the Moonโs shadow is consistent with that expected from the effect of the geomagnetic field. A more elaborate study of the Moonโs shadow using a Monte Carlo technique shows almost same results as those by the experiment (Suga et al. (1999)). Thus, the results obtained by assigning primary energies to the observed events can be directly checked by observing the Moonโs shadow.
The pointing of the array is inferred from the position of the Moonโs shadow by high energy cosmic rays ($`>20`$ TeV) which are negligibly affected by the geomagnetic field. This estimation can also be done by examining the deviations of the Moonโs shadow in the north-south direction, since the effect of the geomagnetic field acts only in the east-west direction. It is then found to be smaller than 0.1 for both arrays.
Figure 2 shows the cumulative deficit counts of the events coming from the direction of the Moon as a function of MJD, obtained with the HD array. The data set used are between February 1997 and August 1997, just corresponding to the observation period of Mrk 501. A linearly increasing of the deficit events may be a sure guarantee against the long-term stability of the array operation. Th meridian zenith angle of the Moon at Yangbajing changes between 12 and 50 every 27.3 days. Naturally the most efficient observations are done when the Moon comes in sight around the smallest zenith angle of about 12 every 27.3 days. This effect will be found as a tier-like structure on the deficit curve, as seen in Fig. 2.
## 4 RESULTS AND DISCUSSIONS
A circular window was used to search for signals and then its size was determined based on the angular resolution estimated by the experiment. The window size is chosen to optimize the significance of signals defined by $`N_s/N_B^{1/2}`$, where $`N_S`$ is the number of signals and $`N_B`$ the number of background events, and to contain more than 50 % of the signals from a source. The radii of search windows used for the events with $`\rho >`$ 15, 50 and 100 were 0.9, 0.8 and 0.5, respectively. The signals were searched for by counting the number of events coming from the on-source window. The background was estimated by averaging over events falling in the ten off-source windows adjacent to the source, but without overlapping each other. The source window traverses a path in local coordinates expressed by the zenith angle and azimuth angle through every day. In order to reduce a strong zenith angle dependence of the background, these off-source windows were taken in the azimuth angle directions with the same zenith angle, except two windows adjacent to the on-source window.
Figure 3 shows the cumulative excess counts for all events as a function of MJD and background, obtained with the HD array. No excess counts were observed until the middle of March 1997. However, excess events rapidly increased in the period from April through June and then it became slightly dull. The operation of the array was stopped on August 25 of 1997 to calibrate the operation system. As discussed in ยง3, one should first note that the observed excess counts are by no means due to some artificial noise or unstable operation of the system. The statistical significance of the excess counts reached a 3.7 $`\sigma `$ level during this period. The excess counts very rapidly increased during the period from April 7 through June 16 and the statistical significance of the excess counts was a 4.7 $`\sigma `$. These observed features are almost consistent with other observations by atmospheric Cerenkov telescopes (Protheroe et al. (1997)).
Shown in Fig. 4 is the contour map of the excess event densities around Mrk 501 for the events with $`\rho >`$ 15 observed between April 7 and June 16 in 1997. This map was obtained using the same method as done for the Moon and Sun shadows (Amenomori et al. (1993)). Mrk 501 is well observed in the right direction by our air shower array.
Figure 5 shows the distribution of the opening angles relative to the Mrk 501 direction for all events with $`\rho >`$ 15 in the HD array. The excess in the small opening angle region (less than 0.5) could be attributed to $`\gamma `$-rays from Mrk 501. The simulation result done for $`\gamma `$-ray events coming from Mrk 501 can well reproduce the experiment as shown in Fig. 5, when we take account of the systematic pointing errors estimated in ยง3. For the observation period from 1997 February to 1997 August, the statistical significances of the excess events with $`\rho >`$ 15, 30 and 50 were 3.7 $`\sigma `$, 2.3 $`\sigma `$ and 1.6 $`\sigma `$, respectively.
We also searched for $`\gamma `$-ray emission using the entire Tibet-II array, but no excess was found in this period and upper limits on the excess number of the events at the 90 % confidence level were obtained.
We estimated the $`\gamma `$-ray spectrum from Mrk 501 by a Monte Carlo simulation (we used a GENAS code by Kasahara & Torii (Kasahara & Torii (1991)).), assuming a differential power-law spectrum with the form $`E^\beta `$ and the cut-off at a certain energy, $`E_c`$, where the cut-off means that the spectral slope steepens by 1.0 at $`E_c`$. The value of $`\beta `$ was changed between 2.4 and 2.7 and also the effect of $`E_c`$ was examined between 7 TeV and 30 TeV. Primary $`\gamma `$-rays with energies between 0.2 TeV and 50 TeV were thrown from the direction of Mrk 501. Observation of simulated events at Yangbajing level was done as in our experiment, estimating the collecting area, trigger efficiency and threshold energy for $`\gamma `$-rays generating the events at observation level. Simulated events in respective size ($`\rho `$) bins were then compared with those by the experiment. The energy of $`\gamma `$-rays was defined as the energy of the maximum flux of simulated events observed in each size bin. These steps were repeated until the observed results are well reproduced. A combination of $`\beta 2.6`$ and $`E_c`$ 20-30 TeV can reproduce the data well. We examined that the absolute flux values except the highest energy bin stay almost unchanged for above trials, but it is of course difficult to settle the spectral slope from this experiment because of very small energy range fitted here. The systematic errors on the flux arise mainly from the event selection procedure, which depends upon the array performance, and from the calculations of the collecting area and the air shower size distribution by the simulation. They are estimated to be 13 % and 8 %, respectively (Amenomori et al. (1999)).
Shown in Figs. 6 and 7 are the energy spectra averaged in the period from February 15 to August 25 in 1997 and from February 15 to June 8 in 1997, respectively. The latter observation time corresponds to that of the Whipple Collaboration (Samuelson et al. (1998)). It is seen that the results reported recently by other experiments (Samuelson et al. (1998) ; Hayashida et al. (1998) & Konopelko et al. (1999)) are almost compatible with ours, although these do not cover the same observation times. It should pay attention, however, that our results were obtained by the continuous observation of Mrk 501 extending February through August in 1997, while those by Cerenkov telescopes were obtained for very limited periods of moonless and cloudless nights.
Mrk 501 and Mrk 421, nearby AGNs, are at almost the same red-shift (0.033 and 0.031, respectively) and have been detected in TeV energies (Ong (1998)). In particular, Mrk 501 during the strong, long-lasting 1997 flare provided a good opportunity to study the energy spectrum of $`\gamma `$-rays from this source in detail (Protheroe et al. (1997)), suggesting a spectral feature different with that of Mrk 421 (Krennrich et al. (1999)). In both sources, it is likely that a synchrotron-inverse Compton picture plays an important part (Protheroe et al. (1997)). Since the attenuation mechanism of TeV $`\gamma `$-rays by intergalactic infra-red photon field is almost the same for both sources, a difference of spectral features, if any, could be attributed to the production mechanism of $`\gamma `$-rays at the sources. Therefore, it is very important to continue the observation of high energy $`\gamma `$-rays from both sources with as small uncertainties as possible.
## 5 SUMMARY
Mrk 501 suddenly came into a very active phase from March in 1997, with several large flares and lasted for $`1/2`$ yr. The maximum $`\gamma `$-ray flux during this period reached about 10 times as high as the Crab Nebula. Following a successful observation of steady emission of multi-TeV $`\gamma `$-rays from the Crab(Amenomori et al. (1999)), we further detected multi-TeV $`\gamma `$-rays from Mrk 501 which was in a high flaring state between March 1997 and August 1997, and estimated the absolute fluxes of $`\gamma `$-rays around multi-TeV region, using the high resolution Tibet air shower array. The detection of a signal from this source was achieved by the improvement of the array performance, which can be directly checked by observing the Moonโs shadow. Monthly observations of the Moonโs shadow provide a direct check of the angular resolution, pointing accuracy, and also the stable operation of the array over a long period. Furthermore, the observation of the displacement of the Moonโs shadow due the effect of the geomagnetic field provides an important check of the results obtained by assigning energies to all the events. This is the first attempt to be done in the air shower experiments, and it suggests that the Moon is a unique cosmic-ray anti-source capable of calibrating the array performance thoroughly. Hence, the results obtained by the Tibet experiment using a different technique will be a great help to understand the possible bias and errors involved in the Cerenkov observations.
The area of the present HD array will be extended by a factor of about five in 1999, while its effective area will be increased by a factor of about seven by the reduction of edge effects. Then, the Tibet array could cover the energy range from 3 TeV to $``$100 TeV with significantly better statistics and angular resolution at high energies. Air shower arrays are wide aperture and high duty cycle instruments, in contrast to atmospheric Cerenkov telescopes with relatively narrow fields of view and small duty cycle of $``$10 %. These features will be indispensable for understanding a time variability of emission of high energy $`\gamma `$-rays from point sources such as AGNs and GRBs (gamma ray bursts). The Tibet experiment, therefore, will have unique capabilities for the discovery of new, relatively bright sources and for a general survey of the overhead sky.
This work is supported in part by Grants-in-Aid for Scientific Research and also for International Science Research from the Ministry of Education, Science, Sports and Culture in Japan and for International Science Research from the Committee of the Natural Science Foundation and the Academy of Sciences in China. |
warning/0002/math0002130.html | ar5iv | text | # 1 Introduction and results
## 1 Introduction and results
All algebraic objects in this paper are defined over the ring of integers $`\mathrm{}`$. Our work was motivated by reconsidering the following classical Basic Perturbation Lemma (see and the historical account there).
Recall that strong deformation retract (SDR) data (also called a contraction) are given by chain complexes $`(M,d_M)`$, $`(N,d_N)`$, chain maps $`๐ต:(M,d_M)(N,d_N)`$, $`๐ถ:(N,d_N)(M,d_M)`$ and a chain homotopy $`๐ท:MM`$ satisfying
$`๐ตd_M`$ $`=`$ $`d_N๐ต,`$
$`๐ถd_N`$ $`=`$ $`d_M๐ถ,`$ (1)
$`\mathrm{๐ถ๐ต}\mathrm{๐ท๐ท}_M`$ $`=`$ $`d_M๐ท+๐ทd_M\text{ and}`$
$`\mathrm{๐ต๐ถ}`$ $`=`$ $`\mathrm{๐ท๐ท}_N.`$
This of course means that $`(N,d_N)`$ is a strong deformation retract of $`(M,d_M)`$. One usually assumes that the following side conditions (also called annihilation properties) hold:
$$\mathrm{๐ท๐ท}=0,\mathrm{๐ท๐ถ}=0\text{ and }\mathrm{๐ต๐ท}=0.$$
(2)
Then the following statement is true.
Basic Perturbation Lemma (BPL). Suppose we are given strong deformation data (1) satisfying (2) and a perturbation $`\stackrel{~}{d}_M`$ of the differential $`d_M`$ on M. Then there are perturbations $`\stackrel{~}{d}_N,\stackrel{~}{๐ต},\stackrel{~}{๐ถ}`$ and $`\stackrel{~}{๐ท}`$ of $`d_N,๐ต,๐ถ`$ and $`๐ท`$ that again form strong deformation data (1),
$$\stackrel{~}{๐ต}\stackrel{~}{d}_M=\stackrel{~}{d}_N\stackrel{~}{๐ต},\stackrel{~}{๐ถ}\stackrel{~}{d}_N=\stackrel{~}{d}_M\stackrel{~}{๐ถ},\stackrel{~}{๐ถ}\stackrel{~}{๐ต}\mathrm{๐ท๐ท}_M=\stackrel{~}{d}_M\stackrel{~}{๐ท}+\stackrel{~}{๐ท}\stackrel{~}{d}_M\text{ and }\stackrel{~}{๐ต}\stackrel{~}{๐ถ}=\mathrm{๐ท๐ท}_N.$$
All notions used in the formulation of the BPL are standard and we believe it is not necessary to repeat their definitions here. Filtered objects and perturbations are treated in Section 3. The perturbation $`(\stackrel{~}{d}_N,\stackrel{~}{๐ต},\stackrel{~}{๐ถ},\stackrel{~}{๐ท})`$ is given by the following explicit formulas (see again ):
$$\begin{array}{ccc}\stackrel{~}{d}_N\hfill & =& \hfill d_N+๐ต(_M+_M๐ท_M+_M๐ท_M๐ท_M+_M๐ท_M๐ท_M๐ท_M+\mathrm{})๐ถ,\\ \stackrel{~}{๐ต}\hfill & =& \hfill ๐ต+๐ต(_M+_M๐ท_M+_M๐ท_M๐ท_M+_M๐ท_M๐ท_M๐ท_M+\mathrm{})๐ท,\\ \stackrel{~}{๐ถ}\hfill & =& \hfill ๐ถ+๐ท(_M+_M๐ท_M+_M๐ท_M๐ท_M+_M๐ท_M๐ท_M๐ท_M+\mathrm{})๐ถ,\\ \stackrel{~}{๐ท}\hfill & =& \hfill ๐ท+๐ท(_M+_M๐ท_M+_M๐ท_M๐ท_M+_M๐ท_M๐ท_M๐ท_M+\mathrm{})๐ท,\end{array}$$
where $`_M:=\stackrel{~}{d}_Md_M`$. The formulas above contain infinite series, so one must assume some conditions assuring that they converge. This is usually achieved by assuming that both $`(M,d_M)`$ and $`(N,d_N)`$ are filtered complete, see again Section 3.
Our original motivation was to understand why there is such a formula and what is the role of side conditions. As usual, the best way to understand a problem is to formulate it in as general a form as possible. So let us consider the following:
Perturbation Problem (PP). Suppose we are given two complete filtered complexes $`M=(N,d_M)`$ and $`N=(N,d_N)`$ and chain maps $`๐ต:MN`$ and $`๐ถ:NM`$ that are chain homotopy inverse to each other, with homotopies $`๐ท:MM`$ and $`๐ป:NN`$, that is
$$๐ตd_M=d_N๐ต,๐ถd_N=d_M๐ถ,\mathrm{๐ถ๐ต}\mathrm{๐ท๐ท}_M=d_M๐ท+๐ทd_M,\mathrm{๐ต๐ถ}\mathrm{๐ท๐ท}_N=d_N๐ป+๐ปd_N.$$
(3)
Given a perturbation $`\stackrel{~}{d}_M`$ of the differential $`d_M`$, find perturbations $`\stackrel{~}{d}_N,\stackrel{~}{๐ต},\stackrel{~}{๐ถ},\stackrel{~}{๐ท}`$ and $`\stackrel{~}{๐ป}`$ of $`d_N,๐ต,๐ถ,๐ท`$ and $`๐ป`$ such that $`\stackrel{~}{๐ต}`$ and $`\stackrel{~}{๐ถ}`$ are chain maps with respect to the perturbed differentials, homotopy inverse to each other, with homotopies $`\stackrel{~}{๐ท}`$ and $`\stackrel{~}{๐ป}`$, that is
$$\stackrel{~}{๐ต}\stackrel{~}{d}_M=\stackrel{~}{d}_N\stackrel{~}{๐ต},\stackrel{~}{๐ถ}\stackrel{~}{d}_N=\stackrel{~}{d}_M\stackrel{~}{๐ถ},\stackrel{~}{๐ถ}\stackrel{~}{๐ต}\mathrm{๐ท๐ท}_M=\stackrel{~}{d}_M\stackrel{~}{๐ท}+\stackrel{~}{๐ท}\stackrel{~}{d}_M,\stackrel{~}{๐ต}\stackrel{~}{๐ถ}\mathrm{๐ท๐ท}_N=\stackrel{~}{d}_N\stackrel{~}{๐ป}+\stackrel{~}{๐ป}\stackrel{~}{d}_N.$$
Observe that in the formulation of the BPL and the PP we consider not only the differentials and the chain maps, but also the homotopies to be a part of the structure which has to be perturbed. Ignoring homotopies leads to the โcrudeโ perturbation lemma formulated at the end of this Introduction.
The fact, both frustrating and provoking, is that the PP has, for general input data, no solution! โ a rigorous formulation and proof of this negative statement is provided by Theorem 16. The reason is that the chain homotopy equivalence $`(๐ต,๐ถ,๐ท,๐ป)`$ of (3) is not a homotopy invariant concept and it must be replaced by a subtler notion of a strong (chain) homotopy equivalence:
###### Definition 1
A strong homotopy equivalence (SHE) consists of degree $`2m`$ maps $`๐ต_{2m}:MN`$, $`๐ถ_{2m}:NM`$ and degree $`2m+1`$ โhomotopiesโ $`๐ท_{2m+1}:MM`$ , $`๐ป_{2m+1}:NN`$, for all $`m0`$, such that
$$\begin{array}{cccccc}\hfill ๐ต_0d_M& =& d_N๐ต_0,\hfill & \hfill ๐ถ_0d_N& =& d_M๐ถ_0,\hfill \\ \hfill ๐ถ_0๐ต_0\mathrm{๐ท๐ท}_M& =& d_M๐ท_1+๐ท_1d_M,\hfill & \hfill ๐ต_0๐ถ_0\mathrm{๐ท๐ท}_N& =& d_N๐ป_0+๐ป_0d_N\hfill \end{array}$$
and that, for each $`m1`$,
$`d_N๐ต_{2m}๐ต_{2m}d_N`$ $`=`$ $`{\displaystyle \underset{0i<m}{}}(๐ต_{2i}๐ท_{2(mi)1}๐ป_{2(mi)1}๐ต_{2i}),`$
$`d_M๐ท_{2m+1}+๐ท_{2m+1}d_M`$ $`=`$ $`{\displaystyle \underset{0jm}{}}๐ถ_{2j}๐ต_{2(mj)}{\displaystyle \underset{0j<m}{}}๐ท_{2j+1}๐ท_{2(mj)1},`$
$`d_M๐ถ_{2m}๐ถ_{2m}d_N`$ $`=`$ $`{\displaystyle \underset{0i<m}{}}(๐ถ_{2i}๐ป_{2(mi)1}๐ท_{2(mi)1}๐ถ_{2i}),`$
$`d_N๐ป_{2m+1}+d_N๐ป_{2m+1}`$ $`=`$ $`{\displaystyle \underset{0jm}{}}๐ต_{2j}๐ถ_{2(mj)}{\displaystyle \underset{0j<m}{}}๐ป_{2j+1}๐ป_{2(mj)1}.`$
See 8.1 where we expanded the axioms above for some small $`m`$. To understand better the meaning of a SHE, we offer the following analogy.
A homotopy associative algebra is a chain complex $`V=(V,d_V)`$ with a homotopy associative multiplication $`\mu :V^2V`$:
$$\mu (\mu 11)\mu (11\mu )0\text{ modulo a chain homotopy }\nu :V^3V.$$
As argued in , a proper homotopy invariant version of this concept is that of a strongly homotopy associative algebra, which is a structure consisting of infinitely many multilinear operations $`\{\mu _n:V^nV\}_{n2}`$ such that the โmultiplicationโ $`\mu _2:V^2V`$ is homotopy associative up to the homotopy $`\mu _3:V^3V`$, and there is, for each $`n4`$, a certain โcoherence relationโ assumed to be zero modulo the homotopy $`\mu _n`$, see . While each strongly homotopy associative algebra defines, by $`\mu :=\mu _2`$ and $`\nu :=\mu _3`$, a homotopy associative one, the converse is not true; there are obstructions for extending a homotopy associative multiplication to a strongly homotopy associative one.
The situation in Definition 1 is similar. While a strong homotopy equivalence defines, by $`๐ต:=๐ต_0`$, $`๐ถ:=๐ถ_0`$, $`๐ท:=๐ท_1`$ and $`๐ป:=๐ป_1`$ an ordinary homotopy equivalence, the converse is not true โ there is a primary obstruction $`[๐ฌ]`$ for extending a homotopy equivalence to a strong one. The surprising Theorem 11 says that vanishing of this primary obstruction already implies the existence of the extension.
A strong homotopy equivalence of $`M`$ and $`N`$ will be denoted as $`(\underset{ยฏ}{๐ต},\underset{ยฏ}{๐ถ},\underset{ยฏ}{๐ท},\underset{ยฏ}{๐ป}):MN`$. Let us formulate our Ideal Perturbation Lemma.
Ideal Perturbation Lemma (IPL). Suppose we are given two complete filtered complexes $`M=(N,d_M)`$ and $`N=(N,d_N)`$ and a strong homotopy equivalence $`(\underset{ยฏ}{๐ต},\underset{ยฏ}{๐ถ},\underset{ยฏ}{๐ท},\underset{ยฏ}{๐ป}):MN`$.
Given a perturbation $`\stackrel{~}{d}_M`$ of the differential $`d_M`$, there exist a perturbation $`\stackrel{~}{d}_N`$ of the differential $`d_N`$ and a perturbation $`(\underset{ยฏ}{\overset{~}{๐ต}},\underset{ยฏ}{\overset{~}{๐ถ}},\underset{ยฏ}{\overset{~}{๐ท}},\underset{ยฏ}{\overset{~}{๐ป}})`$ of $`(\underset{ยฏ}{๐ต},\underset{ยฏ}{๐ถ},\underset{ยฏ}{๐ท},\underset{ยฏ}{๐ป})`$ which is a strong homotopy equivalence of the perturbed complexes $`(M,\stackrel{~}{d}_M)`$ and $`(N,\stackrel{~}{d}_N)`$. Moreover, the perturbations $`\stackrel{~}{d}_M`$ and $`(\underset{ยฏ}{\overset{~}{๐ต}},\underset{ยฏ}{\overset{~}{๐ถ}},\underset{ยฏ}{\overset{~}{๐ท}},\underset{ยฏ}{\overset{~}{๐ป}})`$ depend functorially on $`\stackrel{~}{d}_M`$ and $`(\underset{ยฏ}{๐ต},\underset{ยฏ}{๐ถ},\underset{ยฏ}{๐ท},\underset{ยฏ}{๐ป})`$.
The IPL is proved in Section 6, see also 8.6 for explicit formulas. As most ideal things, the Ideal Perturbation Lemma is almost useless. In practice, the input data are formulated only in terms of an ordinary homotopy equivalence, and the answer is also expected to be a perturbation of this ordinary homotopy equivalence. Here is our mundane version of the Ideal Perturbation Lemma.
###### Theorem 2
Suppose that the obstruction $`[๐ฌ]`$ to the extension of the homotopy equivalence (3) to a strong one vanishes. Then the Perturbation Problem has a solution, functorial up to a choice of the extension of (3) to a strong homotopy equivalence.
The theorem immediately follows from the IPL and the above notes. There are situations when the obstruction $`[๐ฌ]`$ vanishes and when there even exists a functorial extension of the homotopy equivalence (3) to a SHE. This the case of our motivating example of the Basic Perturbation Lemma (the case $`๐ป=0`$). It immediately follows from Theorem 12 that the side conditions (2) guarantee the existence of a functorial extension of (1) to a strong homotopy equivalence. So Theorem 2 implies the BPL.
Another trick that overrides the nonexistence of a solution of the PP is to change the initial data a bit. We show in Theorem 13 that changing in (3) the homotopy $`๐ท`$ to $`๐ท๐ถ(\mathrm{๐ต๐ท}\mathrm{๐ป๐ต})`$ (or, dually, $`๐ป`$ to $`๐ป๐ต(\mathrm{๐ถ๐ป}\mathrm{๐ท๐ถ})`$) annihilates the obstruction $`[๐ฌ]`$ and we reprove the following recent result by J. Huebschmann and T. Kadeishvili .
###### Theorem 3
Let $`M=(M,d_M)`$ and $`N=(N,d_N)`$ be complete filtered chain complexes and $`(๐ต,๐ถ,๐ท,๐ป)`$ a chain homotopy equivalence (3).
Given a perturbation $`\stackrel{~}{d}_M`$ of $`d_M`$, there exist a perturbation $`\stackrel{~}{d}_N`$ of the differential $`d_N`$ and a homotopy equivalence $`(\stackrel{~}{๐ต},\stackrel{~}{๐ถ},\stackrel{~}{๐ท},\stackrel{~}{๐ป})`$ of the perturbed complexes $`(M,\stackrel{~}{d}_M)`$ and $`(N,\stackrel{~}{d}_N)`$ that is a perturbation of $`(๐ต,๐ถ,๐ท๐ถ(\mathrm{๐ต๐ท}\mathrm{๐ป๐ต}),๐ป)`$.
Changing $`๐ป`$ to $`๐ป๐ต(\mathrm{๐ถ๐ป}\mathrm{๐ท๐ถ})`$ and leaving $`๐ท`$ untouched gives the following complement to Theorem 3.
Complement to Theorem 3. Under the assumption of Theorem 3, there exists another perturbation $`\stackrel{~}{d}_N^{}`$ of the differential $`d_N`$ and another homotopy equivalence $`(\stackrel{~}{๐ต}^{},\stackrel{~}{๐ถ}^{},\stackrel{~}{๐ท}^{},\stackrel{~}{๐ป}^{})`$ of the perturbed complexes $`(M,\stackrel{~}{d}_M)`$ and $`(N,\stackrel{~}{d}_N^{})`$ that is a perturbation of $`(๐ต,๐ถ,๐ท,๐ป๐ต(\mathrm{๐ถ๐ป}\mathrm{๐ท๐ถ}))`$.
Ignoring the homotopies in the Perturbation Problem, we get the following
Crude Perturbation Lemma. Suppose we are given two complete filtered complexes $`M=(N,d_M)`$ and $`N=(N,d_N)`$ and chain maps $`๐ต:MN`$ and $`๐ถ:NM`$ that are chain homotopy inverse to each other.
Given a perturbation $`\stackrel{~}{d}_M`$ of the differential $`d_M`$, there are perturbations $`\stackrel{~}{d}_N,\stackrel{~}{๐ต}`$ and $`\stackrel{~}{๐ถ}`$ of $`d_N,๐ต`$ and $`๐ถ`$ such that $`\stackrel{~}{๐ต}`$ and $`\stackrel{~}{๐ถ}`$ are chain maps with respect to the perturbed differentials, homotopy inverse to each other.
A conceptual explanation of these results is given in Section 7.
โ โ โ โ โ
Plan of the paper: In Section 2 we recall colored operads and introduce the operad $`\mathrm{๐ ๐}`$ describing isomorphisms of chain complexes. In Section 3 we repeat necessary facts on filtrations and perturbations and define the filtered operad $`๐\mathrm{๐๐}`$ describing perturbations of differentials. The filtered operad $`_{\mathrm{iso}}`$ that describes strong homotopy equivalences is introduced in Section 4 where we also discuss extensions of a homotopy equivalence to a strong one. In Section 5 we introduce the operad $`\stackrel{~}{}_{\mathrm{iso}}`$ for perturbations of strong homotopy equivalences and construct a retraction $`r`$ that gives the functorial solution to the IPL. Some of the proofs are postponed to Section 6. In Section 7 we give a conceptual explanation of the results. In the Appendix (Section 8) we present some explicit formulas.
## 2 Language of operads
Roughly speaking, operads are objects that describe types of algebraic systems. Colored operads are then objects describing diagrams of algebraic systems. The definition of a (colored) operad is classical (see or ) and we will not repeat it here in its full generality.
By an operad we will always mean an operad in the symmetric monoidal category $`\mathrm{๐ฒ๐๐๐๐}_{\mathrm{}}`$ of differential graded complexes of abelian groups (that is, complexes of $`\mathrm{}`$-modules). Operads in this category behave in many aspects as associative algebras, so we may speak about suboperads, ideals, presentations, resolutions, etc., see .
All algebraic objects in this paper will have only unary operations. Colored operads describing algebraic systems with only unary operations are the same as small additive categories enriched over $`\mathrm{๐ฒ๐๐๐๐}_{\mathrm{}}`$. This means that all hom-sets are chain complexes and composition maps are homomorphisms of chain complexes. All operads in this paper will be of this type.
We will use the โoperadicโ notation and terminology. Thus, for such an operad/category $`๐ซ`$, we call $`:=\mathrm{๐๐}(๐ซ)`$ the set of colors and, for $`c,d`$, we denote
$$๐ซ\left(\text{}\begin{array}{c}d\\ c\end{array}\text{}\right):=\mathrm{๐๐๐}_๐ซ(c,d).$$
We will usually express the fact that $`f๐ซ\left(\text{}\begin{array}{c}d\\ c\end{array}\text{}\right)`$ by writing $`f:cd`$.
In the particular case when $`\mathrm{card}()=1`$, the $``$-colored operads are exactly differential graded associative unital algebras. In this paper, by a colored operad we always mean an operad colored by the two-point set $`=\{๐ฑ,๐\}`$ ($`๐ฑ`$ from black, $`๐`$ from white) or by a set isomorphic to this one.
###### Example 4
Let $`M=(M,d_M)`$ be a chain complex, then the endomorphism operad $`\mathrm{๐๐}_M`$ is defined to be the chain complex $`\mathrm{๐ป๐๐}(M,M)`$ with the operadic structure (which in this particular case is the same as that of an unital associative algebra) given by the composition. An algebra over an operad $`๐ซ`$ is an operadic homomorphism $`A:๐ซ\mathrm{๐๐}_M`$. In this situation we also say that the operad $`๐ซ`$ acts on the chain complex $`M`$.
###### Example 5
This example describes a colored version of the endomorphism operad recalled in Example 4. Let $`M=(M,d_M)`$ and $`N=(N,d_N)`$ be chain complexes. By a colored endomorphism operad $`\mathrm{๐๐}_{M,N}`$ we mean the full subcategory of $`\mathrm{๐ฒ๐๐๐๐}_{\mathrm{}}`$ with objects $`M`$ and $`N`$. If $`๐ซ`$ is a $`\{๐ฑ,๐\}`$-colored operad, then by a $`๐ซ`$-algebra we mean a homomorphism $`A:๐ซ\mathrm{๐๐}_{M,N}`$ such that $`A(๐ฑ)=M`$ and $`A(๐)=N`$.
###### Example 6
Let $`f:๐ฑ๐`$, $`g:๐๐ฑ`$ be two degree-zero generators and denote
$$\mathrm{๐ ๐}:=\left(\frac{(f,g)}{(fg=1_๐,gf=1_๐ฑ)},d=0\right).$$
In the above display, $`(f,g)`$ denotes the free $`\{๐ฑ,๐\}`$-colored operad on the set $`\{f,g\}`$ and $`(fg=1_๐,gf=1_๐ฑ)`$ the operadic ideal generated by $`fg1_๐`$ and $`gf1_๐ฑ`$.
An algebra $`A:\mathrm{๐ ๐}\mathrm{๐๐}_{M,N}`$ consists of two degree zero chain maps $`๐ต:MN`$, $`๐ถ:NM`$ that are inverse to each other. Thus the operad $`\mathrm{๐ ๐}`$ describes isomorphisms of chain complexes, whence its name.
## 3 Filtrations and Perturbations
Let $`M=(M,d_M)`$ be a chain complex. A (descending) filtration on $`M`$ is a descending sequence $`\{F^pM\}_{p0}`$ of subcomplexes of $`M`$. If not stated otherwise, we always assume that the filtration is complete. This, by definition, means that the module $`M`$ is complete in the $`F^p`$-adic topology. This guarantees that each sum $`_{p0}m_p`$ with $`m_pF^pM`$ represents a unique element of $`M`$. A typical example is the module of power series $`\mathrm{}[[h]]`$ with the filtration defined by $`F^p\mathrm{}[[h]]:=h^p\mathrm{}[[h]]`$, $`p0`$.
Morphisms of filtered chain complexes are maps that preserve filtrations. A linear map $`g:MN`$ is a perturbation or deformation of a linear map $`f:MN`$ if
$$(fg)(F^pM)F^{p+1}N\text{ for each }p0\text{.}$$
If $`M`$ and $`N`$ are filtered complexes, then the chain complex $`\mathrm{๐ป๐๐}(M,N)`$ is also filtered, by
$$F^q\mathrm{๐ป๐๐}(M,N):=\{f\mathrm{๐ป๐๐}(M,N);f(F^pM)F^{p+q}N\text{ for each }p\}.$$
(4)
We believe that the notion of a filtered algebra, operad, etc., is clear; we require that all structure operations preserve the filtration.
If $`M=(M,d_M)`$ is a filtered chain complex, then (4) defines a filtration of the endomorphism operad $`\mathrm{๐๐}_M`$. A filtered algebra over a filtered operad $`๐ซ`$ is a homomorphism $`A:๐ซ\mathrm{๐๐}_M`$ of filtered operads. There is an evident colored analog of this notion.
###### Example 7
Let $`\overline{x}`$ be a generator of degree $`1`$ and let
$$\text{pre}๐\mathrm{๐๐}:=((\overline{x}),d),$$
(5)
with $`d`$ the โderivationโ in the operadic sense defined by $`d\overline{x}:=\overline{x}\overline{x}`$. The free operad $`(\overline{x})`$ on $`\overline{x}`$ is the same as the polynomial ring $`\mathrm{}[\overline{x}]`$. We define the filtration
$$F^p\text{pre}๐\mathrm{๐๐}:=\text{ the subspace spanned by monomials in }\overline{x}\text{ of length }p0\text{.}$$
The differential $`d`$ clearly preserves the filtration, as well as does the composition, so the operad $`\text{pre}๐\mathrm{๐๐}`$ is filtered. Let $`๐\mathrm{๐๐}`$ be the completion of $`\text{pre}๐\mathrm{๐๐}`$ with respect to the above filtration; of course, $`๐\mathrm{๐๐}`$ coincides with the algebra of power series $`\mathrm{}[[\overline{x}]]`$.
Filtered $`๐\mathrm{๐๐}`$-algebras $`A:๐\mathrm{๐๐}\mathrm{๐๐}_M`$ on $`M=(M,d_M)`$ correspond to perturbations $`\stackrel{~}{d}_M=d_M+_M`$ of the differential $`d_M`$, the correspondence being given by $`_M:=A(\overline{x})`$. Indeed, $`d\overline{x}=\overline{x}\overline{x}`$ is mapped by $`A`$ to $`_Md_M+d_M_M=_M_M`$, which is the same as $`(d_M+_M)^2=0`$.
###### Proposition 8
The operad $`๐\mathrm{๐๐}`$ of Example 7 is acyclic, that is, $`H_{}(๐\mathrm{๐๐})\mathfrak{1}`$, where $`\mathfrak{1}`$ is the trivial operad.
The proof of the above proposition is easy and we leave it as an exercise. One feels that the proposition must be โphilosophicallyโ true. Algebras over the trivial operad $`\mathfrak{1}`$ are just chain complexes with no additional structure, i.e. with only the structure given by the unperturbed differential. The operad $`๐\mathrm{๐๐}`$ describes perturbations of this differential, so it must be homologically the same as $`\mathfrak{1}`$.
## 4 Strong homotopy equivalences
In Example 6 we introduced a colored operad $`\mathrm{๐ ๐}`$ describing chain maps $`๐ต:MN`$, $`๐ถ:NM`$ such that $`\mathrm{๐ต๐ถ}=\mathrm{๐ท๐ท}_N`$ and $`\mathrm{๐ต๐ถ}=\mathrm{๐ท๐ท}_M`$.
A general belief is that the homotopy analog of this situation is given by a quadruple $`(๐ต,๐ถ,๐ท,๐ป)`$, where $`๐ต:MN`$ and $`๐ถ:NM`$ are degree zero chain maps that are homotopy inverses of each other, with associated homotopies $`๐ท`$ and $`๐ป`$:
$$\mathrm{๐ถ๐ต}\mathrm{๐ท๐ท}_M=d_M๐ท+๐ทd_M,\mathrm{๐ต๐ถ}\mathrm{๐ท๐ท}_N=d_N๐ป+๐ปd_N.$$
(6)
Such a quadruple is clearly an algebra over the operad
$$_{\mathrm{fake}}:=((f_0,g_0,f_1,g_1),d),$$
(7)
where
$$f_0:๐ฑ๐,g_0:๐๐ฑ,f_1:๐ฑ๐ฑ\text{ and }g_1:๐๐$$
are generators with $`\mathrm{๐๐๐}(f_0)=\mathrm{๐๐๐}(g_0)=0`$, $`\mathrm{๐๐๐}(f_1)=\mathrm{๐๐๐}(g_1)=1`$, and the differential $`d`$ is given by
$$d(f_0):=0,d(g_0):=0,d(f_1):=g_0f_01_๐ฑ\text{ and }d(g_1):=f_0g_01_๐.$$
There is a dg operad map $`\alpha _{\mathrm{fake}}:_{\mathrm{fake}}\mathrm{๐ ๐}`$ given by
$$\alpha _{\mathrm{fake}}(f_0):=f,\alpha _{\mathrm{fake}}(g_0):=g,\alpha _{\mathrm{fake}}(f_1):=0\text{ and }\alpha _{\mathrm{fake}}(g_1):=0.$$
The following fact which shows that $`_{\mathrm{fake}}`$ is not an acyclic resolution of the operad $`\mathrm{๐ ๐}`$ is crucial.
Fact. The map $`\alpha _{\mathrm{fake}}`$ is not a homology isomorphism. For instance, $`f_0f_1g_1g_0`$ is a cycle in the kernel of $`\alpha _{\mathrm{fake}}`$ that is not homologous to zero.
A proper resolution of $`\mathrm{๐ ๐}`$ was described in . It is a graded colored differential operad
$$_{\mathrm{iso}}:=((f_0,f_1,\mathrm{};g_0,g_1,\mathrm{}),d),$$
with generators of two types,
$$\begin{array}{cc}\hfill \text{(i)}& \text{generators }\{f_n\}_{n0}\text{}\mathrm{๐๐๐}(f_n)=n\text{,}\{\begin{array}{c}f_n:๐ฑ๐\text{ if }n\text{ is even,}\hfill \\ f_n:๐ฑ๐ฑ\text{ if }n\text{ is odd,}\hfill \end{array}\hfill \\ \hfill \text{(ii)}& \text{generators }\{g_n\}_{n0}\text{}\mathrm{๐๐๐}(g_n)=n\text{,}\{\begin{array}{c}g_n:๐๐ฑ\text{ if }n\text{ is even,}\hfill \\ g_n:๐๐\text{ if }n\text{ is odd.}\hfill \end{array}\hfill \end{array}$$
(8)
The differential $`d`$ is given by
$$\begin{array}{cc}df_0:=0,\hfill & dg_0:=0,\hfill \\ df_1:=g_0f_01,\hfill & dg_1:=f_0g_01\hfill \end{array}$$
and, on remaining generators, by the formula
$`df_{2m}`$ $`:=`$ $`{\displaystyle \underset{0i<m}{}}(f_{2i}f_{2(mi)1}g_{2(mi)1}f_{2i}),m0,`$
$`df_{2m+1}`$ $`:=`$ $`{\displaystyle \underset{0jm}{}}g_{2j}f_{2(mj)}{\displaystyle \underset{0j<m}{}}f_{2j+1}f_{2(mj)1},m1,`$ (9)
$`dg_{2m}`$ $`:=`$ $`{\displaystyle \underset{0i<m}{}}(g_{2i}g_{2(mi)1}f_{2(mi)1}g_{2i}),m0,`$
$`dg_{2m+1}`$ $`:=`$ $`{\displaystyle \underset{0jm}{}}f_{2j}g_{2(mj)}{\displaystyle \underset{0j<m}{}}g_{2j+1}g_{2(mj)1},m1,`$
see also (8.2). The above formulas can be written in a compact form by introducing elements
$$\begin{array}{cc}f_{}:=f_0+f_2+f_4+\mathrm{}:๐ฑ๐,\hfill & h_{}:=f_1+f_3+f_5+\mathrm{}:๐ฑ๐ฑ,\hfill \\ g_{}:=g_0+g_2+g_4+\mathrm{}:๐๐ฑ,\hfill & l_{}:=l_1+l_3+l_5+\mathrm{}:๐๐.\hfill \end{array}$$
(10)
Then $`_{\mathrm{iso}}=(f_{},g_{},h_{},l_{})`$ with the differential given by
$$df_{}=f_{}h_{}l_{}f_{},dh_{}=g_{}f_{}h_{}h_{}1_๐ฑ,dg_{}=g_{}l_{}h_{}g_{}\text{ and }dl_{}=f_{}g_{}l_{}l_{}1_๐.$$
We will use this kind of abbreviation quite often, but we shall always keep in mind that each formula of this type in fact represents infinitely many formulas for homogeneous parts. The operad $`_{\mathrm{iso}}`$ is โtriviallyโ filtered, by
$$F^p_{\mathrm{iso}}:=\{\begin{array}{cc}_{\mathrm{iso}},\hfill & \text{for }p=0\text{, and}\hfill \\ 0,\hfill & \text{for }p>0\text{.}\hfill \end{array}$$
This filtration is obviously complete. Algebras over the operad $`_{\mathrm{iso}}`$ are the strong homotopy equivalences introduced in Definition 1.
The following theorem, formulated without proof in , claims that $`_{\mathrm{iso}}`$ is an acyclic resolution of the operad $`\mathrm{๐ ๐}`$.
###### Theorem 9
The map $`\alpha _{\mathrm{iso}}:_{\mathrm{iso}}\mathrm{๐ ๐}`$ defined by
$$\alpha _{\mathrm{iso}}(f_0):=[f],\alpha _{\mathrm{iso}}(g_0):=[g],\text{ while }\alpha _{\mathrm{iso}}(f_n):=0,\alpha _{\mathrm{iso}}(g_n)=0\text{ for }n1,$$
(11)
is a map of differential graded colored operads that induces an isomorphism of cohomology.
Proof. It is clear that $`\alpha _{\mathrm{iso}}`$ commutes with the differentials and that it induces an isomorphism $`H_0(_{\mathrm{iso}},d)\mathrm{๐ ๐}`$. It thus remains to prove that $`_{\mathrm{iso}}`$ is acyclic in positive dimensions.
The operad $`(f_0,f_1,\mathrm{};g_0,g_1,\mathrm{})`$ is the free abelian group spanned by composable chains of generators. The length of these chains induces another grading, which we call the homogeneity. The differential $`d`$ decomposes as $`d=d_1+d_{+1}`$, where $`d_i`$ raises the homogeneity by $`i=\pm 1`$. Explicitly, $`d_{+1}`$ is given on generators by
$$d_{+1}f_{}=f_{}h_{}l_{}f_{},d_{+1}h_{}=g_{}f_{}h_{}h_{},d_{+1}g_{}=g_{}l_{}h_{}g_{},d_{+1}l_{}=f_{}g_{}l_{}l_{},$$
while $`d_1`$ is given by
$$d_1f_{}=0,d_1h_{}=1_๐ฑ,d_1g_{}=0\text{ and }d_1l_{}=1_๐.$$
We claim that
$$((f_0,f_1,\mathrm{};g_0,g_1,\mathrm{}),d_{+1})\text{ is acyclic in positive degrees.}$$
(12)
We prove (12) by introducing a contracting homotopy
$$\theta :(f_0,f_1,\mathrm{};g_0,g_1,\mathrm{})(f_0,f_1,\mathrm{};g_0,g_1,\mathrm{})$$
as follows. Let $`z_1,z_2,\mathrm{}`$ denote generators of $`(f_0,f_1,\mathrm{};g_0,g_1,\mathrm{})`$, then let, for $`m0`$,
$$R(z_1z_2):=\{\begin{array}{cc}f_{2m+2},\hfill & \text{ if }z_1z_2=f_0f_{2m+1}\text{,}\hfill \\ g_{2m+2},\hfill & \text{ if }z_1z_2=g_0g_{2m+1}\text{,}\hfill \\ g_{2m+1},\hfill & \text{ if }z_1z_2=f_0g_{2m}\text{,}\hfill \\ f_{2m+1},\hfill & \text{ if }z_1z_2=g_0f_{2m}\text{,}\hfill \\ 0\hfill & \text{ if otherwise.}\hfill \end{array}$$
Then the contracting homotopy $`\theta `$ is defined by
$$\theta (z_1z_2\mathrm{}z_t):=\{\begin{array}{cc}R(z_1z_2)z_3\mathrm{}z_t,\hfill & \text{if }t2\text{, and}\hfill \\ 0,\hfill & \text{otherwise.}\hfill \end{array}$$
It is immediate to check that indeed $`\theta d_{+1}(x)+d_{+1}\theta (x)=x`$ whenever $`x`$ has positive degree, which proves (12).
Suppose that $`x(f_0,f_1,\mathrm{};g_0,g_1,\mathrm{})`$ is a $`d`$-cycle of positive degree and let
$$x=x_1+\mathrm{}+x_N,\text{ }x_j\text{ has homogeneity }j\text{}1jN\text{}N>1\text{,}$$
be its decomposition into homogeneity-homogeneous parts. Then clearly $`d_{+1}(x_N)=0`$, thus, by (12), there exists some $`b_{N1}`$ of homogeneity $`N1`$ such that $`x_N=d_{+1}(b_{N1})`$. Then $`xd(b_{N1})`$ is a $`d`$-cycle homologous to $`x`$, whose decomposition contains no terms of homogeneity $`N`$. Repeating this process as many times as necessary, we end up with some $`x^{}`$, homologous to $`x`$, of homogeneity $`1`$, i.e. linear in the generators. An immediate inspection shows that there is no such nontrivial $`x^{}`$ of positive degree, therefore $`x^{}=0`$ which finishes the proof, since $`x^{}`$ is, by construction, homologous to $`x`$. mm
Observe that, in the course of the proof of Theorem 9, we proved the following interesting statement:
###### Proposition 10
The map
$$\alpha _{+1}:((f_0,f_1,\mathrm{};g_0,g_1,\mathrm{}),d_{+1})\frac{(f,g)}{(fg=0,gf=0)}$$
given by
$`\alpha _{+1}(f_0):=[f]`$, $`\alpha _{+1}(g_0):=[g]`$, while $`\alpha _{+1}(f_n)=0`$ and $`\alpha _{+1}(g_n)=0`$ for $`n1`$,
is a homology isomorphism.
In the rest of this section we study when a homotopy equivalence (6) extends to a strong homotopy one. Observe first that (6) induces a โrestrictedโ action $`A_{\mathrm{res}}:_{\mathrm{iso}}\mathrm{๐๐}_{M,N}`$ by
$$A_{\mathrm{res}}(f_0)=๐ต,A_{\mathrm{res}}(g_0)=๐ถ,A_{\mathrm{res}}(f_1)=๐ท\text{ and }A_{\mathrm{res}}(g_1)=๐ป.$$
Related to these data are two obstruction cycles
$$๐ฌ_M:=\mathrm{๐ต๐ท}\mathrm{๐ป๐ต}\mathrm{๐ป๐๐}_1(M,N)\text{ and }๐ฌ_N:=\mathrm{๐ถ๐ป}\mathrm{๐ท๐ถ}\mathrm{๐ป๐๐}_1(N,M).$$
(13)
###### Theorem 11
The obstruction $`[๐ฌ_M]H_1(\mathrm{๐ป๐๐}(M,N))`$ vanishes if and only if the obstruction $`[๐ฌ_N]H_1(\mathrm{๐ป๐๐}(N,M))`$ does.
The restricted action $`A_{\mathrm{res}}`$ can be extended to a full action $`A:_{\mathrm{iso}}\mathrm{๐๐}_{M,N}`$ if and only if one (and hence both) of the above obstructions vanish.
Proof. Let us denote by $`((f_{<n};g_{<n}),d)`$ the suboperad of $`(f_0,f_1,\mathrm{};g_0,g_1,\mathrm{})`$ generated by $`\{f_n,g_n\}_{i<n}`$, with the induced differential. It is clear from the definition that $`df_n,dg_n(f_{<n};g_{<n})`$ for any $`n1`$, thus it makes sense to consider the homology classes $`[df_n]_{n1}`$ and $`[dg_n]_{n1}`$ of these elements in $`H_{n1}((f_{<n};g_{<n}),d)`$. We claim that
$`[g_0][df_{2m}]_{2m1}+[df_{2m}]_{2m1}[f_0]`$ $`=`$ $`0\text{ in }H_{2m1}((f_{<2m};g_{<2m}),d)\text{, and}`$ (14)
$`[f_0][df_{2m+1}]_{2m}+[dg_{2m+1}]_{2m}[f_0]`$ $`=`$ $`0\text{ in }H_{2m}((f_{<2m+1};g_{<2m+1}),d)\text{.}`$ (15)
The first equation follows from the inspection of the degree $`2m+1`$ part of
$$dh_{}^2=d(f_{}g_{})$$
(16)
which is
$$(dh_{}^2)_{2m+1}=f_0(df_{2m})+(dg_{2m})g_0+d(\underset{\begin{array}{c}i+j=m\\ i,j1\end{array}}{}f_{2i}g_{2j});$$
equation (16) can be verified directly. Equation (15) follows in the same manner from
$$d(f_{}h_{}l_{}f_{})=0.$$
which follows from $`d^2=0`$. Observe that (14) gives, for $`m=1`$,
$$[g_0][df_2]_2+[df_2]_2[f_0]=0$$
which is mapped by $`A_{\mathrm{res}}:(f_{<2};g_{<2})\mathrm{๐๐}_{M,N}`$ to
$$[๐ถ][\mathrm{๐ต๐ท}\mathrm{๐ป๐ต}]+[\mathrm{๐ถ๐ป}\mathrm{๐ท๐ถ}][๐ต]=0\text{ in }H_1(\mathrm{๐ป๐๐}(M,N))\text{,}$$
which is of course
$$[๐ถ][๐ฌ_M]+[๐ฌ_N][๐ต]=0.$$
This implies the first part of the statement, since multiplication by the homology class of $`f`$ (resp. of $`g`$) is an isomorphism, as these maps are homotopy invertible.
Let us prove the second part of the theorem. One implication is clear โ if the restricted action $`A_{\mathrm{res}}`$ can be extended to a full one, then obviously both obstructions must vanish.
Suppose that both obstructions vanish. Then the restricted action can be clearly extended to $`f_2`$ and $`g_2`$, i.e. on $`(f_{<3};g_{<3})`$; we denote this extended action by $`A_2`$.
Let us suppose that we have extended $`A_{\mathrm{res}}`$ to some $`A_{n1}:(f_{<n};g_{<n})\mathrm{๐๐}_{M,N}`$, $`n3`$, and try to extend it to $`f_n`$ and $`g_n`$. We must distinguish whether $`n`$ is even or odd; suppose first that $`n=2m`$. The extension clearly exists if and only if
$$\begin{array}{cc}\hfill [A_{2m1}(df_{2m})]=0& \text{in }H_{2m1}(\mathrm{๐ป๐๐}(M,N))\text{, and}\hfill \\ \hfill [A_{2m1}(dg_{2m})]=0& \text{in }H_{2m1}(\mathrm{๐ป๐๐}(N,M))\text{.}\hfill \end{array}$$
(17)
This, unfortunately, need not be true in general, but we can use the following trick. Observe that if we change the definition of $`A_{2m1}(f_{2m1})`$ by adding a cycle $`\varphi \mathrm{๐ป๐๐}_{2m1}(M,M)`$ and $`A_{2m1}(g_{2m1})`$ by adding a cycle $`\psi \mathrm{๐ป๐๐}_{2m1}(M,M)`$, the extension $`A_{2m1}`$ remains well defined. We show that by such a โrecalibration,โ we may always achieve that the elements in (17) vanish. Indeed, it follows from the definition of the differential, from $`A_{2m1}(f_0)=๐ต`$ and $`A_{2m1}(g_0)=๐ถ`$, that (17) changes to
$$[๐ต][\varphi ]+[A_{2m1}(df_{2m})][\psi ][๐ต]=0\text{ and }[๐ถ][\psi ]+[A_{2m1}(dg_{2m})][\varphi ][๐ถ]=0.$$
This system can clearly be solved if and only if
$$[๐ถ][A_{2m1}(df_{2m})]+[A_{2m1}(dg_{2m})][๐ต]=0,$$
which is the image of (14) under $`A_{2m1}`$. The case of odd $`n`$ is discussed in the same manner, using (15) instead of (14). mm
In the light of Theorem 11, we will make no distinction between $`[๐ฌ_M]`$ and $`[๐ฌ_N]`$ and denote both obstructions by $`[๐ฌ]`$. The following statement is a โchain-levelโ version of Theorem 11.
###### Theorem 12
The restricted action $`A_{\mathrm{res}}`$ can be extended to a full action $`A:_{\mathrm{iso}}\mathrm{๐๐}_{M,N}`$ by putting $`A_{\mathrm{res}}(f_n)=0`$ and $`A_{\mathrm{res}}(g_n)=0`$ for $`n2`$ if and only if the obstruction cycles (13) vanish and if $`\mathrm{๐ท๐ท}=0`$ and $`\mathrm{๐ป๐ป}=0`$.
The proof is an easy exercise. In we formulated without proof the following theorem.
###### Theorem 13
Let $`(๐ต,๐ถ,๐ท,๐ป)`$ be a homotopy equivalence (6). By changing either $`๐ท`$ or $`๐ป`$ we may always achieve that the obstruction $`[๐ฌ]`$ vanishes, i.e. that, by Theorem 11, the homotopy equivalence $`(๐ต,๐ถ,๐ท,๐ป)`$ extends to a strong one. Examples of these changes are
$`(๐ต,๐ถ,๐ท,๐ป)`$ $``$ $`(๐ต,๐ถ,๐ท๐ถ(\mathrm{๐ต๐ท}\mathrm{๐ป๐ต}),๐ป),\text{ or}`$
$`(๐ต,๐ถ,๐ท,๐ป)`$ $``$ $`(๐ต,๐ถ,๐ท,๐ป๐ต(\mathrm{๐ถ๐ป}\mathrm{๐ท๐ถ})).`$
Proof. Let us show that the first substitution annihilates the obstructions. Denote for simplicity $`๐ท^{}:=๐ท๐ถ(\mathrm{๐ต๐ท}\mathrm{๐ป๐ต})`$. Then it can be verified directly that
$`๐ฌ_M(๐ต,๐ถ,๐ท^{},๐ป)`$ $`=`$ $`\mathrm{๐ต๐ท}\mathrm{๐ป๐ต}\mathrm{๐ต๐ถ}(\mathrm{๐ต๐ท}\mathrm{๐ป๐ต})=d(๐ป(\mathrm{๐ต๐ท}\mathrm{๐ป๐ต})),\text{ and}`$
$`๐ฌ_N(๐ต,๐ถ,๐ท^{},๐ป)`$ $`=`$ $`\mathrm{๐ถ๐ป}\mathrm{๐ท๐ถ}+๐ถ(\mathrm{๐ต๐ท}\mathrm{๐ป๐ต})g=d(๐ท^2๐ถ+\mathrm{๐ถ๐ป}^2\mathrm{๐ท๐ถ๐ป}),`$
therefore $`[๐ฌ_M(๐ต,๐ถ,๐ท^{},๐ป)]=[๐ฌ_N(๐ต,๐ถ,๐ท^{},๐ป)]=0`$. The discussion of the second substitution is the same. mm
## 5 The retraction
Let us introduce a filtered colored operad $`\stackrel{~}{}_{\mathrm{iso}}`$ describing perturbations of strong homotopy equivalences. It is the completion of the operad $`\mathrm{๐๐๐}\stackrel{~}{}_{\mathrm{iso}}`$ generated by two types of generators:
* generators $`\{f_n\}_{n0}`$ and $`\{g_n\}_{n0}`$ as in (8) for an unperturbed strongly homotopy equivalence, and
* generators for a perturbation, that is, a generator $`\overline{x}`$ for a perturbation of the โblackโ differential, a generator $`\overline{y}`$ for a perturbation of the โwhiteโ differential, and generators $`\overline{f}_n`$ and $`\overline{g}_n`$ for perturbations of $`f_n`$ resp. $`g_n`$, $`n0`$.
For homogeneity of the notation, we will sometimes write $`f_n^0`$ (resp. $`g_n^0`$) instead of $`f_n`$ (resp. $`g_n`$) and $`f_n^1`$ (resp. $`g_n^1`$) instead of $`\overline{f}_n`$ (resp. $`\overline{g}_n`$). With these conventions assumed,
$$\mathrm{๐๐๐}\stackrel{~}{}_{\mathrm{iso}}:=((\overline{x},\overline{y},\{f_n^s\}_{n0}^{s=1,2},\{g_n^s\}_{n0}^{s=1,2}),d)$$
with $`\mathrm{๐๐๐}(\overline{x})=\mathrm{๐๐๐}(\overline{y})=1`$ and $`\mathrm{๐๐๐}(f_n^s)=\mathrm{๐๐๐}(g_n^s)=n`$. The differential $`d`$ will be defined later. To define on $`\mathrm{๐๐๐}\stackrel{~}{}_{\mathrm{iso}}`$ a filtration, we assign to each generator another degree $`\underset{ยฏ}{\mathrm{๐๐}}g`$ by
$$\underset{ยฏ}{\mathrm{๐๐}}g(\overline{x})=\underset{ยฏ}{\mathrm{๐๐}}g(\overline{y})=\underset{ยฏ}{\mathrm{๐๐}}g(\overline{f}_n)=\underset{ยฏ}{\mathrm{๐๐}}g(\overline{g}_n)=1,\underset{ยฏ}{\mathrm{๐๐}}g(f_n)=\underset{ยฏ}{\mathrm{๐๐}}g(g_n)=0,n0.$$
This assignment expresses the fact that overlined generators describe perturbations. The $`\underset{ยฏ}{\mathrm{๐๐}}g`$-grading of generators induces, in the standard way, a grading on $`\mathrm{๐๐๐}\stackrel{~}{}_{\mathrm{iso}}`$ and we define
$$F^p\mathrm{๐๐๐}\stackrel{~}{}_{\mathrm{iso}}:=\{z\mathrm{๐๐๐}\stackrel{~}{}_{\mathrm{iso}};\underset{ยฏ}{\mathrm{๐๐}}g(z)p\},p0.$$
Let us denote by $`\stackrel{~}{}_{\mathrm{iso}}`$ the completion of $`\mathrm{๐๐๐}\stackrel{~}{}_{\mathrm{iso}}`$. A typical element of $`\stackrel{~}{}_{\mathrm{iso}}`$ is a formal sum $`_{i0}m_i`$ with $`m_i\mathrm{๐๐๐}\stackrel{~}{}_{\mathrm{iso}}`$ and $`\underset{ยฏ}{\mathrm{๐๐}}g(m_i)=i`$.
The best way to describe the differential is to introduce a condensed notation (compare (10)):
$$\begin{array}{cc}\stackrel{~}{f}_{}:=_{m0}f_{2m}+\overline{f}_{2m},\hfill & \stackrel{~}{g}_{}:=_{m0}g_{2m}+\overline{g}_{2m},\hfill \\ \stackrel{~}{h}_{}:=_{m0}f_{2m+1}+\overline{f}_{2m+1},\hfill & \stackrel{~}{l}_{}:=_{m0}g_{2m+1}+\overline{g}_{2m+1}.\hfill \end{array}$$
(18)
The differential is given by
$$\begin{array}{cc}d\overline{x}=\overline{x}\overline{x},\hfill & d\overline{y}=\overline{y}\overline{y},\hfill \\ d\stackrel{~}{f}_{}=\stackrel{~}{f}_{}(\overline{x}+\stackrel{~}{h}_{})(\overline{y}+\stackrel{~}{l}_{})\stackrel{~}{f}_{},\hfill & d\stackrel{~}{g}_{}=\stackrel{~}{g}_{}(\overline{y}+\stackrel{~}{l}_{})(\overline{x}+\stackrel{~}{h}_{})\stackrel{~}{g}_{},\hfill \\ d\stackrel{~}{h}_{}=(\stackrel{~}{h}_{}\overline{x}+\overline{x}\stackrel{~}{h}_{})+\stackrel{~}{g}_{}\stackrel{~}{f}_{}\stackrel{~}{h}_{}\stackrel{~}{h}_{}1,\hfill & d\stackrel{~}{l}_{}=(\stackrel{~}{l}_{}\overline{y}+\overline{y}\stackrel{~}{l}_{})+\stackrel{~}{f}_{}\stackrel{~}{g}_{}\stackrel{~}{l}_{}\stackrel{~}{l}_{}1.\hfill \end{array}$$
(19)
A momentโs reflection shows that the differential operad $`\stackrel{~}{}_{\mathrm{iso}}`$ really describes perturbations of strongly homotopy equivalences. Expanding (19) we get more explicit formulas for the differential:
$`df_{2m}^1`$ $`:=`$ $`{\displaystyle \underset{t=1,2}{}}(f_{2m}^t\overline{x}\overline{y}f_{2m}^t)+{\displaystyle \underset{t+r1}{}}({\displaystyle \underset{0i<m}{}}(f_{2i}^tf_{2(mi)1}^rg_{2(mi)1}^tf_{2i}^r)),`$
$`df_{2m+1}^1`$ $`:=`$ $`{\displaystyle \underset{t=1,2}{}}(f_{2m+1}^t\overline{x}+\overline{x}f_{2m+1}^t)+`$
$`+{\displaystyle \underset{t+r1}{}}({\displaystyle \underset{0jm}{}}g_{2j}^tf_{2(mj)}^r{\displaystyle \underset{0j<m}{}}f_{2j+1}^tf_{2(mj)1}^r),`$
$`dg_{2m}^1`$ $`:=`$ $`{\displaystyle \underset{t=1,2}{}}(g_{2m}^t\overline{y}\overline{x}g_{2m}^t)+{\displaystyle \underset{t+r1}{}}({\displaystyle \underset{0i<m}{}}(g_{2i}^tg_{2(mi)1}^rf_{2(mi)1}^tg_{2i}^r)),`$
$`dg_{2m+1}^1`$ $`:=`$ $`{\displaystyle \underset{t=1,2}{}}(g_{2m+1}^t\overline{y}+\overline{y}g_{2m+1}^t)+`$
$`+{\displaystyle \underset{t+r1}{}}({\displaystyle \underset{0jm}{}}f_{2j}^tg_{2(mj)}^r{\displaystyle \underset{0j<m}{}}g_{2j+1}^tg_{2(mj)1}^r).`$
The action of $`d`$ on $`f_n^0=f_n`$ and $`g_n^0=g_n`$ is, of course, the same as in (9). See also 8.3. The following theorem claims that $`\stackrel{~}{}_{\mathrm{iso}}`$ is a resolution of the operad $`\mathrm{๐ ๐}`$ introduced in Example 6.
###### Theorem 14
The map $`\alpha :\stackrel{~}{}_{\mathrm{iso}}\mathrm{๐ ๐}`$ given by
$$\alpha (f_0^0):=[f]\text{ and }\alpha (g_0^0):=[g],$$
(20)
while $`\alpha `$ is zero on the remaining generators, is a map of differential filtered operads that induces an isomorphism of cohomology.
Proof. It is immediate to see that $`\alpha `$ decomposes as $`\alpha =\alpha _{\mathrm{iso}}\stackrel{~}{\alpha }`$, with $`\stackrel{~}{\alpha }:\stackrel{~}{}_{\mathrm{iso}}_{\mathrm{iso}}`$ given by $`\stackrel{~}{\alpha }(f_n^0)=f_n`$, $`\stackrel{~}{\alpha }(g_n^0)=g_n`$ and $`\stackrel{~}{\alpha }`$ trivial on remaining generators.
Since $`\alpha _{\mathrm{iso}}`$ is, by Theorem 9, a homology isomorphism, it is enough to show that $`\stackrel{~}{\alpha }`$ is also a homology isomorphism. This can be done by a spectral sequence argument which we omit, since we will not need the theorem in our proofs.
The philosophical meaning is that a perturbation cannot introduce nontrivial homology classes. mm
Let us consider the free product
$$๐\mathrm{๐๐}_{\mathrm{iso}}=((\overline{x},\{f_n\}_{n0},\{g_n\}_{n0}),d)$$
with the differential given by (5) and (9). It is clear that the map $`\iota :๐\mathrm{๐๐}_{\mathrm{iso}}\stackrel{~}{}_{\mathrm{iso}}`$ defined by
$$\iota (f_n):=f_n^0,\iota (g_n):=g_n^0\text{ and }\iota (\overline{x}):=\overline{x},n0,$$
(21)
is an inclusion of differential filtered colored operads. Let us formulate the main statement of this section.
###### Theorem 15
The operad $`๐\mathrm{๐๐}_{\mathrm{iso}}`$ is a retract of $`\stackrel{~}{}_{\mathrm{iso}}`$, that is, there exists a map $`r:\stackrel{~}{}_{\mathrm{iso}}๐\mathrm{๐๐}_{\mathrm{iso}}`$ of differential filtered colored operads such that $`r\iota =11_{๐\mathrm{๐๐}_{\mathrm{iso}}}`$.
Proof. We construct the retraction $`r`$ explicitly. Let us define, for each odd $`r1`$, a โkernelโ $`_r:๐ฑ๐ฑ`$, $`_r๐\mathrm{๐๐}_{\mathrm{iso}}`$, of degree $`r`$ by the formula
$$_r:=\underset{t0}{}\overline{x}f_{2m_1+1}\overline{x}\mathrm{}\overline{x}f_{2m_t+1}\overline{x}$$
where the summation runs over all $`2(m_1+\mathrm{}+m_t)1=r`$, $`m_10,\mathrm{},m_t0`$. See 8.5 for some explicit formulas. The retraction $`r:\stackrel{~}{}_{\mathrm{iso}}๐\mathrm{๐๐}_{\mathrm{iso}}`$ is then given by the following formulas:
$$\begin{array}{cccccc}\hfill r(\overline{x})& :=& \overline{x},\hfill & \hfill r(\overline{y})& :=& f_0_1g_0,\hfill \\ \hfill r(f_n)& :=& f_n,\hfill & \hfill r(g_n)& :=& g_n,\hfill \\ \hfill r(\overline{f}_{2m})& :=& \underset{a+b+c=m}{}f_{2a}_{2b1}f_{2c+1},\hfill & \hfill r(\overline{g}_{2m})& :=& \underset{a+b+c=m}{}f_{2a+1}_{2b1}g_{2b},\hfill \\ \hfill r(\overline{f}_{2m+1})& :=& \underset{a+b+c=m}{}f_{2a+1}_{2b1}f_{2c+1},\hfill & \hfill r(\overline{g}_{2m+1})& :=& \underset{a+b+c=m+1}{}f_{2a}_{2b1}g_{2b},\hfill \end{array}$$
(22)
where $`m,n0`$ and $`a,b,c`$ are nonnegative integers. In compact notation
$$_{}:=\underset{q0}{}(\overline{x}h_{})^q\overline{x}$$
we can rewrite (22) as
$$\begin{array}{cc}r(\stackrel{~}{f}_{})=f_{}(1+_{}h_{}),\hfill & r(\stackrel{~}{g}_{})=(1+h_{}_{})g_{},\hfill \\ r(\stackrel{~}{h}_{})=h_{}+h_{}_{}h_{},\hfill & r(\overline{y}+\stackrel{~}{l}_{})=l_{}+f_{}_{}g_{},\hfill \end{array}$$
see (18) for the meaning of $`\stackrel{~}{f}_{},\stackrel{~}{g}_{},\stackrel{~}{h}_{}`$ and $`\stackrel{~}{l}_{}`$. It is clear that $`r`$ defined above is a retraction. Let us prove that it commutes with the differentials, that is
$$dr=rd.$$
(23)
It is, of course, enough to prove (23) on generators $`\stackrel{~}{f}_{},\stackrel{~}{g}_{},\stackrel{~}{h}_{}`$ and $`\stackrel{~}{l}_{}`$ of $`\stackrel{~}{}_{\mathrm{iso}}`$. For $`\stackrel{~}{f}_{}`$ we have
$`dr(\stackrel{~}{f}_{})`$ $`=`$ $`d(f_{}(1+_{}h_{}))=(f_{}h_{}l_{}f_{})(1+_{}h_{})+`$
$`+f_{}_{}(g_{}f_{}h_{}h_{})_{}h_{}f_{}_{}(g_{}f_{}h_{}h_{}1),`$
where we used the obvious relation
$$d_{}=_{}(g_{}f_{}h_{}h_{})_{}.$$
On the other hand,
$`rd(\stackrel{~}{f}_{})`$ $`=`$ $`r(\stackrel{~}{f}_{}(\overline{x}+\stackrel{~}{h}_{})(\overline{y}+\stackrel{~}{l}_{})\stackrel{~}{f}_{})=`$
$`=`$ $`f_{}(1+_{}h_{})(\overline{x}+h_{}+h_{}_{}h_{})(l_{}+f_{}_{}g_{})f_{}(1+_{}h_{}).`$
Comparing (5) to (5), using another obvious relation
$$\overline{x}+_{}h_{}\overline{x}=_{},$$
we indeed check that $`rd(\stackrel{~}{f}_{})=dr(\stackrel{~}{f}_{})`$. Equation (23) can be verified on remaining generators by the same direct argument. mm
## 6 Proofs
The initial data of the Perturbation Problem define an algebra $`A_{\mathrm{in}}`$ over the free product
$$๐\mathrm{๐๐}_{\mathrm{fake}}=((\overline{x},f_0,g_0,f_1,g_1),d),$$
of the operad $`๐\mathrm{๐๐}`$ of Example 7 with the operad $`_{\mathrm{fake}}`$ introduced in (7), $`A_{\mathrm{in}}:๐\mathrm{๐๐}_{\mathrm{fake}}\mathrm{๐๐}_{M,N}`$, by
$$A_{\mathrm{in}}(\overline{x}):=_M:=\stackrel{~}{d}_Md_M,A_{\mathrm{in}}(f_0):=๐ต,A_{\mathrm{in}}(g_0):=๐ถ,A_{\mathrm{in}}(f_1):=๐ท,A_{\mathrm{in}}(g_1):=๐ป.$$
We seek a solution of the PP encoded to an algebra $`A_{\mathrm{out}}`$ over the differential filtered suboperad
$$\stackrel{~}{}_{\mathrm{fake}}:=((\overline{x},\overline{y},f_0,f_1,g_0,g_1,\overline{f}_0,\overline{f}_1,\overline{g}_0,\overline{g}_1),d)$$
of the operad $`\stackrel{~}{}_{\mathrm{iso}}`$ introduced in Section 5 as
$`\stackrel{~}{d}_N:=A_{\mathrm{out}}(\overline{y}),\stackrel{~}{๐ต}=A_{\mathrm{out}}(\overline{f}_0)+๐ต,\stackrel{~}{๐ถ}=A_{\mathrm{out}}(\overline{g}_0)+๐ถ,`$
$`\stackrel{~}{๐ท}:=A_{\mathrm{out}}(\overline{f}_1)+๐ท\text{ and }\stackrel{~}{๐ป}=A_{\mathrm{out}}(\overline{g}_1)+๐ป.`$
There is a natural inclusion $`\iota _{\mathrm{fake}}:๐\mathrm{๐๐}_{\mathrm{fake}}\stackrel{~}{}_{\mathrm{fake}}`$ given by
$$\iota _{\mathrm{fake}}(\overline{x}):=\overline{x},\iota _{\mathrm{fake}}(f_0):=f_0,\iota _{\mathrm{fake}}(g_0):=g_0,\iota _{\mathrm{fake}}(f_1):=f_1\text{ and }\iota _{\mathrm{fake}}(g_1):=g_1.$$
A โfunctorialโ solution of the Perturbation Problem means to find a retraction
$$r_{\mathrm{fake}}:\stackrel{~}{}_{\mathrm{fake}}๐\mathrm{๐๐}_{\mathrm{fake}},r_{\mathrm{fake}}\iota _{\mathrm{fake}}=11_{๐\mathrm{๐๐}_{\mathrm{fake}}}.$$
(26)
###### Theorem 16
There is no retraction $`r_{\mathrm{fake}}:\stackrel{~}{}_{\mathrm{fake}}๐\mathrm{๐๐}_{\mathrm{fake}}`$ as in (26).
Proof. The proof is a straightforward obstruction theory, but since the non-existence of the retraction $`r_{\mathrm{fake}}`$ motivated all this work, we reproduce the proof here in its full length. All calculations below are made modulo terms of filtration $`2`$, so we, in fact, work in the associated graded operad. The following equations must be satisfied (see 8.5):
$`dr_{\mathrm{fake}}(\overline{y})`$ $`=`$ $`r_{\mathrm{fake}}(d\overline{y})=0,`$
$`dr_{\mathrm{fake}}(\overline{f}_0)`$ $`=`$ $`r_{\mathrm{fake}}(d\overline{f}_0)=f_0\overline{x}r_{\mathrm{fake}}(\overline{y})f_0,`$ (27)
$`dr_{\mathrm{fake}}(\overline{g}_0)`$ $`=`$ $`r_{\mathrm{fake}}(d\overline{g}_0)=g_0r_{\mathrm{fake}}(\overline{y})\overline{x}g_0,`$ (28)
$`dr_{\mathrm{fake}}(\overline{f}_1)`$ $`=`$ $`r_{\mathrm{fake}}(d\overline{f}_1)=(\overline{x}f_1+f_1\overline{x})+r_{\mathrm{fake}}(\overline{g}_0)f_0+g_0r_{\mathrm{fake}}(\overline{g}_0),\text{ and}`$ (29)
$`dr_{\mathrm{fake}}(\overline{g}_1)`$ $`=`$ $`r_{\mathrm{fake}}(d\overline{g}_1)=(r_{\mathrm{fake}}(\overline{y})g_1+g_1r_{\mathrm{fake}}(\overline{y}))+r_{\mathrm{fake}}(\overline{f}_0)g_0+f_0r_{\mathrm{fake}}(\overline{g}_0).`$ (30)
It follows from (27) and (28) that, for some $`b`$, $`r_{\mathrm{fake}}(\overline{y})=f_0\overline{x}g_0+db`$ and that
$$r_{\mathrm{fake}}(\overline{f}_0)=f_0\overline{x}g_0bf_0+c_1,r_{\mathrm{fake}}(\overline{g}_0)=f_1\overline{x}g_0+g_0b+c_2,$$
for some cycles $`c_1,c_2`$. The right hand side of (29) then becomes
$$(g_0f_01)\overline{x}f_1f_1\overline{x}(1g_0f_0)+c_2f_0+g_0c_1=d(f_1\overline{x}f_1)+c_2f_0+g_0c_1,$$
while the right hand side of (30) becomes
$`f_0\overline{x}(f_1g_0g_0g_1)+(f_0f_1g_1f_0)\overline{x}gbf_0g_0dbg_1+f_0g_0bg_1db+c_1g_0+f_0c_2=`$
$`=`$ $`d(g_1bbg_1)+f_0\{\overline{x}(f_1g_0g_0g_1)+c_2\}+\{(f_0f_1g_1f_0)\overline{x}+c_1\}g_0.`$
From this we see that (29) and (30) can be solved in $`r_{\mathrm{fake}}(\overline{f}_1)`$ and $`r_{\mathrm{fake}}(\overline{g}_1)`$ if and only if
$$f_0\overline{x}(f_1g_0g_0g_1)+(f_0f_1g_1f_0)\overline{x}g_0$$
is homologous to zero. It can be easily seen that this is not true. mm
Proof of the IPL. The initial data of the IPL can be organized into an action $`E_{\mathrm{in}}:๐\mathrm{๐๐}_{\mathrm{iso}}\mathrm{๐๐}_{M,N}`$. Then the action
$$E_{\mathrm{out}}:\stackrel{~}{}_{\mathrm{iso}}\stackrel{r}{}๐\mathrm{๐๐}_{\mathrm{iso}}\stackrel{E_{\mathrm{in}}}{}\mathrm{๐๐}_{M,N},$$
where $`r`$ is the retraction of Theorem 15, clearly solves the IPL. mm
## 7 A conceptual explanation
We believe in the existence of a model category (MC) structure on the category of operads. Let us ignore in this conceptual section the fact that the existence of this structure has been proved only for some special cases and certainly not for the category $`\mathrm{๐ต๐๐๐๐พ๐}_{\mathrm{}}`$ of general filtered colored operads over $`\mathrm{}`$.
Our candidate for cofibrations in $`\mathrm{๐ต๐๐๐๐พ๐}_{\mathrm{}}`$ are maps such that the associated maps of graded operads are cofibrations in the sense of an obvious integral version of \[5, Definition 15\] (or something close to it). Fibrations are then epimorphisms and weak equivalences are homology isomorphisms.
As argued in , homotopy invariant algebras are those over cofibrant operads. By Theorem 10, $`_{\mathrm{iso}}`$ is a cofibrant resolution of the operad $`\mathrm{๐ ๐}`$, that is why strong homotopy equivalences, as algebras over $`_{\mathrm{iso}}`$, are proper homotopy versions of strict isomorphisms.
Let us show how the IPL follows from the properties of the MC structure on $`\mathrm{๐ต๐๐๐๐พ๐}_{\mathrm{}}`$. The situation is summarized in the following diagram.
In the above diagram, $`\alpha `$ is the map from Theorem 14, $`\iota `$ the inclusion (21), $`p:=\alpha \iota `$ and the action $`E_{\mathrm{in}}`$ summarizes the input data of the IPL. The solution of the IPL will then be given by $`E_{\mathrm{out}}:=E_{\mathrm{in}}r`$.
The map is $`p`$ clearly an epimorphism, hence a fibration. It is also a weak equivalence, because, if we ignore the acyclic (by Proposition 8) factor $`๐\mathrm{๐๐}`$, the map $`p`$ is exactly the map $`\alpha _{\mathrm{iso}}`$ of Theorem 9. The map $`\iota `$ is a cofibration, thus the existence of $`r`$ follows from the axioms of a MC structure.
The above is, of course, just an explanation, not a proof, so we had, in the proof of Theorem 15, to construct the retraction $`r`$ by other means.
## 8 Appendix: Explicit Formulas
###### 8.1
Explicit axioms for a SHE:
$`d_NF_0F_0d_M`$ $`=`$ $`0,`$
$`d_MG_0G_0d_N`$ $`=`$ $`0,`$
$`d_MH_1+H_1d_M`$ $`=`$ $`G_0F_011_M,`$
$`d_NL_1+L_1d_N`$ $`=`$ $`F_0G_011_N,`$
$`d_NF_2F_2d_M`$ $`=`$ $`F_0H_1L_1F_0,`$
$`d_MG_2G_2d_N`$ $`=`$ $`G_0L_1H_1G_0,`$
$`d_MH_3+H_3d_M`$ $`=`$ $`G_0F_2H_1H_1+G_2F_0,`$
$`d_NL_3+L_3d_N`$ $`=`$ $`F_0G_2L_1L_1+F_2G_0,`$
$`d_NF_4F_4d_M`$ $`=`$ $`F_0H_3L_1F_2+F_2H_1L_3F_0,`$
$`d_MG_4G_4d_N`$ $`=`$ $`G_0L_3H_1G_2+G_2L_1H_3G_0,`$
$`\mathrm{}`$
###### 8.2
Formulas for the differential of $`_{\mathrm{iso}}`$:
$$\begin{array}{cccccc}\hfill df_0& =& 0,\hfill & \hfill dg_0& =& 0,\hfill \\ \hfill df_1& =& g_0f_01,\hfill & \hfill dg_1& =& f_0g_01,\hfill \\ \hfill df_2& =& f_0f_1g_1f_0,\hfill & \hfill dg_2& =& g_0g_1f_1g_0,\hfill \\ \hfill df_3& =& g_0f_2f_1f_1+g_2f_0,\hfill & \hfill dg_3& =& f_0g_2g_1g_1+f_2g_0\hfill \\ \hfill df_4& =& f_0f_3g_1f_2+f_2f_1g_3f_0,\hfill & \hfill dg_4& =& g_0g_3f_1g_2+g_2g_1f_3g_0,\hfill \\ & \mathrm{}& & & \mathrm{}& \end{array}$$
###### 8.3
Formulas for the differential of $`\stackrel{~}{}_{\mathrm{iso}}`$: the action on $`f_n`$, $`g_n`$ is, for $`n0`$, the same as in 8.2, and
$`d\overline{x}`$ $`=`$ $`\overline{x}\overline{x},d\overline{y}=\overline{y}\overline{y},`$
$`d\overline{f}_0`$ $`=`$ $`f_0\overline{x}\overline{y}f_0+\overline{f}_0\overline{x}\overline{y}\overline{f}_0,`$
$`d\overline{g}_0`$ $`=`$ $`g_0\overline{y}\overline{x}g_0+\overline{g}_0\overline{y}\overline{x}\overline{g}_0,`$
$`d\overline{f}_1`$ $`=`$ $`(\overline{x}f_1+f_1\overline{x})+(\overline{g}_0f_0+g_0\overline{f}_0)(\overline{x}\overline{f}_1+\overline{f}_1\overline{x})+\overline{g}_0\overline{f}_0,`$
$`d\overline{g}_1`$ $`=`$ $`(\overline{y}g_1+g_1\overline{y})+(\overline{f}_0g_0+f_0\overline{g}_0)(\overline{y}\overline{g}_1+\overline{g}_1\overline{y})+\overline{f}_0\overline{g}_0,`$
$`d\overline{f}_2`$ $`=`$ $`(f_2\overline{x}\overline{y}f_2)+(\overline{f}_0f_1+f_0\overline{f}_1)(\overline{g}_1f_0+g_1\overline{f}_0)+(\overline{f}_2\overline{x}\overline{y}\overline{f}_2)+(\overline{f}_1\overline{f}_0\overline{g}_1\overline{g}_0),`$
$`d\overline{g}_2`$ $`=`$ $`(g_2\overline{y}\overline{x}g_2)+(\overline{g}_0g_1+g_0\overline{g}_1)(\overline{f}_1g_0+f_1\overline{g}_0)+(\overline{g}_2\overline{y}\overline{x}\overline{g}_2)+(\overline{g}_1\overline{g}_0\overline{f}_1\overline{f}_0),`$
$`\mathrm{}`$
###### 8.4
Formulas for the kernel $`_n`$:
$`_1`$ $`=`$ $`\overline{x}+\overline{x}f_1\overline{x}+\overline{x}f_1\overline{x}f_1\overline{x}+\overline{x}f_1\overline{x}f_1\overline{x}f_1\overline{x}+\mathrm{},`$
$`_1`$ $`=`$ $`\overline{x}f_3\overline{x}+\overline{x}f_1\overline{x}f_3\overline{x}+\overline{x}f_3\overline{x}f_1\overline{x}+\overline{x}f_1\overline{x}f_1\overline{x}f_3\overline{x}+\overline{x}f_1\overline{x}f_3\overline{x}f_1\overline{x}+\overline{x}f_3\overline{x}f_1\overline{x}f_1\overline{x}+\mathrm{},`$
$`_3`$ $`=`$ $`\overline{x}f_5\overline{x}+\overline{x}f_5\overline{x}f_1\overline{x}+\overline{x}f_3\overline{x}f_3\overline{x}+\overline{x}f_1\overline{x}f_5\overline{x}+\overline{x}f_1\overline{x}f_3\overline{x}f_3\overline{x}+\overline{x}f_3\overline{x}f_1\overline{x}f_3\overline{x}+`$
$`+\overline{x}f_3\overline{x}f_3\overline{x}f_1\overline{x}+\overline{x}f_1\overline{x}f_1\overline{x}f_5\overline{x}+\overline{x}f_1\overline{x}f_5\overline{x}f_1\overline{x}+\overline{x}f_5\overline{x}f_1\overline{x}f_1\overline{x}+\mathrm{},`$
$`\mathrm{}`$
###### 8.5
Formulas for the retraction $`r:\stackrel{~}{}_{\mathrm{iso}}๐\mathrm{๐๐}_{\mathrm{iso}}`$:
$`r(\overline{y})`$ $`:=`$ $`f_0_1g_0`$
$`r(\overline{f}_0)`$ $`:=`$ $`f_0_1f_1`$
$`r(\overline{g}_0)`$ $`:=`$ $`f_1_1g_0`$
$`r(\overline{f}_1)`$ $`:=`$ $`f_1_1f_1`$
$`r(\overline{g}_1)`$ $`:=`$ $`f_0_1g_2+f_2_1g_0+f_0_1g_0`$
$`r(\overline{f}_2)`$ $`:=`$ $`f_0_1f_3+f_2_1f_1+f_0_1f_1`$
$`r(\overline{g}_2)`$ $`:=`$ $`f_1_1g_2+f_3_1g_0+f_1_1g_0`$
$`r(\overline{f}_3)`$ $`:=`$ $`f_3_1f_1+f_1_1f_3+f_1_1f_1`$
$`r(\overline{g}_3)`$ $`:=`$ $`f_4_1g_0+f_2_1g_2+f_0_1g_4+f_2_1g_0+f_0_1g_2+f_0_3g_0,`$
$`\mathrm{}`$
###### 8.6
Formulas for the solution of the IPL:
$`\stackrel{~}{d}_N`$ $`=`$ $`d_N+๐ต_0(_M+_M๐ท_0_M+_M๐ท_0_M๐ท_0_M+_M๐ท_0_M๐ท_0_M๐ท_0_M+\mathrm{})๐ถ_0,`$
$`\stackrel{~}{๐ต}_0`$ $`=`$ $`๐ต_0+๐ต_0(_M+_M๐ท_0_M+_M๐ท_0_M๐ท_0_M+_M๐ท_0_M๐ท_0_M๐ท_0_M+\mathrm{})๐ท_0,`$
$`\stackrel{~}{๐ถ}_0`$ $`=`$ $`๐ถ_0+๐ท_0(_M+_M๐ท_0_M+_M๐ท_0_M๐ท_0_M+_M๐ท_0_M๐ท_0_M๐ท_0_M+\mathrm{})๐ถ_0,`$
$`\stackrel{~}{๐ท}_0`$ $`=`$ $`๐ท_0+๐ท_0(_M+_M๐ท_0_M+_M๐ท_0_M๐ท_0_M+_M๐ท_0_M๐ท_0_M๐ท_0_M+\mathrm{})๐ท_0,`$
$`\stackrel{~}{๐ป}_0`$ $`=`$ $`๐ป_0+๐ต_0(_M+_M๐ท_0_M+_M๐ท_0_M๐ท_0_M+_M๐ท_0_M๐ท_0_M๐ท_0_M+\mathrm{})๐ถ_2`$
$`+๐ต_2(_M+_M๐ท_0_M+_M๐ท_0_M๐ท_0_M+_M๐ท_0_M๐ท_0_M๐ท_0_M+\mathrm{})๐ถ_0`$
$`+๐ต_0(_M๐ต_3_M+_M๐ท_0_M๐ต_3_M+_M๐ต_3_M๐ท_0_M+\mathrm{})๐ถ_0,`$
$`\mathrm{}`$
In the above formulas, $`_M=\stackrel{~}{d}_Md_M`$.
ACKNOWLEDGEMENTS
I would like to express my thanks to Jim Stasheff and Johannes Huebschmann for reading the manuscript and many useful comments, and apologize to the latter for not implementing all his suggestions. |
warning/0002/nucl-th0002045.html | ar5iv | text | # Microscopic study of energy and centrality dependence of transverse collective flow in heavy-ion collisions
## I Introduction
Collective effects, such as the expansion of highly compressed nuclear matter in the direction perpendicular to the beam axis of colliding heavy ions at relativistic energies, are very important for the study of the nuclear equation of state (EOS) and for the search of a predicted transition to the a phase of matter, quark-gluon plasma (QGP). At present the transverse flow of particles is believed to be one of the most clear signals to detect the creation of the QGP in heavy-ion experiments (for recent review, see ). This explains the great interest of both experimentalists and theoreticians in the transverse flow phenomenon (see, e.g., and references therein), which was predicted about 25 years ago in nuclear shock wave model analysis.
Initially, the collective flow has been conventionally subdivided into the radial flow, which is azimuthally symmetric, the bounce-off or directed flow in the reaction plane along the impact parameter axis ($`x`$-axis), and the squeeze-out flow developing out of the reaction plane. The latter two components represent the anisotropic part of the transverse flow and appear only in noncentral heavy-ion collisions. The first observation of the transverse flow was made by the Plastic Ball and the Streamer Chamber collaborations at the BEVALAC energies ($`E_{lab}=100`$A MeV \- 1.8A GeV). Later on the directed flow of charged particles has been detected by E877 collaboration at the AGS energies ($`E_{lab}=10.7`$A GeV) and by NA49 and WA98 collaborations at the SPS energies ($`E_{lab}=158`$ A GeV).
The collective flow is a very suitable observable to characterise the reaction dynamics because it is extremely sensitive to the interactions between the particles. At intermediate (SIS) energies the evolution of flow is mainly governed by the density and momentum dependence of the long-range attractive and short-range repulsive nuclear forces in the medium, i.e., the nuclear mean field . With rising energy (AGS, SPS) the mean field gets less important while new degrees of freedom, strings, come into play. It has been shown also that the transverse flow could carry the primary information about the softening of the EOS due to the QGP creation , including the subsequent hadronization, as well as the relaxation of the excited matter to (local) thermal equilibrium.
The advanced technique for the analysis of the flow at high energies, based on the Fourier expansion of the particle azimuthal distribution, has been developed in Refs. . The distribution of the particles in the azimuthal plane can be presented as
$$\frac{dN}{d\varphi }=a_0\left[1+2\underset{n=1}{\overset{\mathrm{}}{}}v_ncos(n\varphi )\right],$$
(1)
where $`\varphi `$ is the azimuthal angle between the momentum of the particle and the reaction plane. The first two coefficients, $`v_1`$ and $`v_2`$, are the amplitudes of the first and second harmonics in the Fourier expansion of the azimuthal distribution, respectively. The asymmetric fraction of the collective flow is decomposed in this analysis into the directed (bounce-off) flow of particles emitted preferentially along the $`x`$-axis, and the elliptic component, which is developed mostly either along the $`x`$-axis or in the squeeze-out direction. The coefficient $`v_2`$, therefore, characterises the eccentricity of the flow ellipsoid .
The importance of the elliptic flow to study collective effects in heavy-ion collisions was first stressed in . In this paper the rotation of the elliptic flow from the squeeze-out direction to the bounce-off direction with rising projectile energy was discussed as well. The alignment of the elliptic flow in the plane of the directed flow has been experimentally detected in Au+Au collisions at the AGS energies and in Pb+Pb collisions at the SPS energies . In the sensitivity of the elliptic flow to the early pressure was noticed. The elliptic flow seems to be generated only during the very beginning of the collective expansion , while the radial flow is developing almost until the freeze-out. It was also pointed out that the elliptic flow should have a kinky structure if the expanding and cooling fireball undergoes a first-order phase transition from the QGP to hadrons. Therefore, the characteristic features of the plasma hadronization can be traced by the dependence of the elliptic flow on the impact parameter .
Although the collective flow is a unique complex phenomenon, the variety of its signals is very rich. Lacking a first principles theoretical description of heavy-ion collisions, one definitely needs to explore semi-phenomenological models whose numerical predictions can be compared with the experimental data on nuclear collective effects in a wide energy range. These models can be classified in general either as macroscopic models or as microscopic ones. Macroscopic models are based on the hypothesis of (local) thermal equilibrium in the system achieved by the large number of various inelastic and elastic processes in the course of a nuclear collision. The many-body distribution functions, which characterise the nonequilibrium states, are rapidly reduced to the one-particle distribution functions (one for each particle species), and the kinetic stage emerges. At a longer time scale the system can reach the hydrodynamic stage, where its evolution is described in terms of the moments of the one-particle distribution functions, such as average velocities, energies, and number of particles. The evolution of a relativistic perfect fluid obeys the conservation of energy and momentum
$$_\mu T^{\mu \nu }=0,$$
(2)
where
$$T^{\mu \nu }=(\epsilon +P)u^\mu u^\nu +Pg^{\mu \nu }$$
(3)
is the energy-momentum tensor, and $`\epsilon ,P,u^\mu `$ are the energy density, pressure, and local four-velocity, respectively.
Without the EOS, which links the pressure $`P`$ to the energy density $`\epsilon `$, the system of hydrodynamic equations (2) - (3) is incomplete. Usually, the EOS is taken in a form
$$P=a\epsilon c_s^2\epsilon ,$$
(4)
with $`c_s`$ being the speed of sound in the medium. Inserting different equations of state, particularly with and without the plasma EOS, into the one-, two-, or three-fluid hydrodynamic model one can study the properties of the particle collective flow at various incident energies .
The microscopic models, developed to describe heavy-ion collisions in a wide range of bombarding energies, e.g. and others, do not rely on the assumption of thermal equilibrium. They employ a dynamical picture of heavy-ion interactions, in which the parton-, string-, and transport approaches can be relevant. Though these models do not explicitly assume the formation of the QGP, the creation of the field of strongly interacted coloured strings may be considered as a precursor of the quark-gluon plasma. Because of the uncertainties in the description of the early stage of heavy-ion collisions at ultrarelativistic energies, the microscopic and macroscopic models can be merged to implement the phase transition to the deconfined phase directly in the microscopic model .
In the present paper two microscopic models, QMD and QGSM, are employed to study the anisotropic flow components in collisions of light and heavy ions at energies from SIS to SPS. The main goal is to understand to what extent the characteristic signals of the hot nuclear matter can be reproduced without invoking the assumption of QGP creation. In other words, if the experimental data will noticeably diverge from the results of simulations, this can be considered as an indication on new processes not included into the models. The paper is organised as follows. A brief description of the models is given in Sec. II. Sections III and IV present the mass and impact parameter dependence of the simulated directed and elliptic flow, $`v_1`$ and $`v_2`$, at SIS, AGS, and SPS energies. Results obtained are discussed in Sec. V. Finally, conclusions are drawn in Sec. VI.
## II Models
The dynamics of nucleus-nucleus collisions at energies up to $`\sqrt{s}2`$A GeV per nucleon can be described in terms of reactions between hadrons and their excited states, resonances. At higher energies additional degrees of freedom, i.e. strings, should be taken into account to describe correctly the processes of multiparticle production. Therefore, we employ the quantum molecular dynamics (QMD) model at the SIS energies, while at the AGS and SPS energies the quark-gluon string model (QGSM) is applied.
In the QMD approach the particles are propagated according to Hamilton equations of motion until their mutual interactions. Each nucleon is represented by a Gaussian-shaped density in the phase space. The black disk approximation is used to determine the binary collision of hadrons. It implies that two hadrons can collide if the centroids of two Gaussians are closer than the distance $`d_0=\sqrt{\sigma _{\mathrm{tot}}(\sqrt{s})/\pi }`$ during their propagation. The Pauli principle is taken into account by blocking the collision if the final states are already occupied in the phase space by other particles. Among the inelastic channels the $`\mathrm{\Delta }(1232)`$ resonance is the dominant one. Pion production takes place via resonance decay of $`\mathrm{\Delta }(1232)`$ and $`N^{}(1440)`$. Pions, which can propagate freely, i.e. without any mean field, undergo, however, a complex chain of reabsorption and subsequent resonance decay processes before their freeze-out . At SIS energies the reaction dynamics is governed by the interplay between the nuclear mean field and binary collisions collisions, which pay a minor role at low energies $`(100`$A MeV) due to the Pauli-blocking of possible scattering states, but become more and more important with rising incident energy. In the present work we employ Skyrme-type mean field with a density dependence corresponding to a hard EOS ($`K=380`$ MeV) and momentum dependence fitted to the empirical nucleon-nucleus optical potential . This type of interaction has been shown to give a good description of flow data in the considered energy range around 1A GeV . Note, that the in-medium cross-section as well as the mean field can also be based on microscopic many-body approaches like Brรผckner theory which is, however, not the scope of the present investigation. This approach allows to describe heavy-ion collisions at energies up to few A GeV. At higher energies the strings come into play.
The QGSM is based on the $`1/N_c`$ (where $`N_c`$ is the number of quark colours or flavours) topological expansion of the amplitude for processes in quantum chromodynamics and string phenomenology of particle production in inelastic binary collisions of hadrons. The diagrams of various topology, which arose due to the $`1/N_c`$ expansion, correspond at high energies to processes with exchange of Regge singularities in the $`t`$-channel. For instance, planar and cylindrical diagram corresponds to the Reggeon and Pomeron exchange, respectively. Therefore, QGSM treats the elementary hadronic interactions on the basis of the Gribov-Regge theory, similar to the dual parton model and the VENUS model . The model simplifies the nuclear effects and concentrates on hadron rescattering. As independent degrees of freedom QGSM includes octet and nonet vector and pseudoscalar mesons, and octet and decuplet baryons, and their antiparticles.
The formation of the quark-gluon plasma is not assumed in the present version of the model. Thus, the effects similar to softening of the EOS in ultrarelativistic heavy-ion collisions, discussed below, are merely attributed to the dynamics of hadron rescattering and nuclear shadowing. We start from the study of energy and centrality dependence of the directed flow.
## III Directed flow
For the simulations at all three energies, namely 1A GeV, 11.6A GeV, and 160A GeV, we choose light (<sup>32</sup>S+<sup>32</sup>S) and heavy (<sup>197</sup>Au+<sup>197</sup>Au and <sup>208</sup>Pb+<sup>208</sup>Pb) symmetric systems. The directed and elliptic flow of nucleons and pions as a function of rapidity is defined as
$$v_n^i=\mathrm{cos}(n\varphi _i)\mathrm{cos}(n\varphi _i)\frac{d๐ฉ^i}{dy}dy/\frac{d๐ฉ^i}{dy}dy,$$
(5)
where $`n=1,2`$ and $`i=N,\pi `$. The mean directed and elliptic flow integrated over the whole rapidity range is simply
$$v_n^i=\mathrm{cos}(n\varphi _i)\mathrm{cos}(n\varphi _i)\frac{d๐ฉ^i}{dy}๐y/\frac{d๐ฉ^i}{dy}๐y,$$
(6)
To compare different systems colliding at different energies the reduced rapidity $`\stackrel{~}{y}=y/y_{proj}`$ and reduced impact parameter $`\stackrel{~}{b}=b/b_{max}`$ has been used. The maximum impact parameter for a symmetric system is $`b_{max}=2R_A`$. The value of $`\stackrel{~}{b}`$ in the simulations is varying from 0.15 (central collisions) up to 0.9 (most peripheral collisions).
The rapidity distributions of $`v_1`$ at SIS energies are shown in Fig. 1(a) for S+S and in Fig. 1(b) for Au+Au system. The directed flow of nucleons has a characteristic $`S`$-shape attributed to the standard $`p_x/A`$ distribution. Conventionally, we will call this type of flow, for which the slope $`dv_1/d\stackrel{~}{y}`$ is positive, normal flow, in contrast to the antiflow for which $`dv_1/d\stackrel{~}{y}<0`$ in the midrapidity region.
The nucleon flow reaches the maximum at $`\stackrel{~}{b}=0.30.45`$ both in S+S and Au+Au system, and then it drops. In the midrapidity range the flow can be well approximated by a linear dependence. The slope parameters of the $`v_1^N(y)`$ distributions are listed in Table I together with the $`dv_1^N/d\stackrel{~}{y}`$ data at higher energies. Pions at SIS energies show only weak antiflow which reaches a maximum around $`\stackrel{~}{b}=0.450.6`$ for S+S and $`\stackrel{~}{b}=0.60.75`$ for Au+Au collisions, i.e. in more peripheral collisions compared to the maximal nucleon directed flow. This behaviour is understandable since the evolution of a positive nuclear flow due to the (momentum dependent) repulsive $`NN`$-forces requires sufficiently large participant matter, whereas the negative pion flow due to shadowing needs large spectators. The antiflow of pions can also be fitted by a linear dependence; slope parameters are presented in Table I.
The directed flow $`v_1^i(y)`$ calculated for the same systems, S+S and Au+Au, at AGS energies is shown in Fig. 2(a) and Fig. 2(b), respectively. Here the distributions for nucleons differ considerably, especially in light system, from those at 1A GeV. The deviations of $`v_1^N(y)`$ from the straight line in S+S collisions begin noticeable already at $`\stackrel{~}{b}=0.3`$. The nucleon directed flow goes to zero in the midrapidity range with increasing impact parameter. Moreover, even antiflow is developed in very peripheral collisions at $`\stackrel{~}{b}=0.9`$, as seen in Fig. 2(a).
In contrast, in heavy Au+Au collisions at 11.6A GeV there are no singularities in the behaviour of $`v_1^N(y)`$ up to $`\stackrel{~}{b}=0.6`$. The plateau in the midrapidity region seems to build up only at $`\stackrel{~}{b}0.75`$, see Fig. 2(b). Pion directed flow has negative slope in the midrapidity range for both light and heavy colliding system. Values of the slope parameter are listed in Table I. It is worth to mention that $`v_1^\pi `$ increases as the reaction becomes more peripheral, and that both pion and nucleon directed flow does not vanish even at $`\stackrel{~}{b}=0.9`$ compared to the flow at SIS energies.
At the SPS energies the directed flow of nucleons has negative slope in the midrapidity region already in semiperipheral S+S collisions as depicted in Fig. 3(a). With the increase of the impact parameter the nucleon antiflow becomes stronger. The $`y`$-dependence of $`v_1^N`$ in Pb+Pb collisions is shown in Fig. 3(b). Here the deviations from the straight line start to develop at $`\stackrel{~}{b}=0.45`$ in the central rapidity window, $`|\stackrel{~}{y}|0.25`$. It is interesting that the slope of the antiflow of nucleons at $`\stackrel{~}{b}=0.9`$ is similar to that of the pion flow. The latter reaches maximum also at $`\stackrel{~}{b}=0.9`$ in Pb+Pb, as well as in S+S collisions.
The disappearance of the directed flow of hadrons can be regarded as an indication for a softening of the EOS due to a QGP-hadron phase transition. The simple hypothesis would be that, despite the absence of the plasma formation in the microscopic model, the colour field of quark-antiquark and quark-diquark strings can force the softening of the hadronic EOS. This idea explains the disappearance of the directed flow at energies of AGS and higher, but obviously fails to explain why the effect is stronger in peripheral collisions and for light systems like S+S. The correct explanation can be, therefore, that the apparent softening of the equation of state is in fact caused by the nuclear shadowing . The mechanism of the development of nuclear antiflow in the midrapidity range of nuclear peripheral collisions is elaborated in Sec. V.
The mean directed flow $`v_1`$ of pions and nucleons is shown in Fig. 4. Except nearly central events, the pion mean flow is negative for both light and heavy colliding system at all three energies. At the AGS and SPS energies $`v_1^\pi `$ rises steadily as the reaction becomes more peripheral. The mean directed flow of nucleons, which is always positive, seems also to exhibit a similar tendency. The maximum in $`v_1^N(b)`$ distribution is located around $`\stackrel{~}{b}=0.4`$ for Au+Au collisions at 1A GeV. It is shifted to $`\stackrel{~}{b}=0.6`$ at 11.6A GeV, and is completely dissolved at higher energies.
## IV Elliptic flow
Since the elliptic flow develops at the very beginning stage of nuclear collision, it might be even a better tool to probe the nuclear EOS under extreme conditions . Particularly, calculations based on a relativistic hadron transport model indicate a transition of elliptic flow from out-of-plane to in-plane for the case of the QGP formation in Au+Au collisions in the energy range $`111`$A GeV . Recent experimental data confirm the transition from negative to positive elliptic flow at $`E4`$A GeV, which was considered as indication of the softening of nuclear EOS. But can this change in the behaviour of elliptic flow be induced by some other reasons? To answer the question the microscopic study of elliptic flow of nucleons and pions has been performed at energies from 1A GeV to 160A GeV.
Figures 5(a) and 5(b) depict the elliptic flow of pions and nucleons in S+S and Au+Au collisions, respectively, at 1A GeV. The elliptic flow of pions is small and negative at $`\stackrel{~}{b}0.45`$. Nucleon elliptic flow is also negative in peripheral and semiperipheral collisions. The nucleon flow increases to maximum at $`\stackrel{~}{b}=0.75`$ and then drops. For heavy system the pionic flow is negative already at $`\stackrel{~}{b}=0.15`$, while the nucleon flow at $`\stackrel{~}{b}0.45`$ has two positive peaks, centred around $`|\stackrel{~}{y}|1.4`$, and the dip in the midrapidity region, where $`v_2^N`$ is negative. In peripheral collisions the positive peaks in the $`v_2^N(y)`$ distribution vanish, and the negative elliptic flow of nucleons becomes stronger. Generally, the spectator matter at target/projectile rapidities shows in-plane flow $`(v_2^N>0)`$ whereas the participant matter at midrapidity shows preferential out-of-plane emission $`(v_2^N<0)`$.
The elliptic flow of nucleons at the AGS energies, shown in Figs. 6(a) and 6(b), also has a two-hump structure both in S+S and in Au+Au collisions. But in Au+Au interactions the nucleon flow becomes negative in the central part of $`\stackrel{~}{y}`$-distribution only at $`\stackrel{~}{b}0.75`$. The peaks are quite noticeable and shifted closer to the center of the distributions. In Au+Au collisions the elliptic flow of both pions and nucleons is at least twice as large as in S+S interactions. At $`\stackrel{~}{b}0.6`$ the nucleon flow is positive, while the pionic flow becomes negative at midrapidity already at $`\stackrel{~}{b}=0.45`$.
At the SPS energies the elliptic flow of pions in S+S collisions is quite flat and slightly positive, as demonstrated in Fig. 7(a). The flow of nucleons in this reaction is also small and positive except for $`\stackrel{~}{b}0.75`$, where the negative dip around $`\stackrel{~}{y}=0`$ is built up. The negative flow is seen for nucleons in Pb+Pb interactions only at $`\stackrel{~}{b}=0.9`$ in Fig. 7(b). The origin of the negative values of $`v_2^N`$ at midrapidity can be linked to absorption of hadrons, emitted at $`\theta =90^{}`$ angle in the reaction plane, by the dense baryon rich spectators, while hadrons emitted out-of-plane remain almost unaffected. Note also that the positions of positive maxima in $`v_2^N(y)`$ distributions in Pb+Pb reactions are shifted to $`\stackrel{~}{y}0.45`$. Compared to $`v_2^\pi `$ in S+S interactions, the elliptic flow of pions in Pb+Pb collisions is large and positive. It has a noticeable dip at midrapidity only in very peripheral collisions.
The $`\stackrel{~}{b}`$-dependence of elliptic flow integrated over the whole rapidity range is presented in Fig. 8. In S+S collisions the mean elliptic flow of pions is quite weak for all three energies, though it appears to change the sign from negative at 1A GeV to positive at 160A GeV. The nucleon flow in S+S reaction is more distinct. It is negative at the SIS energies, while at both AGS and SPS energies $`v_2^N(\stackrel{~}{b})`$ is positive and almost constant.
The situation is changed drastically with the rise of the mass number of colliding nuclei from $`A=32`$ to $`A=197(208)`$. The mean elliptic flow of pions becomes positive at the AGS energies, in accord with the experimental results . The flow reaches maximum values at 160A GeV. It is easy to see that the strength of $`v_2^\pi (\stackrel{~}{b})`$ increases with $`\stackrel{~}{b}`$ rising to 0.75, which corresponds to $`b=10`$ fm in the calculations, and then drops. At the SIS energies the nucleon mean elliptic flow is positive in semicentral and semiperipheral events with $`\stackrel{~}{b}0.45`$ and negative at higher values of the impact parameter. Thus in the heavy system there appears a transition of in-plane to out-of-plane flow with decreasing centrality of the reaction. A detailed analysis of the EOS dependence on the nuclear mean field at SIS energies will be presented elsewhere. The nucleon flow has the maximum strength at 11.6A GeV, in contrast to the pion mean elliptic flow which rises continuously with increasing incident energy.
## V Discussion of the results
We see in Sec. III that directed flow of nucleons in (semi)peripheral heavy-ion collisions noticeably deviates from the straight line behaviour in the midrapidity range at AGS energies or higher. The comparison between light and heavy systems colliding at the same energy per nucleon shows that the effect is stronger in the light system. The most probable explanation of this phenomenon is nuclear shadowing. To clarify the idea the symmetric system of two nuclei at maximum overlap is shown in Fig. 9 for the three energies under consideration. In addition, Fig. 10 illustrates the development of the antiflow-like behaviour in the midrapidity region. As was discussed in, e.g. , the total flow of hadrons is a result of mutual cancellation of two competitive components, namely, the normal flow which follows the ongoing spectators, and the antiflow which develops towards the baryon dilute areas of collision. The normal flow integrated over the whole rapidity range is always slightly larger than the integrated antiflow. But in the midrapidity window the antiflow can dominate over its normal counterpart. For instance, hadrons with small rapidities, emitted early in the direction of normal flow in heavy-ion collision at 160A GeV, (see Fig. 10) will be absorbed by flying spectators. In contrast, hadrons, emitted in the direction of antiflow even at the angles close to $`\theta =90^{}`$ to the beam axis, propagate freely.
This effect can be reduced by (i) increasing the centrality of the collision and (ii) by decreasing the center-of-mass energy of colliding nuclei (see Fig. 9). In both cases the area where particles can be emitted without shadowing significantly shrinks. It is important to mention here that in heavy-ion collisions at collider energies, RHIC ($`\sqrt{s}=200`$ GeV) and LHC ($`\sqrt{s}=5.5`$ TeV), the disappearance of nucleon directed flow in the midrapidity range should emerge already in semicentral collisions with $`b3`$ fm.
But why the irregularities in $`v_1^N(y)`$-distribution start to develop in light system at smaller impact parameter compared to that of heavy system? To answer this question note that the larger volume of overlapping zone in heavy system leads to the intensive rescattering of baryons and increase of hadron emission from the central fireball. The spectators still absorb several early emitted hadrons, but this process becomes less efficient compared to that of the light system, where the isotropic particle radiation from the central part is not so strong. Since the effect can be misinterpreted as an evidence for the QGP formation, it should be subtracted from the analysis of experimental data.
The presence of spectators, which absorb hadrons early emitted in the direction of normal flow, affects also the development of elliptic flow. Particularly, it leads to the creation of the dip in $`v_2(y)`$ in midrapidity range. As expected from simple geometrical considerations, the effect is stronger in peripheral collisions.
The transition of elliptic flow from the out-of-plane to in-plane direction with the rise of energy from 1A GeV to 11.6A GeV can also be linked to change in geometry of colliding system. The Lorentz-contracted spectators, which rapidly fly away, provide more free space for the in-plane development of the flow than almost noncontracted nuclei, see Figs. 9 and 10.
It is worth to mention that the elliptic flow of nucleons as a function of impact parameter becomes more flat with rising energy of the collision, while the maximum in $`v_2^\pi (b)`$ distribution is shifted to very peripheral events. This tendency is clearly seen in Fig. 8. Figure 11 presents the comparison of the model simulations of the elliptic flow of charged pions in $`3<y<6`$ in Pb+Pb collisions at SPS energies with the experimental data . We see that the QGSM provides a good quantitative agreement with the experiment. Note that the behaviour of the elliptic flow of charged particles is determined by the proton elliptic flow at energies below 11.6A GeV and by the pionic elliptic flow at 160A GeV. It means that although the physics of rescattering in the QGSM is the same in peripheral and central collisions, the nuclear matter undergoes a transition from a baryon dominated to a meson dominated matter with rising energy of colliding nuclei. The transition is similar to the predicted in transition from hadronic to partonic degrees of freedom.
## VI Conclusions
The directed and elliptic flow of hadrons in heavy-ion collisions is very sensitive to the EOS of the nuclear medium. At low and intermediate energies (SIS) hadrons are the relevant degrees of freedom, and the intranuclear interactions, i.e. the mean field, determine the EOS as well as the reaction dynamics. With increasing energy new degrees of freedom are extended, and the formation of small domains of a QGP phase might happen already at the SPS energies or even below. Accompanied by the phase transition to the hadronic phase this enforces a softening of the EOS due to the dropping pressure. Thus, the disappearance of the directed flow in midrapidity range can be considered as an indication of a new state of matter. This conclusion is supported by hydrodynamic simulations. No deviations of the nucleon directed flow from the straight line in $`|\stackrel{~}{y}|1`$ range have been found in the one-fluid calculations with a pure hadronic EOS .
On the other hand, several microscopic models, which do not explicitly imply the QGP formation, predict larger or smaller deviations of the directed flow from the straight line behaviour which is presented at low and intermediate energies. These deviations are attributed to the shadowing effect, which plays a decisive role in the competition between normal flow and antiflow in (semi)peripheral ultrarelativistic collisions of nuclei. Hadrons, emitted with small rapidities at the onset of the collision in the antiflow area can propagate freely, while their counterparts will be absorbed by the flying massive spectators.
The signal becomes stronger with the rise of the impact parameter. In collisions with the same impact parameter the antiflow starts to dominate over the normal flow in the midrapidity range as the reaction becomes more energetic, i.e. the spectators are more Lorentz-contracted and more hadrons can be emitted unscreened with small rapidities in the direction of antiflow. Therefore, this effect should appear in (semi)central collisions with $`b3`$ fm at RHIC energies, and can imitate the softening of the EOS of hot and dense nuclear matter. However, the disappearance of directed flow due to shadowing is more distinct for light systems, like S+S or Ca+Ca, colliding with the same reduced impact parameter. In the case of a plasma creation the effect should be more pronounced in large systems like Pb+Pb. Thus, one can distinguish between the two phenomena, shadowing and quark-hadron phase transition, by the comparison of the directed flow of nucleons in the midrapidity range in light and heavy-ion collisions.
The elliptic flow of nucleons and pions is found to change its orientation from out-of-plane at 1A GeV to in-plane at 11.6A GeV. Since the dynamics of rescattering is the same, the effect can be explained by purely geometric reasons, such as stronger Lorentz-contraction of colliding nuclei. At higher colliding energies the contracted spectators leave the reaction zone faster, thus giving space for the growth of elliptic flow in the reaction plane.
Results of the simulations appear to favour a similarity of hadron rescattering in central and peripheral heavy-ion collisions at energies up to 160A GeV. QGSM predicts that the $`v_2^\pi (b)`$-distribution in Pb+Pb collisions at SPS energies increases as the reaction becomes more peripheral, in accord with the experimental data. The elliptic flow of pions in this reaction drops only for highly peripheral collisions somewhere at $`b12`$ fm. However, if the data will show the further rise of elliptic flow even at such impact parameters, this can be taken as an indication for new processes not included in present version of the model. The situation awaits better data on both directed and elliptic flow in the midrapidity range and in very peripheral collisions of light and heavy nuclei at ultrarelativistic energies.
Acknowledgements. We are thankful to L. Csernai, E. Shuryak, H. Sorge, H. Stรถcker, S. Voloshin, and Nu Xu for the fruitful discussions and comments. This work was supported in part by the Bundesministerium fรผr Bildung und Forschung (BMBF) under contract 06Tร887. |
warning/0002/cond-mat0002202.html | ar5iv | text | # Liouvillian Approach to the Integer Quantum Hall Effect Transition
## I Introduction
The metal-insulator transition in the integer quantum Hall effect (IQHE) is a reentrant zero temperature quantum phase transition in which the sample goes from an insulating phase with longitudinal conductivity $`\sigma _{\mathrm{xx}}=0`$ to another insulating phase by crossing a conducting critical point ($`\sigma _{\mathrm{xx}}0`$) as the magnetic field is varied. The critical point occurs between the plateaus of the Hall conductivity $`\sigma _{\mathrm{xy}}`$ and corresponds to the instance when the Fermi energy is at a critical energy located in the middle of one of the disorder broadened Landau levels.
In general, the disorder induced metal-insulator transition is a transition in the nature of the states (whether they are localized or delocalized) at the Fermi energy and it does not manifest itself in the density of states which remains smooth across the mobility edge. According to the one-parameter theory of scaling, the states of a two-dimensional noninteracting electron gas are all localized in the presence of arbitrary weak disorder. In the IQHE however, the presence of the strong magnetic field pointing perpendicular to the plane drastically changes the nature of the states near the middle of the Landau bands. In the noninteracting picture of the IQHE these states are characterized by a localization length
$`\xi (E)\xi _0\left|{\displaystyle \frac{EE_c^i}{E_0}}\right|^\nu `$ (1)
which determines the extent to which the eigenstates of energy $`E`$ are delocalized. Here $`\xi _0`$ denotes a characteristic length scale of the system, e.g. the magnetic length $`\mathrm{}`$ (see below) and $`E_0`$ a characteristic energy scale, e.g. the bandwidth or disorder strength. The critical energy $`E_c^i`$ is located in the middle of the $`i`$-th Landau band and, in an infinite size system, it is the only energy at which the one-particle eigenstates are delocalized within this Landau band. As the Fermi energy (or magnetic field) is varied, the conductivity $`\sigma _{\mathrm{xx}}`$ will change according to the nature of the states at that energy and sharp peaks in the longitudinal conductivity will be observed.
When studying the IQHE the interaction between the electrons is usually ignored and only the disorder is considered to be responsible for the localization of the single particle states. This assumption must be checked by comparing the predictions of the noninteracting theory to experimental results and the outcome of numerical calculations which include the interactions. The universal localization exponent $`\nu =2.34\pm 0.04`$ numerically obtained within a noninteracting theory is in excellent agreement with experimental measurements of $`\nu `$, but it remains a mystery why the strong interactions, which do affect the dynamical exponent $`z`$, does not seem to affect $`\nu `$. Here we adopt the noninteracting picture. We furthermore assume a strong magnetic field and a Zeeman splitting, which is much larger than the width of each disorder broadened Landau level. We can then focus on the transition within the lowest Landau level (LLL) and neglect the spin degree of freedom of the electrons.
It has been shown numerically that for a finite system delocalized one-particle wave functions near $`E_c`$ show multifractal properties characterized by a set of generalized fractal dimensions $`D_q`$. Also, dynamical studies have shown anomalous slow diffusion of wave packets constructed from these multifractal states. Diffusion can be studied using the spectral function of the disorder averaged retarded density-density correlation function. For the problem considered here the spectral function is given by (after dividing by $`\pi \mathrm{}\omega `$)
$`\overline{S}(r,\omega ;E)`$ $``$ $`<<{\displaystyle \underset{i,j}{}}\delta (E\mathrm{}\omega /2E_i)\delta (E+\mathrm{}\omega /2E_j)`$ (3)
$`\times \psi _i(0)\psi _i^{}(๐ซ)\psi _j(๐ซ)\psi _j^{}(0)>>.`$
Here the $`\psi _i(๐ซ)`$ denote one-particle eigenfunctions and $`E_i`$ the respective eigenenergies for an electron of a two-dimensional spinless electron gas which is subject to a perpendicular magnetic field and a disorder potential. $`\mathrm{}`$ indicates the ensemble average over the disorder. After taking the disorder average, translational invariance is restored and $`\overline{S}`$ only depends on the distance $`r|๐ซ|`$ from the origin of the plane. Assuming that the eigenstates which contribute in Eq. (3) for $`EE_c`$ are of multifractal character, it has been argued that $`\overline{S}`$ decays algebraically
$$\overline{S}(r,\omega 0;EE_c)\left(\frac{r}{\xi (E)}\right)^\eta ,$$
(4)
for $`\xi _0r\xi (E)`$. The anomalous diffusion exponent $`\eta `$ is related to the generalized fractal dimension via $`D_2=2\eta `$.
Assuming a generalized nonlocal (in time and space) relation between the current and the gradient of the density and using the continuity equation, the spectral function in momentum space $`S(q,\omega ;E)`$ at small $`q|๐ช|`$ and $`\omega `$, can be rewritten in terms of a generalized diffusion โcoefficientโ $`D(q,\omega )`$ for $`EE_c`$
$$S(q,\omega ;E)=\frac{\rho (E)}{\pi }\frac{\mathrm{}q^2D(q,\omega )}{[\mathrm{}\omega ]^2+[\mathrm{}q^2D(q,\omega )]^2},$$
(5)
where $`\rho (E)`$ is the density of states per unit area. In the limit of $`\omega ,q0`$ and for large enough system sizes, $`D(q,\omega )`$ is only a function of $`qL_\omega `$, where
$`L_\omega [\rho (E_c)\mathrm{}\omega ]^{1/2}.`$ (6)
Through numerical diagonalization and using Eq. (5) Chalker and Daniell have shown that $`D(q,\omega )`$ approaches a constant $`D_0`$ for small $`qL_\omega `$. The precise value of $`D_0`$ is important since the longitudinal conductivity at the critical point is given by the Einstein relation $`\sigma _{\mathrm{xx}}=e^2\rho (E_c)D_0`$ and is expected to be universal. The $`qL_\omega 0`$ limit of $`D(q,\omega )`$ has later been reinvestigated in an extended numerical study. For $`qL_\omega 1`$, but still in the limit of $`q,\omega 0`$, $`D(q,\omega )`$ decays as
$$D(q,\omega )D_0(qL_\omega )^\eta .$$
(7)
For the anomalous diffusion exponent $`\eta `$ Chalker and Daniell obtain the numerical value $`\eta =0.38\pm 0.04`$, indicating that the delocalized states near the critical energy indeed have multifractal properties. This value for $`\eta `$ has later been confirmed in other numerical studies. For energies $`E`$ away from $`E_c`$, $`S(q,\omega ;E)`$ vanishes in the small $`q`$ and $`\omega `$ limit independent of the order in which the limits are taken due to exponential localization of the states.
Most of the progress in the theoretical understanding of the localization-delocalization transition considered here has been through numerical calculations. Although a field theory has been proposed some time ago by Pruisken and co-workers, up to now no quantitative results such as the critical exponents of the transition have been obtained within this description. More recent studies have introduced alternative field theories. Within the framework of these theories it might in the future be possible to analytically determine critical exponents as has been recently successfully achieved for the SU(2) version of the network model. In this paper we present a novel approach to the transition which may prove more tractable. Although thus far we have not been able to analytically calculate the spectral function Eq. (3) beyond the self-consistent Born approximation, we have numerically verified the possibility of obtaining the critical exponent $`\nu `$ using this approach.
We start by defining the density correlation function at zero temperature as
$$\stackrel{~}{\mathrm{\Pi }}(q,t;E)\frac{i\theta (t)}{N\mathrm{}\mathrm{}^2}\mathrm{Tr}\{\overline{\rho }_๐ช(t)\overline{\rho }_๐ช(0)\delta (EH)\},$$
(8)
with the one-particle Hamiltonian $`H=H_0+H_D`$. Here $`H_0`$ denotes the kinetic energy of a spinless electron moving in the plane in the presence of a perpendicular magnetic field and $`H_D`$ is the potential energy for a fixed realization of the disorder potential $`V(๐ซ)`$. $`\mathrm{}`$ is the magnetic length given by $`\mathrm{}^2=\mathrm{}c/(eB)`$, where $`B`$ is the strength of the magnetic field, and $`N=L^2/(2\pi \mathrm{}^2)`$ is the number of states in the LLL. We consider a square sample of area $`L^2`$. By projecting the one-particle density operator $`\rho _๐ช\mathrm{exp}(i๐ช๐ซ)`$ onto the LLL, denoting the projected density by $`\overline{\rho }_๐ช`$ (see Sec. II), and taking the one-particle trace $`\mathrm{Tr}`$ over the states in the LLL, we restrict our considerations to the transition in the LLL. It will turn out that the equation of motion for the density operators restricted to the LLL can be solved formally in this case. In the small $`\omega `$ limit we have
$`S(q,\omega ;E)={\displaystyle \frac{1}{2\pi ^2}}\mathrm{Im}\mathrm{\Pi }(q,\omega ;E).`$ (9)
Instead of dealing with $`\stackrel{~}{\mathrm{\Pi }}(q,t;EE_c)`$ directly, we will integrate $`\stackrel{~}{\mathrm{\Pi }}(q,t;E)`$ over all energies $`E`$ and focus our attention on
$$\stackrel{~}{\mathrm{\Pi }}(q,t)i\frac{\theta (t)}{N\mathrm{}\mathrm{}^2}\mathrm{Tr}\{\overline{\rho }_๐ช(t)\overline{\rho }_๐ช(0)\}.$$
(10)
Since the localization length only diverges at $`E_c`$, the energetically unconstrained diffusion problem considered by investigating $`\mathrm{Im}\mathrm{\Pi }(q,\omega )`$ still contains useful information about critical exponents. For instance, let us suppose that at time $`t=0`$ we create a wave packet localized at the origin constructed from all the states of the system (localized as well as delocalized states). For large $`t`$ only โdelocalizedโ states with $`\xi (E)/r>1`$ can contribute to the probability amplitude of the wave packet at a distance $`r/\xi _01`$ far away from the origin. This implies that in the limit of small $`q\xi _0`$ only states with $`q\xi (E)>1`$ and thus $`|EE_c|<E_0(q\xi _0)^{1/\nu }`$ contribute to the right hand side of Eq. (10). Hence for $`q\xi _00`$ only a fraction $`(q\xi _0)^{1/\nu }`$ of the states in the LLL contributes and we expect from Eq. (5), that for small $`qL_\omega `$
$$\mathrm{}\omega \mathrm{}^2\mathrm{Im}\mathrm{\Pi }(q,\omega )(q\xi _0)^{1/\nu }\frac{(D_0q^2/\omega )}{1+(D_0q^2/\omega )^2},$$
(11)
where the diffusion parameter is a constant $`D_0`$. The above argument, which we confirm numerically in Sec. V, gives a strong indication that some useful information about the quantum phase transition can be extracted from $`\mathrm{\Pi }(q,\omega )`$. We again emphasize that this is so because the delocalization only occurs at a single critical energy, a characteristic unique to the IQHE where the extended states have zero measure in the energy spectrum. We also point out the importance of the order of limits in obtaining Eq. (11). The limit of $`q`$,$`\omega 0`$ is taken by having $`q`$ approach zero faster than $`\omega `$ so as to obtain a finite diffusion constant. In contrast to the usual approach, in which information about the anomalous diffusion exponent $`\eta `$ is extracted from the spectral function $`S(q,\omega ;EE_c)`$ we will be able to extract information about the localization exponent $`\nu `$ using the same spectral function but integrated over all energies $`E`$.
We will show that $`\mathrm{\Pi }(q,\omega )`$, which is an inherent fermionic disorder averaged two-particle correlation function, can be re-expressed as the single particle correlation function of an interacting (after the disorder average has been performed) dynamical system with an unusual action. Therefore, in order to extract the dynamical behavior of the original problem, one simply has to study the disorder averaged density of states of this new action.
The rest of this paper is organized as follows. In Sec. II we introduce the model and mapping of the problem to the new โHamiltonianโ. In Sec. III we calculate $`\mathrm{\Pi }(q,\omega )`$ within the self-consistent Born approximation. It displays normal diffusion at this level of approximation. In Sec. IV we introduce the field theoretical approach to the disorder averaging. In Sec. V we will demonstrate numerically the validity of the scaling hypothesis stated in Eq. (11), and finally in Sec. VI we present our conclusions.
## II Model and mapping
We consider the two-dimensional spinless electron gas lying in the $`x`$-$`y`$ plane which is subject to a perpendicular magnetic field $`๐=B\widehat{๐ณ}`$ and an external potential $`V(๐ซ)`$. $`\widehat{๐ณ}`$ denotes the unit vector in the $`z`$ direction. In the symmetric gauge the vector potential is given by $`๐=\frac{1}{2}๐ซ\times ๐`$ and the one-particle Hamiltonian reads
$`H`$ $`=`$ $`H_0+H_D`$ (12)
$`=`$ $`{\displaystyle \frac{1}{2m}}\left[๐ฉ+{\displaystyle \frac{e}{c}}๐\right]^2+V(๐ซ).`$ (13)
We restrict our investigations to the LLL and thus project the Hamiltonian onto the states in the LLL. The kinetic energy of all the LLL states is the same and after projecting leads to a constant which we will neglect in what follows. Writing the potential energy in Fourier space the Hamiltonian simplifies to
$$H=\underset{๐ช}{}v(๐ช)\overline{\rho }_๐ช,$$
(14)
where $`v(๐ช)`$ is the Fourier transform of the disorder potential. The projected density operator is given by
$$\overline{\rho }_๐ชe^{\frac{1}{4}\mathrm{}^2q^2}\tau _๐ช,$$
(15)
with $`\tau _๐ช`$ being the unitary magnetic translation operator which translates the electron a distance $`(๐ช\times \widehat{๐ณ})\mathrm{}^2`$. The formalism needed to project the density operator $`\rho _๐ช=e^{i๐ช๐ซ}`$ onto the LLL was developed elsewhere.
The magnetic translation operators have the following special property:
$$\tau _๐ช\tau _๐ฉ=\mathrm{exp}\left(\frac{i\mathrm{}^2}{2}qp\right)\tau _{๐ช+๐ฉ},$$
(16)
where $`qp\left(๐ช\times ๐ฉ\right)\widehat{๐ณ}`$. Hence their commutation relation defines a closed Lie algebra:
$$[\tau _๐ช,\tau _๐ฉ]=2i\mathrm{sin}\left(\frac{\mathrm{}^2}{2}qp\right)\tau _{๐ช+๐ฉ}.$$
(17)
Also we have
$$\mathrm{Tr}\left\{\tau _๐ช\right\}=N\delta _{๐ช,0}.$$
(18)
The latter can be proved by noting that the left hand side is proportional to the one-particle trace of $`\overline{\rho }_๐ช`$. Since the trace is taken over states in the LLL, the projection is unnecessary and we have
$$\mathrm{Tr}\left\{\overline{\rho }_๐ช\right\}=\mathrm{Tr}\left\{e^{i๐ช๐ซ}\right\},$$
(19)
which vanishes unless $`๐ช=0`$.
If there are $`N`$ states in the Hilbert space, there are $`N^2`$ independent operators on the space. However there are exactly $`N^2`$ different wave vectors on the torus, so the set of operators $`\overline{\rho }_๐ช`$ is โcompleteโ; it spans the set of all operators. The Hamiltonian can be expressed in terms of the $`\overline{\rho }_๐ช`$ and the Heisenberg equation of motion of the $`\overline{\rho }_๐ช`$ is closed. This allows us to define the quantum โLiouvillianโ matrix by
$$\dot{\tau }_๐ช(t)=i\underset{๐ช^{}}{}_{\mathrm{๐ช๐ช}^{}}\tau _๐ช^{}(t).$$
(20)
From the simple commutation properties Eq. (17) of the $`\tau _๐ช`$ it readily follows that
$$_{\mathrm{๐ช๐ช}^{}}\frac{2i}{\mathrm{}}v(๐ช๐ช^{})e^{\frac{1}{4}\mathrm{}^2|๐ช^{}๐ช|^2}\mathrm{sin}\left(\frac{\mathrm{}^2}{2}q^{}q\right).$$
(21)
Using the Liouvillian matrix we can immediately write down the formal solution of the equation of motion Eq. (20) for $`\tau _๐ช(t)`$
$$\tau _๐ช(t)=\underset{๐ช^{}}{}\left(e^{it}\right)_{\mathrm{๐ช๐ช}^{}}\tau _๐ช^{}(0).$$
(22)
This leads to a simple expression for the density-density correlation function defined in Eq. (10)
$$\stackrel{~}{\mathrm{\Pi }}(q,t)=i\frac{\theta (t)}{\mathrm{}\mathrm{}^2}e^{\frac{1}{2}\mathrm{}^2q^2}\left(e^{it}\right)_{\mathrm{๐ช๐ช}}.$$
(23)
We can define an $`N^2`$ element operator โsuperspaceโ and view $``$ as the โHamiltonianโ. From this point of view, finding $`\mathrm{Im}\mathrm{\Pi }(q,\omega )`$ is the same as finding the one-particle density of states for a system with Hamiltonian $``$:
$`\mathrm{\Pi }(q,\omega )`$ $`=`$ $`{\displaystyle \frac{i}{\mathrm{}\mathrm{}^2}}e^{\frac{1}{2}\mathrm{}^2q^2}{\displaystyle _0^{\mathrm{}}}๐te^{i(\omega +i\delta )t}๐ช\left|e^{it}\right|๐ช`$ (24)
$`=`$ $`{\displaystyle \frac{1}{\mathrm{}\mathrm{}^2}}e^{\frac{1}{2}\mathrm{}^2q^2}<๐ช|{\displaystyle \frac{1}{\omega +i\delta }}|๐ช>,`$ (25)
where we have introduced states $`|๐ช`$ with $`๐ช\left|\right|๐ช^{}_{\mathrm{๐ช๐ช}^{}}`$ and $`\delta `$ is an infinitesimal small positive number.
This remarkable formula is our central result. Let us now try to understand its import. In a crude sense it represents a kind of bosonization of the problem. Ordinarily in an interacting many-body system the equations of motion for the density are not closed but rather involve a hierarchy of additional operators. However for the special case of one-dimension and a linear dispersion relation (the Tomonaga-Luttinger model) the equations of motion are closed and the density fluctuations become free bosons even though the underlying particles are interacting. In the present problem (without electron-electron interactions) the equations for the density operators close after projection onto a single Landau level (which for simplicity we have taken to be the lowest). This has several advantages. First we do not have to work separately with retarded and advanced one-particle Greenโs functions and their products. Secondly we note that there are no problems with gauge invariance and conserving approximations. This is because the Liouvillian matrix elements $`_{\mathrm{๐ช๐ช}^{}}`$ vanish if either $`๐ช`$ or $`๐ช^{}`$ vanish. Thus the total charge in the system is automatically conserved. Finally this representation allows us to establish a hierarchy of length and time scales which should be suitable for renormalization group (RG) analysis. Because the kinetic energy has been quenched, high momentum of a particle is not associated with high energy. Since the Liouvillian vanishes at small wavevectors, it naturally organizes the decay rates of density fluctuations into short time scales at large wavevectors and long time scales at small wavevectors. As we comment further below however, there are technical obstacles to be overcome before this RG can be carried out.
We take the disorder to be gaussian distributed, but not necessarily white noise, i.e. possibly smoothed. We then have
$$v(๐ช)=0$$
(26)
and
$$v(๐ช)v(๐ช^{})=\frac{2\pi \alpha ^2v^2}{L^2}e^{\frac{1}{2}\mathrm{}^2q^2(\alpha ^21)}\delta _{๐ช+๐ช^{},0},$$
(27)
which in real space translates into
$$V(๐ซ)V(๐ซ^{})=\frac{\alpha ^2v^2}{\mathrm{}^2(\alpha ^21)}\mathrm{exp}\left[\frac{|๐ซ๐ซ^{}|^2}{2\mathrm{}^2(\alpha ^21)}\right].$$
(28)
Here $`v`$ denotes the strength of the disorder potential and $`\alpha `$ is a dimensionless smoothness parameter. In the limit of a distribution which is extremely smooth ($`\alpha \mathrm{}`$), the one-particle electronic density of states approaches a gaussian
$$\rho _{\alpha =\mathrm{}}(E)=\frac{1}{(2\pi )^{3/2}\mathrm{}^2v}\mathrm{exp}\left[\frac{1}{2v^2}(EE_c)^2\right].$$
(29)
An integration over all energies $`E`$ gives the number of states in the LLL divided by the sample area $`N/L^2`$ which is $`1/(2\pi \mathrm{}^2)`$. For $`\alpha =1`$ the disorder distribution goes over to the uncorrelated white noise distribution for which Wegner has determined the density of states. At $`E=E_c`$ it is given by
$$\rho _{\alpha =1}(E_c)=\frac{\sqrt{2}}{\pi ^2\mathrm{}^2v}.$$
(30)
## III Self-consistent Born approximation
We next calculate $`\mathrm{\Pi }(q,\omega )`$ Eq. (25) in the self-consistent Born approximation. We define the complex self-energy $`\mathrm{\Sigma }(q,\omega )=\mathrm{\Sigma }_R(q,\omega )+i\mathrm{\Sigma }_I(q,\omega )`$ for the propagator
$$\widehat{\mathrm{\Pi }}(q,\omega )\mathrm{}\mathrm{}^2e^{\frac{1}{2}\mathrm{}^2q^2}\mathrm{\Pi }(q,\omega ).$$
(31)
by setting
$$\widehat{\mathrm{\Pi }}(q,\omega )=\frac{1}{\omega +i\delta \mathrm{\Sigma }(q,\omega )}.$$
(32)
Within the self-consistent Born approximation the self-energy is given by the expression
$`\mathrm{\Sigma }^B(q,\omega )={\displaystyle \underset{๐ฉ}{}}_{๐ช,๐ช+๐ฉ}_{๐ช+๐ฉ,๐ช}\widehat{\mathrm{\Pi }}^B(|๐ช+๐ฉ|,\omega ).`$ (33)
In contrast to standard many-body perturbation theory the right hand side of this expression does not contain an energy sum. In this approximation all non-crossing diagrams for the propagator $`\widehat{\mathrm{\Pi }}(q,\omega )`$ are summed, as shown in Fig. 1. In this figure a thick solid line stands for $`\widehat{\mathrm{\Pi }}^B(q,\omega )`$ and a thin solid line indicates the โnoninteractingโ propagator $`\widehat{\mathrm{\Pi }}^0(q,\omega )`$, which is given by Eq. (32) with $`\mathrm{\Sigma }(q,\omega )0`$. $`\widehat{\mathrm{\Pi }}^0`$ is independent of $`q`$. The consequences of this for a perturbative treatment will be discussed in the next section. In Fig. 1 the vertex with an incoming and an outgoing solid line and a dashed line stands for a matrix element $`_{\mathrm{๐ช๐ช}^{}}`$ of the Liouvillian. The disorder average introduces โcontractionsโ, i.e. connections, between the dashed lines. In general the Hartree terms are included in the partial sum, but as indicated in Fig. 1, they vanish because of the $`qp`$ term in the matrix elements of the Liouvillian.
Using the distribution introduced in the last section \[see Eq. (27)\] and the definition of the Liouvillian matrix Eq. (21) we obtain the self-consistency equation
$$\mathrm{\Sigma }^B(q,\omega )=\frac{2\pi \alpha ^2v^2}{\mathrm{}^2L^2}\underset{๐ฉ}{}\frac{e^{\frac{1}{2}\mathrm{}^2\alpha ^2|๐ช๐ฉ|^2}4\mathrm{sin}^2(\frac{\mathrm{}^2}{2}qp)}{\omega +i\delta \mathrm{\Sigma }^B(p,\omega )}.$$
(34)
The strength of the disorder $`v`$ can be scaled out of this equation by replacing $`\mathrm{\Sigma }^B\mathrm{}\mathrm{\Sigma }^B/v`$ and $`\omega \mathrm{}\omega /v`$.
As explained in the introduction the diffusive properties can be read off from the small $`q`$ and $`\omega `$ limit of the imaginary part of $`\mathrm{\Pi }`$. For $`q0`$ we have $`\widehat{\mathrm{\Pi }}(q,\omega )\mathrm{}\mathrm{}^2\mathrm{\Pi }(q,\omega )`$ and can thus write
$`\mathrm{}\mathrm{}^2\mathrm{Im}\mathrm{\Pi }^B(q,\omega )={\displaystyle \frac{\mathrm{\Sigma }_I^B(q,\omega )}{\left[\omega \mathrm{\Sigma }_R^B(q,\omega )\right]^2+\left[\mathrm{\Sigma }_I^B(q,\omega )\right]^2}}.`$ (35)
Eq. (34) can be solved numerically by iteration. Following Eq. (11) the best way to extract the diffusive properties is a โscalingโ plot in which (for fixed $`v`$ and $`\alpha `$) $`\mathrm{}\omega \mathrm{}^2\mathrm{Im}\mathrm{\Pi }^B(q,\omega )`$ is plotted as a function of $`qL_\omega `$ for different small $`q`$ and $`\omega `$. Such an evaluation shows that on this level of approximation $`\mathrm{}\omega \mathrm{}^2\mathrm{Im}\mathrm{\Pi }^B(q,\omega )`$ is a function of $`qL_\omega `$ only and thus does not display a sign of the prefactor $`(q\xi _0)^{1/\nu }`$ discussed in connection with Eq. (11). Furthermore $`\mathrm{Im}\mathrm{\Pi }^B(q,\omega )`$ only shows normal diffusion with a diffusion constant $`D_0`$ which for $`q0`$ and $`\omega 0`$ is independent of $`qL_\omega `$. We thus conclude that (as expected) the occurrence of the critical exponents $`\nu `$ and $`\eta `$ is a higher order fluctuation effect. For $`\omega 0`$, $`\mathrm{\Sigma }_R^B(q,\omega )`$ goes to zero for all $`q`$. Thus $`D_0`$ is given by
$`D_0=\underset{\omega 0}{lim}\underset{q0}{lim}\mathrm{\Sigma }_I^B(q,\omega )/q^2.`$ (36)
Because of the scaling property discussed following Eq. (34) $`D_0`$ is proportional to $`v`$. As shown in Fig. 2 $`D_0`$ also depends on the smoothness $`\alpha `$ of the disorder. Between $`\alpha =1`$ (white noise) and $`\alpha =2`$, $`D_0`$ changes by approximately ten percent. For $`\alpha >2`$ the $`\alpha `$ dependence is extremely weak and for $`\alpha \mathrm{}`$, $`D_0`$ saturates at $`D_0^{\alpha =\mathrm{}}0.828v\mathrm{}^2/\mathrm{}`$. For $`\alpha =1`$ we find $`D_0^{\alpha =1}0.965v\mathrm{}^2/\mathrm{}`$.
Using the Einstein relation for the conductivity and Eqs. (29) and (30) we obtain
$`\sigma _{\mathrm{xx}}^{\alpha =\mathrm{}}0.330{\displaystyle \frac{e^2}{h}}`$ (37)
and
$`\sigma _{\mathrm{xx}}^{\alpha =1}0.869{\displaystyle \frac{e^2}{h}}.`$ (38)
If one is interested in the large $`\alpha `$ limit it might be tempting to expand the sine in Eq. (34), as only small $`p`$ contribute to the sum due to the exponential function. Anticipating that for small $`q`$ the self-energy is quadratic in $`q`$ the ansatz $`\mathrm{\Sigma }^B(q,\omega )=iq^2\stackrel{~}{D}_0`$ seems to be plausible. Then the self-consistency equation can be solved analytically leading to $`\stackrel{~}{D}_0=(1/\sqrt{2})v\mathrm{}^2/\mathrm{}0.707v\mathrm{}^2/\mathrm{}`$. A comparison with $`D_0^{\alpha =\mathrm{}}`$ discussed above shows that this procedure does not give the correct large $`\alpha `$ value for $`D_0`$. This is due to the fact that in the exact solution of Eq. (34) the range of $`q`$ values over which $`\mathrm{\Sigma }^B(q,\omega )`$ can be approximated by a purely quadratic function in $`q`$ shrinks as $`1/\alpha `$. Thus in the limit $`\alpha \mathrm{}`$ it would be necessary to include higher order terms in the expansion of $`\mathrm{\Sigma }^B(q,\omega )`$ in order to reproduce the numerical result in Eq. (37). Note that $`\sigma _{\mathrm{xx}}^{\alpha 1}`$ obtained above is independent of the correlation parameter $`\alpha `$ as it should in the limit $`\alpha 1`$. Since the exact conductivity is universal, the present result is a considerable improvement over the traditional self-consistent Born approximation result for which the conductivity vanishes like $`\alpha ^1`$ in this limit.
In a previous numerical study it was found that $`\sigma _{\mathrm{xx}}=(0.54\pm 0.04)e^2/h`$ independent of the smoothness of the disorder. The results for $`\sigma _{\mathrm{xx}}`$ obtained within our approach are of the same order of magnitude as the one calculated using purely numerical methods but in contrast to this one our results depend on $`\alpha `$. This is due to the fact that we have calculated $`D_0`$ within the self-consistent Born approximation but included in the Einstein relation the exact density of states at the critical energy.
Using our approach of calculating the disorder averaged one-particle correlation function for the dynamical system described by the Liouvillian, we observe normal diffusion already at the level of the self-consistent Born approximation. In the usual fermionic picture of noninteracting electrons in the presence of disorder and a magnetic field, much more elaborate techniques, as e.g. Borel resummation, instanton methods, the replica trick, and the supersymmetry method, are used to obtain similar results. In particular, in the more traditional approaches, diffusion is not obtained at the saddle point level and it is necessary to include gaussian fluctuations (i.e. sum ladder diagrams) to obtain diffusion. Because we deal directly with the density itself, we obtain diffusion even at the saddle point level.
## IV Field theoretical approach
To go beyond the self-consistent Born approximation it might prove advantageous to bring our approach into a field theoretical framework. This is what we will do in this section. In reformulating $`\mathrm{\Pi }(q,\omega )`$ using field theoretical methods we use the gaussian integral identity
$$i\overline{\psi }_๐ช\psi _๐ช=<๐ช\left|\frac{1}{\omega +i\delta }\right|๐ช>,$$
(39)
where
$$\overline{\psi }_๐ช\psi _๐ช\frac{1}{Z}๐\overline{\psi }๐\psi e^{S_\psi }\overline{\psi }_๐ช\psi _๐ช,$$
(40)
and
$$S_\psi i\underset{๐ค,๐ค^{}}{}\overline{\psi }_๐ค[\omega +i\delta ]_{๐ค,๐ค^{}}\psi _๐ค^{}.$$
(41)
The $`\psi _๐ช`$ denote complex (bosonic) fields and $`Z`$ is given by
$$Z๐\overline{\psi }๐\psi e^{S_\psi }.$$
(42)
In order to ensemble average over the disorder we introduce additional Grassmann variables to represent $`1/Z`$ as a path integral
$$\frac{1}{Z}=๐\overline{\xi }๐\xi e^{S_\xi },$$
(43)
where
$$S_\xi i\underset{๐ค,๐ค^{}}{}\overline{\xi }_๐ค[\omega +i\stackrel{~}{\delta }]_{๐ค,๐ค^{}}\xi _๐ค^{}.$$
(44)
One can then carry out the ensemble average over the gaussian distributed disorder and obtains the generalized functional
$$\overline{Z}(\omega )=๐\overline{\xi }๐\xi ๐\overline{\psi }๐\psi e^{F(\omega )},$$
(45)
where
$`F(\omega )`$ $``$ $`{\displaystyle \underset{๐ค}{}}\left[(i\omega +\delta _๐ค)\overline{\psi }_๐ค\psi _๐ค+(i\omega +\stackrel{~}{\delta }_๐ค)\overline{\xi }_๐ค\xi _๐ค\right]`$ (48)
$`+{\displaystyle \underset{๐ค,๐ค^{}}{}}{\displaystyle \underset{๐ฉ,๐ฉ^{}}{}}_{\mathrm{๐ค๐ค}^{}}_{\mathrm{๐ฉ๐ฉ}^{}}[\overline{\psi }_๐ค\overline{\psi }_๐ฉ\psi _๐ฉ^{}\psi _๐ค^{}`$
$`+2\overline{\psi }_๐ค\psi _๐ค^{}\overline{\xi }_๐ฉ\xi _๐ฉ^{}+\overline{\xi }_๐ค\overline{\xi }_๐ฉ\xi _๐ฉ^{}\xi _๐ค^{}].`$
Here we have let $`\delta \delta _๐ค`$ so that we can generate the correlation functions by
$$\overline{\xi }_๐ช\xi _๐ช_\omega =\overline{\psi }_๐ช\psi _๐ช_\omega =\frac{\overline{Z}(\omega )}{\stackrel{~}{\delta }_๐ช}=\frac{\overline{Z}(\omega )}{\delta _๐ช}.$$
(49)
Once the disorder averaging is done we finally obtain
$`F(\omega )`$ $`=`$ $`i{\displaystyle \underset{๐ช}{}}\left[(\omega +i\delta _๐ช)\overline{\psi }_๐ช\psi _๐ช+(\omega +i\stackrel{~}{\delta }_๐ช)\overline{\xi }_๐ช\xi _๐ช\right]`$ (52)
$`+{\displaystyle \underset{๐ช_1,๐ช_2,๐ช_3,๐ช_4}{}}f(1,2,3,4)[\overline{\psi }_{๐ช_1}\overline{\psi }_{๐ช_2}\psi _{๐ช_3}\psi _{๐ช_4}`$
$`+2\overline{\psi }_{๐ช_1}\psi _{๐ช_4}\overline{\xi }_{๐ช_2}\xi _{๐ช_3}+\overline{\xi }_{๐ช_1}\overline{\xi }_{๐ช_2}\xi _{๐ช_3}\xi _{๐ช_4}],`$
with
$`f(1,2,3,4)`$ $`=`$ $`{\displaystyle \frac{\pi \alpha ^2v^2}{\mathrm{}^2L^2}}e^{\frac{1}{2}\mathrm{}^2\alpha ^2|๐ช_1๐ช_4|^2}4\mathrm{sin}\left({\displaystyle \frac{\mathrm{}^2}{2}}q_1q_4\right)`$ (54)
$`\times \mathrm{sin}\left({\displaystyle \frac{\mathrm{}^2}{2}}q_2q_3\right)\delta _{๐ช_1+๐ช_2,๐ช_3+๐ช_4}.`$
In contrast to standard many-body theory the action Eq. (52) does not contain a sum over the frequency. $`\omega `$ only enters this equation as an external parameter. As already discussed in the last section the noninteracting propagator ($`v=0`$) is given by $`(\omega +i\delta )^1`$ and does not depend on $`q`$. Thus a perturbation theory or RG procedure can only be set up after a $`q`$ dependent propagator has been generated by self-consistently summing up an entire class of diagrams, as e.g. the non-crossing diagrams in Sec. III. Furthermore the interaction $`f`$ in Eq. (54) has an unusual momentum dependence compared to standard standard $`\varphi ^4`$ theory of critical phenomena: It vanishes if one of the $`๐ช_i`$ goes to zero and is periodic in the momenta.
Using the field theoretical approach we can reproduce the approximation discussed in Sec. III, which is usually called self-consistent mean-field or saddle point approximation in the present context. In the absence of symmetry breaking, the middle of the three quartic terms in the action cannot contribute to the saddle point solution since its coefficient vanishes for $`๐ช_1=๐ช_4`$ and $`๐ช_2=๐ช_3`$. Hence we can deal separately with the bosonic and the fermionic variables when discussing the saddle point solution. By performing the usual pairing of the fields in the quartic interaction term at the mean-field level we have
$`\overline{\psi }_{๐ช_1}\overline{\psi }_{๐ช_2}\psi _{๐ช_3}\psi _{๐ช_4}`$ $`=`$ $`i\widehat{\mathrm{\Pi }}^{\mathrm{MF}}(q_1,\omega )\delta _{๐ช_1,๐ช_3}\overline{\psi }_{๐ช_2}\psi _{๐ช_4}`$ (56)
$`+i\widehat{\mathrm{\Pi }}^{\mathrm{MF}}(q_2,\omega )\delta _{๐ช_2,๐ช_4}\overline{\psi }_{๐ช_1}\psi _{๐ช_3}.`$
Thus we can write
$$F^{\mathrm{MF}}(\omega )=\underset{๐ช}{}\overline{\psi }_๐ช\left[i\omega +\delta +i\mathrm{\Sigma }^{\mathrm{MF}}(q,\omega )\right]\psi _๐ช,$$
(57)
and use this in calculating
$`i\widehat{\mathrm{\Pi }}^{\mathrm{MF}}(q,\omega )`$ $``$ $`\overline{\psi }_๐ช\psi _๐ช_\omega ^{\mathrm{MF}}={\displaystyle \frac{๐\overline{\psi }๐\psi e^{F^{\mathrm{MF}}(\omega )}\overline{\psi }_๐ช\psi _๐ช}{๐\overline{\psi }๐\psi e^{F^{\mathrm{MF}}(\omega )}}}`$ (58)
$`=`$ $`{\displaystyle \frac{i}{\omega +i\delta \mathrm{\Sigma }^{\mathrm{MF}}(q,\omega )}}=i\widehat{\mathrm{\Pi }}^\mathrm{B}(q,\omega ),`$ (59)
which reproduces the self-consistency Eq. (34) for the self-energy.
At present we do not know how to evaluate the correlation function beyond the self-consistent mean-field approximation in a controlled way. However, we hope that in the future it will be possible to analytically extend our results.
## V Numerical results
In this section we will numerically calculate $`\mathrm{Im}\mathrm{\Pi }(q,\omega )`$ by exact diagonalization and verify the scaling hypothesis stated in Eq. (11). We closely follow the procedure and notation used by Chalker and Daniell. Motivated by Eq. (9) we define
$`\overline{S}(r,\omega ){\displaystyle \frac{1}{2\pi ^2}}\mathrm{Im}\overline{\mathrm{\Pi }}(r,\omega )`$ (60)
$`=<<{\displaystyle \underset{i,j}{}}\delta (\mathrm{}\omega +E_iE_j)\psi _i(0)\psi _i^{}(๐ซ)\psi _j(๐ซ)\psi _j^{}(0)>>.`$ (61)
The single particle wave functions $`\psi _i(๐ซ)`$ can be expanded in the basis of the elliptical theta functions $`\varphi _m(๐ซ)`$
$`\psi _i(๐ซ)={\displaystyle \underset{m=1}{\overset{N}{}}}a_i(m)\varphi _m(๐ซ),`$ (62)
where
$`\varphi _m(x,y)`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{L\mathrm{}\pi ^{1/2}}}}{\displaystyle \underset{s=\mathrm{}}{\overset{\mathrm{}}{}}}\mathrm{exp}\left(iX_{m,s}y/\mathrm{}^2\right)`$ (64)
$`\times \mathrm{exp}\left[(xX_{m,s})^2/(2\mathrm{}^2)\right],`$
and
$`X_{m,s}=m{\displaystyle \frac{2\pi }{L}}\mathrm{}^2+sL.`$ (65)
Then the Fourier transform of Eq. (61) can be written as
$`S(q,\omega )`$ $`=`$ $`{\displaystyle \frac{1}{2\pi \mathrm{}^2N^2}}e^{\frac{1}{2}\mathrm{}^2q^2}`$ (67)
$`\times {\displaystyle \underset{i,j}{}}\delta (\mathrm{}\omega +E_iE_j)Q_{ij}(k,l),`$
where
$`Q_{ij}(k,l)`$ $`=`$ $`N|{\displaystyle \underset{m=1}{\overset{N}{}}}a_i(m)a_j^{}([ml])`$ (69)
$`\times \mathrm{exp}\left(i2\pi k{\displaystyle \frac{m}{N}}\right)|^2,`$
and $`๐ช=(2\pi /L)(k,l)=\sqrt{\frac{2\pi }{\mathrm{}^2N}}(k,l)`$, with $`k`$, $`l`$ integer. In Eq. (69) $`[m+l]`$ is defined as being $`m+l`$ for $`1m+lN`$ and $`m+l\pm N`$ otherwise such that $`1|m+l\pm N|N`$. In the numerical calculation we replace the delta function in Eq. (67) by a sharply peaked gaussian $`\delta _\gamma (x)\mathrm{exp}[x^2/\gamma ^2]`$ with a broadening $`\gamma =0.64v/N`$ which is of the order of the level spacing. We then have
$$S(q,\omega )=\frac{1}{2\pi \mathrm{}^2N^2}e^{\frac{1}{2}\mathrm{}^2q^2}K(q,\omega ),$$
(70)
with
$$K(q,\omega )=\frac{_{ij}\delta _\gamma (\mathrm{}\omega +E_iE_j)Q_{ij}(k,l)}{_{ij}\delta _\gamma (\mathrm{}\omega +E_iE_j)}.$$
(71)
This function is suitable for a numerical investigation. We restrict ourselves to a white noise disorder distribution ($`\alpha =1`$). We calculate $`K(q,\omega )`$ for values of $`2k^2+l^225`$ and $`\mathrm{}\omega =\gamma n`$, with $`3n23`$, where the limits have been chosen such that $`L^1<q<\mathrm{}^1`$ and $`\mathrm{}\omega v`$ but $`\mathrm{}\omega `$ greater than the level spacing of the finite size system. The system sizes range from $`N=200`$ to $`N=2000`$, and the number of disorder realizations are 500 or 100 depending on the system size. All values of $`K(q,\omega )`$ were determined to an accuracy better than 1% in the disorder averaging.
For a fixed and small value of $`qL_\omega `$ (so that we are in the range of normal diffusion) and $`q,\omega 0`$ we expect from Eq. (11) that $`\mathrm{}\omega K(q,\omega )`$ scales as $`(\mathrm{}\omega /v)^{\frac{1}{2\nu }}(q\mathrm{})^{\frac{1}{\nu }}`$. The scaling hypothesis is illustrated in Fig. 3, where we plot $`A(qL_\omega )\mathrm{}\omega K(q,\omega )`$ for fixed ratios of $`(qL_\omega )^2(k^2+l^2)/n`$ as a function of $`n\sqrt{1+(k^2+l^2)^2/n^2}/N`$ on a log-log scale. Here each curve is multiplied by a constant factor $`A(qL_\omega )`$ \[different for each $`(k^2+l^2)/n`$ ratio\] to make the comparison of the different lines easier. Also the factor $`\sqrt{1+(k^2+l^2)^2/n^2}`$ multiplying $`n/N\omega `$ is used such that the curves line up horizontally. The fact that data calculated for different system sizes fall onto the same curve \[for a fixed ratio of $`(k^2+l^2)/n`$)\] indicates that the limits chosen above for $`k,l`$, and $`n`$ do avoid large finite size effects. On the log-log scale the different data sets fall onto straight lines and can be fitted by power-laws (solid lines in Fig. 3).
The localization exponent $`\nu `$ extracted from the slope of the lines in Fig. 3 is shown as a function of $`(k^2+l^2)/n(qL_\omega )^2`$ in Fig. 4. Within our error bars and for the $`qL_\omega `$ considered, $`\nu `$ is a constant. Its value $`\nu =2.33\pm 0.05`$ is in excellent agreement with previous finite-size scaling studies and strongly supports the scaling hypothesis Eq. (11). The fact that the lowest $`(k^2+l^2)/n`$ points seem to be moving upwards in Fig. 4 is an indication that there are still some finite-size effects for the low values of $`(k^2+l^2)`$. In contrast to previous numerical studies we are able to obtain information about the critical exponent $`\nu `$ from systems of finite size without doing finite-size scaling.
## VI Conclusion
We have presented a new analytical and numerical approach to the localization-delocalization transition in the LLL of the IQHE. By using the closed Lie algebra of the density operators in the LLL we are able to write the equation of motion for the densities in a closed form which can be solved formally. Using the solution of the equation of motion for the projected densities we can express the integrated spectral function $`๐ES(q,\omega ;E)`$ as the disorder averaged density of states of a dynamical system with a novel action. We show analytically that the self-consistent mean-field approximation of the integrated spectral function yields normal diffusion but it misses the critical scaling. However, it is encouraging to note that even at this level of approximation the longitudinal conductivity is in approximate agreement with previous numerical studies. Finally, using exact diagonalization, we are able to extract the localization critical exponent $`\nu `$ from the integrated spectral function by using the scaling hypothesis Eq. (11), without having to do finite-size scaling. We obtain $`\nu =2.33\pm 0.05`$ in excellent agreement with previous studies. We hope that in the future it will be possible to extend our approach beyond the self-consistent mean-field level and analytically extract information about the critical exponent $`\nu `$.
The authors would like to thank A. Zee, A.H. MacDonald, T. Brandes, K. Schรถnhammer, and J.T. Chalker for helpful discussions. V.M. is grateful to the Deutsche Forschungsgemeinschaft for financial support during his stay at Indiana University. This work was furthermore supported by the NSF Grant No. DMR-9714055 and the NSF Grant No. DMR-9820816. |
warning/0002/hep-th0002004.html | ar5iv | text | # Collapsing D-branes in one-parameter models and small/large radius duality
## Introduction
Recent work on โD-brane geometryโ has lead to renewed interest in the quantum analogue of the notion of โsizeโ. A necessary preliminary of analyses such as is the identification of those D-brane states which become massless at special points in the moduli space of a type II compactification on a Calabi-Yau manifold, which in the language of amounts to identifying the cycles which acquire zero quantum volume at such a point. Many basic questions in quantum geometry still await an answer, one of the most important among these being the central issue of marginal stability and its implications for the extension of mirror symmetry to the the D-brane sector of compactified string theory. Such an extension holds promise of providing a tool for understanding quantum corrections to the moduli space of type IIA BPS saturated D-branes. Progress along these lines should enable us to understand the tantalizing conjectures of and .
One of the obstacles to a detailed and reasonably general investigation of D-brane effects in $`N=2`$ string compactifications is the difficulty of performing computations of a basis of periods of the holomorphic 3-form throughout the complex structure moduli space of a given Calabi-Yau manifold. In , we took a few steps towards removing this obstacle, at least in the one-parameter case, by showing how the largely overlooked <sup>1</sup><sup>1</sup>1An example in which Meijer functions were used for performing the analytic continuation of periods can be found in . We thank Erik Zaslow for bringing this reference to our attention. but classical technique of Meijer functions can be used to give a systematic approach to the problem. In fact, this technique allows us to reduce most one-parameter models to four classes, each of which allows for universal expressions of a special set of periods introduced in . Determining the analytic continuation of periods for all such classes amounts to a complete solution of the problem โ given a one parameter model, all that remains to be done is to substitute in these expressions for the specific values of the hypergeometric parameters. In we made use of this approach in order to undertake a systematic study of quantum volumes in one parameter models and along a special sub-locus of a two-parameter example. Considerations of space prevented us from giving a complete discussion of all classes of one-parameter models. The present paper remedies this lack of completeness by carrying through a similar analysis for the last two classes of this hierarchy, which in a certain sense are the most degenerate situations. This allows us to bring further evidence that the phenomenon noticed in of collapsing 6-branes at the mirror of the conifold point is generic in one-parameter models, and not limited to the case of the quintic , where it was first observed.
The last part of the paper is concerned with a special example which exhibits some rather exotic features. This is a one-parameter family of Calabi-Yau complete intersections in seven-dimensional projective space, which belongs to the most โdegenerateโ family in our classification. As noticed a while ago , the moduli space of this model admits a $`\text{}_2`$ symmetry which identifies the small and large radius limits. This lead to suspicions that the model provides a Calabi-Yau example of small-large radius duality. This would give an example of a โT-dualโ string compactification with reduced (N=2) supersymmetry, with potentially interesting implications for phenomenology. Our knowledge of a basis of periods allows us to address some of the puzzles concerning this model. While doing so in Section 3, we will meet with some surprises. Indeed, we will be able to confirm the suspicions of , but in a rather unexpected way: while small-large radius duality is indeed an exact feature of the model, its realization involves a certain rotation in the space of states, as well as a (less surprising) rescaling of the correlation functions. This conclusion, which can be extracted from the direct computation of periods by a a careful consideration of branch cuts, has some interesting implications for the action of the symmetry on the D-brane states. In particular, the duality exchanges D2 and D4-branes in the mirror, type IIA compactification. In Section 5, we propose an explanation of this phenomenon by making use of the ideas of Strominger, Yau and Zaslow . We will argue that the model admits two $`T^3`$ fibrations, which are interchanged by our symmetry. The nontrivial action on D2/D4 branes appears as a consequence of the fact that the dimension of the holomorphic cycle wrapped by the mirror of a given type IIB D-brane depends on the position of the original special Lagrangian cycle with respect to the fibration: when changing the fibration, the interpretation of the mirror state is modified. The existence of this symmetry has other interesting implications for the D-brane physics of this model. In particular, there exists a two-dimensional space of D-brane states which vanish at the mirror of the โconifoldโ point (modulo issues of marginal stability). Such states can be interpreted as composites of D4 and D6 branes.
## 1 Quantum notions of โsizeโ
The problem of understanding the correct string-theoretic generalization of the notion of size was considered in (see also for a review). The best framework for addressing this issue is that of type II string compactifications on Calabi-Yau manifolds, which have the advantage of allowing for exact computations of stringy corrections while at the same time displaying nontrivial quantum effects. This problem can be approached by considering a type IIA compactification on a Calabi-Yau manifold $`X`$ and its dual, type IIB compactification on the mirror $`Y`$ of $`X`$. The quantum corrections to the notion of size appear in the vector multiplet moduli space, which corresponds to the Kahler moduli of the IIA compactification and to the complex structure moduli of its IIB dual. Since the latter does not suffer quantum corrections , one can use mirror symmetry in order to transport the results accessible on this side to the IIA compactification, thereby extracting exact information about the stringy corrections to the Kahler moduli space of $`X`$. Hence mirror symmetry identifies the quantum-corrected complexified Kahler moduli space of $`X`$ with the complex structure moduli space $``$ of $`Y`$.
The first question one encounters in this framework is that of introducing a physically meaningful parameterization of the corrected complexified Kahler moduli space, which will allow us to measure โquantum areasโ on $`X`$. In this paper, we will follow the proposal of , which consists of using the value of the complexified Kahler class dictated by the mirror map:
$$k(z)=(B+iJ)(z)=\frac{_{\gamma _1}\mathrm{\Omega }(z)}{_{\gamma _0}\mathrm{\Omega }(z)},$$
(1)
where $`z`$ is a coordinate<sup>2</sup><sup>2</sup>2We restrict to one-parameter models for simplicity. on $``$, $`\mathrm{\Omega }`$ is the holomorphic 3-form of $`Y`$ and $`\gamma _0`$, $`\gamma _1`$ are certain 3-cycles in $`Y`$ which can be identified in the manner discussed in . Hence (1) defines a specific class in $`H^2(X,\text{})`$ at each point $`z`$, which is identified as the correct quantum counterpart of the complexified Kahler class at that point. The imaginary part $`J`$ of this class defines the so-called โnonlinear sigma model measureโ on $``$. More precisely, writing:
$$k(z)=t(z)e,$$
(2)
where $`e`$ is the generator of $`H^2(X,\text{})`$ defines a special coordinate on $``$ (in the sense of special geometry). Then the nonlinear sigma model measure is defined by the imaginary part of $`t(z)`$.
An โintermediateโ parameterization of $``$ is given by the so-called โalgebraic coordinateโ, which is defined through:
$$k_{alg}=(b+is)(z)=\frac{1}{2\pi i}\mathrm{log}(\kappa z)e,$$
(3)
where $`\kappa `$ is a certain constant which is determined by the monomial -divisor mirror map of . Measuring distances with $`k_{alg}`$ amounts to using the semiclassical notion of size (which is, strictly speaking, only valid in the large radius limit of $`X`$) throughout the entire moduli space $``$.
An important point, first noticed in and discussed in full generality in is that the classical geometric relation:
$$\mathrm{vol}(\mathrm{\Sigma }_{2p})_{\mathrm{\Sigma }_{2p}}k^p$$
(4)
(with $`\mathrm{\Sigma }_{2p}`$ some 2p-cycle in $`X`$) does not admit a natural generalization to the quantum level. This follows by noticing that the most natural extension of the notion of volume to the quantum setting is to identify the โquantum volumeโ of $`\mathrm{\Sigma }_{2p}`$ with the mass of a $`D_{2p}`$ brane wrapping this cycle (divided by the associated D-brane tension). This can be computed via mirror symmetry techniques in the BPS case (when $`\mathrm{\Sigma }_{2p}`$ is a holomorphic cycle and hence the associated D-brane state is BPS), since the mass of the mirror state (a type IIB $`D3`$-brane wrapping a special Lagrangian 3-cycle $`C`$ mirror to $`\mathrm{\Sigma }_{2p}`$) is given by the exact formula:
$$m(C)=\frac{|_C\mathrm{\Omega }|}{|_C\overline{\mathrm{\Omega }}\mathrm{\Omega }|^{1/2}}=m(\mathrm{\Sigma }_{2g}).$$
(5)
The disagreement between quantum volumes measured in this way and those given by (4) is due to open string instanton corrections<sup>3</sup><sup>3</sup>3These are induced by open strings whose endpoints are constrained to lie in $`\mathrm{\Sigma }_{2p}`$. to the mass of the corresponding $`D_{2p}`$ brane . In fact, using the semiclassical relation (4) amounts to substituting the correct quantum Kahler class into the classical relation for volumesโa procedure somewhat akin to using the algebraic measure (3) instead of the correct, nonlinear sigma model measure.
An important question raised by these considerations is to what extent this notion of quantum volume behaves like its geometric counterpart. Since the definition discussed above includes nontrivial quantum corrections from open string instantons, it is natural to expect that the two quantities will diverge as we move away from the large radius limit of $`X`$ into regions of the moduli space where such corrections are important. In fact, instanton corrections are especially strong in the vicinity of conifold points, so one expects that the most pronounced difference will be manifest there. This suspicion is confirmed by the observation of that the quantum volume of IIA $`D2`$ and $`D4`$ branes on the quintic remains nonzero at the mirror of the conifold point, while the quantum volume of a $`D6`$-brane vanishes. In , we presented evidence that this is a widespread phenomenon in Calabi-Yau compactifications, and not a peculiarity of the quintic. However, the analysis of was limited to only two of the four hypergeometric families of one parameter models. The purpose of next three sections is to complete this argument, by showing that the same behaviour occurs in the remaining families, thus providing more evidence that this is a generic feature of one-parameter compactifications.
Most of the results of the next three sections are of a somewhat technical nature and represent a direct extension of the work of . The reader mainly interested in the discussion of Calabi-Yau small-large radius duality can proceed directly to Section 5.
## 2 Universal results for one-parameter models
This section reviews and completes some results obtained in , which will be used intensively below. These rest on the theory of Meijer functions , a brief account of which can be found in .
### 2.1 Review of large radius results
Let us start by summarizing some material presented in . Following the discussion of that paper, we focus on one-parameter models whose hypergeometric symbol has the form $`\left(\begin{array}{c}\alpha _1,\alpha _2,\alpha _3,\alpha _4\\ 1,1,1\end{array}\right)`$ with $`\alpha _j`$ some rational numbers contained in the interval $`[0,1]`$. In this case, the associated Picard-Fuchs equation has the hypergeometric form:
$$\left[\delta ^4z(\delta +\alpha _1)(\delta +\alpha _2)(\delta +\alpha _3)(\delta +\alpha _4)\right]u=0$$
(6)
(where $`\delta :=z\frac{d}{dz}`$). By using the theory of Meijer functions, it was shown in that an especially convenient basis of periods (called Meijer periods) is given by the integral representations:
$$U_j(z)=\frac{1}{2\pi i}_\gamma ๐s\varphi _j(s),$$
(7)
where:
$`\varphi _j\left(s\right):={\displaystyle \frac{1}{_{i=1\mathrm{}4}\mathrm{\Gamma }\left(\alpha _i\right)}}{\displaystyle \frac{\mathrm{\Gamma }\left(s\right)^{j+1}_{i=1\mathrm{}4}\mathrm{\Gamma }\left(s+\alpha _i\right)}{\mathrm{\Gamma }\left(s+1\right)^{3j}}}\left(\left(1\right)^{j+1}z\right)^s.`$ (8)
In these expressions, the contour $`\gamma `$ is chosen as shown in Figure 1.
Figure 1. The defining contour for the Meijer periods.
The expansions of these periods in the large and small radius regions of the moduli space follow by closing the contour to the right or left, which is allowed for $`|z|<1`$ and $`|z|>1`$ respectively. The expansions for $`|z|<1`$ were computed in and are given by the universal expression:
$$U_j(z)=\frac{(1)^j}{j!}\underset{n=0}{\overset{\mathrm{}}{}}\frac{(\alpha _1)_n(\alpha _2)_n(\alpha _3)_n(\alpha _4)_n}{n!^4}\nu _j(n,z)z^n,$$
(9)
where:
$`\nu _0`$ $`=`$ $`1`$
$`\nu _1(n,z)=g_1^{}(n,z)`$ $`=`$ $`\eta _1\left(n\right)+\mathrm{log}\left(z\right)`$
$`\nu _2(n,z)=g_2^{\prime \prime }(n,z)+\left[g_2^{}(n,z)\right]^2]`$ $`=`$ $`\eta _2^{}\left(n\right)+\left(\eta _2\left(n\right)+\mathrm{log}\left(z\right)\right)^2`$ (10)
$`\nu _3(n,z)=g_3^{\prime \prime \prime }(n,z)+3g_3^{\prime \prime }(n,z)g_3^{}(n,z)+g_3^{}(n,z)^3`$ $`=`$ $`\eta _3^{\prime \prime }\left(n\right)+3\eta _3^{}\left(n\right)\left(\eta _3\left(n\right)+\mathrm{log}z\right)+\left(\eta _3\left(n\right)+\mathrm{log}z\right)^3,`$
with:
$$\eta _j^{\left(i\right)}\left(n\right)=\underset{k=1}{\overset{4}{}}\psi ^{\left(i\right)}\left(n+\alpha _k\right)\left(3j\right)\psi ^{\left(i\right)}\left(n+1\right)\left(1\right)^i\left(j+1\right)\left[\psi ^{\left(i\right)}\left(1\right)+i!\underset{l=1}{\overset{n}{}}\frac{1}{l^{i+1}}\right],$$
(11)
for $`i=0,1,2`$. In , we also computed the monodromy matrix of the Meijer periods about the large complex structure point $`z=0`$, with the result:
$$T\left[0\right]=\left[\begin{array}{cccc}1& 0& 0& 0\\ 2i\pi & 1& 0& 0\\ 4\pi ^2& 2i\pi & 1& 0\\ 0& 0& 2i\pi & 1\end{array}\right].$$
(12)
### 2.2 The special coordinate on the moduli space
For later use, let us derive a universal expression for the special coordinate $`t`$ on the moduli space. As explained in , this is given by a certain ratio of a linear combination of $`\mathrm{log}^0`$ and $`\mathrm{log}^1`$ periods to a $`\mathrm{log}^0`$ period (the latter is, of course, determined up to a global factor). The correct linear combination appearing in the numerator is fixed by the requirement that the asymptotic form of $`t`$ in the large complex structure limit be given by:
$$t_{as}=\frac{1}{2\pi i}\mathrm{log}w,$$
(13)
where $`w=\kappa z`$, with $`\kappa =e^{_{k=1}^4\psi (\alpha _k)4\psi (1)}`$, is a coordinate on the moduli space determined by the monomial-divisor mirror map of . The asymptotic form of the Meijer periods at large complex structure can be easily extracted from the expansions in terms of $`w`$ computed in . Indeed, it was shown there that (9) can be rewritten as:
$$U_j(w)=\underset{s=0}{\overset{j}{}}\stackrel{~}{q}_{sj}(w)(\mathrm{log}w)^s,$$
(14)
where:
$$\stackrel{~}{q}_{sj}(w):=\frac{(1)^j}{j!}\underset{n=0}{\overset{\mathrm{}}{}}\frac{(\alpha _1)_n(\alpha _2)_n(\alpha _3)_n(\alpha _4)_n}{n!^4}\stackrel{~}{v}_{sj}(n)\left(\frac{w}{\kappa }\right)^n,$$
(15)
with $`\stackrel{~}{v}_{sj}(n)`$ some quantities whose explicit form is listed in Subsection 4.2.1 of . Since the matrix $`\stackrel{~}{q}(0):=(q_{sj}(0))_{s,j=\mathrm{0..3}}`$ has a finite limit at $`w=0`$,
$$\stackrel{~}{q}\left(0\right):=\left[\begin{array}{cccc}1& 0& \frac{1}{2}\left(\eta _2^{}\left(0\right)\pi ^2\right)& \frac{1}{6}\eta _3^{\prime \prime }\left(0\right)\\ 0& 1& i\pi & \frac{1}{2}\eta _3^{}\left(0\right)\\ 0& 0& \frac{1}{2}& 0\\ 0& 0& 0& \frac{1}{6}\end{array}\right],$$
(16)
it follows that the leading terms in the large radius expansions of the Meijer periods are:
$$U_j^{as}(w)=\underset{s=0}{\overset{j}{}}\stackrel{~}{q}_{sj}(0)(\mathrm{log}w)^s.$$
(17)
In particular, we obtain:
$$U_0^{as}(w)=1,U_1^{as}(w)=\mathrm{log}w$$
(18)
and since the periods $`U_j`$ are adapted to the monodromy weight filtration of the model we immediately deduce that the special coordinate has the simple universal form:
$$t=\frac{1}{2\pi i}\frac{U_1}{U_0}.$$
(19)
Substituting expansion (9) in this formula leads to a general expression for the special coordinate in the large radius region $`|z|<1`$ (which can be used, in particular, to extract universal expressions for the Gromov-Witten invariants of this class of models as functions of the parameters $`\alpha _k`$). On the other hand, the analytic continuations of $`U_0`$ and $`U_1`$ allow us to compute $`t`$ as a function of $`z`$ (or $`w`$) throughout the moduli space.
### 2.3 The hypergeometric hierarchy
As discussed in , the nature of the small radius expansions of the Meijer periods, and hence the nature of the small radius point of the model, depend on the relative values of the parameters $`\alpha _i`$. From an abstract point of view, this leads to a hierarchy of models characterized (up to permutations of $`\alpha _i`$) by one of the conditions:
$`(0)`$ all $`\alpha _i`$ are distinct
$`(1)`$ three of the parameters $`\alpha _i`$ are distinct
$`(2)`$ $`\alpha _1=\alpha _2`$ and $`\alpha _3=\alpha _4`$ but $`\alpha _1\alpha _3`$
$`(3)`$ $`\alpha _1=\alpha _3=\alpha _3=\alpha _4`$.
$`(4)`$ $`\alpha _1=\alpha _2=\alpha _3\alpha _4`$
Only levels $`(0),(1),(2)`$ and $`(3)`$ of this hierarchy are realized through one-parameter complete intersections in projective spaces, as well as through many one-parameter complete intersections in weighted projective spaces and more general toric varieties (see for a discussion of toric geometry). Level $`(4)`$ does not seem to be realized<sup>4</sup><sup>4</sup>4This follows from the results of . through compact one-parameter complete intersections in toric varieties, though it could be realized through more general constructions. Since we are mostly interested in the toric case, we will limit ourselves to the families (0), (1), (2) and (3). A few examples of models belonging to these classes are listed in Table 1.
$`Family`$ $`Model`$ $`(\alpha _1,\alpha _2,\alpha _3,\alpha _4)`$ $`0`$ $`\text{}^4[5]`$ $`(1/5,2/5,3/5,4/5)`$ $`0`$ $`\text{๐}\text{}^{2,1,1,1,1}[6]`$ $`(1/3,2/3,1/6,5/6)`$ $`0`$ $`\text{๐}\text{}^{4,1,1,1,1}[8]`$ $`(1/8,3/8,5/8,7/8)`$ $`0`$ $`\text{๐}\text{}^{5,2,1,1,1}[10]`$ $`(1/10,3/10,7/10,9/10)`$ $`0`$ $`\text{๐}\text{}^{2,1,1,1,1,1}[3,4]`$ $`(1/3,2/3,1/4,3/4)`$ $`0`$ $`\text{๐}\text{}^{3,2,2,1,1,1}[4,6]`$ $`(1/6,1/4,3/4,5/6)`$ $`1`$ $`\text{}^5[2,4]`$ $`(1/2,1/2,1/4,3/4)`$ $`1`$ $`\text{}^6[2,2,3]`$ $`(1/2,1/2,1/3,2/3)`$ $`1`$ $`\text{๐}\text{}^{3,1,1,1,1,1}[2,6]`$ $`(1/2,1/2,1/6,5/6)`$ $`2`$ $`\text{}^5[3,3]`$ $`(1/3,1/3,2/3,2/3)`$ $`2`$ $`\text{๐}\text{}^{2,2,1,1,1,1}[4,4]`$ $`(1/4,1/4,3/4,3/4)`$ $`2`$ $`\text{๐}\text{}^{3,3,2,2,1,1}`$ $`(1/6,1/6,5/6,5/6)`$ $`3`$ $`\text{}^7[2,2,2,2]`$ $`(1/2,1/2,1/2,1/2)`$
Table 1. Some examples of models belonging to various hypergeometric families.
In , we studied only the families $`(0)`$ and $`(1)`$. Here we consider the more degenerate families $`(2)`$ and $`(3)`$. As we show below, these models can also be approached efficiently by the general methods developed in . The highly degenerate family $`(3)`$ displays some surprising, which we discuss in detail in Section 5.
### 2.4 The choice of branch-cuts
Let us clarify the choice of branch-cuts used in the present paper and implicitly in . Our convention is that we start from the large complex structure region $`|z|<1`$ and perform the analytic continuation through the sector $`\mathrm{arg}(z)(\pi ,0)`$, i.e. through the lower half of the unit circle in the complex plane (see Figure 2). Moreover, we will pick the branch-cut of all periods to lie along the negative real axis $`(\mathrm{},0)`$. With this convention, expressions such as $`\mathrm{log}(z)`$, $`\mathrm{log}(1/z)`$ and $`\mathrm{log}(z)`$ are always understood to have the cut on the negative real axis, so that we can write:
$$\mathrm{log}(z)=\mathrm{log}(z)+i\pi ,\mathrm{log}(1/z)=\mathrm{log}(z)i\pi .$$
(20)
A similar convention is used for power functions with non-integral exponents. In particular, we have $`(z)^\alpha =z^\alpha e^{i\pi \alpha }`$ for any real constant $`\alpha `$. For the โgenericโ model considered in the branch-cut along $`(\mathrm{},0)`$ suffices for all periods. For the other families (and in particular for all models discussed in the present paper), the situation is slightly different since in these cases the analytic continuation of the fundamental period $`U_0`$ displays logarithmic behaviour in the region $`|z|>1`$, even though it is regular in the unit disk, $`|z|<1`$. This requires that we enlarge the associated branch-cut in a way consistent with this behaviour, and we shall do so by adding the upper half of the unit circle to the common cut along the negative real axis.
Figure 2. Our choice of branch-cuts for the analytic continuation of periods. The upper half of the unit circle is added only for the fundamental period $`U_0`$, in all cases when this period displays logarithmic behaviour in the region $`|z|>1`$.
## 3 The family $`\alpha _1=\alpha _2`$$`\alpha _3=\alpha _4`$
Consider first the family $`(2)`$, which corresponds to the hypergeometric symbol $`\left(\begin{array}{c}\alpha ,\alpha ,\beta ,\beta \\ 1,1,1\end{array}\right)`$, i.e. to the parameters $`\alpha _1=\alpha _2:=\alpha `$, $`\alpha _3=\alpha _4:=\beta `$ with $`\alpha \beta `$, where we take $`0<\alpha ,\beta <1`$.
### 3.1 The Meijer periods
The expansion of the Meijer periods for $`|z|<1`$ follows by substituting our particular values for $`\alpha _i`$ in the general formula (9). The expansions for $`|z|>1`$ follow by closing the contour to the right, which gives contributions from the B-type poles:
$`(B_1)`$ $`s=n\alpha `$
$`(B_2)`$ $`s=n\beta `$
(with $`n`$ a nonnegative integer). Noticing that all such poles are double, a straightforward residue computation yields:
$`U_j\left(z\right)=\left({\displaystyle \frac{\mathrm{sin}\pi \alpha }{\pi }}\right)^{3j}\left((1)^{j+1}z\right)^\alpha {\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\mathrm{\Gamma }\left(n+\alpha \right)^4\mathrm{\Gamma }\left(n+\beta \alpha \right)^2}{\mathrm{\Gamma }\left(\alpha \right)^2\mathrm{\Gamma }\left(\beta \right)^2n!^2}}z^n\times `$
$`\left[2\psi \left(1\right)+2\psi \left(n+\beta \alpha \right)\left(j+1\right)\psi \left(n+\alpha \right)\left(3j\right)\psi \left(n\alpha +1\right)+2{\displaystyle \underset{k=1}{\overset{n}{}}}{\displaystyle \frac{1}{k}}+\mathrm{log}\left(\left(1\right)^{j+1}z\right)\right]`$ (21)
$`+(\alpha \beta ).`$
### 3.2 Meijer monodromies
The monodromy of the Meijer basis about $`z=0`$ follows by applying the results reviewed above, while the monodromy about $`z=\mathrm{}`$ can be computed by making use of the general techniques developed in . Following that procedure, we first determine the canonical and Jordan forms of the matrix $`R[\mathrm{}]`$:
$`R_{can}\left[\mathrm{}\right]=\left[\begin{array}{cccc}0& 1& 0& 0\\ 0& 0& 1& 0\\ 0& 0& 0& 1\\ \alpha ^2\beta ^2& 2\alpha ^2\beta 2\alpha \beta ^2& \alpha ^2+4\alpha \beta +\beta ^2& 2\alpha 2\beta \end{array}\right],R_J\left[\mathrm{}\right]=\left[\begin{array}{cccc}\beta & 1& 0& 0\\ 0& \beta & 0& 0\\ 0& 0& \alpha & 1\\ 0& 0& 0& \alpha \end{array}\right].`$ (30)
The relation $`R_{can}[\mathrm{}]=PR_J[\mathrm{}]P^1`$ allows us to determine a choice for the transition matrix $`P`$ from a Jordan basis to the canonical basis:
$`P=\left[\begin{array}{cccc}\frac{\alpha ^2\beta }{\alpha ^22\alpha \beta +\beta ^2}& \frac{\left(\alpha 3\beta \right)\alpha ^2}{\alpha ^33\alpha ^2\beta +3\alpha \beta ^2\beta ^3}& \frac{\alpha \beta ^2}{\alpha ^22\alpha \beta +\beta ^2}& \frac{\left(3\alpha \beta \right)\beta ^2}{\alpha ^33\alpha ^2\beta +3\alpha \beta ^2\beta ^3}\\ \frac{\alpha ^2\beta ^2}{\alpha ^22\alpha \beta +\beta ^2}& 2\frac{\alpha ^2\beta ^2}{\alpha ^33\alpha ^2\beta +3\alpha \beta ^2\beta ^3}& \frac{\alpha ^2\beta ^2}{\alpha ^22\alpha \beta +\beta ^2}& 2\frac{\alpha ^2\beta ^2}{\alpha ^33\alpha ^2\beta +3\alpha \beta ^2\beta ^3}\\ \frac{\beta ^3\alpha ^2}{\alpha ^22\alpha \beta +\beta ^2}& \frac{\alpha ^2\beta ^2\left(\alpha +\beta \right)}{\alpha ^33\alpha ^2\beta +3\alpha \beta ^2\beta ^3}& \frac{\alpha ^4\beta ^2}{\alpha ^22\alpha \beta +\beta ^2}& 2\frac{\alpha ^3\beta ^3}{\alpha ^33\alpha ^2\beta +3\alpha \beta ^2\beta ^3}\end{array}\right].`$ (34)
In the present case, the singular content of the periods around $`z=\mathrm{}`$ can be extracted by writing $`U^t(z)=Z(z)q(z)`$, where $`Z\left(z\right)=\left[\begin{array}{cccc}z^\alpha & z^\alpha \mathrm{log}z& z^\beta & z^\beta \mathrm{log}z\end{array}\right]`$ and $`q(z)=(q_{sj}(z))_{s,j=\mathrm{0..3}}`$, with:
$`q_{0j}\left(z\right)=(\delta _{j,odd}+\delta _{j,even}e^{i\pi \alpha })\left({\displaystyle \frac{\mathrm{sin}\pi \alpha }{\pi }}\right)^{3j}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\mathrm{\Gamma }\left(n+\alpha \right)^4\mathrm{\Gamma }\left(n+\beta \alpha \right)^2}{\mathrm{\Gamma }\left(\alpha \right)^2\mathrm{\Gamma }\left(\beta \right)^2n!^2}}z^n\times `$
$`[2\psi \left(1\right)+2\psi (n+\beta \alpha )(j+1)\psi (n+\alpha )(3j)\psi (n\alpha +1)+2{\displaystyle \underset{k=1}{\overset{n}{}}}{\displaystyle \frac{1}{k}}+i\pi \delta _{j,even})]`$
$`q_{1j}\left(z\right)=\left(\delta _{j,odd}+\delta _{j,even}e^{i\pi \alpha }\right)\left({\displaystyle \frac{\mathrm{sin}\pi \alpha }{\pi }}\right)^{3j}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\mathrm{\Gamma }\left(n+\alpha \right)^4\mathrm{\Gamma }\left(n+\beta \alpha \right)^2}{\mathrm{\Gamma }\left(\alpha \right)^2\mathrm{\Gamma }\left(\beta \right)^2n!^2}}z^n,`$
and $`q_{2j}(z)=q_{0j}(z)|_{\alpha \beta }`$, $`q_{3j}(z)=q_{1j}(z)|_{\alpha \beta }`$.
On the other hand, the matrix $`z^{R_J[\mathrm{}]}`$ has the simple form:
$$z^{R_J\left[\mathrm{}\right]}=\left[\begin{array}{cccc}z^\beta & \mathrm{ln}\left(z\right)z^\beta & 0& 0\\ 0& z^\beta & 0& 0\\ 0& 0& z^\alpha & \mathrm{ln}\left(z\right)z^\alpha \\ 0& 0& 0& z^\alpha \end{array}\right].$$
(35)
This allows us to find the matrix $`q_J(z)`$ which satisfies $`U_J^t(\mathrm{})=Z(z)q_J(z)`$:
$$q_J\left(z\right)=\left[\begin{array}{cccc}0& 0& S_{1,3}\left(z\right)& S_{1,4}\left(z\right)\\ 0& 0& 0& S_{1,3}\left(z\right)\\ S_{1,1}\left(z\right)& S_{1,2}\left(z\right)& 0& 0\\ 0& S_{1,1}\left(z\right)& 0& 0\end{array}\right].$$
(36)
In this expression, $`S_{ij}(z)`$ are the entries of the matrix $`S(z)`$ which defines the nilpotent orbit of the fundamental system $`\mathrm{\Phi }_J(z)`$ associated with the Jordan basis $`U_J(z)`$:
$$\mathrm{\Phi }_J(z)=S(z)z^{R_J}.$$
(37)
Since $`S(\mathrm{})=P`$, we obtain:
$$q_J\left(\mathrm{}\right)=\left[\begin{array}{cccc}0& 0& \frac{\alpha \beta ^2}{\left(\beta +\alpha \right)^2}& \frac{\left(3\alpha \beta \right)\beta ^2}{\left(\beta +\alpha \right)^3}\\ 0& 0& 0& \frac{\alpha \beta ^2}{\left(\beta +\alpha \right)^2}\\ \frac{\alpha ^2\beta }{\left(\beta +\alpha \right)^2}& \frac{\left(\alpha 3\beta \right)\alpha ^2}{\left(\beta +\alpha \right)^3}& 0& 0\\ 0& \frac{\alpha ^2\beta }{\left(\beta +\alpha \right)^2}& 0& 0\end{array}\right].$$
(38)
We can now compute the matrix $`M=q(\mathrm{})^tq_J(\mathrm{})^t`$ and the Meijer monodromy about the small radius point:
$$T[\mathrm{}]=MT_J[\mathrm{}]M^1,$$
(39)
where:
$$T_J\left[\mathrm{}\right]=e^{2\pi iR_J\left[\mathrm{}\right]^t}=\left[\begin{array}{cccc}e^{2i\pi \alpha }& 0& 0& 0\\ 2i\pi e^{2i\pi \alpha }& e^{2i\pi \alpha }& 0& 0\\ 0& 0& e^{2i\pi \beta }& 0\\ 0& 0& 2i\pi e^{2i\pi \beta }& e^{2i\pi \beta }\end{array}\right].$$
(40)
### 3.3 The model $`\text{}^5[3,3]`$
The mirror $`Y`$ of this model can be realized as an orbifold <sup>5</sup><sup>5</sup>5We refer the reader to for details about the orbifold action in this case. of a complete intersection $`p_1=p_2=0`$ of two cubics in $`\text{}^5`$:
$`p_1`$ $`=`$ $`x_1^3+x_2^3+x_3^33\psi x_4x_5x_6`$
$`p_2`$ $`=`$ $`x_4^3+x_5^3+x_6^33\psi x_1x_2x_3.`$
The fundamental period and special coordinate in this example are discussed in . Reference also discusses the counting of holomorphic curves for this model.
In this example, the hypergeometric coordinate is given by $`z=\frac{1}{\psi ^6}`$. The matrices $`R_{can}[\mathrm{}],R_J[\mathrm{}]`$ and a choice for the matrix $`P`$ are given in Appendix A. This data allows us to compute the Meijer monodromies:
$$T\left[0\right]=\left[\begin{array}{cccc}1& 0& 0& 0\\ 2i\pi & 1& 0& 0\\ 4\pi ^2& 2i\pi & 1& 0\\ 0& 0& 2i\pi & 1\end{array}\right],T\left[\mathrm{}\right]=\left[\begin{array}{cccc}1& 0& 0& 0\\ 5& 3\frac{i}{\pi }& 9/4\pi ^2& \frac{9}{8}\frac{i}{\pi ^3}\\ 10& 15/2\frac{i}{\pi }& \frac{27}{4}\pi ^2& 9/2\frac{i}{\pi ^3}\\ 8& 9\frac{i}{\pi }& \frac{27}{4}\pi ^2& \frac{27}{4}\frac{i}{\pi ^3}\end{array}\right]$$
(41)
and $`T[1]=T[0]^1T[\mathrm{}]`$. These monodromy matrices satisfy:
$$(T[0]I)^4=0,(T[1]I)^2=0,(T[\mathrm{}]^3I)^2=0.$$
A set of periods associated with with a basis of a full sublattice of the integral lattice $`H_3(Y,\text{})`$ (up to a common factor) is given by:
$$U_E\left(z\right)=EU\left(z\right),\text{with }E=\left[\begin{array}{cccc}1& 0& 0& 0\\ 5& 3\frac{i}{\pi }& 9/4\pi ^2& \frac{9}{8}\frac{i}{\pi ^3}\\ 10& 15/2\frac{i}{\pi }& \frac{27}{4}\pi ^2& 9/2\frac{i}{\pi ^3}\\ 8& 9\frac{i}{\pi }& \frac{27}{4}\pi ^2& \frac{27}{4}\frac{i}{\pi ^3}\end{array}\right].$$
(42)
It is also easy to check that the period $`U_v(z)=\frac{3}{\pi ^3}\left[\frac{3}{8}U_3\pi ^2U_1\right]`$ vanishes at $`z=1`$. This period is weakly integral since:
$$U_v\left(z\right)=i[3,3,2,1]U_E\left(z\right).$$
(43)
The relation $`U_v(1)=0`$ is equivalent with an arithmetic identity which we write down in the Appendix.
In this case, the constant $`\kappa =e^{2\psi (\alpha )+2\psi (\beta )4\psi (1)}=\frac{1}{729}=3^6`$ and the imaginary part of the algebraic coordinate on the moduli space is $`s=\frac{1}{2\pi }\mathrm{log}(\kappa |z|)=\frac{3}{\pi }\mathrm{log}(\frac{3}{|z|})`$. Figure 3 displays the values of $`|U_v|`$ versus $`s`$. The point $`z=1`$ corresponds to $`s=\frac{3\mathrm{log}3}{\pi }1.049`$. For comparison, we also display the absolute values of the weakly integral period $`\frac{9}{4\pi ^2}U_2`$ and of the special coordinate $`t`$. Figure 4 shows the absolute value of the special coordinate $`t=\frac{1}{2\pi i}\frac{U_1}{U_0}`$ as a function of $`s`$, for $`s[6,2]`$. The asymptotic form of $`t`$ in the small radius limit $`z\mathrm{}`$ can be easily computed from the small radius expansions given above:
$`t{\displaystyle \frac{2\left(i\sqrt{3}\mathrm{log}\left(z\right)+\left(1+i\right)\pi \right)}{\sqrt{3}\left(1+i\sqrt{3}\right)\mathrm{log}z}}+O\left(\left(\mathrm{log}z\right)^2\right)={\displaystyle \frac{9\mathrm{log}36\pi s3i\sqrt{3}\mathrm{log}3+i\pi \left(\sqrt{3}s2\right)}{6\left(3\mathrm{log}3+\pi s\right)}}+O\left(s^2\right).`$
In particular, the value of $`t`$ in the limit $`z=\mathrm{}`$ is:
$$t_{lim}=\frac{1}{2}+i\frac{\sqrt{3}}{6}|t_{lim}|=\frac{1}{\sqrt{3}}.577.$$
(44)
$`\begin{array}{cc}\begin{array}{c}\text{}\end{array}& \begin{array}{c}\text{}\end{array}\\ \begin{array}{c}\\ \text{Figure 3. }\text{Graph of }|U_v|,\frac{9}{4\pi ^2}|U_2|\text{ and }|t|\text{ versus the imaginary }\\ \text{part }s\text{ of the algebraic coordinate for }s[0,2]\text{. The }\\ \text{point }z=1\text{ corresponds to }s=\frac{3\mathrm{l}\mathrm{o}\mathrm{g}3}{\pi }1.049\text{}\end{array}& \begin{array}{c}\\ \\ \text{Figure 4. }\text{Graph of }|t|\text{ versus }s\\ \text{for }s[6,2]\text{.}\end{array}\end{array}`$
## 4 The family $`\alpha _1=\alpha _2=\alpha _3=\alpha _4`$
In this section we consider the family $`(3)`$, associated with the hypergeometric symbol $`\left(\begin{array}{c}\alpha ,\alpha ,\alpha ,\alpha \\ 1,1,1\end{array}\right)`$, i.e. to the parameters $`\alpha _1=\alpha _2=\alpha _3=\alpha _4=\alpha `$, where we take $`0<\alpha <1`$ for simplicity. The hypergeometric equation of this family has the form:
$$\left[\delta ^4z(\delta +\alpha )^4\right]u=0.$$
(45)
### 4.1 The Meijer periods
The expansion of the Meijer periods for $`|z|<1`$ follows from the general results of Section 2, while the expansion for $`|z|>1`$ is obtained by closing the contour to the right and applying the residue theorem. This brings contributions from the quadruple poles $`s=n\alpha `$, with $`n`$ a nonnegative integer. In this case, the computation is rather similar to that leading to the large radius expansions (9), giving the result:
$`U_j\left(z\right)={\displaystyle \frac{1}{6}}\left({\displaystyle \frac{\mathrm{sin}\pi \alpha }{\pi }}\right)^{3j}\left(\left(1\right)^{j+1}z\right)^\alpha {\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}\left[{\displaystyle \frac{\left(\alpha \right)_n}{n!}}\right]^4\mu _j(n,z)z^n,`$ (46)
where the quantities $`\mu _j(n,z)`$ are defined through:
$$\mu _j(n,z)=\xi _j^{^{\prime \prime }}\left(n\alpha \right)+3\xi _j^{^{}}\left(n\alpha \right)\left(\xi _j\left(n\alpha \right)+i\pi \delta _{j,even}+\mathrm{log}z\right)+\left(\xi _j\left(n\alpha \right)+i\pi \delta _{j,even}+\mathrm{log}z\right)^3,$$
(47)
with:
$$\xi _j^{\left(i\right)}\left(n\alpha \right)=4\left[\psi ^{\left(i\right)}\left(1\right)+i!\underset{k=1}{\overset{n}{}}\frac{1}{k^{i+1}}\right]\left(1\right)^i\left(j+1\right)\psi ^{\left(i\right)}\left(n+\alpha \right)\left(3j\right)\psi ^{\left(i\right)}\left(1n\alpha \right)$$
(48)
for $`i=0,1,2`$.
### 4.2 Meijer monodromies
The monodromies of the Meijer basis can be extracted by a procedure very similar to the one employed above. For the benefit of the reader interested in reproducing our computations, let us mention that in this case the correct row vector needed for extracting the singular behaviour around $`z=\mathrm{}`$ is:
$$Z=\left[\begin{array}{cccc}z^\alpha & z^\alpha \mathrm{log}\left(z\right)& z^\alpha \mathrm{log}\left(z\right)^2& z^\alpha \mathrm{log}\left(z\right)^3\end{array}\right]$$
(49)
and that writing $`U^t(z)=Z(z)q(z)`$ produces a regular matrix function $`q(z)`$ whose value $`q(\mathrm{})`$ at the point of interest has entries:
$$q_{ij}\left(\mathrm{}\right)=\frac{1}{6}\left(\frac{\mathrm{sin}\pi \alpha }{\pi }\right)^{3j}\left(\delta _{j,odd}+\delta _{j,even}e^{i\pi \alpha }\right)v_{ij}\left(\mathrm{}\right),$$
(50)
where:
$`v_{0j}\left(\mathrm{}\right)=\xi _j^{^{\prime \prime }}\left(\alpha \right)+3\xi _j^{^{}}\left(\alpha \right)\left(\xi _j\left(\alpha \right)+i\pi \delta _{j,even}\right)+\left(\xi _j\left(\alpha \right)+i\pi \delta _{j,odd}\right)^3,v_{3j}\left(\mathrm{}\right)=1.`$
$`v_{1j}\left(\mathrm{}\right)=3\xi _j^{^{}}(\alpha )+3(\xi _j(\alpha )+i\pi \delta _{j,even})^2,v_{2j}\left(\mathrm{}\right)=3(\xi _j(\alpha )+i\pi \delta _{j,odd}).`$
(Here $`\xi _j^{(i)}(\alpha )`$ are obtained from (48) by setting $`n=0`$.)
The canonical and Jordan forms of the matrix $`R[\mathrm{}]`$ are:
$$R_{can}\left[\mathrm{}\right]=\left[\begin{array}{cccc}0& 1& 0& 0\\ 0& 0& 1& 0\\ 0& 0& 0& 1\\ \alpha ^4& 4\alpha ^3& 6\alpha ^2& 4\alpha \end{array}\right],R_J\left[\mathrm{}\right]=\left[\begin{array}{cccc}\alpha & 1& 0& 0\\ 0& \alpha & 1& 0\\ 0& 0& \alpha & 1\\ 0& 0& 0& \alpha \end{array}\right]$$
(51)
while a choice for the matrix $`P`$ which defines a Jordan basis is:
$$P=\left[\begin{array}{cccc}\alpha ^3& \alpha ^2& \alpha & 1\\ \alpha ^4& 0& 0& 0\\ \alpha ^5& \alpha ^4& 0& 0\\ \alpha ^6& 2\alpha ^5& \alpha ^4& 0\end{array}\right].$$
(52)
The matrix $`q_J(z)`$ is expressed in terms of the nilpotent orbit $`S(z)`$ of $`\mathrm{\Phi }_J(z)`$ via:
$`q_J\left(z\right)=\left[\begin{array}{cccc}S_{1,1}\left(z\right)& S_{1,2}\left(z\right)& S_{1,3}\left(z\right)& S_{1,4}\left(z\right)\\ 0& S_{1,1}\left(z\right)& S_{1,2}\left(z\right)& S_{1,3}\left(z\right)\\ 0& 0& 1/2S_{1,1}\left(z\right)& 1/2S_{1,2}\left(z\right)\\ 0& 0& 0& 1/6S_{1,1}\left(z\right)\end{array}\right]\stackrel{\left(S\left(\mathrm{}\right)=P\right)}{}q_J\left(0\right)=\left[\begin{array}{cccc}\alpha ^3& \alpha ^2& \alpha & 1\\ 0& \alpha ^3& \alpha ^2& \alpha \\ 0& 0& 1/2\alpha ^3& 1/2\alpha ^2\\ 0& 0& 0& 1/6\alpha ^3\end{array}\right].`$ (61)
Finally, the Meijer monodromy about $`z=\mathrm{}`$ can be computed as $`T[\mathrm{}]=MT_J[\mathrm{}]M^1`$, where $`M=q(0)^tq_J(0)^t`$ and:
$$T_J\left[\mathrm{}\right]=\left[\begin{array}{cccc}e^{2\pi i\alpha }& 0& 0& 0\\ 2\pi ie^{2\pi i\alpha }& e^{2\pi i\alpha }& 0& 0\\ 2\pi ^2e^{2\pi i\alpha }& 2i\pi e^{2\pi i\alpha }& e^{2\pi i\alpha }& 0\\ \frac{4i\pi ^3}{3}e^{2\pi i\alpha }& 2\pi ^2e^{2\pi i\alpha }& 2\pi ie^{2\pi i\alpha }& e^{2\pi i\alpha }\end{array}\right].$$
The expression of the special coordinate (19) as a function of $`z`$ follows easily from the small and large radius expansions of the Meijer periods. For later reference, we write down the form of $`t(z)`$ in the region $`|z|>1`$:
$$t\left(z\right)=\frac{ie^{i\pi \alpha }}{2\mathrm{sin}\pi \alpha }\frac{_{n=0}^{\mathrm{}}\left[\frac{\left(\alpha \right)_n}{n!}\right]^4\mu _1(n,z)z^n}{_{n=0}^{\mathrm{}}\left[\frac{\left(\alpha \right)_n}{n!}\right]^4\mu _0(n,z)z^n}.$$
(62)
This allows us to extract the asymptotic form of $`t`$ for $`z\mathrm{}`$:
$$t_{as}=\frac{ie^{i\pi \alpha }}{2\mathrm{sin}\pi \alpha }\frac{\mu _1(0,z)}{\mu _0(0,z)}=\frac{ie^{i\pi \alpha }}{2\mathrm{sin}\pi \alpha }\left[1+\frac{3\left(\psi \left(\alpha \right)\psi \left(1\alpha \right)i\pi \right)}{\mathrm{log}z}\right]+O\left(\left(\mathrm{log}z\right)^2\right),$$
(63)
where we used the relations:
$`\xi _0\left(\alpha \right)=4\psi \left(1\right)\psi \left(\alpha \right)3\psi \left(1\alpha \right),\xi _1\left(\alpha \right)=4\psi \left(1\right)2\psi \left(\alpha \right)2\psi \left(1\alpha \right).`$
### 4.3 The model $`\text{}^7[2,2,2,2]`$
The mirror of this model is given by an orbifold $`Y`$ of the complete intersection $`\{p_1=p_2=p_3=p_4=0\}`$, where:
$`p_1`$ $`=`$ $`x_1^2+x_2^22\psi x_3x_4`$
$`p_2`$ $`=`$ $`x_3^2+x_4^22\psi x_5x_6`$
$`p_3`$ $`=`$ $`x_5^2+x_6^22\psi x_7x_8`$ (64)
$`p_4`$ $`=`$ $`x_7^2+x_8^22\psi x_1x_2.`$
The fundamental period of this example was determined in ( see also ), while the semiclassical structure of the Kรคhler moduli space was analyzed in detail in by making use of the linear sigma model technology of . Our techniques allow us to go further and perform a systematic analysis of all periods. In Section 5, we will use the results derived below in order to address certain puzzles about the small radius limit of this model.
In this example, $`\psi `$ is related to the hypergeometric coordinate through $`z=\psi ^8`$. The associated hypergeometric symbol is $`\left(\begin{array}{c}1/2,1/2,1/2,1/2\\ 1,1,1\end{array}\right)`$, so the model fits into the scheme discussed above for the particular value $`\alpha =1/2`$. The canonical and Jordan forms of the monodromy about $`z=\mathrm{}`$, as well as a choice for the matrix $`P`$ are given in Appendix A, while the Meijer monodromies are given by:
$$T\left[0\right]=\left[\begin{array}{cccc}1& 0& 0& 0\\ 2i\pi & 1& 0& 0\\ 4\pi ^2& 2i\pi & 1& 0\\ 0& 0& 2i\pi & 1\end{array}\right],T\left[\mathrm{}\right]=\left[\begin{array}{cccc}7& 4\frac{i}{\pi }& 4\pi ^2& 2\frac{i}{\pi ^3}\\ 2i\pi & 1& 0& 0\\ 4\pi ^2& 2i\pi & 1& 0\\ 0& 0& 2i\pi & 1\end{array}\right],$$
(65)
and $`T[1]=T[0]^1T[\mathrm{}]`$. These matrices satisfy:
$$(T[0]I)^4=0,(T[1]I)^2=0,(T[\mathrm{}]^2I)^4=0.$$
Note that the matrix $`T[\mathrm{}]`$ is not maximally unipotent.
Partial information about the integral structure is provided by a set of periods associated (up a common factor) with a basis of a full sublattice of $`H_3(Y,\text{})`$:
$$U_E\left(z\right)=EU\left(z\right),\text{with }E=\left[\begin{array}{cccc}1& 0& 0& 0\\ 7& 4\frac{i}{\pi }& 4\pi ^2& 2\frac{i}{\pi ^3}\\ 25& 16\frac{i}{\pi }& 20\pi ^2& 12\frac{i}{\pi ^3}\\ 63& 44\frac{i}{\pi }& 56\pi ^2& 38\frac{i}{\pi ^3}\end{array}\right].$$
(66)
In this case, one obtains two weakly integral periods vanishing at $`z=1`$:
$`U_{v1}`$ $`=`$ $`{\displaystyle \frac{2i}{\pi ^3}}\left[U_32\pi ^2U_1\right]=[5,6,4,1]U_E,`$ (67)
$`U_{v2}`$ $`=`$ $`{\displaystyle \frac{8}{\pi ^2}}\left[U_2+i\pi U_12\pi ^2U_0\right]=[15,11,5,1]U_E.`$
In the mirror picture, these correspond to a $`D6`$ and a $`D4`$-brane which become massless at $`z=1`$. In fact, any linear combination of these periods will also vanish there, so we can for example also consider the vanishing period $`U_{v1}+U_{v2}`$, which in the mirror picture also corresponds to a collapsing $`D6`$-brane. This situation will be discussed in more detail in Section 5.
In this example, the constant $`\kappa =e^{4(\psi (\alpha )\psi (1))}`$ has the value $`2^8=\frac{1}{256}`$. Figure 5 displays the absolute values of the special coordinate $`t`$ and of the weakly integral periods $`U_{v1},U_{v2}`$ as functions of the imaginary part $`s=\frac{4}{\pi }\mathrm{log}\frac{2}{|z|}`$ of the algebraic coordinate on the moduli space. In Figure 6 we plot the absolute value of the special coordinate $`t`$ as a function of $`s`$, including the region $`s<0`$ of the moduli space, which has no classical analogue. In this example, we have $`\psi (\alpha )=\psi (1\alpha )=\psi (1/2)`$, so the asymptotic form of $`t`$ for $`|z|\mathrm{}`$ is:
$$t_{as}=\frac{1}{2}+\frac{3i\pi }{2\mathrm{log}z}.$$
(68)
In particular, $`J=\mathrm{Im}(t)\frac{3}{2}\frac{\pi }{\mathrm{log}|z|}`$ remains nonnegative for $`|z|>>1`$, as pointed out <sup>6</sup><sup>6</sup>6The reader should note that the variable $`z`$ used in equation (37) of is the inverse of the hypergeometric coordinate $`z`$ used in the present paper. in .
$`\begin{array}{cc}\begin{array}{c}\text{}\end{array}& \begin{array}{c}\text{}\end{array}\\ \begin{array}{c}\\ \text{Figure 5. }\text{Graph of }|U_{v1}|,|U_{v2}|\text{ and }|t|\text{ versus the imaginary}\\ \text{part }s\text{ of the algebraic coordinate for }s[0,1.2]\text{. The }\\ \text{point }z=1\text{ corresponds to }s=\frac{4\mathrm{l}\mathrm{o}\mathrm{g}2}{\pi }0.882\text{.}\end{array}& \begin{array}{c}\\ \\ \text{Figure 6. }\text{Graph of }|t|\text{ versus }s\\ \text{for }s[4,1.2]\text{.}\end{array}\end{array}`$
## 5 Small/large radius duality
### 5.1 Basic considerations
The results we have obtained for the model $`\text{}^7[2,2,2,2]`$ apparently preclude us from interpreting $`z=\mathrm{}`$ as a large complex structure point. Indeed, the associated monodromy matrix is not maximally unipotent, but rather satisfies $`(T[\mathrm{}]^2I)^4=0`$. This behaviour is due to the factors of $`z^\alpha =z^{1/2}`$ in the expansions (46) of the periods for $`|z|>1`$. One may be tempted to interpret this result as showing that the limit $`z\mathrm{}`$ of the model does not admit a standard geometric (i.e. Calabi-Yau) description . However, the form of the monodromy about $`z=\mathrm{}`$ is tantalizing close to that of the monodromy about a large complex structure point, which is an indication that something more interesting may be going on.
Indeed, it was noticed in that the moduli space of this model admits a symmetry $`z1/z`$. This follows by replacing $`\psi `$ with its inverse and performing the change of coordinates:
$`\begin{array}{ccccc}x_1& =& y_1+iy_2,x_2& =& iy_1+y_2\\ x_3& =& y_7+iy_8,x_4& =& iy_7+y_8\\ x_5& =& y_5+iy_6,x_6& =& iy_5+y_6\\ x_7& =& y_3+iy_4,x_8& =& iy_3+y_4\end{array},`$ (73)
which preserves the form of the defining equations (4.3). Hence the manifolds $`Y_\psi `$ and $`Y_{\frac{1}{\psi }}`$ described by (4.3) for the parameters $`\psi `$ and $`\frac{1}{\psi }`$ are isomorphic, which implies that the nature of the points $`z=0`$ and $`z=\mathrm{}`$ is identical. Indeed, the isomorphism between $`Y_\psi `$ and $`Y_{\frac{1}{\psi }}`$ forces us to conclude that the points $`z=0`$ and $`z=\mathrm{}`$ are physically indistinguishable โ this is an exact statement in the full IIB string theory on $`Y`$, since its vector multiplet moduli space does not receive quantum corrections . This, however, seems to be at odds with the different behaviour of the Meijer periods in the two limits $`z=0`$ and $`z=\mathrm{}`$.
In order to clarify the situation, let us consider the effect of the change of variable $`zs:=1/z`$ on the hypergeometric equation (45). Under this operation, the equation is transformed into:
$$\left[s\delta ^4(\delta ^{}\alpha )^4\right]\stackrel{~}{u}(s)=0,$$
(74)
where $`\delta ^{}=s\frac{d}{ds}=z\frac{d}{dz}`$ and $`\stackrel{~}{u}(s):=u(1/s)`$. Thus (45) is not invariant under this symmetry. However, it is not hard to see that the form of (45) is preserved under the combined change of variable and function:
$`z`$ $`s:={\displaystyle \frac{1}{z}}`$
$`u`$ $`u^{}:=z^\alpha u,`$
i.e. $`u(z)z^\alpha u(1/z)`$. Since $`u(z)=_\gamma \mathrm{\Omega }(z)`$ is the period of the holomorphic 3-form $`\mathrm{\Omega }`$ on a 3-cycle $`\gamma H_3(Y,\text{})`$, it follows that the implementation of the symmetry $`z\frac{1}{z}`$ requires a rescaling of $`\mathrm{\Omega }`$:
$$\mathrm{\Omega }(z)z^\alpha \mathrm{\Omega }(1/z).$$
(76)
What, then, is the correct interpretation of the point $`z=\mathrm{}`$ ? The answer follows by recalling that the moduli space of the closed conformal field theory on $`Y`$ is built by considering marginal deformations, a process which is analytic in the deformation parameter $`z`$. This forces the periods to have different behavior in the regions $`|z|<1`$ and $`|z|>1`$. There is, however, a basic point to take into account: when performing marginal deformations one must specify a starting point ! In fact, one could as well choose this point to be $`z=\mathrm{}`$ and use the periods $`\stackrel{~}{U}_j(z)=U_j(1/z)`$ instead of $`U_j(z)`$. Therefore, the interpretation of $`z=0`$ and $`z=\mathrm{}`$ as โlargeโ and โsmallโ radius points is indeed conventional and can be reversed, even though the analytic continuations of the associated periods do not coincide. In fact, interchanging these points corresponds to starting on different branches of a double cover of the moduli space. This follows by noticing that, since $`Y_z`$ and $`Y_{\frac{1}{z}}`$ are isomorphic, the complex structure moduli space of $`Y`$ is not the copy of $`\text{}^1`$ parameterized by $`z`$, but rather its quotient $``$ via this identification. This quotient is again a $`\text{}^1`$, which can be parameterized, for example, by the variable:
$$x=\frac{2z}{z^2+1}.$$
(77)
The map $`zu`$ gives a double cover of $``$, branched over the points $`x=+1`$ and $`x=1`$, which are the images of $`z=1`$ and $`z=1`$, respectively. The unit circle $`|z|=1`$ is mapped into the region $`x[\mathrm{},1][1,\mathrm{}]`$, which represents a segment on the associated Riemann sphere (see Figure 7). The points $`z=0`$ and $`z=\mathrm{}`$ are both mapped into the point $`x=0`$. Since these points lie on different branches of our double cover, picking one of them as the large complex structure limit amounts to choosing a particular realization of the model.
Figure 7. The coordinate $`z`$ parameterizes a double cover of the complex structure moduli space of $`Y`$. The figure shows the topology of the restriction of this cover above the circle $`\mathrm{Im}(x)=0`$ on the Riemann sphere of $`x`$.
This situation is similar to the standard interpretation of T-duality for the conformal field theory on a circle. In that case, marginal deformations starting from a point $`R>R_0`$ build the continuation of the theory through the self-dual point $`R_0=\sqrt{\alpha ^{}}`$, into the region $`R<R_0`$. The duality $`R\stackrel{~}{R}=\alpha ^{}/R`$ identifies this continued theory with its form at a radius $`\stackrel{~}{R}>R_0`$, but this discrete identification is not captured by the marginal deformations. Just as in the case of T-duality, the global identification $`z1/z`$ in our model is โaccidentalโ in the sense that it is not captured by marginal deformations associated with the $`(c,c)`$ ring.
In order to make this more precise, let us compute the action of our symmetry on $`H^3(Y,\text{})`$. Consider acting with the transformation (5.1) on the Meijer periods:
$$U_j(z)U_j^{}(z)=z^{1/2}U_j(1/z).$$
(78)
Since (5.1) is a symmetry of (45), it follows that both $`(U_j)_{j=\mathrm{0..3}}`$ and $`(U_j^{})_{j=\mathrm{0..3}}`$ give a basis of solutions. Hence there must exist a constant matrix $`C=(c_{ij})_{i,j=\mathrm{0..3}}`$ such that $`U_i^{}(z)=c_{ij}U_j(z)`$. In fact, this conclusion is a bit too quick, since the functions $`U_j`$ are multi-valued, so we must be careful to take the branch-cuts into account. The correct statement is that such a relation must hold on every open and connected subset $`V`$ of the moduli space which does not intersect the cuts. In fact, the matrix $`C`$ can depend on $`V`$ <sup>7</sup><sup>7</sup>7In mathematical parlance, $`C`$ is a locally constant matrix -valued function defined on the moduli space with the branch-cuts removed.. Since the set $`V_0=\{|z|<1\}`$ does not contain any cuts (see Figure 2), it suffices to start by considering our relation in this region. Hence we define $`C`$ to be the matrix associated with $`V_0`$. Then performing a transformation $`z1/z`$ shows that the matrix associated with the region $`V_1=\{|z|>1\}`$ is the inverse of $`C`$. Thus, we expect the relations:
$`z^{1/2}U(1/z)`$ $`=`$ $`CU(z)\text{, if }|z|<1`$ (79)
$`z^{1/2}U(1/z)`$ $`=`$ $`C^1U(z)\text{, if }|z|>1.`$
In order to check these equalities and determine the matrix $`C`$, let us take $`z`$ to be such that $`|z|<1`$. Then $`|\frac{1}{z}|>1`$ and we have:
$`U_j(z)`$ $`=`$ $`{\displaystyle \frac{(1)^j}{j!}}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}\left[{\displaystyle \frac{(\frac{1}{2})_n}{n!}}\right]^4\nu _j(n,z)z^n`$ (80)
$`U_i^{}(z)`$ $`=`$ $`{\displaystyle \frac{(\delta _{i,odd}i\delta _{i,even})}{6\pi ^{3i}}}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}\left[{\displaystyle \frac{(\frac{1}{2})_n}{n!}}\right]^4\mu _i(n,1/z)z^n.`$ (81)
Defining $`b_{ij}`$ via:
$$c_{ij}=(1)^jj!\frac{(\delta _{i,odd}i\delta _{i,even})}{6\pi ^{3i}}b_{ij},$$
(82)
it suffices to compute the matrix $`B=(b_{ij})_{i,j=\mathrm{0..3}}`$, which satisfies:
$$\mu _i(n,1/z)=b_{ij}\nu _j(n,z).$$
(83)
Using the explicit form of these sequences given in (2.1,11) and (47,48), it is not very hard to show that the required matrix has the form:
$$B=\left[\begin{array}{cccc}6i\pi ^3& 18\pi ^2& 3i\pi & 1\\ 0& 6\pi ^2& 0& 1\\ 0& 12\pi ^2& 3i\pi & 1\\ 0& 0& 0& 1\end{array}\right],$$
(84)
which finally leads to the transition matrix of interest:
$$C=\left[\begin{array}{cccc}1& 3\frac{i}{\pi }& \pi ^2& \frac{i}{\pi ^3}\\ 0& 1& 0& \pi ^2\\ 0& 2i\pi & 1& \frac{i}{\pi }\\ 0& 0& 0& 1\end{array}\right]C^1=\left[\begin{array}{cccc}1& \frac{i}{\pi }& \pi ^2& \frac{i}{\pi ^3}\\ 0& 1& 0& \pi ^2\\ 0& 2i\pi & 1& \frac{i}{\pi }\\ 0& 0& 0& 1\end{array}\right].$$
(85)
The geometric interpretation of these results follows by writing the Meijer periods in the form:
$$U_j(z)=_{g_j(z)}\mathrm{\Omega }(z),$$
(86)
where $`g_j(z)(j=\mathrm{0..3})`$ is a basis of $`H_3(Y_z,\text{})`$. Following the general theory of variations of Hodge structure (see for a review in the context of its applications to mirror symmetry), we take the classes $`g_j`$ to be flat with respect to the Gauss-Manin connection on the moduli space <sup>8</sup><sup>8</sup>8Usually one takes this connection to act on cohomology, but here we use Poincare duality to transport the local system from $`H^3(Y)`$ to $`H_3(Y)`$. Hence we think of $`g_i(z)`$ as being flat sections of a bundle with fiber $`H_3(Y_z)`$. Then $`g_i(z)`$ will be multivalued due to the nontrivial holonomy of the connection.. Here $`\mathrm{\Omega }`$ is normalized such that $`lim_{z0}U_0(z)=1`$. On the other hand, using the variable $`s=1/z`$ and starting with $`s=0z=\mathrm{}`$ as the large complex structure point (in which case the Picard Fuchs equation coincides with the equation obtained from (45) by substituting $`s`$ for $`z`$) gives Meijer periods $`\stackrel{~}{U}_j(s)=U_j(1/s)`$, which can also be written in the form:
$$\stackrel{~}{U}_j(s)=_{\stackrel{~}{g}_j(s)}\stackrel{~}{\mathrm{\Omega }}(s),$$
(87)
where $`\stackrel{~}{g}_j(s)`$ is a flat basis of $`H_3(Y,\text{})`$ while $`\stackrel{~}{\mathrm{\Omega }}`$ is the holomorphic 3-form on $`Y`$ normalized via $`lim_{s0}\stackrel{~}{U}_0(s)=1`$. Then (79) shows that:
$`\stackrel{~}{g}_i(z)`$ $``$ $`C_{ij}g_j(z)`$ (88)
$`\stackrel{~}{\mathrm{\Omega }}(z)`$ $``$ $`z^{1/2}\mathrm{\Omega }(z),`$ (89)
for $`|z|<1`$. It follows that $`C`$ encodes the relation between the Meijer bases $`g_i`$ and $`\stackrel{~}{g}_i`$ of $`H_3(Y,\text{})`$ associated with the points $`z=0`$ and $`z=\mathrm{}`$, while the rescaling by $`z^{1/2}`$ reflects the different normalization of $`\mathrm{\Omega }`$ required by their interpretation as large complex structure points.
We can now shed more light on the vanishing periods at $`z=1`$. Indeed, applying (79) at that point shows that the vector $`U\left(1\right)=\left[\begin{array}{c}U_0\left(1\right)\\ U_1\left(1\right)\\ U_2\left(1\right)\\ U_3\left(1\right)\end{array}\right]`$ is an eigenvector of $`C`$ with eigenvalue one:
$$(CI)U(1)=0.$$
(90)
The kernel of the matrix $`(CI)`$ is a two-dimensional subspace spanned by the row vectors <sup>9</sup><sup>9</sup>9The matrix $`C`$ has eigenvalues $`1`$ and $`+1`$, each of which have multiplicity two. However, it is easy to check that $`C`$ is not diagonalizable. The reader may wonder why we do not apply relation (79) to the other fixed point $`z=1`$ and try to obtain vanishing periods there via a similar argument. The reason is, of course, that $`1`$ is a branch point for our analytic continuations, so that the limit of $`U(z)`$ at this point is not well-defined. While (79) holds in a directional limiting sense at $`z=1`$ (no matter from what direction in the complement of the cut we approach that point), this does not imply vanishing of a period there since the limits of $`U(z)`$ and $`U(1/z)`$ are different as $`z`$ approaches the value $`1`$ (note that $`z`$ and $`1/z`$ lie on different sides of the cut).:
$`[0,2\pi ^2,0,1],[2\pi ^2,i\pi ,1,0]`$
associated with the vanishing periods (67). This reproduces the result of Section 4 that this model admits a two dimensional subspace of periods which vanish at $`z=1`$. It also shows that this somewhat unusual situation is a consequence of the โaccidentalโ symmetry $`z1/z`$.
### 5.2 Physical interpretation
#### 5.2.1 The closed string sector
Let us consider the implications of these results for the bulk conformal field theory associated with our compactification. The B-model defined by $`Y`$ contains chiral primary operators $`๐ช^{p,p}`$ which are in one to one correspondence with generators of the Hodge groups $`H^{p,3p}(Y)`$. When computing correlators, we can replace $`H^{p,3p}(Y)`$ with their holomorphic counterparts $`^{p,3p}=^{3p}๐ฒ_p`$, where:
$$0^0^1^2^3=H^3(X)$$
(91)
is the Hodge filtration and:
$$0๐ฒ_0๐ฒ_1๐ฒ_2๐ฒ_3=H^3(Y)$$
(92)
is the โreducedโ monodromy weight filtration associated with a large complex structure point (see Appendix A of for a short explanation of this concept). Roughly, $`๐ฒ_j`$ is the space of those periods which have $`\mathrm{log}^j`$ leading behaviour near that point <sup>10</sup><sup>10</sup>10The vector space $`^3(Y_z)`$ can be identified with the space spanned by the vectors $`w_j:=\left[\begin{array}{c}U_j\left(z\right)\\ \delta U_j\left(z\right)\\ \delta ^2U_j\left(z\right)\\ \delta ^3U_j\left(z\right)\end{array}\right]`$. Viewing $`w_j`$ as a set of initial conditions for the Picard Fuchs equation at the point $`z`$ further identifies this space with the space of solutions to (45). . The monodromy filtrations can be easily determined by making use of the special logarithmic behaviour of the Meijer periods (see the expansions (9)): if $`z=0`$ is treated as a large complex structure point, then we obtain a filtration $`๐ฒ`$ which can be identified with the spaces of periods spanned by:
$$๐ฒ_0=<U_0>,๐ฒ_1=<U_0,U_1>,๐ฒ_2=<U_0,U_1,U_2>,๐ฒ_3=<U_0,U_1,U_2,U_3>.$$
(93)
On the other hand, treating $`z=\mathrm{}`$ as a large complex structure point gives:
$$\stackrel{~}{๐ฒ}_0=<\stackrel{~}{U}_0>,\stackrel{~}{๐ฒ}_1=<\stackrel{~}{U}_0,\stackrel{~}{U}_1>,\stackrel{~}{๐ฒ}_2=<\stackrel{~}{U}_0,\stackrel{~}{U}_1,\stackrel{~}{U}_2>,\stackrel{~}{๐ฒ}_3=<\stackrel{~}{U}_0,\stackrel{~}{U}_1,\stackrel{~}{U}_2,\stackrel{~}{U}_3>.$$
(94)
Hence (79) implies a nontrivial relation between $`(\stackrel{~}{๐ฒ})`$ and $`(๐ฒ)`$ and thus between $`\stackrel{~}{}^{p,3p}(Y_z)`$ and $`^{p,3p}(Y_z)`$. It follows that our symmetry involves a โrotationโ of the chiral primary operators $`๐ช^{p,3p}`$.
#### 5.2.2 The D-brane sector
The (BPS saturated) D-brane sector of our compactification can be realized by considering the open conformal field theory or, equivalently, by including boundary states. In the large complex structure limit of the IIB theory on $`Y`$, these correspond to special Lagrangian cycles $`C`$ in $`Y`$. Hence given a boundary state we can associate to it the homology class $`\gamma =[C]`$ of the associated cycle and hence the corresponding period $`_\gamma \mathrm{\Omega }`$. As we move away from this limit, the correspondence may be destroyed for some boundary states, due to the fact that the path we use for performing the marginal deformations could cross a marginal stability line . On such a line, the associated special Lagrangian cycle is expected to suffer a splitting transition of the type discussed in . Since we do not have a proper understanding of marginal stability lines in this model, the conclusions we can derive regarding the behaviour of D-brane states are only tentative.
The most basic question about such states concerns the dimensionality of the type IIA D-brane on $`X`$ mirror to a given IIB D-brane on $`Y`$. As discussed in , this is determined by the order of the logarithmic behaviour of the associated period in the large complex structure limit, i.e. by the smallest component of the monodromy weight filtration which contains that period. In our model, we have two points which can play the role of large complex structure points, and hence two monodromy weight filtrations $`(๐ฒ)`$ and $`(\stackrel{~}{๐ฒ})`$. Thus the correspondence between the mirror D-brane states (and even their dimension) involves a nontrivial rotation of $`H_{even}(X)`$.
A rather dramatic effect of this type can be observed as follows. Suppose that we define the large radius/large complex structure limit to correspond to $`z=0`$. Then consider a $`D2`$-brane in the large radius limit on $`X`$, whose mirror D3-brane is associated (up to a factor) with the period $`U_1`$. Note that this period is weakly integral (i.e. proportional with the period of $`\mathrm{\Omega }`$ over an integral homology class of $`Y`$). Now perform marginal deformations until we cross the circle $`|z|=1`$, reaching a point $`z_0`$ which lies outside the unit disk. At this point, we have a boundary state (the deformation of the original D-brane state) in the conformal field theory associated with $`z_0`$. Performing a duality transformation maps this theory into an equivalent conformal theory for which the large radius point correspond to $`z=\mathrm{}`$; this transformation will modify the associated period through the action of $`C^1`$ (and rescaling by $`z_0^{1/2}`$). Inspection of the matrix $`C^1`$ (equation (85)) shows that the associated period has $`\mathrm{log}^2`$ behaviour around $`z=\mathrm{}`$, and hence the mirror boundary state corresponds to a D4-brane! In fact, choosing $`z_0`$ to be far away in the $`z`$-plane assures that we are in the large radius limit of the dual model, and hence the associated $`D4`$-brane must correspond to a holomorphic 4-cycle on $`X`$. In other words, our duality seems to identify some $`D2`$-brane states on $`X`$ with $`D4`$-branes. Of course, this surprising conclusion may be avoided if the path used for analytic continuation crosses a marginal stability curve, or if the homology class under consideration does not actually contain a special Lagrangian cycle.
### 5.3 Small versus large size
What is the action of our symmetry on the size of $`X`$ ? The answer to this question depends on the precise definition of โsizeโ. Let us first consider the nonlinear sigma model measure of , which was shortly reviewed in the introduction. Following , we can start with $`z=0`$ as the large radius point and measure size by using the analytic continuation of the special coordinate $`t(z)`$. Then the symmetry $`z1/z`$ identifies $`t(z)`$ with $`\stackrel{~}{t}(z):=t(1/z)`$. Eliminating $`z`$ defines a map $`\stackrel{~}{t}=f(t)`$, which can be determined numerically and is plotted in Figure 8 for $`|z|`$ belonging to the interval $`(1,10^4)`$ ($`\stackrel{~}{J}`$ remains positive in this range, even though this is not obvious at the scale and from the viewing angle of this figure). We see that the duality indeed maps small into large distances โ a conclusion which is now established at the quantum level. Figure 9 displays the values of $`t(z)`$ for $`|z|(1,10^8)`$. Note that $`J=\mathrm{Im}(t)`$ remains positive when $`\mathrm{Im}(z)0`$.
$`\begin{array}{cc}\begin{array}{c}\text{}\end{array}& \begin{array}{c}\text{}\end{array}\\ \begin{array}{c}\\ \text{Figure 8. }\text{Graph of }\stackrel{~}{J}=Im(\stackrel{~}{t})\text{ vs }t=B+iJ\text{}\end{array}& \begin{array}{c}\\ \\ \text{Figure 9. Values of }t(z)\text{ for }1<|z|<10^8\text{}\end{array}\end{array}`$
What about the quantum volume of $`X`$ ? As discussed in the introduction, this is measured by the mass of a D6-brane wrapped over $`X`$, and it is natural to pick the D6-brane state whose mass vanishes at $`z=1`$, which is plotted in Figure 5. There is no positive lower bound for the (quantum) volume of $`X`$ โ string theory allows the entire manifold to shrink to zero size.
### 5.4 Phases
The semiclassical Kahler moduli space of $`X`$ was studied in , where it was shown that the model admits two phases, one of which is a large radius Calabi-Yau phase. The other phase can be analyzed via the linear sigma model techniques of , with the result that it is a hybrid phase which can be roughly described as a fibration of a $`\text{}_2`$ Landau-Ginzburg orbifold over a $`\text{}^3`$. This picture is tantalizingly close to a purely geometric description of that phase (say, in terms of a nonlinear sigma model having $`\text{}^3`$ or a closely related space as a target, for example through a construction along the lines of ) but, as pointed out in , the semiclassical picture provided by the linear sigma model is affected by strong quantum corrections which have the potential to seriously modify the discussion, thus making this geometric interpretation inconclusive. Our results allow us to make a precise statement about the effect of these corrections: they modify the theory in such a way that it becomes equivalent with its large radius incarnation ! In fact, once quantum corrections have been taken into account, there is no physical difference between the two limits and the model has a single phase (see Figure 10).
Figure 10. The effect of quantum corrections on the phase diagram of the IIA compactification on $`X`$.
### 5.5 Interpretation via special Lagrangian fibrations
How can we understand the behaviour of this model from the point of view of the SYZ conjecture ?. Since both $`z=0`$ and $`z=\mathrm{}`$ can be viewed as large complex structure points, the natural expectation is that $`Y_z`$ should admit two special Lagrangian fibrations, well-defined on some vicinities of the points $`z=0`$ and $`z=\mathrm{}`$, and related by the transformation (73). It was shown in that the monodromy weight filtration is determined by the fibration. Hence using one or the other of these fibrations corresponds to declaring $`z=0`$ or $`z=\mathrm{}`$ to be the large complex structure point. Then our small-large radius duality appears as a consequence of the fact that the two fibrations are isomorphic.
The techniques for constructing special Lagrangian fibrations of Calabi-Yau manifolds are not yet fully developed (see for partial results in this direction), so it is premature to attempt a complete analysis along these lines. However, simple and powerful methods are available in the large complex structure limit , where the problem can be reduced to one of toric geometry and hence can be approached with the machinery available in such situations . Appendix B uses a simple generalization of these techniques in order to identify the topology of the relevant fibrations. As in the hypersurface case, the base of each fibration turns out to be a 3-sphere.
The SYZ picture provides a natural interpretation of the nontrivial action of the duality on D-brane states: since mirror symmetry amounts to T-duality along the $`T^3`$ fibers, the dimension of the mirror D-brane depends on the relative position of a given IIB D3-brane with respect to the fibration of interest. Changing the fibration modifies this relative position, and hence can modify the dimension of the mirror holomorphic cycle. This is just the familiar fact that the dimension of a D-brane increases or decreases when performing T-duality along a direction orthogonal or parallel with its volume.
## 6 Conclusions
We completed the study of the hierarchy of one-parameter models introduced in , providing more evidence that the phenomena discussed in that paper are generic: in a typical IIA compactification on a one-parameter Calabi-Yau manifold, the non-perturbative state which becomes massless at the mirror of the conifold point is associated with a D6-brane. The general results derived in and in the present paper should open the way for extensions of the work of to more general Calabi-Yau compactifications, as well as providing a convenient framework for a systematic study of issues of marginal stability (see for a few steps in this direction) through the effective field theory methods of .
From a methodological point of view, our results show that most one-parameter models fit into a hypergeometric hierarchy, which allows for a very systematic approach to the computation of all periods. This should help prepare the ground for further investigations of D-brane effects in Calabi-Yau compactifications. The universal large radius expansions we have obtained should also help clarify some of the arithmetic properties of the mirror map when combined with the work of and .
We also performed a detailed study of a special one-parameter example, which displays some unusual features. In particular, we were able to bring some detailed evidence that this model realizes a Calabi-Yau version of large-small radius duality, thus confirming the suspicions of . We also presented evidence that, in the framework of , this duality is realized through the existence of two special Lagrangian fibrations โ a feature which has interesting implications for the physics of D-branes in the associated string theory compactification. It would be interesting to investigate this phenomenon further, as well as its implications for the problem of marginal stability of D-brane states. Since the duality exchanges the small and radius points, it should be possible to use it in this model in order to extract strong results regarding this issue.
Another interesting question is to what extent these phenomena generalize. Multiple large complex structure points are common in multi-parameter models (any model admitting topology-changing transitions possesses at least two such points). It would be interesting to see if similar discrete identifications occur in such models, and what can be learned from this about quantum corrections to the Kahler moduli space.
## Appendix A Some intermediate results for the models $`\text{}^5[3,3]`$ and $`\text{}^7[2,2,2,2]`$
### A.1 The model $`\text{}^5[3,3]`$
The canonical and Jordan form of the matrix $`R[\mathrm{}]`$, as well as a choice for the matrix $`P`$ are given below:
$`\begin{array}{c}R_{can}\left[\mathrm{}\right]=\left[\begin{array}{cccc}0& 1& 0& 0\\ 0& 0& 1& 0\\ 0& 0& 0& 1\\ \frac{4}{81}& 4/9& \frac{13}{9}& 2\end{array}\right],R_J\left[\mathrm{}\right]=\left[\begin{array}{cccc}2/3& 1& 0& 0\\ 0& 2/3& 0& 0\\ 0& 0& 1/3& 1\\ 0& 0& 0& 1/3\end{array}\right]\\ P=\left[\begin{array}{cccc}2/3& 5& 4/3& 4\\ 4/9& 8/3& 4/9& 8/3\\ \frac{8}{27}& 4/3& \frac{4}{27}& 4/3\\ \frac{16}{81}& \frac{16}{27}& \frac{4}{81}& \frac{16}{27}\end{array}\right]\end{array}`$ (109)
The small radius arithmetic identity associated to the collapsing period at $`z=1`$ is: $`_{n=0}^{\mathrm{}}\frac{a_n}{n!^2}=0`$, where
$`a_n=\mathrm{\Gamma }\left(n+1/3\right)^4\mathrm{\Gamma }\left(n+1/3\right)^2\psi \left(n+1/3\right)+\mathrm{\Gamma }\left(n+1/3\right)^4\mathrm{\Gamma }\left(n+1/3\right)^2\psi \left(n+1\right)+`$
$`\mathrm{\Gamma }\left(n+2/3\right)^4\mathrm{\Gamma }\left(n1/3\right)^2\psi \left(n1/3\right)+\mathrm{\Gamma }\left(n+2/3\right)^4\mathrm{\Gamma }\left(n1/3\right)^2\psi \left(n+1\right)`$
$`2\mathrm{\Gamma }\left(n+1/3\right)^4\mathrm{\Gamma }\left(n+1/3\right)^2\psi \left(n+2/3\right)2\mathrm{\Gamma }\left(n+2/3\right)^4\mathrm{\Gamma }\left(n1/3\right)^2\psi \left(n+1/3\right).`$
A pair identity follows from the large radius expansions.
### A.2 The model $`\text{}^7[2,2,2,2]`$
In this case, we have:
$`\begin{array}{c}R_{can}\left[\mathrm{}\right]=\left[\begin{array}{cccc}0& 1& 0& 0\\ 0& 0& 1& 0\\ 0& 0& 0& 1\\ 1/16& 1/2& 3/2& 2\end{array}\right],R_J\left[\mathrm{}\right]=\left[\begin{array}{cccc}1/2& 1& 0& 0\\ 0& 1/2& 1& 0\\ 0& 0& 1/2& 1\\ 0& 0& 0& 1/2\end{array}\right]\\ P=\left[\begin{array}{cccc}1/8& 1/4& 1/2& 1\\ 1/16& 0& 0& 0\\ 1/32& 1/16& 0& 0\\ \frac{1}{64}& 1/16& 1/16& 0\end{array}\right]\end{array}`$ (124)
## Appendix B Special Lagrangian fibrations of $`Y`$
A topological $`T^3`$ fibration in the large complex structure limit can be obtained by the methods of . This fibration is believed to admit a deformation to a special Lagrangian fibration of $`Y`$ as we move away from the large complex structure point. While the arguments discussed in those papers are restricted to hypersurfaces in toric varieties, our model $`\text{}^7[2,2,2,2]`$ is more general since it is a complete intersection. Assuming that some generalization of those arguments goes through in our case, we can attempt to construct our fibration along the same lines.
For this, let us first consider the point $`z=0\psi =\mathrm{}`$. In this limit, the defining equations (4.3) become:
$`x_1x_2=0`$ , $`x_3x_4=0`$
$`x_5x_6=0`$ , $`x_7x_8=0,`$
so that $`Y`$ reduces to a union of $`16`$ copies of $`\text{}^3`$ intersecting with normal crossings<sup>11</sup><sup>11</sup>11In this appendix, $`Y_{\mathrm{}}`$ and $`Y_0`$ mean $`Y_{\psi =\mathrm{}}`$ and $`Y_{\psi =0}`$, respectively.:
$$Y_{\mathrm{}}=\underset{u_1,u_2,u_3,u_4\text{}_2}{}Z_{u_1,u_2,u_3,u_4},$$
(125)
where $`Z_{u_1,u_2,u_3,u_4}=\{x=[x_1\mathrm{}x_8]\text{}^7|x_{1+u_1}=x_{3+u_2}=x_{5+u_3}=x_{7+u_4}=0\}\text{}^3`$. Following the procedure of , we consider the map $`\mu :\text{}^7\text{}^7`$ given by:
$`\mu (x)={\displaystyle \frac{_{k=1}^8|x_k|^2P_k}{_{k=1}^8|x_k|^2}},`$ (126)
with $`P_1\mathrm{}P_8`$ some points in general position in $`\text{}^7`$. The convex hull of these points defines a 7-simplex denoted by $`\mathrm{\Delta }`$, which clearly coincides with the image of $`\mu `$. According to the discussion of , a candidate for the desired $`T^3`$ fibration of $`Y`$ in the large radius limit is given by the restriction of $`\mu `$ to $`Y_{\mathrm{}}`$:
$$\mu _0:=\mu |_Y_{\mathrm{}}:Y_{\mathrm{}}\mathrm{im}(\mu _0)\mathrm{\Delta }.$$
(127)
Indeed, it is easy to see that the generic fiber of this map is a 3-torus. In the hypersurface case considered in , the image of $`\mu _0`$ coincides with the boundary of $`\mathrm{\Delta }`$ (which is topologically a 3-sphere, since in the hypersurface case $`\mathrm{\Delta }`$ has dimension 4), but for our complete intersection the situation is different. Indeed, it is easy to see that the image of each of the components $`Z`$ is a three-dimensional face of $`\mathrm{\Delta }`$. For example, we have:
$$\mu (Z_{1,1,1,1})=\{\frac{|x_2|^2P_2+|x_4|^2P_4+|x_6|^2P_6+|x_8|^2P_8}{|x_2|^2+|x_4|^2+|x_6|^2+|x_8|^2}|(x_2,x_4,x_6,x_8)\text{}^3\},$$
(128)
which coincides with the three dimensional face $`<P_2,P_4,P_6,P_8>`$ spanned by the vertices $`P_2,P_4,P_6`$ and $`P_8`$. Hence the base of our fibration coincides with the union $`\mathrm{im}(\mu _0)=\mathrm{\Delta }_0`$ of 16 three-dimensional tetrahedra <sup>12</sup><sup>12</sup>12$`\mathrm{\Delta }_0`$ is a subset of (but does not coincide with) the $`3`$-skeleton of $`\mathrm{\Delta }`$ (i.e. the union of all of its three-dimensional faces).. These tetrahedra intersect along common vertices, edges and facets, and the fibers of $`\mu _0`$ degenerate at the points of intersection.
What is the topology of the base $`\mathrm{\Delta }_0`$ ? To answer this question, note that the 16 tetrahedra composing the base are spanned by the vertices:
$`\begin{array}{ccccccc}<2,4,5,8>& ,& <1,4,5,8>& ,& <1,4,6,8>& ,& <2,4,6,8>\\ <2,4,5,7>& ,& <1,4,5,7>& ,& <1,4,6,7>& ,& <2,4,6,7>\\ <2,3,5,7>& ,& <1,3,5,7>& ,& <1,3,6,7>& ,& <2,3,6,7>\\ <2,3,5,8>& ,& <1,3,5,8>& ,& <1,3,6,8>& ,& <2,3,6,8>\end{array},`$ (133)
and $`\mathrm{\Delta }_0`$ is obtained by gluing these along their common faces. Then a momentโs thought shows that the resulting body is a 3-sphere (see Figures 11 and 12). Thus, just as in the hypersurface case, $`Y_{\mathrm{}}`$ is a $`T^3`$ fibration over $`S^3`$.
Figure 11. Arrangement of the 16 tetrahedra which form the base $`\mathrm{\Delta }_0`$. The points $`8`$ and $`8^{}`$ are identified, together will all identifications of edges and facets implied by this.
Figure 12. The identifications in Figure 11 can be performed in two steps. First, identify the edges starting from the point $`P_8P_8^{}`$; we represent this by introducing $`6`$ copies of that point. This shows that the topology of the base with the point $`P_8`$ removed is that of $`\text{}^3`$. Identifying the $`6`$ copies of $`P_8`$ amounts to adding a point to $`\text{}^3`$, which can be thought of as โthe point at infinityโ. This produces a 3-sphere.
Let us now consider the limit $`z\mathrm{}\psi 0`$. In this limit, the defining equations reduce to:
$`x_1^2+x_2^2=0`$ , $`x_3^2+x_4^2=0`$
$`x_5^2+x_6^2=0`$ , $`x_7^2+x_8^2=0,`$
which, via the transformation (73) are equivalent with:
$`y_1y_2=0`$ , $`y_3y_4=0`$
$`y_5y_6=0`$ , $`y_7y_8=0.`$
Hence $`Y_0`$ reduces once again to 16 copies of $`\text{}^3`$ intersecting transversely, as should be expected from the fact that $`Y_\psi `$ and $`Y_{1/\psi }`$ are isomorphic as complex manifolds. Since the form of (B) is the same as above, we can once again use the map:
$`\stackrel{~}{\mu }(y)={\displaystyle \frac{_{k=1}^8|y_k|^2Q_k}{_{k=1}^8|y_k|^2}}`$ (134)
(with $`Q_k`$ some points in general position in $`\text{}^7`$) in order to produce a $`T^3`$-fibration $`\stackrel{~}{\mu }_0`$ of $`Y_0`$ whose basis is a 3-sphere.
The fibrations $`\mu _0`$ and $`\stackrel{~}{\mu }_0`$ are related through the biholomorphic map $`\varphi :Y_{\mathrm{}}Y_0`$ which identifies the complex structures $`J_{\mathrm{}}`$ and $`J_0`$ of $`Y_{\mathrm{}}`$ and $`Y_0`$:
$$J_0=d\varphi J_{\mathrm{}}(d\varphi )^1.$$
(135)
We may hope that some appropriate deformations of the fibrations $`\mu _0,\stackrel{~}{\mu }_0`$ are special Lagrangian with respect to $`J_{\mathrm{}},J_0`$ and the associated metrics. |
warning/0002/hep-ex0002009.html | ar5iv | text | # BNLโ67140 February 3, 2000 Experimental Results on Radiative Kaon Decays11footnote 1 To be published in the Proceedings of the 3๐โข๐ Workshop on PHYSICS AND DETECTORS FOR DAฮฆNE, Frascati, Italy, Nov. 16โ19, 1999 (Frascati Physics Series Vol. XVI, 2000)
## 1 Introduction
Radiative kaon decays provide a testing ground for Chiral Perturbation Theory (ChPT). ChPT provides a framework for calculating the decay rates for several modes, either directly or relative to other measured modes. The radiative modes are important for determining long distance contributions to other decays of interest: the two-photon contribution to $`K_L^{}\mu ^+\mu ^{}`$, and the CP-conserving and indirect CP-violating contributions to $`K_L^{}\pi ^{}e^+e^{}`$ and $`K_L^{}\pi ^{}\mu ^+\mu ^{}`$. They are also important as backgrounds to other modes (e.g. the $`K_L^{}e^+e^{}\gamma \gamma `$ background to $`K_L^{}\pi ^{}e^+e^{}`$).
A number of recent results have been reported in the literature, as well as in several recent conferences.
## 2 Radiative K<sub>ฯ2</sub> Decays
The radiative K<sub>ฯ2</sub> decays: $`K^+\pi ^+\pi ^{}\gamma `$, $`K_L^{}\pi ^+\pi ^{}\gamma `$ and $`K_\mathrm{S}^{}\pi ^+\pi ^{}\gamma `$ have two contributions. One is inner bremsstrahlung (IB) radiation from one of the charged particles. The second is direct emission (DE) from the vertex. The branching ratio of the IB contribution scales with the underlying K<sub>ฯ2</sub> decay rate. Whereas, the rate for direct emission is expected to be roughly comparable for all three modes.
A new result for $`K_L^{}\pi ^+\pi ^{}\gamma `$ from KTeV is shown in Fig. 1. The energy of the photon is shown, along with the contributions from IB and DE.
The DE component is modified by a โ$`\rho `$-propagatorโ that serves to soften the DE spectrum. The branching ratio for the direct emission component (see eq.1) is
$$\mathrm{BR}(\mathrm{K}_\mathrm{L}^{}\pi ^+\pi ^{}\gamma ;\mathrm{DE})=(3.70\pm 0.10)\times 10^5(\mathrm{E}_\gamma ^{}>20\mathrm{M}\mathrm{e}\mathrm{V})$$
(1)
The ratio of direct emission to DE+IB is (see eq.2)
$$\mathrm{DE}/(\mathrm{DE}+\mathrm{IB})=0.685\pm 0.009\pm 0.017$$
(2)
This result is based on $``$5% of the total KTeV data for this mode.
There are new results from E787 in the charged decay mode ($`K^+\pi ^+\pi ^{}\gamma `$) as well. This result is striking, in that the branching ratio is a factor of 4 lower than the previous value. The data is traditionally expressed in terms of the variable W, which is defined as:
$`\mathrm{W}^2`$ $``$ $`(\mathrm{p}\mathrm{q})/\mathrm{m}_{\mathrm{K}^+}^2\times (\mathrm{p}_+\mathrm{q})/\mathrm{m}_{\pi ^+}^2`$
$`=`$ $`\mathrm{E}_\gamma ^2\times (\mathrm{E}_{\pi ^+}\mathrm{P}_{\pi ^+}\times \mathrm{cos}\theta _{\pi ^+\gamma })/(\mathrm{m}_{\mathrm{K}^+}^2\times \mathrm{m}_{\pi ^+}^2)`$
The new result from E787, shown in Figure 2,
has about 8 times higher statistics than the old one. The branching ratio for the direct emission component, from a fit to IB and DE (see eq.4) is
$$\mathrm{BR}(\mathrm{K}^+\pi ^+\pi ^{}\gamma ;\mathrm{DE})=(4.72\pm 0.77)\times 10^6(55<\mathrm{T}_{\pi ^+}<90\mathrm{M}\mathrm{e}\mathrm{V})$$
(4)
This represents half of the E787 data that is currently on tape. The interference term is small, $`(0.4\pm 1.6)`$% and the direct emission is $`(1.85\pm 0.30)`$%. The decay rate, corrected to full phase space<sup>2</sup><sup>2</sup>2This correction assumes that the form factor has no energy dependence., is now measured to be similar to that for $`K_L`$: $`\mathrm{\Gamma }(K^+\pi ^+\pi ^{}\gamma ;DE)=808\pm 132s^1`$ vs. $`\mathrm{\Gamma }(K_L^{}\pi ^+\pi ^{}\gamma ;DE)=617\pm 18s^1`$.
KTeV also has new results on $`K_L^{}\pi ^+\pi ^{}e^+e^{}`$, where the photon has internally converted to two electrons. In addition to measuring the branching ratio, a T-odd observable in the angular distribution of the plane of the $`\pi `$-pair vs. the plane of the electron pair is observed. This data represents one quarter of the final KTeV sample.
A summary of the current experimental status of radiative K<sub>ฯ2</sub> decays is shown in Table 1.
## 3 $`K\pi \gamma \gamma `$ Decays
The decay $`K_L^{}\pi ^{}\gamma \gamma `$ is very interesting, since to $`๐ช(p^4)`$ of ChPT the decay rate and spectral shape are completely determined, without any free parameters. The prediction of the spectral shape is a striking success of ChPT; however, the decay rate is a factor of 3 too small. To match the experimental number a model dependent contribution from $`๐ช(p^6)`$ is needed, which is usually parameterized with a constant $`a_V`$. The CP-conserving contribution to $`K_L^{}\pi ^{}e^+e^{}`$ depends on the value of $`a_V`$. Based on half of the total data sample, KTeV has recently measured $`a_V=0.72\pm 0.05\pm 0.06`$, implying a contribution of $`12\times 10^{12}`$.
The charged mode $`K^+\pi ^+\gamma \gamma `$ is more complicated, requiring an unknown parameter, $`\widehat{c}`$, even at $`๐ช(p^4)`$. Both the decay rate and spectral shape are predicted with this single parameter. E787 has measured $`\widehat{c}=1.8\pm 0.6`$.
The experimental measurements of $`K\pi \gamma \gamma `$ are summarized in Table 2.
The KTeV measurement of $`K_L^{}\pi ^{}e^+e^{}\gamma `$ should improve by $`\times `$3; the measurements of $`K_L^{}\pi ^{}\mathrm{}^+\mathrm{}^{}`$ are background limited, and will improve by $`\sqrt{3}`$.
## 4 K<sup>0</sup> to Two Real or Off-shell Photons
The decay $`K_\mathrm{S}^{}\gamma \gamma `$ is predicted in $`๐ช(p^4)`$ of ChPT, without any free parameters, to occur with BR($`K_\mathrm{S}^{}\gamma \gamma `$) = $`2.0\times 10^6`$. This is in good agreement with the experimental value (see Table 3), although the experimental errors need to be reduced.
The decay $`K_L^{}\gamma \gamma `$ is of interest for its importance in interpreting the measurement of $`K_L^{}\mu ^+\mu ^{}`$. The decay $`K_L^{}\mu ^+\mu ^{}`$ is sensitive to internal top quark loops, that would allow a determination of the fundamental SM parameter $`\rho `$. The decay is, however, dominated by the decay $`K_L^{}\gamma \gamma `$ with the photons converting to a $`\mu ^\pm `$ pair. For this reason a precise measure of $`K_L^{}\gamma \gamma `$ is needed. With the improved precision on $`K_L^{}\mu ^+\mu ^{}`$ from E871, the uncertainties on $`K_L^{}\gamma \gamma `$ and $`\frac{K_\mathrm{L}^{}\pi ^{}\pi ^{}}{K_\mathrm{S}^{}\pi ^+\pi ^{}}`$ are now contributing significantly to the uncertainty on the ratio
$`{\displaystyle \frac{\mathrm{\Gamma }(K_L^{}\mu ^+\mu ^{})}{\mathrm{\Gamma }(K_L^{}\gamma \gamma )}}`$ $`=`$ $`\left[\begin{array}{c}B(K_L^{}\mu ^+\mu ^{})\\ \\ B(K_L^{}\pi ^+\pi ^{})\end{array}\right]\times `$
$`\left[\left|\begin{array}{c}\eta _+\\ \\ \eta _{}\end{array}\right|\begin{array}{c}B(K_S^{}\pi ^+\pi ^{})\\ \\ B(K_\mathrm{S}^{}\pi ^{}\pi ^{})\end{array}\right]\times \left[\begin{array}{c}B(K_\mathrm{L}^{}\pi ^{}\pi ^{})\\ \\ B(K_L^{}\gamma \gamma )\end{array}\right]`$
$`[1.55\%][(0.23\%)(1.28\%)][1.42\%]`$
$`=`$ $`(1.213\pm 0.030)\times 10^5`$
KLOE should be able to contribute to improving both of these measurements. Finally there is a long distance dispersive contribution, from two off-shell photons, for which additional input from ChPT and measurements of the decays $`K_L^{}e^+e^{}\gamma `$, $`K_L^{}\mu ^+\mu ^{}e^+e^{}`$ and $`K_L^{}e^+e^{}e^+e^{}`$ are needed.
Results of kaon decays to two real or off-shell photons are summarized in Table 3.
The KTeV measurements of $`K_L^{}e^+e^{}e^+e^{}`$ and $`K_L^{}\mu ^+\mu ^{}e^+e^{}`$ should improve by $`\times `$4 and the modes $`K_L^{}e^+e^{}\gamma \gamma `$ and $`K_L^{}\mu ^+\mu ^{}\gamma \gamma `$ should improve by $`\times `$3 with the final KTeV data set. The $`K_L^{}e^+e^{}\gamma `$ should improve by $`\times `$20 and $`K_L^{}\mu ^+\mu ^{}\gamma `$ should improve by $`\times `$30. The $`K_S`$ modes may be improved by NA48 in a special run, after $`ฯต^{}/ฯต`$. The $`K_L^{}\gamma \gamma `$ and $`K_S^{}\gamma \gamma `$ as well as several other modes will be improved by KLOE. There is no improvement in the foreseeable future for $`K_L^{}e^+e^{}`$ or $`K_L^{}\mu ^+\mu ^{}`$.
## 5 Radiative K<sub>โ2</sub> Decays
The form factors in the decays $`K^+\mathrm{}^+\nu _{\mathrm{}}\gamma `$, A and V, and, R, in the decays $`K^+\mathrm{}^+\nu _{\mathrm{}}\mathrm{}^+\mathrm{}^{}`$, are predicted by ChPT. Recent measurements should allow precise experimental determinations of all three parameters. The most recent determination of $`|F_V+F_A|=0.165\pm 0.007\pm 0.011`$ from the E787 measurement of the direct emission component of $`K^+\mu ^+\nu \gamma `$, usually called Structure Dependent (SD<sup>+</sup>) radiation, is consistent with the previous determination of $`|F_V+F_A|=0.148\pm 0.010`$ from $`K^+e^+\nu \gamma `$. A limit of $`0.25<F_VF_A<0.07`$ is derived from the $`K^+\mu ^+\nu \gamma (SD^+)`$. An improved measure of $`F_VF_A`$ along with a measure of $`R`$ should be available soon from E865.
A summary of the recent radiative K<sub>โ2</sub> results is presented in Table 4.
## 6 Other Radiative Kaon Decays
The experimental sensitivity for the other radiative kaon decays K<sub>ฯ3ฮณ</sub>, K<sub>โ3ฮณ</sub> and K<sub>ฯ4ฮณ</sub> are such as to only be sensitive to IB contributions. All of these measurements are consistent with theoretical predictions. A summary of the results is given in Table 5.
A couple of modes should be seen for the first time in existing data, $`K^+\pi ^{}\mu ^+\nu _\mu \gamma `$ (E787) and $`K^+\pi ^+\pi ^{}e^+\nu _e\gamma `$ (E865). Improvements in other modes may be possible, particularly at IHEP.
## 7 Conclusions
Several new results are expected from KTeV and NA48, as well as a few more from E787,E865 and E871. With the turn on of DA$`\mathrm{\Phi }`$NE and KLOE, which is well equipped for the radiative modes, we can expect another round of new measurements. Finally, the next generation of rare kaon experiments, designed to fully constrain the CKM unitarity triangle, by measuring the โGolden modesโ $`K^+\pi ^+\nu \overline{\nu }`$ and $`K_L^{}\pi ^{}\nu \overline{\nu }`$ , are under construction (E391a, E949) or being designed (KOPIO, CKM, KAMI). These experiments will provide even more precise measurements of several radiative modes.
## Acknowledgements
I would like to thank several people from various experiments for providing data and discussions for this talk, including: Hong Ma, Mike Zeller, Bob Tschirhart, John Belz, Lutz Koepke, Bill Molzon, Leonid Landsberg, Takeshi Komatsubara, Takashi Nakano and Laurie Littenberg. This work was supported under U.S. Department of Energy contract #DE-AC02-98CH10886. |
warning/0002/hep-th0002225.html | ar5iv | text | # Universal Ratios in the 2-D Tricritical Ising Model
## Abstract
We consider the universality class of the twoโdimensional Tricritical Ising Model. The scaling form of the freeโenergy naturally leads to the definition of universal ratios of critical amplitudes which may have experimental relevance. We compute these universal ratios by a combined use of results coming from Perturbed Conformal Field Theory, Integrable Quantum Field Theory and numerical methods.
An unifying principle in the study of critical phenomena goes under the name of universality . In the vicinity of a phase transition, when the correlation length is much larger than any microscopic scale, one can assign each system to a universality class, which is identified by its dimensionality $`D`$, the symmetry properties of the order parameters and the number of relevant fields. The first characteristic of a given universality class is the set of critical exponents, expressed in terms of algebraic expressions of the conformal dimensions of the relevant fields. Additional data of a universality class may be derived by the scaling properties of the freeโenergy alone. These data โ called universal ratios โ are pure numbers, obtained by taking particular combinations of various thermodynamical amplitudes in such a way to cancel any dependence on the microscopic scales. Together with critical exponents, universal ratios are ideal fingerprints of the universality classes. From an experimental point of view, there is by now a large literature on universal ratio measurements of various systems extending from binary fluids to magnetic systems and polymer conformations (for an extensive review on the subject, see ).
In recent years, due to the theoretical progress achieved in the study of twoโdimensional models (at criticality by the methods of Conformal Field Theory (CFT) , and away from criticality by the approach of Perturbed Conformal Theories ), several universal quantities have been computed by different techniques for a large variety of bidimensional systems, such as the selfโavoiding walks , the Ising model , the q-state Potts model , to name few. In this letter we will focus on the first determination of some universal ratios relative to the class of universality of the $`2`$$`D`$ Tricritical Ising Model (TIM) for which very few universal quantities are known (see ). Whereas the $`3`$$`D`$ TIM describes, for instance, the universality class of an antiโferromagnet with strong uniaxial anisotropy like FeCl<sub>2</sub>, its $`2D`$ version can describe the tricritical behaviour of a binary mixture of thin films of $`He^3He^4`$ or orderโdisorder transitions in absorbed systems (for a review on the theory of tricritical points, see ). Hence there is an obvious interest in computing the amplest set of universal data for this universality class and in testing the theoretical predictions versus their experimental determinations.
In a continuum version of the TIM (which is, after all, a particular representative of this universality class), it is convenient to adopt a LandauโGinzburg (LG) formulation based on a scalar field $`\mathrm{\Phi }(x)`$ with $`\mathrm{\Phi }^6`$ interaction. The LG approach permits to have a clear bookkeeping of the symmetry properties of each order parameter and to easily understand the phase diagram of the model, at least qualitatively. The class of universality of the TIM is then described by the LG euclidean action
$`๐={\displaystyle d^Dx\left[\frac{1}{2}(_\mu \mathrm{\Phi })^2+g_1\mathrm{\Phi }+g_2\mathrm{\Phi }^2+g_3\mathrm{\Phi }^3+g_4\mathrm{\Phi }^4+\mathrm{\Phi }^6\right]}`$
with the tricritical point identified by the bare conditions $`g_1=g_2=g_3=g_4=0`$. Adopting a magnetic terminology, the statistical interpretation of the coupling constants is as follows: $`g_1`$ plays the role of an external magnetic field $`h`$, $`g_2`$ measures the displacement of the temperature from its critical value $`(TT_c)`$, $`g_3`$ may be regarded as a staggered magnetic field $`h^{}`$ and finally $`g_4`$ may be thought as a chemical potential $`\mu `$ for the vacancy density. Dimensional analysis shows that the upper critical dimension of the model is $`D=3`$, where tricritical exponents are expected to have their classical values (apart logarithmic corrections). In two dimensions, although the mean field solution of the model cannot be trusted for the strong fluctuations of the order parameters, an exact solution at criticality is provided by CFT. In fact, the TIM is described by the second model of the unitary minimal series of CFT , with central charge equal to $`C=\frac{7}{10}`$. There are six primary fields, identified with the normal ordered composite LG fields , which close an algebra under the Operator Product Expansion (OPE). Only four of them are relevant (i.e. with conformal dimension $`\mathrm{\Delta }<1`$): $`\sigma =\phi _1\mathrm{\Phi }`$ ($`\mathrm{\Delta }_1=\frac{3}{80}`$), $`\epsilon =\phi _2:\mathrm{\Phi }^2:`$ ($`\mathrm{\Delta }_2=\frac{1}{10}`$), $`\sigma ^{}=\phi _3:\mathrm{\Phi }^3:`$ ($`\mathrm{\Delta }_3=\frac{7}{16}`$) and $`t=\phi _4:\mathrm{\Phi }^4:`$ ($`\mathrm{\Delta }_4=\frac{3}{5}`$). The fields $`\epsilon `$ and $`t`$ are even under the $`Z_2`$ spinโsymmetry whereas $`\sigma `$ and $`\sigma ^{}`$ are odd. There is another $`Z_2`$ symmetry of the model (related to its self-duality), under which $`D\epsilon D^1=\epsilon `$, $`DtD^1=t`$, whereas the magnetic order parameters are mapped onto their corresponding disorder parameters. Each of the above relevant fields can be used to move the TIM away from criticality (the resulting phases of the model are discussed in ).
In order to derive the scaling form of the free energy and the set of universal ratios for the $`2`$$`D`$ TIM, let us first normalise the two-point functions of the fields as $`\phi _i(r)\phi _i(0)\frac{A_i}{r^{4\mathrm{\Delta }_i}}`$ when $`r0`$ (in the perturbed CFT approach to the model, $`A_i=1`$). When the TIM is moved away from criticality by means of one (or several) of its relevant fields, with the resulting action $`๐=๐_{CFT}+_pg_pd^2x\phi _p(x)`$, a finite correlation length $`\xi `$ generally appears. Its scaling form may be written in four possible equivalent ways, according to which coupling constant is selected out as a prefactor
$$\xi =a(K_ig_i)^{\frac{1}{22\mathrm{\Delta }_i}}_i\left(\frac{K_jg_j}{(K_ig_i)^{\varphi _{ji}}}\right),$$
(1)
where $`a`$ is some microscopic length scale, $`\varphi _{ji}\frac{1\mathrm{\Delta }_j}{1\mathrm{\Delta }_i}`$, and $`_i`$ are universal homegeneous scaling functions of the ratios $`\frac{K_jg_j}{(K_ig_i)^{\varphi _{ji}}}`$. The terms $`K_i`$ are nonโuniversal metric factors which depend on the unit chosen for measuring the external source $`g_i`$, alias on the particular realization of the universality class. Let $`f[g_1,g_2,g_3,g_4]`$ be the singular part of the freeโenergy (per unit volume). According to which coupling constant is selected out as a prefactor, it can be parameterised in four possible equivalent ways as:
$$f[g_1,g_2,g_3,g_4]\left(K_ig_i\right)^{\frac{1}{1\mathrm{\Delta }_i}}_i\left(\frac{K_jg_j}{(K_ig_i)^{\varphi _{ji}}}\right),$$
(2)
where $`_i`$ are scaling functions. For the Vacuum Expectation Value (VEV) of the fields $`\phi _j`$ in the $`i^{\mathrm{th}}`$ direction (i.e. for the offโcritical theory finally obtained by $`g_i0`$, $`g_k=0`$, $`ki`$), we have
$$\phi _j_i=\frac{f}{g_j}|{}_{g_k=0}{}^{}B_{ji}g_i^{\frac{\mathrm{\Delta }_j}{1\mathrm{\Delta }_i}},$$
(3)
where, from (2), $`B_{ji}K_jK_i^{\frac{\mathrm{\Delta }_j}{1\mathrm{\Delta }_i}}`$. In a similar manner, for the generalized susceptibilities we have
$$\widehat{\mathrm{\Gamma }}_{jl}^i=\frac{^2f}{g_lg_j}|{}_{g_k=0}{}^{}\mathrm{\Gamma }_{jl}^ig_i^{\frac{\mathrm{\Delta }_j+\mathrm{\Delta }_l1}{1\mathrm{\Delta }_i}},$$
(4)
where, from (2), $`\mathrm{\Gamma }_{jl}^iK_jK_lK_i^{\frac{\mathrm{\Delta }_j+\mathrm{\Delta }_l1}{1\mathrm{\Delta }_i}}`$. These quantities are obviously symmetric in the lower indices ($`\widehat{\mathrm{\Gamma }}_{22}^i`$ and $`\widehat{\mathrm{\Gamma }}_{11}^i`$ are respectively the usual specific heat and magnetic susceptibility in the $`i^{\mathrm{th}}`$ direction). Similarly, for the correlation length we have $`\xi _i=a\xi _0g_i^{\frac{1}{22\mathrm{\Delta }_i}}`$, with $`\xi _0K_i^{\frac{1}{22\mathrm{\Delta }_i}}`$. From the above formulas, appropriate combinations can be found such that the nonโuniversal metric factors $`K_i`$ cancel out. Some of the $`2`$$`D`$ universal ratios are:
$`(R_c)_{jk}^i`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Gamma }_{ii}^i\mathrm{\Gamma }_{jk}^i}{B_{ji}B_{ki}}}`$ (5)
$`(R_\chi )_j^i`$ $`=`$ $`\mathrm{\Gamma }_{jj}^iB_{jj}^{\frac{\mathrm{\Delta }_j1}{\mathrm{\Delta }_j}}B_{ji}^{\frac{12\mathrm{\Delta }_j}{\mathrm{\Delta }_j}}`$ (6)
$`R_\xi ^i`$ $`=`$ $`\left(\mathrm{\Gamma }_{ii}^i\right)^{1/2}\xi _i^0`$ (7)
$`(R_A)_j^i`$ $`=`$ $`\mathrm{\Gamma }_{jj}^iB_{ii}^{\frac{2\mathrm{\Delta }_j\mathrm{\Delta }_i+2}{\mathrm{\Delta }_i}}B_{ij}^{\frac{2\mathrm{\Delta }_j2}{\mathrm{\Delta }_i}}`$ (8)
$`(Q_2)_{jk}^i`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Gamma }_{jj}^i}{\mathrm{\Gamma }_{jj}^k}}\left({\displaystyle \frac{\xi _k^0}{\xi _j^0}}\right)^{24\mathrm{\Delta }_j}.`$ (9)
In this letter we only consider the case $`i=1,2`$, which correspond to the most important physical deformations of the model (the magnetic and the thermal ones), i.e. those which are most accessible from an experimental point of view. For both the magnetic and thermal deformation there are no mixing among the conformal fields due to ultraviolet renormalization . A complete analysis relative to all deformations of the TIM and the theoretical details of our approach will be published elsewhere .
The $`ฯต`$ perturbation around the critical TIM is integrable and its behavior is governed by the $`E_7`$ algebra . Therefore the $`B_{j2}`$โs in eq. (3) have been computed exactly in . On the other hand, the $`\sigma `$ perturbation is non-integrable (numerical indications were discussed in ). In this case, the $`B_{j1}`$โs have been numerically evaluated in by using the soโcalled Truncated Conformal Space Approach (TCSA) . This method consists in diagonalizing the offโcritical Hamiltonian on a cylinder in a truncated conformal basis of the critical TIM such that an extimation of $`\phi _j_1`$ can be obtained from the knowledge of the eigenvectors (only the ground state eigenvector is needed for the VEV). All these calculations can be easily performed by means of the numerical program of ref. .
In order to estimate the universal ratios, it is still necessary to calculate the $`\mathrm{\Gamma }_{jk}^i`$โs. Their values can be extracted in two different ways. The first method is purely numerical and of immediate use, since it consists in employing the TCSA to compute numerically the derivative $`\frac{}{g_k}\phi _j_i`$ (details will be found in ). The second method is based on the fluctuationโdissipation theorem which permits to express the generalised susceptibilities as
$$\widehat{\mathrm{\Gamma }}_{jk}^i=d^2x\phi _j(x)\phi _k_i^c,$$
(10)
where $`\mathrm{}^c`$ indicates the connected correlator. Therefore in this second approach we first need to evaluate the $`2`$point correlation functions and then to perform the integration. For our calculation of the universal ratios, we have employed both methods, finding an agreement in their final outputs. Let us briefly discuss the second method. First of all, write the integral (10) in polar coordinate as $`\widehat{\mathrm{\Gamma }}_{jk}^i=2\pi ๐rr\phi _j(r)\phi _k_i^c`$. Secondly, decompose the integral over $`r`$ into two integrals over the regions $`0<r<R`$ and $`rR`$ with $`R\xi `$. When $`r<R`$, the correlation function $`\phi _j(r)\phi _k(0)_i`$ can be efficiently evaluated by using a shortโdistance expansion
$$\phi _j(r)\phi _k(0)_i=\underset{l}{}C_{jk}^l(r)\phi _l_i$$
(11)
where the nonโanalytic dependence on the coupling constant is completly encoded into the VEVโs, whereas the structure constants $`C_{jk}^l(r)`$ can be evaluated perturbatively in $`g`$
$$C_{jk}^l(r)=r^{2(\mathrm{\Delta }_l\mathrm{\Delta }_j\mathrm{\Delta }_k)}\underset{n=0}{\overset{\mathrm{}}{}}C_{jk}^{l(n)}\left(g_ir^{22\mathrm{\Delta }_i}\right)^n.$$
(12)
For the TIM, $`C_{jk}^{l(0)}`$ have been computed in whereas their first correction can be obtained by the formula
$$C_{jk}^{l(1)}=^{^{}}d^2z\phi _l(\mathrm{})\phi _i(z)\phi _k(1)\phi _j(0)_{CFT},$$
(13)
where the prime indicates a suitable infrared regularization of the integral. As shown in , an efficient way to compute the regularised integrals is through a Mellin transformation. Hence, the calculation of the above integral (13) on the conformal functions plus the knowledge of the various expectation values $`\phi _l_i`$ enables us to reach a quite accurate approximation of $`\phi _j(r)\phi _k(0)_i`$ in the ultraviolet limit, i.e. for $`r<R`$. By choosing $`R\xi `$, one can obtain an overlap between the ultraviolet and the infrared representations of the correlation functions. The latter is expressed by means of the spectral series of the correlators on the massive states $`A_k(\theta )`$ of the offโcritical theory
$$\phi _j(x)\phi _k(0)_c=\underset{n=1}{\overset{\mathrm{}}{}}g_n(r),$$
(14)
where
$`g_n(r)={\displaystyle }{\displaystyle \frac{d\theta _1}{2\pi }}\mathrm{}{\displaystyle \frac{d\theta _n}{2\pi }}0|\phi _j(0)|A_{a_1}(\theta _1)\mathrm{}A_{a_n}(\theta _n)\times `$
$`A_{a_1}(\theta _1)\mathrm{}A_{a_n}(\theta _n)|\phi _k(0)|0e^{r_{k=1}^nm_k\mathrm{cosh}\theta _k}.`$
As tested in several examples (see, for instance ), the above series (14) converges very fast even for $`r\xi `$ so that its truncation to the lowest terms is able to capture the correct behaviour of the correlator in the interval $`r\xi `$. For the integrable theory defined by the thermal deformation of the TIM, one can truncated the series up to the lowest $`2`$โparticle states, with the relative matrix elements computed along the lines of the refs. . For the nonโintegrable theory defined by the magnetic deformation of the TIM, it is hard to go beyond the oneโparticle matrix elements and one has to be satisfied with the estimate of the correlators obtained by the oneโparticle contributions only: since this theory has two lowest masses with mass ratio $`m_22m_1\mathrm{cos}\frac{\pi }{5}`$ , in this case we have
$`\phi _j(r)\phi _k(0)_i{\displaystyle \frac{1}{\pi }}\left(f_j^1f_k^1K_0(m_1r)+f_j^2f_k^2K_0(m_2r)\right)`$
where $`K_0(x)`$ is the modified Bessel function and the indices $`1,2`$ refer to the first and second massive states. The oneโparticle matrix elements of this model $`f_j^k=0|\phi _j(0)|A_k`$ can be also computed numerically by using the TCSA .
Once an overlap of the short and large distance expansions of the correlators in the region $`r\xi `$ has been checked, a numerical integration of the correlators provides the $`\mathrm{\Gamma }_{jk}^i`$โs. An explicit test of the validity of the above method (with a corresponding estimate of its errors) is provided by the comparison of the values of $`\mathrm{\Gamma }_{ik}^i`$ (obtained by the numerical integration) with their exact determination extracted by the $`\mathrm{\Delta }`$theorem sum rule, when this theorem applies :
$$\mathrm{\Gamma }_{ik}^i=\frac{\mathrm{\Delta }_k}{1\mathrm{\Delta }_k}B_{ki}.$$
(15)
This check shows that the uncertanties for $`\mathrm{\Gamma }_{jk}^i`$ is at worst about $`5\%`$, better for the strongest relevant operators. Gathering all these results, a set of universal ratios for the TIM have been obtained. Some of them are exact, like $`(R_c)_{1,k}^1=\frac{240}{5929}\mathrm{\Delta }_k`$, $`(R_c)_{2,k}^2=\frac{10}{81}\mathrm{\Delta }_k`$ ($`k=1,\mathrm{},4`$). We have also computed those relative to the low and high temperature phase of the model (Table 1). An interesting universal ratio is provided in this case by the correlation length prefactors $`\xi _0^{}`$, below and above the critical temperature (as extracted from the correlation function of the magnetic operator using its duality properties)
$$\frac{\xi _0^{}}{\xi _0^+}=2\mathrm{cos}\left(\frac{5\pi }{18}\right)1.28557\mathrm{}$$
(16)
which can be inferred by the exact mass spectrum of the model and the parity properties of the excitations .
In summary, we have combined techniques coming from CFT, integrable models and numerical methods to obtain for the first time a set of universal quantities for the class of universality of the $`2D`$ Tricritical Ising Model. It would be interesting to have an experimental determination of these quantities and a comparison with the theoretical predictions presented here.
This work done under partial support of the EC TMR Programme Integrability, non-perturbative effects and symmetry in Quantum Field Theories. We would like also to thank A.B. and Al.B. Zamolodchikov for useful discussions. We are also grateful to V. Rittenberg and M. den Nijs for suggestions and for their interest in this work.
Table 1: Amplitude ratios $`R_{jk}^2=\frac{\mathrm{\Gamma }_{jk}^{2+}}{\mathrm{\Gamma }_{jk}^2}`$.
| $`R_{11}^2`$ | = | $`3.54`$ | $`R_{13}^2`$ | = | $`2.06`$ |
| --- | --- | --- | --- | --- | --- |
| $`R_{22}^2`$ | = | $`1`$ | $`R_{24}^2`$ | = | $`1`$ |
| $`R_{33}^2`$ | = | $`1.30`$ | $`R_{44}^2`$ | = | $`1`$ |
Table 2: Universal ratios $`(R_c)_{jk}^1`$ and $`(R_c)_{jk}^2`$.
| $`(R_c)_{22}^1`$ | = | $`1.0510^2`$ | $`(R_c)_{23}^1`$ | = | $`4.8510^2`$ |
| --- | --- | --- | --- | --- | --- |
| $`(R_c)_{24}^1`$ | = | $`6.710^2`$ | $`(R_c)_{33}^1`$ | = | $`3.810^1`$ |
| $`(R_c)_{11}^2`$ | = | $`2.010^3`$ | $`(R_c)_{14}^2`$ | = | $`2.3410^2`$ |
| $`(R_c)_{13}^2`$ | = | $`1.7910^2`$ | $`(R_c)_{33}^2`$ | = | $`3.410^1`$ |
Table 3: Universal ratio $`(R_\chi )_j^i`$ for $`i,j=1,2`$.
| $`(R_\chi )_1^1`$ | = | $`3.89710^2`$ | $`(R_\chi )_2^{2+}`$ | = | $`0.1111`$ |
| --- | --- | --- | --- | --- | --- |
| $`(R_\chi )_2^1`$ | = | $`0.116`$ | $`(R_\chi )_1^2`$ | = | $`0.040`$ |
| $`(R_\chi )_1^{2+}`$ | = | $`0`$ | $`(R_\chi )_2^2`$ | = | $`0.1111`$ |
Table 4: Universal ratios $`R_\xi ^i`$ and $`(R_A)_j^i`$ for $`i,j=1,2^{},2^+`$.
| $`R_\xi ^1`$ | = | $`7.55710^2`$ | | | |
| --- | --- | --- | --- | --- | --- |
| $`R_\xi ^{2+}`$ | = | $`1.078410^1`$ | $`R_\xi ^2`$ | = | $`8.38910^1`$ |
| $`(R_A)_{2+}^1`$ | = | $`0`$ | $`(R_A)_2^1`$ | = | $`3.91810^2`$ |
| $`(R_A)_1^{2+}`$ | = | $`2.95810^1`$ | $`(R_A)_1^2`$ | = | $`8.26010^1`$ |
Table 5: Universal ratios $`(Q_2)_{jk}^i`$ for $`i,j,k=1,2^+,2^{}`$.
| $`(Q_2)_{2^+1}^1`$ | = | $`1.260`$ | $`(Q_2)_{2^{}1}^1`$ | = | $`1.884`$ |
| --- | --- | --- | --- | --- | --- |
| $`(Q_2)_{2^+2^+}^1`$ | = | $`1.973`$ | $`(Q_2)_{2^+2^{}}^1`$ | = | $`1.320`$ |
| $`(Q_2)_{11}^{2+}`$ | = | $`1.56`$ | $`(Q_2)_{11}^2`$ | = | $`0.442`$ |
| $`(Q_2)_{12^{}}^{2+}`$ | = | $`1.70`$ | | | | |
warning/0002/astro-ph0002197.html | ar5iv | text | # First VLT spectra of white dwarfs in a globular cluster
## 1 Introduction
White dwarfs are the final stage of all low-mass stars and therefore all single stars in a globular cluster that currently finish their nuclear-burning lifetimes are expected to evolve into white dwarfs. As this has been the situation for many billions of years globular clusters should contain many white dwarfs. However, these stars managed to evade detection until photometric white dwarf sequences in globular clusters were discovered recently by observations with the Hubble Space Telescope (HST) (Paresce et al. pade95 (1995), Richer et al. rifa95 (1995), rifa97 (1997), Cool et al. copi96 (1996), Renzini et al. rebr96 (1996)). Photometric observations contain only a limited amount of information: The two chemically distinct white dwarfs sequences (hydrogen-rich DAโs and helium-rich DBโs) in principle can be distinguished by their photometric properties alone in the temperature range $`10,000KT_{\mathrm{eff}}15,000`$ K (see Bergeron et al. 1995a ). Renzini et al. (rebr96 (1996)) classified two white dwarfs in NGC 6752 as DBโs by this method. However, without a spectral classification, both stars can also be explained as high mass DA white dwarfs, possibly a product of merging. Richer et al. (rifa97 (1997)) speculate that the brightest white dwarf in M 4 (V=22.08) might be a hot (27,000K) DB star.
The location of the white dwarf cooling sequence (and thus the brightness of the white dwarfs) is also sensitive to the white dwarf mass. Renzini et al. (rebr96 (1996)) argued that the white dwarf masses in globular clusters are constrained to the narrow range 0.51$`\mathrm{M}_{}`$$`\mathrm{M}_{\mathrm{WD}}`$ 0.55$`\mathrm{M}_{}`$, but some systematic differences between clusters are obvious: At a given metallicity some globular clusters (e.g. NGC 6752) possess very blue horizontal branches (HBโs) with HB star masses as low as 0.50$`\mathrm{M}_{}`$. Such extreme HB stars evolve directly to low-mass C/O white dwarfs (bypassing the AGB), shifting the mean white dwarf mass closer to 0.51$`\mathrm{M}_{}`$. Other clusters show only red HB stars, which will evolve to the AGB and form preferably white dwarfs with masses of $``$0.55$`\mathrm{M}_{}`$.
Low-mass white dwarfs (M$`<`$0.45$`\mathrm{M}_{}`$) with a degenerate He core are produced if the red giant branch evolution is terminated by binary interaction before the helium core exceeds the minimum mass for helium burning. Recently, Cool et al. (cogr98 (1998)) found 3 faint UV-bright stars in NGC 6397 which they suggest could be He white dwarfs (supported by Edmonds et al. edgr99 (1999)). Massive white dwarfs may be produced from blue stragglers or by collisions of white dwarf-binaries with subsequent merging (e.g. Marsh et al. madh95 (1995)).
Only a detailed spectroscopic investigation can provide masses and absolute luminosities of the individual globular cluster white dwarfs. This is also very important for the use of white dwarfs as standard candles to derive distances to globular clusters (Renzini et al. rebr96 (1996)): The basic idea is to fit the white dwarf cooling sequence of a globular cluster to an appropriate empirical cooling sequence of local white dwarfs with well determined trigonometric parallaxes. The procedure is analogous to the classical main sequence fitting but avoids the complications with metallicity โ white dwarfs have virtually metal free atmospheres. In addition they are locally much more abundant than metal-poor subdwarfs. The arrival of the Hipparcos results as well as new metallicity determinations have rekindled the debate on globular cluster distances (see the review by Reid reid99 (1999) and references therein). A further check on the distance is therefore urgently needed.
We started an observing programme at the ESO Very Large Telescope (VLT) to obtain spectra of white dwarfs in globular clusters. The programme consists of two parts: First, low S/N ($`10`$) spectra of the white dwarf candidates are obtained to verify their spectral type and estimate their effective temperatures. In a second run we plan to observe higher S/N ($`30`$) spectra that will allow to derive $`\mathrm{log}g`$ with an internal error of $`0.1`$dex. Here we report on the very first results for NGC 6397.
## 2 Observations and Data Reduction
Cool et al. (copi96 (1996)) discovered the white dwarfs using the Wide Field and Planetary Camera 2 (WFPC2) onboard the HST to observe the globular cluster NGC 6397. From the improved colour-magnitude diagram of King et al. (kian98 (1998)) targets brighter than $`V25^\mathrm{m}`$ were selected (see Fig. 1). The WFPC2 images were convolved to a seeing of 0$`.\mathrm{}`$5 to select targets that are sufficiently uncrowded to be observable from the ground (see Table 1 and Fig. 2). The stars were observed with the FOcal Reducer/low dispersion Spectrograph (FORS) at Unit Telescope 1 of the VLT using the high resolution collimator (0$`.\mathrm{}`$1/pixel) to allow better extraction of the spectra and get a better handle on cosmic rays. We used the multi-object spectroscopy (MOS) mode with the grism 300V and a slit width of 0$`.\mathrm{}`$8. The slit width was chosen to be larger than the required seeing to avoid slit losses due to imperfect pointing of the telescope. The data were obtained in service mode under excellent conditions (seeing below 0$`.\mathrm{}`$55, no moon) with a total exposure time of 90 minutes. The final resolution as judged from a wavelength calibration spectrum obtained with a 0$`.\mathrm{}`$5 slit is $``$11.5 ร
. A trace along the spatial axis of the slitlets at about 4550 ร
is plotted in Fig. 2. Unfortunately WF4-205 lies so close to a bright star that even at this excellent seeing its spectrum could not be extracted.
Due to the use of slit blades instead of fibers or masks the MOS slitlets are very well defined and can be treated like long slits. The spectra were therefore corrected for bias, flat-fielded, wavelength calibrated, and extracted as described by Moehler et al. (mohe97 (1997)). We find only a diffuse and rather low sky background without any strong sky lines below 5150 ร
. The spectra were relatively flux calibrated using LTT 7987 (Hamuy et al. hawa92 (1992)) and are plotted in Fig. 3. All four stars display only strong broad Balmer lines, which is characteristic for hydrogen-rich white dwarfs (DA stars).
## 3 Atmospheric parameters
Although the white dwarf spectra have low signal-to-noise, they are sufficient for rough parameter estimates. The atmospheric parameters are obtained by simultaneously fitting profiles of the observed Balmer lines with model spectra using the least-square algorithm of Bergeron et al. (besa92 (1992); see Napiwotzki et al. 1999 for minor modifications). Analyses were performed with Koesterโs LTE models as described in Finley et al. (fiko97 (1997)). As a check we repeated the analysis of the hottest star in our sample (WF4-358) with the non-LTE grid described in Napiwotzki et al. (nagr99 (1999)). Since the non-LTE code does not treat convection and ignores molecular opacities reliable atmospheric models cannot be calculated for the three cooler white dwarfs.
Fitting the lines H<sub>ฮฒ</sub> to H<sub>ฯต</sub> for WF4-358 (see Fig. 4) gives 18,200$`\pm `$1300 K and 7.30$`\pm `$0.36 for $`T_{\mathrm{eff}}`$ and $`\mathrm{log}g`$, respectively ($`\chi ^2`$ = 0.93). The errors given here are 1$`\sigma `$ errors obtained from the $`\chi ^2`$ fit. Omitting H<sub>ฯต</sub> from the fit results in 17,800 K and 7.19 ($`\chi ^2`$ = 1.02) with more or less unchanged errors. The results of the non-LTE analysis are essentially identical to those obtained with Koesterโs models, differing only by small fractions of the formal errors ($`\mathrm{\Delta }`$$`T_{\mathrm{eff}}`$$``$500 K, $`\mathrm{\Delta }`$$`\mathrm{log}g`$$``$ 0.07 dex). The surface gravity is surprisingly low and suggests that WF4-358 could be a bright ($`M_V`$=$`9\stackrel{\mathrm{m}}{.}7`$) helium white dwarf of (0.36$`\pm `$0.12)$`\mathrm{M}_{}`$. Within the error bars, however, the derived parameters are also consistent with a low-mass C/O white dwarf. For the remaining three stars the S/N is too low to determine $`T_{\mathrm{eff}}`$ and $`\mathrm{log}g`$ simultaneously. We thus fitted H<sub>ฮด</sub>, H<sub>ฮณ</sub>, H<sub>ฮฒ</sub> (H<sub>ฯต</sub> being too noisy) for WF2-51, WF2-479, and WF2-846 for three fixed values of $`\mathrm{log}g`$ (8.0, 7.7, 7.5, see Fig. 4 for an example). These $`\mathrm{log}g`$ values correspond to C/O white dwarfs of $``$ 0.6$`\mathrm{M}_{}`$, low-mass C/O white dwarfs of $``$0.5$`\mathrm{M}_{}`$, and He white dwarfs of $``$0.4$`\mathrm{M}_{}`$, respectively (see below). The formal errors are $``$550 K (WF2-479), $``$650 K (WF2-846, WF4-358), and $``$790 K (WF2-51). The errors for the cooler stars are relatively small despite their low S/N because โ at fixed $`\mathrm{log}g`$ โ the line profiles are much more sensitive to temperature variations at $`T_{\mathrm{eff}}`$$``$11,000 K than at $`T_{\mathrm{eff}}`$$``$18,000 K. The relatively large $`\chi ^2`$ value for WF2-51 suggests that either the noise in this spectrum has been underestimated or that the spectrum contains additional features that are not well described by the model spectra. The temperatures derived from the Balmer lines agree quite well with those obtained from $`(VI)_0`$ using the theoretical colours of Bergeron et al. (1995a , $`\mathrm{log}g`$ = 8.0).
The masses given in Table 2 were derived by interpolation between the evolutionary tracks of C/O white dwarfs calculated by Blรถcker (bloe95 (1995)) and the He white dwarf tracks of Driebe et al. (drsch98 (1998)). Finally, absolute magnitudes $`M_V`$ were calculated for each parameter set with the photometric calibration of Bergeron et al. (1995a ).
## 4 The distance to NGC 6397
The old distance modulus to NGC 6397 was $`(mM)_0`$ = $`11\stackrel{\mathrm{m}}{.}71`$ with a reddening of $`E_{BV}`$ of $`0\stackrel{\mathrm{m}}{.}18`$ (Djorgovski djor93 (1993)). Using local metal-poor subdwarfs to fit the main sequence of NGC 6397 Reid & Gizis (regi98 (1998)) obtained a mean distance modulus of $`(mM)_0`$ = $`12\stackrel{\mathrm{m}}{.}20`$$`\pm `$$`0\stackrel{\mathrm{m}}{.}15`$ for $`E_{BV}`$ = $`0\stackrel{\mathrm{m}}{.}19`$. Thus NGC 6397 is a good example for the large differences between old and new distances to globular clusters. The absolute magnitudes given in Table 2 for $`\mathrm{log}g`$ = 7.5, 7.7, 8.0 (M<sub>WD</sub> = 0.4$`\mathrm{M}_{}`$, 0.5$`\mathrm{M}_{}`$, 0.6$`\mathrm{M}_{}`$) yield mean true distance moduli (for $`E_{BV}`$ = $`0\stackrel{\mathrm{m}}{.}18`$) $`(mM)_0`$ of $`12\stackrel{\mathrm{m}}{.}3`$, $`12\stackrel{\mathrm{m}}{.}0`$, and $`11\stackrel{\mathrm{m}}{.}6`$, respectively, with an r.m.s. error of $`0\stackrel{\mathrm{m}}{.}17`$. Considering the error bars of the various distance determinations the long distance scale would be more consistent with white dwarf masses $``$0.5$`\mathrm{M}_{}`$ and the short distance scale with masses $`>`$0.5$`\mathrm{M}_{}`$. The longer distance moduli obtained for low-mass white dwarfs also result in masses for blue HB stars (Heber et al. hemo97 (1997)) and a hot post-AGB star (ROB 162, Heber & Kudritzki heku86 (1986)) that agree with canonical evolutionary theory.
The distance moduli derived for a given $`\mathrm{log}g`$ from Tables 1 and 2 show a systematic variation with the brightest star (WF4-358) yielding the smallest distance and the faintest star (WF2-846) giving the largest distance. This variation could reflect the fact that the stars may not all have the same surface gravity: From their different apparent magnitudes (i.e. different absolute magnitudes) it is plausible that WF4-358 has the smallest $`\mathrm{log}g`$ and WF2-846 the largest. The quality of the current data, however, does not allow to verify this idea.
## 5 Conclusions
Using VLT-FORS1 multi object spectroscopy we have confirmed four white dwarf candidates to be hydrogen-rich DA white dwarfs. The gravity determined for the brightest star, WF4-358, suggests that it could be a He white dwarf with a mass of (0.36$`\pm `$0.12)$`\mathrm{M}_{}`$, although the error bars are large enough to also accommodate a C/O white dwarf. Temperatures derived for the three cooler and fainter stars for fixed $`\mathrm{log}g`$ would put them near the red edge of the ZZ Ceti instability strip, for which Bergeron et al. (1995b ) determine a temperature range of 11,160 K to 12,460 K using their preferred ML2/$`\alpha `$=0.6 prescription for the treatment of convection. Therefore, a search for photometric variability of WF2-479 and WF2-51 โ if successful โ could place important additional constraints on these stars. The systematic variation of the distance moduli derived for a given $`\mathrm{log}g`$ shows that the assumption of a constant mass for all white dwarfs in a globular cluster may bias a distance determination. However, due to the low S/N of the current data, these results are preliminary. Once higher quality spectra are available, which will allow more accurate parameter (and thus mass) determinations, analyses of white dwarfs in globular clusters will become a powerful tool for independent distance estimates.
###### Acknowledgements.
We highly appreciate the work performed by the FORS team in building an excellent instrument and the efforts of the staff at the ESO Paranal observatory and ESO Garching that made these observations possible. We thank Dr. A. Cool for providing us with the photometry and coordinates of the white dwarf candidates and an anonymous referee for valuable comments. S.M. acknowledges financial support from the DARA under grant 50 OR 96029-ZA. |
warning/0002/gr-qc0002059.html | ar5iv | text | # References
I. Introduction
Hamiltonian formulations, when consistently established, not only guarantee that field quantities have a well defined time evolution, but also allow us to understand physical theories from a different perspective. We have learned from the work of Arnowitt, Deser and Misner (ADM) that the Hamiltonian analysis of Einsteinโs general relativity reveals the intrinsic structure of the theory: the time evolution of field quantities is determined by the Hamiltonian and vector constraints. Thus four of the ten Einsteinโs equations acquire a prominent status in the Hamiltonian framework. Ultimately this is an essential feature for the canonical approach to the quantum theory of gravity.
It is the case in general relativity that two distinct Lagrangian formulations that yield Einsteinโs equations lead to completely different Hamiltonian constructions. An important example in this respect is the reformulation of the ordinary variational principle, based on the Hilbert-Einstein action, in terms of self-dual connections that define Ashtekar variables. Under a Palatini type variation of the action integral constructed out of these field quantities one obtains precisely Einsteinโs equations. Interesting features of this approach reside in the Hamiltonian domain.
Einsteinโs general relativity can also be reformulated in the context of the teleparallel (Weitzenbรถck) geometry. In this geometrical setting the dynamical field quantities correspond to orthornormal tetrad fields $`e_\mu ^a`$ ($`a,\mu `$ are SO(3,1) and space-time indices, respectively). These fields allow the construction of the Lagrangian density of the teleparallel equivalent of general relativity (TEGR) , which offers an alternative geometrical framework for Einsteinโs equations. The Lagrangian density for the tetrad field in the TEGR is given by a sum of quadratic terms in the torsion tensor $`T_{\mu \nu }^a=_\mu e_\nu ^a_\nu e_\mu ^a`$, which is related to the anti-symmetric part of Cartanโs connection $`\mathrm{\Gamma }_{\mu \nu }^\lambda =e^{a\lambda }_\mu e_{a\nu }`$. The curvature tensor constructed out of the latter vanishes identically. This connection defines a space with teleparallelism, or absolute parallelism.
In a space-time with an underlying tetrad field two vectors at distant points are called parallel if they have identical components with respect to the local tetrads at the points considered. Thus consider a vector field $`V^\mu (x)`$. At the point $`x^\lambda `$ its tetrad components are given by $`V^a(x)=e_\mu ^a(x)V^\mu (x)`$. For the tetrad components $`V^a(x+dx)`$ it is easy to show that $`V^a(x+dx)=V^a(x)+DV^a(x)`$, where $`DV^a(x)=e_\mu ^a(_\lambda V^\mu )dx^\lambda `$. The covariant derivative $``$ is constructed out of Cartanโs connection $`\mathrm{\Gamma }_{\mu \nu }^\lambda =e^{a\lambda }_\mu e_{a\nu }`$. Therefore the vanishing of such covariant derivative defines a condition for absolute parallelism in space-time. Hence in the teleparallel geometry tetrad fields transform under the global SO(3,1) group. Teleparallel geometry is less restrictive than Riemannian geometry. For a given Riemaniann geometry there are many ways to construct the teleparallel geometry, since one Riemaniann geometry corresponds to a whole equivalence class of teleparallel geometries.
In the framework of the TEGR it is possible to make definite statements about the energy and momentum of the gravitational field. This fact constitutes the major motivation for considering this theory. In the 3+1 formulation of the TEGR, and by imposing Schwingerโs time gauge condition, we find that the Hamiltonian and vector constraints contain each one a divergence in the form of scalar and vector densities, respectively, that can be identified with the energy and momentum densities of the gravitational field.
In this paper we carry out the Hamiltonian formulation of the TEGR without imposing the time gauge condition, by rigorously performing the Legendre transform. We have not found it necessary to establish a 3+1 decomposition for the tetrad field. We only assume $`g^{00}0`$, a condition that ensures that $`t=constant`$ hypersurfaces are spacelike. The Lagrange multipliers are given by the zero components of the tetrads, $`e_{a0}`$. The constraints corresponding to the Hamiltonian ($`H_0`$) and vector ($`H_i`$) constraints are obtained in the form $`C^a=0`$. The dynamical evolution of the field quantities is completely determined by $`H_0`$ and by a set of primary constraints $`\mathrm{\Gamma }^{ik}`$ and $`\mathrm{\Gamma }^k`$, as we will show. The surprising feature is that if $`H_0=0`$ in the subspace of the phase space determined by $`\mathrm{\Gamma }^{ik}=\mathrm{\Gamma }^k=0`$, then it follows that $`H_i=0`$. As we will see, $`H_i`$ can be obtained from the very definition of $`H_0`$. Furthermore by calculating Poisson brackets we show that the constraints constitute a first class set. Hence the theory is well defined regarding time evolution.
As a consequence of this analysis, we arrive at a scalar density that transforms as a four-vector in the SO(3,1) space, again arising in the expression of the constraints of the theory, and whose zero component is related to the energy of the gravitational field. In analogy with previous investigations, we interpret the constraint equations $`C^a=0`$ as energy-momentum equations for the gravitational field.
The analysis developed here is similar to that developed in Ref. , in which the Hamiltonian formulation of the TEGR in null surfaces was established. The 3+1 formulation of the TEGR has already been considered in Ref. . There are several differences between the latter and the present analysis. The investigation in Ref. has not pointed out neither the emergence of the scalar densities mentioned above nor the relationship between $`H_0`$ and $`H_i`$. Our approach is different and allowed us to proceed further in the understanding of the constraint structure of the theory.
Notation: spacetime indices $`\mu ,\nu ,\mathrm{}`$ and SO(3,1) indices $`a,b,\mathrm{}`$ run from 0 to 3. Time and space indices are indicated according to $`\mu =0,i,a=(0),(i)`$. The tetrad field $`e_\mu ^a`$ yields the definition of the torsion tensor: $`T_{\mu \nu }^a=_\mu e_\nu ^a_\nu e_\mu ^a`$. The flat, Minkowski spacetime metric is fixed by $`\eta _{ab}=e_{a\mu }e_{b\nu }g^{\mu \nu }=(+++)`$.
II. Lagrangian formulation
In order to carry out the 3+1 decomposition we need a first order differential formulation of the Lagrangian density of the TEGR. For this purpose we introduce an auxiliary field quantity $`\varphi _{abc}=\varphi _{acb}`$ that will be related to the torsion tensor. The first order differential Lagrangian formulation in empty space-time reads
$$L(e,\varphi )=ke\mathrm{\Lambda }^{abc}(\varphi _{abc}\mathrm{\hspace{0.17em}2}T_{abc}),$$
$`(1)`$
where $`T_{abc}=e_b^\mu e_c^\nu T_{a\mu \nu }`$. $`\mathrm{\Lambda }^{abc}`$ is defined by
$$\mathrm{\Lambda }^{abc}=\frac{1}{4}(\varphi ^{abc}+\varphi ^{bac}\varphi ^{cab})+\frac{1}{2}(\eta ^{ac}\varphi ^b\eta ^{ab}\varphi ^c),$$
$`(2)`$
and $`\varphi _b=\varphi _{ab}^a`$. The Lagrangian density (1) is invariant under coordinate and global SO(3,1) transformations.
Variation of the action constructed out of (1) with respect to $`\varphi ^{abc}`$ yields an equation that can be reduced to $`\varphi _{abc}=T_{abc}`$. This equation can be split into two equations:
$$\varphi _{a0k}=T_{a0k}=_0e_{ak}_ke_{a0},$$
$`(3a)`$
$$\varphi _{aik}=T_{aik}=_ie_{ak}_ke_{ai}.$$
$`(3b)`$
The variation of the action integral with respect to $`e_{a\mu }`$ yields the field equation
$$\frac{\delta L}{\delta e^{a\mu }}=e_{a\lambda }e_{b\mu }_\nu (e\mathrm{\Sigma }^{b\lambda \nu })e\left(\mathrm{\Sigma }_a^{b\nu }T_{b\nu \mu }\frac{1}{4}e_{a\mu }T_{bcd}\mathrm{\Sigma }^{bcd}\right)=\mathrm{\hspace{0.33em}0}.$$
$`(4)`$
The tensor $`\mathrm{\Sigma }^{abc}`$ is defined in terms of $`T^{abc}`$ exactly like $`\mathrm{\Lambda }^{abc}`$ in terms of $`\varphi ^{abc}`$. By explicit calculations it is verified that these equations are equivalent to Einsteinโs equations in tetrad form:
$$\frac{\delta L}{\delta e^{a\mu }}\frac{1}{2}e\left\{R_{a\mu }(e)\frac{1}{2}e_{a\mu }R(e)\right\}.$$
We note finally that by substituting (3a,b) into (1) the Lagrangian density reduces to
$$L(e_{a\mu })=ke\mathrm{\Sigma }^{abc}T_{abc}=ke(\frac{1}{4}T^{abc}T_{abc}+\frac{1}{2}T^{abc}T_{bac}T^aT_a).$$
III. Legendre transform and the 3+1 decomposition
The Hamiltonian density will be obtained by the standard prescription $`L=p\dot{q}H_0`$ and by properly identifying primary constraints. We have not found it necessary to establish any kind of 3+1 decomposition for the tetrad fields. Therefore in the following both $`e_{a\mu }`$ and $`g_{\mu \nu }`$ are space-time fields. We will follow here the procedure presented in .
Lagrangian density (1) can be expressed as
$$L(e,\varphi )=4ke\mathrm{\Lambda }^{a0k}\dot{e}_{ak}+4ke\mathrm{\Lambda }^{a0k}_ke_{a0}2ke\mathrm{\Lambda }^{aij}T_{aij}+ke\mathrm{\Lambda }^{abc}\varphi _{abc},$$
$`(5)`$
where the dot indicates time derivative, and $`\mathrm{\Lambda }^{a0k}=\mathrm{\Lambda }^{abc}e_b^0e_c^k`$, $`\mathrm{\Lambda }^{aij}=\mathrm{\Lambda }^{abc}e_b^ie_c^j`$.
Therefore the momentum canonically conjugated to $`e_{ak}`$ is given by
$$\mathrm{\Pi }^{ak}=4ke\mathrm{\Lambda }^{a0k},$$
$`(6)`$
In terms of (6) expression (5) reads
$$L=\mathrm{\Pi }^{ak}\dot{e}_{ak}\mathrm{\Pi }^{ak}_ke_{a0}2ke\mathrm{\Lambda }^{aij}T_{aij}+ke\mathrm{\Lambda }^{abc}\varphi _{abc}$$
$$=\mathrm{\Pi }^{ak}\dot{e}_{ak}\mathrm{\Pi }^{ak}_ke_{a0}ke\mathrm{\Lambda }^{aij}(2T_{aij}\varphi _{aij})+2ke\mathrm{\Lambda }^{a0k}\varphi _{a0k}.$$
$`(7)`$
The last term on the right hand side of equation (7) is identified as $`2ke\mathrm{\Lambda }^{a0k}\varphi _{a0k}=\frac{1}{2}\mathrm{\Pi }^{ak}\varphi _{a0k}`$.
The Hamiltonian formulation is established once we rewrite the Lagrangian density (7) in terms of $`e_{ak}`$, $`\mathrm{\Pi }^{ak}`$ and further nondynamical field quantities. It is carried out in two steps. First, we take into account equation (3b) in (7) so that half of the auxiliary fields, $`\varphi _{aij}`$, are eliminated from the Lagrangian by means of the identification
$$\varphi _{aij}=T_{aij}.$$
As a consequence we have
$$ke\mathrm{\Lambda }^{aij}(2T_{aij}\varphi _{aij})=ke\mathrm{\Lambda }^{aij}T_{aij}$$
$$=ke\left(\frac{1}{4}g^{im}g^{nj}T_{mn}^aT_{aij}+\frac{1}{2}g^{nj}T_{mn}^iT_{ij}^mg^{ik}T_{ji}^jT_{nk}^n\right)$$
$$+ke\left(\frac{1}{2}g^{0i}g^{jk}\varphi _{0k}^aT_{aij}\frac{1}{2}g^{jk}\varphi _{0k}^iT_{ij}^0+\frac{1}{2}g^{0j}\varphi _{0k}^iT_{ij}^kg^{0k}\varphi _{0j}^jT_{ik}^i+g^{ik}\varphi _{0i}^0T_{jk}^j\right).$$
The last five terms of the expression above may be rewritten as
$$\frac{1}{2}ke\varphi _{a0k}\left[g^{0i}g^{kj}T_{ij}^ae^{ai}(g^{0j}T_{ij}^kg^{kj}T_{ij}^0)+2(e^{ak}g^{0i}e^{a0}g^{ki})T_{ji}^j\right].$$
Therefore we have
$$L(e_{ak},\mathrm{\Pi }^{ak},e_{a0},\varphi _{a0k})=\mathrm{\Pi }^{ak}\dot{e}_{ak}+e_{a0}_k\mathrm{\Pi }^{ak}_k(e_{a0}\mathrm{\Pi }^{ak})$$
$$ke\left(\frac{1}{4}g^{im}g^{nj}T_{mn}^aT_{aij}+\frac{1}{2}g^{nj}T_{mn}^iT_{ij}^mg^{ik}T_{ji}^jT_{nk}^n\right)$$
$$\frac{1}{2}\varphi _{a0k}\left\{\mathrm{\Pi }^{ak}+ke\left[g^{0i}g^{kj}T_{ij}^ae^{ai}(g^{0j}T_{ij}^kg^{kj}T_{ij}^0)+2(e^{ak}g^{0i}e^{a0}g^{ki})T_{ji}^j\right]\right\}.$$
$`(8)`$
The second step consists of expressing the remaining auxiliary field quantities, the โvelocitiesโ $`\varphi _{a0k}`$, in terms of the momenta $`\mathrm{\Pi }^{ak}`$. This is the nontrivial step of the Legendre transform.
We need to consider the full expression of $`\mathrm{\Pi }^{ak}`$. It is given by equation (6),
$$\mathrm{\Pi }^{ak}=ke\{g^{00}(g^{kj}\varphi _{0j}^ae^{aj}\varphi _{0j}^k+2e^{ak}\varphi _{0j}^j)$$
$$+g^{0k}(g^{0j}\varphi _{0j}^a+e^{aj}\varphi _{0j}^0)+e^{a0}(g^{0j}\varphi _{0j}^k+g^{kj}\varphi _{0j}^0)2(e^{a0}g^{0k}\varphi _{0j}^j+e^{ak}g^{0j}\varphi _{0j}^0)$$
$$g^{0i}g^{kj}T_{ij}^a+e^{ai}(g^{0j}T_{ij}^kg^{kj}T_{ij}^0)2(g^{0i}e^{ak}g^{ik}e^{a0})T_{ji}^j\},$$
$`(9)`$
where we have already identified $`\varphi _{aij}=T_{aij}`$. Denoting $`(..)`$ and $`[..]`$ as the symmetric and anti-symmetric parts of field quantities, respectively, we decompose $`\mathrm{\Pi }^{ak}`$ into irreducible components:
$$\mathrm{\Pi }^{ak}=e_i^a\mathrm{\Pi }^{(ik)}+e_i^a\mathrm{\Pi }^{[ik]}+e_0^a\mathrm{\Pi }^{0k},$$
$`(10)`$
where
$$\mathrm{\Pi }^{(ik)}=ke\{g^{00}(g^{kj}\varphi _{0j}^ig^{ij}\varphi _{0j}^k+2g^{ik}\varphi _{0j}^j)+g^{0k}(g^{0j}\varphi _{0j}^i+g^{ij}\varphi _{0j}^0g^{0i}\varphi _{0j}^j)$$
$$+g^{0i}(g^{0j}\varphi _{0j}^k+g^{kj}\varphi _{0j}^0g^{0k}\varphi _{0j}^j)2g^{ik}g^{0j}\varphi _{0j}^0+\mathrm{\Delta }^{ik}\},$$
$`(11a)`$
$$\mathrm{\Delta }^{ik}=g^{0m}(g^{kj}T_{mj}^i+g^{ij}T_{mj}^k2g^{ik}T_{mj}^j)(g^{km}g^{0i}+g^{im}g^{0k})T_{mj}^j,$$
$`(11b)`$
$$\mathrm{\Pi }^{[ik]}=ke\left\{g^{im}g^{kj}T_{mj}^0+(g^{im}g^{0k}g^{km}g^{0i})T_{mj}^j\right\},$$
$`(12)`$
$$\mathrm{\Pi }^{0k}=2ke(g^{kj}g^{0i}T_{ij}^0g^{0k}g^{0i}T_{ij}^j+g^{00}g^{ik}T_{ij}^j).$$
$`(13)`$
The crucial point in this analysis is that only the symmetrical components $`\mathrm{\Pi }^{(ij)}`$ depend on the โvelocitiesโ $`\varphi _{a0k}`$. The other six components, $`\mathrm{\Pi }^{[ij]}`$ and $`\mathrm{\Pi }^{0k}`$ depend solely on $`T_{aij}`$. Therefore we can express only six of the โvelocityโ fields $`\varphi _{a0k}`$ in terms of the components $`\mathrm{\Pi }^{(ij)}`$. With the purpose of finding out which components of $`\varphi _{a0k}`$ can be inverted in terms of the momenta we decompose $`\varphi _{a0k}`$ identically as
$$\varphi _{0j}^a=e^{ai}\psi _{ij}+e^{ai}\sigma _{ij}+e^{a0}\lambda _j,$$
$`(14)`$
where $`\psi _{ij}=\frac{1}{2}(\varphi _{i0j}+\varphi _{j0i})`$, $`\sigma _{ij}=\frac{1}{2}(\varphi _{i0j}\varphi _{j0i})`$, $`\lambda _j=\varphi _{00j}`$, and $`\varphi _{\mu 0j}=e_\mu ^a\varphi _{a0j}`$ (like $`\varphi _{abc}`$, the components $`\psi _{ij}`$, $`\sigma _{ij}`$ and $`\lambda _j`$ are also auxiliary field quantities). Next we substitute (14) in (11a). By defining
$$P^{ik}=\frac{1}{ke}\mathrm{\Pi }^{(ik)}\mathrm{\Delta }^{ik},$$
$`(15)`$
we find that $`P^{ik}`$ depends only on $`\psi _{ij}`$:
$$P^{ik}=2g^{00}(g^{im}g^{kj}\psi _{mj}g^{ik}\psi )$$
$$+2(g^{0i}g^{km}g^{0j}+g^{0k}g^{im}g^{0j})\psi _{mj}2(g^{ik}g^{0m}g^{0j}\psi _{mj}+g^{0i}g^{0k}\psi ),$$
$`(16)`$
where $`\psi =g^{mn}\psi _{mn}`$.
We can now invert $`\psi _{mj}`$ in terms of $`P^{ik}`$. After a number of manipulations we arrive at
$$\psi _{mj}=\frac{1}{2g^{00}}\left(g_{im}g_{kj}P^{ik}\frac{1}{2}g_{mj}P\right),$$
$`(17)`$
where $`P=g_{ik}P^{ik}`$.
At last we need to rewrite the third line of the Lagrangian density (8) in terms of canonical variables. By making use of (9), (14) and (17) we can rewrite
$$\frac{1}{2}\varphi _{a0k}\left\{\mathrm{\Pi }^{ak}+ke\left[g^{0i}g^{kj}T_{ij}^ae^{ai}(g^{0j}T_{ij}^kg^{kj}T_{ij}^0)+2(e^{ak}g^{0i}e^{a0}g^{ki})T_{ji}^j\right]\right\}$$
in the form
$$\frac{1}{4g^{00}}ke\left(g_{ik}g_{jl}P^{ij}P^{kl}\frac{1}{2}P^2\right).$$
Thus we finally obtain the primary Hamiltonian density $`H_0=\mathrm{\Pi }^{ak}\dot{e}_{ak}L`$,
$$H_0(e_{ak},\mathrm{\Pi }^{ak},e_{a0})=e_{a0}_k\mathrm{\Pi }^{ak}\frac{1}{4g^{00}}ke\left(g_{ik}g_{jl}P^{ij}P^{kl}\frac{1}{2}P^2\right)$$
$$+ke\left(\frac{1}{4}g^{im}g^{nj}T_{mn}^aT_{aij}+\frac{1}{2}g^{nj}T_{mn}^iT_{ij}^mg^{ik}T_{ji}^jT_{nk}^n\right).$$
$`(18)`$
We may now write the total Hamiltonian density. For this purpose we have to identify the primary constraints. They are given by expressions (12) and (13), which represent relations between $`e_{ak}`$ and the momenta $`\mathrm{\Pi }^{ak}`$. Thus we define
$$\mathrm{\Gamma }^{ik}=\mathrm{\Gamma }^{ki}=\mathrm{\Pi }^{[ik]}ke\left\{g^{im}g^{kj}T_{mj}^0+(g^{im}g^{0k}g^{km}g^{0i})T_{mj}^j\right\},$$
$`(19)`$
$$\mathrm{\Gamma }^k=\mathrm{\Pi }^{0k}+\mathrm{\hspace{0.33em}2}ke(g^{kj}g^{0i}T_{ij}^0g^{0k}g^{0i}T_{ij}^j+g^{00}g^{ik}T_{ij}^j).$$
$`(20)`$
Therefore the total Hamiltonian density is given by
$$H(e_{ak},\mathrm{\Pi }^{ak},e_{a0},\alpha _{ik},\beta _k)=H_0+\alpha _{ik}\mathrm{\Gamma }^{ik}+\beta _k\mathrm{\Gamma }^k+_k(e_{a0}\mathrm{\Pi }^{ak}),$$
$`(21)`$
where $`\alpha _{ik}`$ and $`\beta _k`$ are Lagrange multipliers.
IV. Secondary constraints
Since the momenta $`\{\mathrm{\Pi }^{a0}\}`$ vanish identically they also constitute primary constraints that induce the secondary constraints
$$C^a\frac{\delta H}{\delta e_{a0}}=0.$$
$`(22)`$
In order to obtain the expression of $`C^a`$ we have only to vary $`H_0`$ with respect to $`e_{a0}`$, because variations of $`\mathrm{\Gamma }^{ik}`$ and $`\mathrm{\Gamma }^k`$ with respect to $`e_{a0}`$ yield the constraints themselves:
$$\frac{\delta \mathrm{\Gamma }^{ik}}{\delta e_{a0}}=\frac{1}{2}(e^{ai}\mathrm{\Gamma }^ke^{ak}\mathrm{\Gamma }^i),$$
$`(23a)`$
$$\frac{\delta \mathrm{\Gamma }^k}{\delta e_{a0}}=e^{a0}\mathrm{\Gamma }^k.$$
$`(23b)`$
In (23a,b) we have made use of variations like $`\delta e^{b\mu }/\delta e_{a0}=e^{a\mu }e^{b0}.`$ In the process of obtaining $`C^a`$ we need the variation of $`P^{ij}`$ with respect to $`e_{a0}`$. It reads
$$\frac{\delta P^{ij}}{\delta e_{a0}}=e^{a0}P^{ij}+\gamma ^{aij},$$
with $`\gamma ^{aij}`$ defined by
$$\gamma ^{aij}=\frac{1}{2ke}(e^{ai}\mathrm{\Gamma }^j+e^{aj}\mathrm{\Gamma }^i)e^{ak}[g^{00}(g^{jm}T_{km}^i+g^{im}T_{km}^j+2g^{ij}T_{mk}^m)$$
$$+g^{0m}(g^{0j}T_{mk}^i+g^{0i}T_{mk}^j)2g^{0i}g^{0j}T_{mk}^m+(g^{jm}g^{0i}+g^{im}g^{0j}2g^{ij}g^{0m})T_{mk}^0].$$
$`(24)`$
Note that $`\gamma ^{aij}`$ satisfies $`e_{a0}\gamma ^{aij}=0`$.
After a long calculation we arrive at the expression of $`C^a`$:
$$C^a=_k\mathrm{\Pi }^{ak}+e^{a0}[\frac{1}{4g^{00}}ke(g_{ik}g_{jl}P^{ij}P^{kl}\frac{1}{2}P^2)$$
$$+ke(\frac{1}{4}g^{im}g^{nj}T_{mn}^bT_{bij}+\frac{1}{2}g^{nj}T_{mn}^iT_{ij}^mg^{ik}T_{mi}^mT_{nk}^n)]$$
$$\frac{1}{2g^{00}}ke(g_{ik}g_{jl}\gamma ^{aij}P^{kl}\frac{1}{2}g_{ij}\gamma ^{aij}P)kee^{ai}(g^{0m}g^{nj}T_{ij}^bT_{bmn}$$
$$+g^{nj}T_{mn}^0T_{ij}^m+g^{0j}T_{mj}^nT_{ni}^m2g^{0k}T_{mk}^mT_{ni}^n2g^{jk}T_{ij}^0T_{nk}^n).$$
$`(25)`$
Inspite of the fact that expression above is somehow intricate, we immediately notice that
$$e_{a0}C^a=H_0.$$
$`(26)`$
Therefore the total Hamiltonian becomes
$$H(e_{ak},\mathrm{\Pi }^{ak},e_{a0},\alpha _{ik},\beta _k)=e_{a0}C^a+\alpha _{ik}\mathrm{\Gamma }^{ik}+\beta _k\mathrm{\Gamma }^k+_k(e_{a0}\mathrm{\Pi }^{ak}).$$
$`(27)`$
We observe that $`\{e_{a0}\}`$ arise as Lagrange multipliers (see equation (50) ahead).
Before closing this section we remark that the Hamiltonian formulation described here is different from that developed in Ref. , the difference residing in the definition of the canonical momentum. In the latter reference the canonical momentum is not defined by taking the variation of $`L`$ with respect to $`\dot{e}_{ak}`$. Instead, it is defined by
$$\pi _a^k=\frac{\delta L}{\delta (N^{}T_k^a)}=\frac{\delta L}{\delta (T_{0k}^aN^iT_{ik}^a)},$$
where $`N^{}`$ and $`N^i`$ are the usual lapse and shift functions. As a consequence, three of the six primary constraints of Ref. are different from the corresponding constraints obtained here. The expression of the components $`\tau ^{[ik]}`$ and $`\tau _{}^k`$ of Ref. , equivalent to $`\mathrm{\Pi }^{[ik]}`$ and $`\mathrm{\Pi }^{0k}`$, respectively, given by (12) and (13), read in our notation
$$\tau ^{[ik]}=e\left\{g^{im}g^{kj}T_{ij}^0+N^j(g^{im}g^{0k}g^{km}g^{0i})T_{mj}^0\right\},$$
$$\tau _{}^k=\frac{1}{2k}N^{}\mathrm{\Pi }^{0k}.$$
The Hamiltonian and vector constraints of the above mentioned reference are parametrized in terms of the lapse and shift functions. In the present work we have parametrized the set of four constraints according to equation (26), and identified the Lagrange multipliers as $`e_{a0}`$. The final expression of $`C^a`$ acquires the total divergence $`_k\mathrm{\Pi }^{ak}`$. This divergence is different from the one that appears in the expression of the total Hamiltonian density of gravitational fields for asymptotically flat space-times, either in the metric or in the tetrad formulation (see, for example, Eq. (3.17) of Ref. or Eq. (27) above; it is possible to show that the latter expressions are exactly the same field quantities). We finally notice that the constraint algebra to be presented in the coming section has not been evaluated in Ref. .
V. Simplification of the constraints and Poisson brackets
The first two terms of the expression of $`C^a`$ yield the primary Hamiltonian in the form $`e^{a0}H_0`$. This fact can be easily verified by expressing the first term of (25) as
$$_k\mathrm{\Pi }^{ak}=e^{a0}(e_{b0}_k\mathrm{\Pi }^{bk})+e^{aj}(e_{bj}_k\mathrm{\Pi }^{bk}).$$
The second term considered above is the collection of terms in (25) multiplied by $`e^{a0}`$. Substituting definitions (11b) and (24) for $`\mathrm{\Delta }^{ij}`$ and $`\gamma ^{aij}`$, respectively, into (25) we obtain after a long calculation a simplified form for $`C^a`$,
$$C^a=e^{a0}H_0+e^{ai}F_i,$$
$`(28)`$
with the following definitions:
$$F_i=H_i+\mathrm{\Gamma }^mT_{0mi}+\mathrm{\Gamma }^{lm}T_{lmi}+\frac{1}{2g^{00}}(g_{ik}g_{jl}P^{kl}\frac{1}{2}g_{ij}P)\mathrm{\Gamma }^j,$$
$`(29)`$
$$H_i=e_{bi}_k\mathrm{\Pi }^{bk}\mathrm{\Pi }^{bk}T_{bki}.$$
$`(30)`$
We denote $`H_0`$ the Hamiltonian constraint. $`H_i`$ is the vector constraint. It amounts to a SO(3,1) version of the vector constraint of Ref. . The true constraints of the theory are $`C^a`$, $`\mathrm{\Gamma }^{ik}`$ and $`\mathrm{\Gamma }^k`$. Dispensing with the surface term the total Hamiltonian reads
$$H=e_{a0}C^a+\alpha _{ik}\mathrm{\Gamma }^{ik}+\beta _k\mathrm{\Gamma }^k.$$
$`(31)`$
The Poisson bracket between two quantites $`F`$ and $`G`$ is defined by
$$\{F,G\}=d^3x\left(\frac{\delta F}{\delta e_{ai}(x)}\frac{\delta G}{\delta \mathrm{\Pi }^{ai}(x)}\frac{\delta F}{\delta \mathrm{\Pi }^{ai}(x)}\frac{\delta G}{\delta e_{ai}(x)}\right),$$
by means of which we can write down the evolution equations. The first set of Hamiltonโs equations is given by
$$\dot{e}_{aj}(x)=\{e_{aj}(x),๐\}=d^3y\frac{\delta }{\delta \mathrm{\Pi }^{aj}(x)}\left(H_0(y)+\alpha _{ik}(y)\mathrm{\Gamma }^{ik}(y)+\beta _k(y)\mathrm{\Gamma }^k(y)\right),$$
$`(32)`$
where H is the total Hamiltonian. This equation can be worked to yield
$$T_{a0j}=\frac{1}{2g^{00}}e_a^k(g_{ik}g_{jm}P^{im}\frac{1}{2}g_{kj}P)+e_a^i\alpha _{ij}+e_a^0\beta _j,$$
$`(33)`$
from which we obtain
$$\frac{1}{2}(T_{i0j}+T_{j0i})=\psi _{ij}=\frac{1}{2g^{00}}(g_{ik}g_{mj}P^{km}\frac{1}{2}g_{ij}P),$$
$`(34a)`$
$$\frac{1}{2}(T_{i0j}T_{j0i})=\sigma _{ij}=\alpha _{ij},$$
$`(34b)`$
$$T_{00j}=\lambda _j=\beta _j,$$
$`(34c)`$
according to the definitions in equation (14). Thus the Lagrange multipliers in (31) acquire a well defined meaning. Expression (34a) is in total agreement with (17). Consequently we can obtain an expression for $`\mathrm{\Pi }^{(ij)}`$ in terms of velocities via equations (15) and (16). The dynamical evolution of the field quantities is completed with Hamiltonโs equations for $`\mathrm{\Pi }^{(ij)}`$,
$$\dot{\mathrm{\Pi }}^{(ij)}(x)=\{\mathrm{\Pi }^{(ij)}(x),๐\}=d^3y\left(\frac{\delta \mathrm{\Pi }^{(ij)}(x)}{\delta e_{ak}(y)}\frac{\delta ๐}{\delta \mathrm{\Pi }^{ak}(y)}\frac{\delta \mathrm{\Pi }^{(ij)}(x)}{\delta \mathrm{\Pi }^{ak}(y)}\frac{\delta ๐}{\delta e_{ak}(y)}\right),$$
$`(35)`$
together with
$$\mathrm{\Gamma }^{ik}=\mathrm{\Gamma }^k=0.$$
$`(36)`$
The calculations of the Poisson brackets between these constraints are exceedingly complicated. Here we will just present the results. Instead of considering $`C^a(x)`$ in the calculations below, we found it more appropriate to consider $`H_0(x)`$ and $`H_i(x)`$. The constraint algebra is given by
$$\{H_0(x),H_0(y)\}=0,$$
$`(37)`$
$$\{H_0(x),H_i(y)\}=H_0(x)\frac{}{y^i}\delta (xy)$$
$$H_0e^{a0}_ie_{a0}\delta (xy)F_je^{aj}_ie_{a0}\delta (xy),$$
$`(38)`$
$$\{H_j(x),H_k(y)\}=H_k(x)\frac{}{x^j}\delta (xy)H_j(y)\frac{}{y^k}\delta (xy),$$
$`(39)`$
$$\{\mathrm{\Gamma }^i(x),\mathrm{\Gamma }^j(y)\}=0,$$
$`(40)`$
$$\{\mathrm{\Gamma }^{ij}(x),\mathrm{\Gamma }^k(y)\}=(g^{0j}\mathrm{\Gamma }^{ki}g^{0i}\mathrm{\Gamma }^{kj})\delta (xy),$$
$`(41)`$
$$\{\mathrm{\Gamma }^{ij}(x),\mathrm{\Gamma }^{kl}(y)\}=\frac{1}{2}\left(g^{il}\mathrm{\Gamma }^{jk}+g^{jk}\mathrm{\Gamma }^{il}g^{ik}\mathrm{\Gamma }^{jl}g^{jl}\mathrm{\Gamma }^{ik}\right)\delta (xy),$$
$`(42)`$
$$\{H_0(x),\mathrm{\Gamma }^{ij}(y)\}=[\frac{1}{2g^{00}}P^{kl}(\frac{1}{2}g_{kl}g_{mn}g_{km}g_{nl})(g^{mi}\mathrm{\Gamma }^{nj}g^{mj}\mathrm{\Gamma }^{ni})+$$
$$+\frac{1}{2}(\mathrm{\Gamma }^{nj}e^{ai}\mathrm{\Gamma }^{ni}e^{aj})_ne_{a0}]\delta (xy),$$
$`(43)`$
$$\{H_0(x),\mathrm{\Gamma }^i(y)\}=[g^{0i}H_0+\frac{1}{g^{00}}P^{kl}(\frac{1}{2}g_{kl}g_{jm}g_{kj}g_{ml})g^{0j}\mathrm{\Gamma }^{mi}$$
$$+\left(\mathrm{\Gamma }^{ni}e^{a0}+\mathrm{\Gamma }^ne^{ai}\right)_ne_{a0}+\frac{1}{2}\mathrm{\Gamma }^{mn}T_{nm}^i$$
$$+2_n\mathrm{\Gamma }^{ni}+g^{in}(H_n\mathrm{\Gamma }^jT_{0nj}\mathrm{\Gamma }^{mj}T_{mnj})]\delta (xy)$$
$$+\mathrm{\Gamma }^{ni}(x)\frac{}{x^n}\delta (xy),$$
$`(44)`$
$$\{H_i(x),\mathrm{\Gamma }^j(y)\}=\delta _i^j\mathrm{\Gamma }^n(y)\frac{}{y^n}\delta (xy)+\mathrm{\Gamma }^j(x)\frac{}{x^i}\delta (xy)\mathrm{\Gamma }^je^{a0}_ie_{a0}\delta (xy),$$
$`(45)`$
$$\{H_k(x),\mathrm{\Gamma }^{ij}(y)=\mathrm{\Gamma }^{ij}(x)\frac{}{x^k}\delta (xy)+(\delta _k^j\mathrm{\Gamma }^{ni}(y)\delta _k^i\mathrm{\Gamma }^{nj}(y))\frac{}{x^n}\delta (xy)$$
$$+\frac{1}{2}\left(e^{aj}(x)\mathrm{\Gamma }^i(x)e^{ai}(x)\mathrm{\Gamma }^j(x)\right)\frac{}{x^k}e_{a0}(x)\delta (xy).$$
$`(46)`$
It is clear from the constraint algebra above that $`H_0`$, $`H_i`$, $`\mathrm{\Gamma }^{ik}`$ and $`\mathrm{\Gamma }^k`$ constitute a set of first class constraints. Now it is easy to conclude that $`C^a`$, $`\mathrm{\Gamma }^{ik}`$ and $`\mathrm{\Gamma }^k`$ also constitute a first class set. By means of equation (28) we have $`\{C^a(x),C^b(y)\}=e^{a0}(x)\{H_0(x),H_0(y)\}e^{b0}(y)+`$ $`H_0(x)\{e^{a0}(x),H_0(y)\}e^{b0}(y)+\mathrm{}`$. On the right hand side of this Poisson bracket as well as of the brackets $`\{C^a(x),\mathrm{\Gamma }^{ik}(y)\}`$ and $`\{C^a(x),\mathrm{\Gamma }^k(y)\}`$ there will always appear a combination of the constraints $`H_0=e_{a0}C^a`$, $`\mathrm{\Gamma }^{ik}`$, $`\mathrm{\Gamma }^k`$ and
$$H_i=e_{ai}C^a\mathrm{\Gamma }^mT_{0mi}\mathrm{\Gamma }^{lm}T_{lmi}\frac{1}{2g^{00}}(g_{ik}g_{jl}P^{kl}+\frac{1}{2}g_{ij}P)\mathrm{\Gamma }^j.$$
$`(47)`$
The expression above follows from equation (29). All constraints of the theory are first class, and therefore the theory is well defined regarding time evolution.
The Hamiltonian density (31) determines the time evolution of any field quantity $`f(x)`$:
$$\dot{f}(x)=d^3y\{f(x),H(y)\}|_{\mathrm{\Gamma }^{ik}=\mathrm{\Gamma }^k=0}.$$
$`(48)`$
Physical quantities take values in the subspace of the phase space $`๐_\mathrm{\Gamma }`$ defined by (36). In this subspace the constraints $`C^a`$ become
$$C^a=e^{a0}H_0+e^{ai}H_i.$$
$`(49)`$
Restricting considerations to $`๐_\mathrm{\Gamma }`$ we note that if $`H_0`$ vanishes, then $`e_{a0}C^a`$ also vanishes. Since $`\{e_{a0}\}`$ are arbitrary, it follows that $`C^a=0`$. In order to arrive at this conclusion we note that the constraints $`C^a`$ are independent of $`e_{a0}`$. From the orthogonality relation $`e_{a\mu }e^{a\lambda }=\delta _\mu ^\lambda `$ we obtain $`\delta e^{b\mu }/\delta e_{a0}=e^{a\mu }e^{b0}`$. Using this variational relation and equations (22) and (49) it is possible to show that
$$\frac{\delta C^a}{\delta e_{b0}}=\frac{\delta }{\delta e_{b0}}\left(e^{a0}H_0+e^{ai}H_i\right)=e^{b0}e^{a0}H_0+e^{a0}\frac{\delta H_0}{\delta e_{b0}}e^{bi}e^{a0}H_i$$
$$=e^{b0}e^{a0}H_0+e^{a0}(e^{b0}H_0+e^{bi}H_i)e^{bi}e^{a0}H_i=0.$$
$`(50)`$
$`H_i`$ does not depend explicitly or implicitly on $`e_{a0}`$. We remark that by taking the variation with respect to $`e_{b0}`$ of both sides of equation (26), $`H_0=e_{a0}C^a`$, we arrive at
$$C^b=C^b+e_{a0}\frac{\delta C^a}{\delta e_{b0}},$$
from what follows the general result $`e_{a0}(\delta C^a/\delta e_{b0})=0`$. Taking into account the arbitrariness of $`e_{a0}`$ in the latter equation we are led to equation (50).
Therefore the vanishing of the Hamiltonian constraint $`H_0`$ implies the vanishing of $`C^a`$, and ultimately of the vector constraint $`H_i`$. Moreover we observe from (47) and (49) that $`H_i`$ can be obtained from $`H_0`$ in $`๐_\mathrm{\Gamma }`$ according to
$$e_{ai}\frac{\delta }{\delta e_{a0}}H_0=e_{ai}C^a=H_i.$$
$`(51)`$
Thus $`H_i`$ is derived from $`H_0`$. In the complete phase space the vanishing of $`H_i`$ is a consequence of the vanishing of $`H_0`$, $`\mathrm{\Gamma }^{ik}`$ and $`\mathrm{\Gamma }^k`$.
Finally we would like to remark that the Hamiltonian formulation of the theory can be described more succintly in terms of the constraints $`H_0`$, $`\mathrm{\Gamma }^{ik}`$ and $`\mathrm{\Gamma }^k`$, by the Hamiltonian density in the form
$$H(e_{ak},\mathrm{\Pi }^{ak},e_{a0},\alpha _{ik},\beta _k)=H_0+\alpha _{ik}\mathrm{\Gamma }^{ik}+\beta _k\mathrm{\Gamma }^k.$$
$`(52)`$
The Poisson brackets between these constraints are given by equations (37), (40-44). They constitute a first class set except for the fact that on the right hand side of (44) there appears the constraint $`H_i`$. However it poses no problem for the consistency of the constraints provided $`H_0`$, $`\mathrm{\Gamma }^{ik}`$ and $`\mathrm{\Gamma }^k`$ are taken to vanish at the intial time $`t=t_0`$. Let $`\varphi (x^i,t)`$ represent any of the latter constraints. At the initial time we have $`\varphi (x^i,t_0)=0`$. At $`t_0+\delta t`$ we find $`\varphi (x^i,t_0+\delta t)=\varphi (x^i,t_0)+\dot{\varphi }(x^i,t_0)\delta t`$ such that $`\dot{\varphi }(x^i,t_0)=\{\varphi (x^i,t_0),๐\}`$. Since the vanishing of $`H_i`$ at an instant of time is a consequence of the vanishing of $`H_0`$, $`\mathrm{\Gamma }^{ik}`$ and $`\mathrm{\Gamma }^k`$ at the same time, the consistency of the constraints is guaranteed at any $`t>t_0`$.
VI. Discussion
The Weitzenbรถck space-time allows a consistent description of the Hamiltonian formulation of the gravitational field. Although the underlying geometry is not Riemannian, the Lagrangian field equations (4) assure that the theory determined by (1) is equivalent to Einsteinโs general relativity. To our knowledge there does not exist any impediment based on experimental facts that rules out the teleparallel geometry in favour of the Riemannian geometry for the description of the physical space-time. The natural geometrical setting for teleparallel gravity is the teleparallel geometry. The Hamiltonian formulation of the TEGR in the Riemannian geometry, with local SO(3,1) symmetry, requires the introduction of a large number of field variables that renders an intricate constraint structure.
We have shown that the vector constraint $`H_i`$ can be obtained from the Hamiltonian constraint $`H_0`$ by means of a functional derivative of $`H_0`$, making use of the orthogonality properties of the tetrads in the reduced phase space $`๐_\mathrm{\Gamma }`$. However, it is an independent constraint. In contrast, in the ADM formulation the Hamiltonian and vector constraints are not mutually related, and in practice one has to consider both constraints for the dynamical evolution via Hamilton equations.
The number of degrees of freedom may be counted as the total number of canonical variables, $`e_{ak}`$ and $`\mathrm{\Pi }^{ak}`$, minus twice the number of first class constraints. Therefore we have $`2420=4`$ degrees of freedom in the phase space, as expected. Since the constraints $`\mathrm{\Gamma }^{ik}`$ and $`\mathrm{\Gamma }^k`$ are first class they act on $`e_{ak}`$, and $`\mathrm{\Pi }^{ak}`$ and generate symmetry transformations. In particular, for $`e_{a\mu }`$ we have
$$\delta e_{ak}(x)=d^3z\left[\epsilon _{ij}(z)\{e_{ak}(x),\mathrm{\Gamma }^{ij}(z)\}+\epsilon _j(z)\{e_{ak}(x),\mathrm{\Gamma }^j(z)\}\right]$$
$$=d^3z\left[\epsilon _{ij}(z)\frac{\delta \mathrm{\Gamma }^{ij}(z)}{\delta \mathrm{\Pi }^{ak}(x)}+\epsilon _j(z)\frac{\delta \mathrm{\Gamma }^j(z)}{\delta \mathrm{\Pi }^{ak}(x)}\right]=\epsilon _{ik}e_a^i+\epsilon _ke_a^0,$$
$`(53)`$
where $`\epsilon _{ij}(x)=\epsilon _{ji}(x)`$ and $`\epsilon _j(x)`$ are infinitesimal parameters. Note that these transformations do not act on $`e_{a0}`$. This issue has not been completely analyzed. The physical implications of these symmetries to the theory are currently under investigation.
In the analysis of a theory described by a Lagrangian density similar to (1), Mรธller pointed out that some supplementary conditions on the tetrads are needed. He suggested these conditions to arise from suitable boundary conditions for the field equations, possibly in the form of an anti-symmetric tensor. These supplementary conditions would uniquely determine a tetrad lattice, apart from a constant rotation of the tetrads in the lattice The problem of consistently defining these supplementary conditions is likely to be related to the symmetry transformation determined by (53).
The Hamiltonian density (52) determines the time evolution of field quantities via equation (48), and in particular of the metric tensor $`g_{ij}`$ of three-dimensional spacelike hypersurfaces. This property might simplify approaches to a canonical, nonperturbative quantization of gravity provided we manage to construct the reduced phase space determined by (36).
After implementing the primary constraints via equations (36), the first term of $`C^a`$ is given by $`_i\mathrm{\Pi }^{ai}`$, with $`\mathrm{\Pi }^{ai}`$ defined by (9). From our previous experience (cf. ref. ) we are led to conclude that this term is related to energy and momentum of the gravitational field. In the present case we also interpret equations $`C^a=0`$ as energy-momentum equations for the gravitational field. According to this interpretation, the integral form of the constraint equation $`C^{(0)}=0`$ can be written in the form $`E=0`$. Integration of $`_i\mathrm{\Pi }^{ai}`$ over the whole three-dimensional space yields the ADM energy. A complete analysis of this issue will be presented elsewhere.
Acknowledgements
J. F. R. N. is supported by CAPES, Brazil. |
warning/0002/hep-lat0002006.html | ar5iv | text | # A Quantum Perfect Lattice Action for Monopoles and Strings
## Abstract
A quantum perfect lattice action in four dimensions can be derived analytically as a renormalized trajectory when we perform a block spin transformation of monopole currents in a simple but non-trivial case of quadratic monopole interactions. The spectrum of the lattice theory is identical to that of the continuum theory. The perfect monopole action is transformed exactly into a lattice action of a string model. A perfect operator evaluating a static potential between electric charges is also derived explicitly. If the monopole interactions are weak as in the case of infrared $`SU(2)`$ QCD, the string interactions become strong. The static potential and the string tension is estimated analytically by the use of the strong coupling expansion and the continuum rotational invariance is restored completely.
$`PACS`$: 12.38.Gc, 11.15.Ha
$`Keywords`$: lattice QCD; perfect action;block spin transformation; monopole; dual transformation; string model
preprint: KANAZAWA 99-09, June 1999
To obtain the continuum limit is crucial in the framework of lattice field theories. A block spin transformation which is one of renormalization group transformations is adopted usually on the lattice. A quantum perfect lattice action is an action on the renormalized trajectory on which one can take the continuum limit.
In principle, we obtain the renormalized trajectory when we perform infinite steps (the $`n\mathrm{}`$ limit) of block spin transformations for fixed physical length $`b=na`$ where $`a`$ is the lattice constant. But this is actually impossible in ordinary cases. What we can do in actual simulations is to approach the renormalized trajectory carrying out as many steps of block spin transformations as possible on a finite lattice $`N^4`$. If the effective action $`S(n,a,N)`$ obtained satisfies well the two conditions that (1) $`S(n,a,N)`$ is a function of $`b=na`$ alone (the scaling behavior and volume independence) and (2) the continuum rotational invariance is satisfied, then the effective action could be regarded as a good approximation of the renormalized trajectory. The first condition can be checked when we compare $`S(n,a,N)`$ themselves for various $`a`$, $`n`$ and $`N^4`$. But to test the rotational invariance, we have to determine the correct form of physical operators on the blocked lattice. It is the perfect lattice operator on the renormalized trajectory which reproduces the continuum rotational invariance. To find the perfect lattice operator is highly nontrivial, too.
The purpose of this note is to give a simple but a non-trivial lattice model composed of general monopole quadratic interactions alone with which a block spin transformation can be done analytically. The renormalized trajectory and the perfect operator corresponding to a potential between static electric charges can be derived analytically. This is similar to the blocking from the continuum theory as developed by Bietenholz and Wiese. The spectrum of the lattice theory is the same as in the continuum theory. The continuum rotational invariance is shown exactly with the operator. In addition, this model is very interesting, since the effective monopole action obtained after an abelian projection of pure $`SU(2)`$ lattice QCD is known to be well dominated by such quadratic monopole interactions alone in the infrared region .
Let us start from the following action composed of quadratic interactions between magnetic monopole currents. It is formulated on an infinite lattice with very small lattice constant $`a`$:
$`S[k]={\displaystyle \underset{s,s^{},\mu }{}}k_\mu (s)D_0(ss^{})k_\mu (s^{}).`$ (1)
Since we are starting from the region very near to the continuum limit, it is natural to assume the direction independence of $`D_0(ss^{})`$. Also we have adopted only parallel interactions, since we can avoid perpendicular interactions from short-distance terms using the current conservation. Moreover, for simplicity, we adopt only the first three Laurent expansions, i.e., Coulomb, self and nearest-neighbor interactions. Explicitly, $`D_0(ss^{})`$ is expressed as $`\beta \mathrm{\Delta }_L^1(ss^{})+\alpha \delta _{s,s^{}}+\gamma \mathrm{\Delta }_L(ss^{})`$. Here $`\mathrm{\Delta }_L(ss^{})=_\mu _\mu _\mu ^{}\delta _{s,s^{}}`$ and $`(^{})`$ is the forward (backward) difference. Including more complicated quadratic interactions is not difficult.
How to evaluate a static potential between electrically charged particles is a problem. It is known that the theory with the above action (1) is equivalent to an abelian gauge theory of the Villain form. In this model, it is natural to use an abelian Wilson loop $`W(๐)=\mathrm{exp}i_๐(\theta _\mu (s),J_\mu (s))`$, where $`\theta _\mu (s)`$ is an abelian angle variable of the modified Villain action. Also the theory with the above action (1) can be rewritten in the lattice form of the modified London limit of the dual abelian Higgs model. The static potential is evaluated by a โtHooft operator in the model. However the expectation values of both operators are not completely equivalent, although the term of the area law is the same . When use is made of BKT transformation, we see that the area law term is given correctly also by the following operator in the monopole action <sup>*</sup><sup>*</sup>*Using the monopole definition $`\stackrel{`}{\text{a}}`$ la DeGrand-Toussaint, we can prove it also directly from abelian Wilson loops.:
$`W_m(๐)`$ $`=`$ $`\mathrm{exp}\left(2\pi i{\displaystyle \underset{s,\mu }{}}N_\mu (s,S^J)k_\mu (s)\right),`$ (2)
$`N_\mu (s,S_J)`$ $`=`$ $`{\displaystyle \underset{s^{}}{}}\mathrm{\Delta }_L^1(ss^{}){\displaystyle \frac{1}{2}}ฯต_{\mu \alpha \beta \gamma }_\alpha S_{\beta \gamma }^J(s^{}+\widehat{\mu }),`$ (3)
where $`S_{\beta \gamma }^J(s^{}+\widehat{\mu })`$ is a plaquette variable satisfying $`_\beta ^{}S_{\beta \gamma }^J(s)=J_\gamma (s)`$ and the coordinate displacement $`\widehat{\mu }`$ is due to the interaction between dual variables. It is possible to prove that any choice of $`S_{\beta \gamma }^J(s)`$ for fixed electric currents $`J_\gamma (s)`$ gives the same value in the continuum limit $`a0`$, since the difference is a closed surface and the exponent of $`W_m(C)`$ for the difference is just the four-dimensional linking number times $`2\pi i`$. Hence we take a flat surface for $`S_{\beta \gamma }^J(s)`$ in the following. Since the area law term is the same, let us consider only the operator (2) in the following. The details of the definition of the operator evaluating the static potential are discussed in Ref..
Now let us define a blocked monopole current:
$`K_\mu (s^{(n)})`$ $`=`$ $`{\displaystyle \underset{i,j,l=0}{\overset{n1}{}}}k_\mu (ns^{(n)}+(n1)\widehat{\mu }+i\widehat{\nu }+j\widehat{\rho }+l\widehat{\sigma })`$ (4)
$``$ $`_{k_\mu }(s^{(n)}).`$ (5)
With this definition, the current $`K_\mu (s^{(n)})`$ on the coarse lattice with a lattice distance $`b=na`$ satisfies the current conservation $`_\mu ^{}K_\mu (s^{(n)})=_\mu (K_\mu (s^{(n)})K_\mu (s^{(n)}b\widehat{\mu }))=0`$. The block spin transformation is expressed as
$`Z[K,J]`$ $`=`$ $`{\displaystyle \underset{\genfrac{}{}{0pt}{}{k_\mu =\mathrm{}}{_\mu ^{}k_\mu =0}}{\overset{\mathrm{}}{}}}\mathrm{exp}\left\{{\displaystyle \underset{s,s^{},\mu }{}}k_\mu (s)D_0(ss^{})k_\mu (s^{})+2\pi i{\displaystyle \underset{s,\mu }{}}N_\mu (s)k_\mu (s)\right\}`$ (7)
$`\times \delta \left(K_\mu (s^{(n)})_{k_\mu }(s^{(n)})\right).`$
The vacuum expectation value of the Wilson loop (2) is written in terms of $`K_\mu (s^{(n)})`$ as follows:
$`W_m(๐)={\displaystyle \underset{\genfrac{}{}{0pt}{}{K_\mu =\mathrm{}}{_\mu ^{}K_\mu =0}}{\overset{\mathrm{}}{}}}Z[K,J]/{\displaystyle \underset{\genfrac{}{}{0pt}{}{K_\mu =\mathrm{}}{_\mu ^{}K_\mu =0}}{\overset{\mathrm{}}{}}}Z[K,0].`$ (8)
Introducing auxiliary field $`\varphi `$ and $`\gamma `$, we rewrite the constraints $`_\mu ^{}k_\mu =0`$ and $`K_\mu (s^{(n)})=_{k_\mu }(s^{(n)})`$. Then we change the integral region of $`\gamma `$ and $`\varphi `$ from the first Brillouin zone to the infinite region, since the monopole currents take integer values. Making use of the Poisson sum rule and recovering dimensional lattice constants $`a`$ and $`b=na`$, we get
$`Z[K,J]`$ $``$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}๐\gamma \mathrm{exp}\left\{ib^4{\displaystyle \underset{x^{(n)},\mu }{}}\gamma _\mu (bs^{(n)})K_\mu (bs^{(n)})\right\}`$ (11)
$`\times {\displaystyle \underset{\genfrac{}{}{0pt}{}{al=\mathrm{}}{alZ}}{\overset{\mathrm{}}{}}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}๐\varphi {\displaystyle _{\mathrm{}}^{\mathrm{}}}๐F\mathrm{exp}\{a^8{\displaystyle \underset{s,s^{},\mu }{}}F_\mu (as)D_0(asas^{})F_\mu (as^{})`$
$`+ia^4{\displaystyle \underset{s,\mu }{}}[2\pi N_\mu (as)+\varphi (as)_\mu ^{}+2\pi l_\mu (as)]F_\mu (as)a^4{\displaystyle \underset{s^{(n)},\mu }{}}n\gamma _\mu (nas^{(n)})_{F_\mu }(s^{(n)})\}.`$
Since we take the continuum limit $`a0`$ finally, $`l=0`$ alone may remains in the sum with respect to $`alZ`$. Carrying out explicitly integrals with respect to $`F`$, $`\varphi `$, $`\gamma `$ and taking the continuum limit $`a0`$, we obtain the expectation value of the operator and the effective action on the coarse lattice:
$`W(๐)`$ $`=`$ $`\mathrm{exp}\{\pi ^2{\displaystyle _{\mathrm{}}^{\mathrm{}}}d^4xd^4y{\displaystyle \underset{\mu }{}}N_\mu (x)D_0^1(xy)N_\mu (y)`$ (16)
$`+\pi ^2b^8{\displaystyle \underset{\genfrac{}{}{0pt}{}{s^{(n)},s^{(n)^{}}}{\mu ,\nu }}{}}B_\mu (bs^{(n)})A_{\mu \nu }^{\mathrm{GF}1}(bs^{(n)}bs^{(n)^{}})B_\nu (bs^{(n)^{}})\}`$
$`\times {\displaystyle \underset{\genfrac{}{}{0pt}{}{b^3K_\mu (bs)=\mathrm{}}{_\mu ^{}K_\mu =0}}{\overset{\mathrm{}}{}}}\mathrm{exp}\{b^8{\displaystyle \underset{\genfrac{}{}{0pt}{}{s^{(n)},s^{(n)^{}}}{\mu ,\nu }}{}}K_\mu (bs^{(n)})A_{\mu \nu }^{\mathrm{GF}1}(bs^{(n)}bs^{(n)^{}})K_\nu (bs^{(n)^{}})`$
$`+2\pi ib^8{\displaystyle \underset{\genfrac{}{}{0pt}{}{s^{(n)},s^{(n)^{}}}{\mu ,\nu }}{}}B_\mu (bs^{(n)})A_{\mu \nu }^{\mathrm{GF}1}(bs^{(n)}bs^{(n)^{}})K_\nu (bs^{(n)^{}})\}/{\displaystyle }_{\genfrac{}{}{0pt}{}{b^3K_\mu (bs)=\mathrm{}}{_\mu ^{}K_\mu =0}}^{\mathrm{}}Z[K,0],`$
where
$`B_\mu (bs^{(n)})`$ $``$ $`\underset{\genfrac{}{}{0pt}{}{a0}{n\mathrm{}}}{lim}a^8{\displaystyle \underset{s,s^{},\nu }{}}\mathrm{\Pi }_{\neg \mu }(bs^{(n)}as)A_{\mu \nu }(asas^{})N_\nu (as^{}),`$ (17)
$`\mathrm{\Pi }_{\neg \mu }(bs^nas)`$ $``$ $`{\displaystyle \frac{1}{n^3}}\delta \left(nas_\mu ^{(n)}+(n1)aas_\mu \right)\times {\displaystyle \underset{i(\mu )}{}}\left({\displaystyle \underset{I=0}{\overset{n1}{}}}\delta \left(nas_i^{(n)}+Iaas_i\right)\right),`$ (18)
$`A_{\mu \nu }(asas^{})`$ $``$ $`\left\{\delta _{\mu \nu }{\displaystyle \frac{_\mu _\nu ^{}}{_\rho _\rho _\rho ^{}}}\right\}D_0^1(asas^{}).`$ (19)
$`A_{\mu \nu }^{\mathrm{GF}1}(bs^{(n)}bs^{(n)^{}})`$ is a gauge-fixed inverse of the following operator:
$`A_{\mu \nu }^{}(bs^{(n)}bs^{(n)^{}})a^8{\displaystyle \underset{s,s^{}}{}}\mathrm{\Pi }_{\neg \mu }(bs^{(n)}as)A_{\mu \nu }(asas^{})\mathrm{\Pi }_{\neg \nu }(bs^{(n)^{}}as^{}).`$ (20)
Here we have used $`_\mu _\mu N_\mu (s)=0`$ and have adopted a gauge including $`\lambda \{_\mu \gamma _\mu (bs^{(n)})\}^2`$ in the integral with respect to $`\gamma `$.
The momentum representation of the gauge fixed propagator is given explicitly by
$`A_{\mu \nu }^{\mathrm{GF}}(p)=\left\{A_{\mu \nu }^{}(p)+\lambda \widehat{p}_\mu \widehat{p}_\nu \right\}e^{i(p_\mu p_\nu )/2},`$ (21)
where $`\widehat{p}_\mu =2\mathrm{sin}(p_\mu /2)`$ and $`A_{\mu \nu }^{}(p)`$ is written as
$`A_{\mu \nu }^{}(p)`$ $``$ $`({\displaystyle \underset{i=1}{\overset{4}{}}}{\displaystyle \underset{l_i=\mathrm{}}{\overset{\mathrm{}}{}}})\{D_0^1(p+2\pi l)\left[\delta _{\mu \nu }{\displaystyle \frac{(p+2\pi l)_\mu (p+2\pi l)_\nu }{_i(p+2\pi l)_i^2}}\right]{\displaystyle \frac{(p+2\pi l)_\mu (p+2\pi l)_\nu }{_i(p+2\pi l)_i^2}}\}`$ (23)
$`\times {\displaystyle \frac{\left(_{i=1}^4\widehat{p}_i\right)^2}{\widehat{p}_\mu \widehat{p}_\nu }},`$
From the gauge invariance condition $`_\mu _\mu ^{}A_{\mu \nu }^{}(ss^{})=0`$, we get $`_\mu \widehat{p}_\mu A_{\mu \nu }^{}(p)=0`$. The inverse of $`A_{\mu \nu }^{\mathrm{GF}}(p)`$ is as follows:
$`A_{\mu \nu }^{\mathrm{GF}1}(p)`$ $`=`$ $`D_{\mu \nu }(p)+{\displaystyle \frac{1}{\lambda }}{\displaystyle \frac{\widehat{p}_\mu \widehat{p}_\nu }{(_i\widehat{p}_i^2)^2}}e^{i(p_\mu p_\nu )/2},`$ (24)
where $`D_{\mu \nu }(p)`$ is the $`\lambda `$ independent part of $`A_{\mu \nu }^{\mathrm{GF}1}(p))`$. To derive the explicit form of $`D_{\mu \nu }(p)`$ is not so easy. However, evaluating the determinant and cofactors of the matrix $`A_{\mu \nu }^{\mathrm{GF}}(p)`$, we can get
$`D_{\mu \nu }`$ $`=`$ $`{\displaystyle \frac{3}{\left\{\left(_aA_{aa}^{}\right)^33\left(_aA_{aa}^{}\right)\left(_{ab}A_{ab}^{}A_{ba}^{}\right)+2\left(_{abc}A_{ab}^{}A_{bc}^{}A_{ca}^{}\right)\right\}}}`$ (27)
$`\times [\{\left({\displaystyle \underset{a}{}}A_{aa}^{}\right)^2\left({\displaystyle \underset{ab}{}}A_{ab}^{}A_{ba}^{}\right)\}\delta _{\mu \nu }2\left({\displaystyle \underset{a}{}}A_{aa}^{}\right)A_{\mu \nu }^{}+2{\displaystyle \underset{a}{}}A_{\mu a}^{}A_{a\nu }^{}`$
$`\{\left({\displaystyle \underset{a}{}}A_{aa}^{}\right)^2\left({\displaystyle \underset{ab}{}}A_{ab}^{}A_{ba}^{}\right)\}{\displaystyle \frac{\widehat{p}_\mu \widehat{p}_\nu }{_a\widehat{p}_a^2}}]e^{i(p_\mu p_\nu )/2}.`$
Since $`_\mu _\mu ^{}K_\mu (bs^{(n)})=_\mu _\mu ^{}B_\mu (bs^{(n)})=0`$, Eq.(16) is independent of the gauge parameter $`\lambda `$.
Now let us evaluate the spectrum of the monopole current $`K_\mu (s^{(n)})`$ on the coarse lattice. Define an operator with definite spatial momentum $`\stackrel{}{p}`$:
$`K_i(\stackrel{}{p})_{x_4}={\displaystyle _\pi ^{+\pi }}{\displaystyle \frac{dp_4}{2\pi }}K_i(\stackrel{}{p},p_4)e^{ip_4x_4}.`$
Then the correlation function is written as
$`K_i(\stackrel{}{x},0)K_i(\stackrel{}{p})_{x_4}`$ $`=`$ $`{\displaystyle _\pi ^{+\pi }}{\displaystyle \frac{d^4k}{(2\pi )^4}}{\displaystyle _\pi ^{+\pi }}{\displaystyle \frac{dp_4}{2\pi }}e^{i\stackrel{}{k}\stackrel{}{x}}e^{ip_4x_4}K_i(\stackrel{}{k},k_4)K_i(\stackrel{}{p},p_4).`$ (28)
Since the monopole action on the b-lattice is written as $`_{\genfrac{}{}{0pt}{}{s,s^{}}{\mu ,\nu }}K_\mu (s)D_{\mu \nu }(ss^{})K_\nu (s^{})`$, we see that the spectrum is essentially fixed by the gauge invariant part of the inverse of $`D_{\mu \nu }(ss^{})`$. In (28), it is $`\delta ^4(k+p)A_{ii}^{}(\stackrel{}{p},p_4)`$, where
$`A_{ii}^{}(\stackrel{}{p},p_4)`$ $`=`$ $`{\displaystyle \frac{\left(_{i=1}^4\widehat{p}_i\right)^2}{\widehat{p}_i\widehat{p}_i}}{\displaystyle \underset{i=1}{\overset{4}{}}}({\displaystyle \underset{l_i=\mathrm{}}{\overset{\mathrm{}}{}}})\{D_0^1(p+2\pi l)[\delta _{ii}{\displaystyle \frac{(p+2\pi l)_i(p+2\pi l)_i}{_i(p+2\pi l)_i^2}}]`$
$`{\displaystyle \frac{(p+2\pi l)_i(p+2\pi l)_i}{_i(p+2\pi l)_i^2}}\}.`$
The integral
$`{\displaystyle _\pi ^{+\pi }}{\displaystyle \frac{dp_4}{2\pi }}e^{i\stackrel{}{k}\stackrel{}{x}}e^{ip_4x_4}A_{ii}^{}(\stackrel{}{p},p_4)`$
can be performed when we change $`p_4+2\pi l_4p_4`$ and $`_{l_4}_{\pi +2\pi l_4}^{\pi +2\pi l_4}_{\mathrm{}}^{\mathrm{}}`$. Here
$`D_0^1(p)=\kappa \left({\displaystyle \frac{m_1^2}{p^2+m_1^2}}{\displaystyle \frac{m_2^2}{p^2+m_2^2}}\right),`$
where $`\kappa `$, $`m_1`$ and $`m_2`$ are expressed by the original couplings in Eq.(1) as $`\kappa (m_1^2m_2^2)=\gamma ^1`$, $`m_1^2+m_2^2=\alpha /\gamma `$ and $`m_1^2m_2^2=\beta /\gamma `$. Hence the $`p_4`$ integral gives us
$`e^{E_i(\stackrel{}{p}+2\pi \stackrel{}{l})x_4},`$
where $`E_i(\stackrel{}{p}+2\pi \stackrel{}{l})^2=p_4^2=(\stackrel{}{p}+2\pi \stackrel{}{l})^2+m_i^2`$. The spectrum is identical with the one of the continuum.
The above monopole action can be transformed exactly into that of the string model . When use is made of the Poisson sum rule, we write the monopole part of Eq.(16) as
$`{\displaystyle \underset{\genfrac{}{}{0pt}{}{K_\mu =\mathrm{}}{_\mu ^{}K_\mu =0}}{\overset{\mathrm{}}{}}}\mathrm{exp}\left\{{\displaystyle \underset{\genfrac{}{}{0pt}{}{s,s^{}}{\mu ,\nu }}{}}K_\mu (s)D_{\mu \nu }(ss^{})K_\nu (s^{})+2\pi i{\displaystyle \underset{\genfrac{}{}{0pt}{}{s,s^{}}{\mu ,\nu }}{}}B_\mu (s)D_{\mu \nu }(ss^{})K_\nu (s^{})\right\}`$ (31)
$`={\displaystyle _{\mathrm{}}^+\mathrm{}}๐F_\mu (s){\displaystyle _\pi ^{+\pi }}๐\varphi (s){\displaystyle \underset{K_\mu (s)=\mathrm{}}{\overset{\mathrm{}}{}}}\mathrm{exp}\{{\displaystyle \underset{\genfrac{}{}{0pt}{}{s,s^{}}{\mu ,\nu }}{}}F_\mu (s)D_{\mu \nu }(ss^{})F_\nu (s^{})`$
$`+i{\displaystyle \underset{s,\mu }{}}F_\mu (s)[_\mu \varphi (s)+2\pi {\displaystyle \underset{s^{},\mu ,\nu }{}}D_{\mu \nu }(ss^{})B_\nu (s^{})]\},`$
$`=`$ $`{\displaystyle _\pi ^{+\pi }}๐\varphi (s){\displaystyle \underset{l_\mu (s)=\mathrm{}}{\overset{\mathrm{}}{}}}\mathrm{exp}\{{\displaystyle \frac{1}{4}}{\displaystyle \underset{\genfrac{}{}{0pt}{}{s,s^{}}{\mu ,\nu }}{}}[_\mu \varphi (s)+2\pi l_\mu (s)+2\pi {\displaystyle \underset{s_1,\mu ,\alpha }{}}D_{\mu \alpha }(ss_1)B_\alpha (s_1)]`$ (33)
$`D_{\mu \nu }^1(ss^{})[_\nu \varphi (s^{})+2\pi l_\nu (s^{})+2\pi {\displaystyle \underset{s_2,\nu ,\beta }{}}D_{\nu \beta }(s^{}s_2)B_\beta (s_2)]\}.`$
Performing BKT transformation and Hodge decomposition, we obtain
$`l_\mu (s)`$ $`=`$ $`s_\mu (s)+_\mu r(s)`$ (34)
$`=`$ $`_\mu \left\{{\displaystyle \underset{s^{}}{}}\mathrm{\Delta }_{Ls,s^{}}^1_\nu ^{}s_\nu (s^{})+r_\mu (s^{})\right\}+{\displaystyle \underset{s^{}}{}}_\nu ^{}\mathrm{\Delta }_{Ls,s^{}}^1\sigma _{\nu \mu }(s^{}),`$ (35)
where $`\sigma _{\nu \mu }(s)_{[\mu }s_{\nu ]}`$ is the closed string variable satisfying the conservation rule
$`_{[\alpha }\sigma _{\mu \nu ]}=_\alpha \sigma _{\mu \nu }+_\mu \sigma _{\nu \alpha }+_\nu \sigma _{\alpha \mu }=0.`$ (36)
The compact field $`\varphi (s)`$ is absorbed into a non-compact field $`\varphi _{NC}(s)`$. Integrating out the auxiliary non-compact field, we see
$`(\text{33})`$ $`=`$ $`{\displaystyle \underset{\genfrac{}{}{0pt}{}{\sigma _{\mu \nu }(s)=\mathrm{}}{_{[\alpha }\sigma _{\mu \nu ]}(s)=0}}{\overset{\mathrm{}}{}}}\mathrm{exp}\{\pi ^2{\displaystyle \underset{\genfrac{}{}{0pt}{}{s,s^{}}{\genfrac{}{}{0pt}{}{\mu \alpha }{\nu \beta }}}{}}\sigma _{\mu \alpha }(s)_\alpha _\beta ^{}D_{\mu \nu }^1(ss_1)\mathrm{\Delta }_L^2(s_1s^{})\sigma _{\nu \beta }(s^{})`$ (38)
$`2\pi ^2{\displaystyle \underset{\genfrac{}{}{0pt}{}{s,s^{}}{\mu ,\nu }}{}}\sigma _{\mu \nu }(s)_\mu \mathrm{\Delta }_L^1(ss^{})B_\nu (s^{})\pi ^2{\displaystyle \underset{\genfrac{}{}{0pt}{}{s,s^{}}{\mu ,\nu }}{}}B_\mu (s)D_{\mu \nu }(ss^{})B_\nu (s^{})\}.`$
The term independent of the string variable exactly cancels the second classical term of Eq.(16). We find finally
$`W_m(๐)`$ $`=`$ $`{\displaystyle \frac{1}{Z}}\mathrm{exp}\left\{\pi ^2{\displaystyle _{\mathrm{}}^{\mathrm{}}}d^4xd^4y{\displaystyle \underset{\mu }{}}N_\mu (x)D_0^1(xy)N_\mu (y)\right\}`$ (41)
$`\times {\displaystyle \underset{\genfrac{}{}{0pt}{}{\sigma _{\mu \nu }(s)=\mathrm{}}{_{[\alpha }\sigma _{\mu \nu ]}(s)=0}}{\overset{\mathrm{}}{}}}\mathrm{exp}\{\pi ^2{\displaystyle \underset{\genfrac{}{}{0pt}{}{s,s^{}}{\genfrac{}{}{0pt}{}{\mu \alpha }{\nu \beta }}}{}}\sigma _{\mu \alpha }(s)_\alpha _\beta ^{}D_{\mu \nu }^1(ss_1)\mathrm{\Delta }_L^2(s_1s^{})\sigma _{\nu \beta }(s^{})`$
$`2\pi ^2{\displaystyle \underset{\genfrac{}{}{0pt}{}{s,s^{}}{\mu ,\nu }}{}}\sigma _{\mu \nu }(s)_\mu \mathrm{\Delta }_L^1(ss^{})B_\nu (s^{})\}.`$
It is very interesting that we can evaluate analytically the potential between the static electric charges when the monopole action on the dual lattice is in the weak coupling region for large $`b`$ as realized in the infrared region of pure $`SU(2)`$ and $`SU(3)`$ QCD. Then the string model on the original lattice is in the strong coupling region. As shown later explicitly, the potential between the static electric charges is then evaluated mainly by the first classical part of Eq.(41) alone. Hence let us evaluate first the classical part. Since the classical part is written in the continuum form, the continuum rotational invariance for any b lattice site is trivial.
The plaquette variable $`S_{\alpha \beta }`$ in Eq.(3) for the static potential $`V(Ib,0,0)`$ is expressed by
$`S_{\alpha \beta }(z)`$ $`=`$ $`\delta _{\alpha 1}\delta _{\beta 4}\delta (z_2)\delta (z_3)\theta (z_1)\theta (Ibz_1)\theta (z_4)\theta (Tbz_4).`$ (42)
Also the variable $`S_{\alpha \beta }`$ for the static potential $`V(Ib,Ib,0)`$ is given by
$`S_{\alpha \beta }(z)`$ $`=`$ $`\left(\delta _{\alpha 1}\delta _{\beta 4}+\delta _{\alpha 2}\delta _{\beta 4}\right)\delta (z_3)\theta (z_4)\theta (Tbz_4)`$ (44)
$`\times \theta (z_1)\theta (Ibz_1)\theta (z_2)\theta (Ibz_2)\delta (z_1z_2).`$
Let us evaluate $`V(Ib,0,0)`$ as an example. The Fourier transform of $`S_{\alpha \beta }(z)`$ in this case is
$`S_{\alpha \beta }(p)`$ $`=`$ $`\left(\delta _{\alpha 4}\delta _{\beta 1}\delta _{\alpha 1}\delta _{\beta 4}\right){\displaystyle _0^{Ib}}๐z_1e^{ip_1z_1}{\displaystyle _0^{Tb}}๐z_4e^{ip_4z_4},`$ (45)
$`=`$ $`\left(\delta _{\alpha 4}\delta _{\beta 1}\delta _{\alpha 1}\delta _{\beta 4}\right)\left({\displaystyle \frac{2}{p_1}}\right)e^{i\frac{p_1bI}{2}}\mathrm{sin}({\displaystyle \frac{p_1Ib}{2}})\left({\displaystyle \frac{2}{p_4}}\right)e^{i\frac{p_4bT}{2}}\mathrm{sin}({\displaystyle \frac{p_4Tb}{2}}).`$ (46)
Since we study large $`T`$ and large $`b`$ behaviors, we use the following formula:
$`\underset{T\mathrm{}}{lim}\left({\displaystyle \frac{\mathrm{sin}\alpha T}{\alpha }}\right)^2`$ $`=`$ $`\pi T\delta (\alpha ).`$ (47)
We get
$`{\displaystyle \genfrac{}{}{0pt}{}{W(Ib,0,0,Tb)}{}}{\displaystyle \genfrac{}{}{0pt}{}{}{\genfrac{}{}{0pt}{}{T\mathrm{}}{b\mathrm{}}}}{\displaystyle \genfrac{}{}{0pt}{}{\mathrm{exp}\left\{\pi ^2(TIb^2){\displaystyle \frac{d^2p}{(2\pi )^2}\left[\frac{1}{\mathrm{\Delta }D_0}\right](0,p_2,p_3,0)}\right\}}{}}.`$ (48)
Similarly we can evaluate $`W(Ib,Ib,0,Tb)`$ from the classical term. The static potentials $`V(Ib,0,0)`$ and $`V(Ib,Ib,0)`$ can be written as
$`V(Ib,0,0)`$ $`=`$ $`\pi ^2(Ib){\displaystyle \frac{d^2p}{(2\pi )^2}\left[\frac{1}{\mathrm{\Delta }D_0}\right](0,p_2,p_3,0)},`$ (49)
$`=`$ $`{\displaystyle \frac{\pi \kappa Ib}{2}}\mathrm{ln}{\displaystyle \frac{m_1}{m_2}},`$ (50)
$`V(Ib,Ib,0)`$ $`=`$ $`{\displaystyle \frac{\sqrt{2}\pi \kappa Ib}{2}}\mathrm{ln}{\displaystyle \frac{m_1}{m_2}}.`$ (51)
The potentials from the classical part take only the linear form and the rotational invariance is recovered completely even for the nearest $`I=1`$ sites. The string tension from the classical part is evaluated as
$`\sigma _{cl}={\displaystyle \frac{\pi \kappa }{2}}\mathrm{ln}{\displaystyle \frac{m_1}{m_2}}.`$ (52)
This is consistent with the analytical results in Type-2 superconductor. The two constants $`m_1`$ and $`m_2`$ may be regarded as the coherence and the penetration lengths.
Next let us evaluate the quantum fluctuation coming from the interaction of the string variable and the classical source. Since we have introduced the source term corresponding to the Wilson loop on the fine $`a`$ lattice, the recovery of the rotational invariance of the static potential is naturally expected also for the quantum fluctuation. Hence here we evaluate the quantum fluctuation for the flat Wilson loop $`W(Ib,0,0,Tb)`$. Then it is to be emphasized that the same static potential for the flat Wilson loop can be obtained for $`I,T\mathrm{}`$ when we consider the naive Wilson loop operator on the course $`b`$ lattice instead of that on the fine lattice (2):
$`\stackrel{~}{W}_m(๐)`$ $`=`$ $`\mathrm{exp}\left(2\pi i{\displaystyle \underset{s,\mu }{}}\stackrel{~}{N}_\mu (s,S^J)K_\mu (s)\right),`$ (53)
$`\stackrel{~}{N}_\mu (s,S_J)`$ $`=`$ $`{\displaystyle \underset{s^{}}{}}\mathrm{\Delta }_L^1(ss^{}){\displaystyle \frac{1}{2}}ฯต_{\mu \alpha \beta \gamma }_\alpha \stackrel{~}{S}_{\beta \gamma }^J(s^{}+\widehat{\mu }),`$ (54)
where $`\stackrel{~}{S}_{\beta \gamma }^J(s^{}+\widehat{\mu })`$ is a flat plaquette variable satisfying $`_\beta ^{}\stackrel{~}{S}_{\beta \gamma }^J(s)=\stackrel{~}{J}_\gamma (s)`$ and $`\stackrel{~}{J}_\gamma (s)`$ is the electric current on the course lattice.
Similar arguments as above shows that the expectation value of the operator $`\stackrel{~}{W}_m(๐)`$ is expressed as follows:
$`\stackrel{~}{W}_m(๐)`$ $`=`$ $`{\displaystyle \frac{1}{Z}}\mathrm{exp}\left\{\pi ^2{\displaystyle \underset{\genfrac{}{}{0pt}{}{s,s^{}}{\mu ,\nu }}{}}{\displaystyle \underset{\mu }{}}\stackrel{~}{N}_\mu (s)D_{\mu \nu }^1(ss^{})\stackrel{~}{N}_\nu (s^{})\right\}`$ (57)
$`\times {\displaystyle \underset{\genfrac{}{}{0pt}{}{\sigma _{\mu \nu }(s)=\mathrm{}}{_{[\alpha }\sigma _{\mu \nu ]}(s)=0}}{\overset{\mathrm{}}{}}}\mathrm{exp}\{\pi ^2{\displaystyle \underset{\genfrac{}{}{0pt}{}{s,s^{}}{\genfrac{}{}{0pt}{}{\mu \alpha }{\nu \beta }}}{}}\sigma _{\mu \alpha }(s)_\alpha _\beta ^{}D_{\mu \nu }^1(ss_1)\mathrm{\Delta }_L^2(s_1s^{})\sigma _{\nu \beta }(s^{})`$
$`2\pi ^2{\displaystyle \underset{\genfrac{}{}{0pt}{}{s,s^{}}{\mu ,\nu }}{}}\sigma _{\mu \nu }(s)_\mu \mathrm{\Delta }_L^1(ss_1)D_{\nu \alpha }(s_1s^{})\stackrel{~}{N}_\alpha (s^{})\}.`$
Here we note that $`D_{\mu \nu }^1(ss^{})`$ is given by Eq.(20) and that the Fourier transform $`\stackrel{~}{S}_{\beta \gamma }(k)`$ of $`\stackrel{~}{S}_{\beta \gamma }^J(s)`$ is similar to Eq.(46) with the momentum defined on the course lattice. Hence the first classical part becomes
$`\pi ^2{\displaystyle \underset{\genfrac{}{}{0pt}{}{s,s^{}}{\mu ,\nu }}{}}N_\mu (s)D_{\mu \nu }^1(ss^{})N_\nu (s^{})`$ (58)
$`=`$ $`4\pi ^2ฯต_{\mu \alpha 14}ฯต_{\nu \beta 14}{\displaystyle _\pi ^\pi }{\displaystyle \frac{d^4k}{(2\pi )^4}}\mathrm{\Delta }_L^2(k)\mathrm{sin}\left({\displaystyle \frac{k_\alpha }{2}}\right)\mathrm{sin}\left({\displaystyle \frac{k_\beta }{2}}\right)e^{i(\frac{k_\alpha }{2}\frac{k_\beta }{2})}`$ (61)
$`\times \stackrel{~}{S}_{14}(k)\stackrel{~}{S}_{14}(k){\displaystyle \underset{l}{}}\left(\delta _{\mu \nu }{\displaystyle \frac{(k+2\pi l)_\mu (k+2\pi l)_\nu }{(k+2\pi l)^2}}\right)D_0^1(k+2\pi l)`$
$`\times \mathrm{\Pi }_{\genfrac{}{}{0pt}{}{i\mu }{j\nu }}\pi _i(k+2\pi l)\pi _j^{}(k+2\pi l),`$
where
$`\pi _i(k)={\displaystyle \frac{\mathrm{sin}(k_i/2)}{k_i/2}}e^{ik_i/2}.`$ (62)
Changing the integral variable as $`k+2\pi lk`$, we can absorb the summation with respect to $`l`$ using the integral over the infinite momentum range. When we use Eq.(47) for large $`I`$ and $`T`$, we find the classical part (61) agrees exactly with (48).
Similarly
$`{\displaystyle \underset{s^{}}{}}D_{\mu \nu }^1(ss^{})N_\nu (s^{})`$ (63)
$`=`$ $`{\displaystyle _\pi ^\pi }{\displaystyle \frac{d^4k}{(2\pi )^4}}{\displaystyle \underset{l}{}}\left(\delta _{\mu \nu }{\displaystyle \frac{(k+2\pi l)_\mu (k+2\pi l)_\nu }{(k+2\pi l)^2}}\right)D_0^1(k+2\pi l)`$ (64)
$`\times `$ $`\mathrm{\Pi }_{\genfrac{}{}{0pt}{}{i\mu }{j\nu }}\pi _i^{}(k+2\pi l)\pi _j(k+2\pi l)e^{ik_\mu }2iฯต_{\nu \alpha 14}\mathrm{\Delta }_L^1(k)\mathrm{sin}\left({\displaystyle \frac{k_\alpha }{2}}\right)e^{i\frac{k_\alpha }{2}}{\displaystyle \frac{\stackrel{~}{\pi }_1(k)\stackrel{~}{\pi }_4(k)}{\pi _1(k)\pi _4(k)}}e^{iks},`$ (65)
where
$`\stackrel{~}{\pi }_1(k)`$ $`=`$ $`{\displaystyle \frac{\mathrm{sin}(k_1I/2)}{k_1/2}}e^{ik_1I/2},\stackrel{~}{\pi }_4(k)={\displaystyle \frac{\mathrm{sin}(k_4T/2)}{k_4/2}}e^{ik_4T/2}.`$ (66)
When use is made of
$`\mathrm{sin}\left({\displaystyle \frac{k_1I}{2}}\right)\mathrm{sin}\left({\displaystyle \frac{k_1}{2}}\right)=\mathrm{sin}^2\left({\displaystyle \frac{k_1(I+1)}{4}}\right)\mathrm{sin}^2\left({\displaystyle \frac{k_1(I1)}{4}}\right),`$ (67)
we get for large $`I`$ and $`T`$
$`\pi _1^{}(k)\stackrel{~}{\pi }_1(k)\pi _4^{}(k)\stackrel{~}{\pi }_4(k)=(2\pi )^2\delta (k_1)\delta (k_4)e^{i(1I)k_1/2}e^{i(1T)k_4/2}.`$ (68)
We find for large $`I`$ and $`T`$ that $`_s^{}D_{\mu \nu }^1(ss^{})N_\nu (s^{})`$ is equivalent to $`B_\mu (s)`$ in (17).
Hence the expectation value of the naive Wilson loop $`\stackrel{~}{W}_m(๐)`$ in (57) coincides with that of the perfect operator in (41). Now we introduce the dual string variable $`{}_{}{}^{}\sigma `$ as follows:
$`\sigma _{\mu \nu }(s)`$ $`=`$ $`{\displaystyle \frac{1}{2}}ฯต_{\mu \nu \alpha \beta }{}_{}{}^{}\sigma _{\alpha \beta }^{}(s+\widehat{\alpha }+\widehat{\beta }),`$ (69)
$`{}_{}{}^{}\overline{\sigma }_{\alpha \beta }^{}(s)`$ $``$ $`{}_{}{}^{}\sigma _{\alpha \beta }^{}+S_{\alpha \beta }(s).`$ (70)
Then the expectation value (57) can be expressed simply as
$`\stackrel{~}{W}_m(๐)`$ $`=`$ $`{\displaystyle \frac{_{{}_{}{}^{}\overline{\sigma }_{\mu \nu }^{}(s)=\mathrm{}}^{\mathrm{}}\mathrm{\Pi }_{\genfrac{}{}{0pt}{}{s}{\mu ,\nu }}\delta (_\mu ^{}{}_{}{}^{}\overline{\sigma }_{\mu \nu }^{}(s)\stackrel{~}{J}_\nu (s))\mathrm{exp}\{S({}_{}{}^{}\overline{\sigma })\}}{_{{}_{}{}^{}\overline{\sigma }_{\mu \nu }^{}(s)=\mathrm{}}^{\mathrm{}}\mathrm{\Pi }_{\genfrac{}{}{0pt}{}{s}{\mu ,\nu }}\delta (_\mu ^{}{}_{}{}^{}\overline{\sigma }_{\mu \nu }^{}(s))\mathrm{exp}\{S({}_{}{}^{}\overline{\sigma })\}}},`$ (71)
where
$`S({}_{}{}^{}\overline{\sigma })`$ $`=`$ $`\pi ^2{\displaystyle \underset{\genfrac{}{}{0pt}{}{s,s^{}}{\mu ,\nu }}{}}(_\alpha \mathrm{\Delta }_L^1(ss_1){\displaystyle \frac{1}{2}}ฯต_{\alpha \mu \beta \gamma }{}_{}{}^{}\overline{\sigma }_{\beta \gamma }^{}(s_1+\widehat{\alpha }+\widehat{\mu }))D_{\mu \nu }^1(ss^{})`$ (73)
$`\times (_\beta \mathrm{\Delta }_L^1(s^{}s_2){\displaystyle \frac{1}{2}}ฯต_{\beta \nu \eta \delta }{}_{}{}^{}\overline{\sigma }_{\eta \delta }^{}(s_2+\widehat{\beta }+\widehat{\nu })).`$
The strong coupling expansion can be shown in Fig.1. The leading term is the same as the classical contribution in (61). The next-leading term is a house-type diagram with one $`1\times 1`$ cube attached on the flat surface. Then the open surface variable $`{}_{}{}^{}\overline{\sigma }`$ has four more plaquettes than the leading one. If the self coupling term between the string variables $`{}_{}{}^{}\overline{\sigma }`$ is dominant as in $`SU(2)`$ QCD, the next to leading term is estimated as $`\mathrm{exp}(\sigma _{cl}ITb^24\mathrm{\Pi }(0)b^2)`$, where $`\sigma _{cl}`$ is the string tension from the classical part (52) and $`\mathrm{\Pi }(0)`$ is the self coupling constant. Considering the entropy factor $`4IT`$, we get
$`\stackrel{~}{W}_m(๐)`$ $`=`$ $`\mathrm{exp}\{\sigma _{cl}ITb^2\}+4IT\mathrm{exp}\{\sigma _{cl}ITb^24\mathrm{\Pi }(0)b^2\}+\mathrm{},`$ (74)
$``$ $`\mathrm{exp}\{\sigma _{cl}ITb^2\}\mathrm{exp}\{4ITe^{4\mathrm{\Pi }(0)b^2}\},`$ (75)
$``$ $`\mathrm{exp}\{\sigma ITb^2\}.`$ (76)
Hence the string tension becomes
$`\sigma =\sigma _{cl}{\displaystyle \frac{4}{b^2}}e^{4\mathrm{\Pi }(0)b^2}.`$ (77)
Applications to actual pure $`SU(2)`$ and $`SU(3)`$ QCD will be published elsewhere and the quantum fluctuation term will be found to be very small there.
###### Acknowledgements.
T.S. acknowledges the financial support from JSPS Grant-in Aid for Scientific Research (B) (No.10440073 and No.11695029). |
warning/0002/math0002143.html | ar5iv | text | # Reidemeister Torsion of 3-Dimensional Euler Structures with Simple Boundary Tangency and Pseudo-Legendrian Knots
## Introduction
Reidemeister torsion is a classical yet very vital topic in 3-dimensional topology, and it was recently used in a variety of important developments. To mention a few, torsion is a fundamental ingredient of the Casson-Walker-Lescop invariants (see e.g. ), and more generally of the perturbative approach to quantum invariants (see e.g. ). Relations have been pointed out between torsion and hyperbolic geometry . Turaevโs torsion of Euler structures has recently been recognized by Turaev himself () to have deep connections with the Seiberg-Witten invariants of $`\mathrm{Spin}^\mathrm{c}`$-structures on 3-manifolds, after the proof of Meng and Taubes that a suitable combination of these invariants can be identified with the classical Milnor torsion.
Turaevโs theory actually exists in all dimensions. We quickly review it before proceeding. A smooth Euler structure $`\xi `$ on a compact oriented manifold $`M`$, possibly with $`M=\mathrm{}`$, is a non-singular vector field on $`M`$ viewed up to local modifications in $`\mathrm{Int}(M)`$ and homotopy relative to $`M`$. Orientability of $`M`$ is not strictly necessary, but we find it convenient to assume it. Turaev allows only โmonochromaticโ boundary components, i.e. black ones (on which the field points outwards) and white ones (on which it points inwards). This implies the constraint that $`\chi (M,W)=0`$, where $`W`$ is the white portion of $`M`$, but in and Turaev only focuses on the more specialized case where $`M`$ is 3-dimensional and closed or bounded by tori. In all dimensions, the set $`\mathrm{Eul}^\mathrm{s}(M,W)`$ of smooth Euler structures compatible with $`(M,W)`$ is an affine space over $`H_1(M;\text{})`$. The two main ingredients of Turaevโs theory are as follows. First, he defines a certain set of $`1`$-chains, called the space $`\mathrm{Eul}^\mathrm{c}(M,W)`$ of combinatorial Euler structures compatible with $`(M,W)`$, he shows that this is again affine over $`H_1(M;\text{})`$, and he describes an $`H_1(M;\text{})`$-equivariant bijection $`\mathrm{\Psi }:\mathrm{Eul}^\mathrm{c}(M,W)\mathrm{Eul}^\mathrm{s}(M,W)`$ called the reconstruction map. Second, for $`\xi \mathrm{Eul}^\mathrm{c}(M,W)`$ and for any representation $`\phi `$ of $`\pi _1(M)`$ into the units of a suitable ring $`\mathrm{\Lambda }`$ he defines a torsion invariant $`\tau ^\phi (M,\xi )`$, or more generally $`\tau ^\phi (M,\xi ,\text{h})`$, with values in $`K_1(\mathrm{\Lambda })/(\pm 1)`$. This invariant is by definition a lifting of the classical Reidemeister torsion (see ) $`\tau ^\phi (M)K_1(\mathrm{\Lambda })/(\pm \phi (\pi _1(M)))`$, and it satisfies the $`H_1(M;\text{})`$-equivariance formula
$$\tau ^\phi (M,\xi ^{},\text{h})=\tau ^\phi (M,\xi ,\text{h})\phi (\xi ^{}\xi )$$
(1)
where $`\xi ^{}\xi H_1(M;\text{})`$. For $`\xi \mathrm{Eul}^\mathrm{s}(M,W)`$ one defines $`\tau ^\phi (M,\xi )`$ as $`\tau ^\phi (M,\mathrm{\Psi }^1(\xi ))`$, and the $`H_1(M;\text{})`$-equivariance of the reconstruction map $`\mathrm{\Psi }`$ implies that formula (1) holds also for smooth structures. We emphasize that the definition of $`\mathrm{\Psi }`$ is based on an explicit geometric construction, but its bijectivity is only established through $`H_1(M;\text{})`$-equivariance. This makes the definition of torsion for smooth structures somewhat implicit.
In the present paper, and in other papers in preparation, we are concerned with generalizations and improvements of Turaevโs theory. Here we consider 3-manifolds. This work had two main initial aims. Our first aim was to find a geometric description of the map $`\mathrm{\Psi }^1`$, and hence to turn the computation of Turaevโs torsion into a more effective procedure, using our encoding of non-singular vector fields up to homotopy (also called โcombingsโ) in terms of branched standard spines. Our second aim was to define torsion invariants of pseudo-Legendrian pairs $`(v,L)`$, consisting of a link $`L`$ transversal to a non-singular vector field $`v`$, viewed up to pseudo-Legendrian isotopy, namely transversality-preserving simultaneous isotopy of $`L`$ and homotopy of $`v`$. For a given $`v`$ we will just say that $`L`$ is pseudo-Legendrian in $`(M,v)`$, or just in $`v`$. Note that an ordinary Legendrian link in a given oriented contact structure $`\xi `$ is pseudo-Legendrian in $`\xi ^{}`$, and Legendrian isotopy implies pseudo-Legendrian isotopy. A specific motivation to look for invariants of pseudo-Legendrian links comes from the remarkable relation recently discovered by Fintushel and Stern between the Alexander polynomial (i.e. Milnor torsion) of a knot $`KS^3`$ and (a suitable combination of) the Seiberg-Witten invariants of the โsurgeredโ $`4`$-manifold $`X_K`$ obtained using $`K`$ (and a suitable base $`4`$-manifold $`X`$). Both our initial aims lead us to consider Euler structures on $`3`$-manifolds $`M`$ (without restrictions on $`M`$) allowing simple tangency circles to $`M`$ of concave type (see Fig. 1 below). On the other hand it turns out that, to define torsion, the natural objects to deal with are Euler structures with convex tangency circles. It is a fortunate fact, peculiar of dimension 3, that there is a canonical way to associate a convex field to any simple (i.e. mixed concave and convex) one. This allows to define torsion for all smooth simple Euler structures, and eventually to achieve both the objectives we had in mind.
Let us now summarize the contents of this paper. The foundational part of our work consists in extending to the context of Euler structures with simple tangency the notions of combinatorial structure $`\mathrm{Eul}^\mathrm{c}`$ and reconstruction map $`\mathrm{\Psi }`$. This part follows the same scheme as and relies on technical results of Turaev. Our main contribution here is the proof that the natural transformations of a concave structure into a convex one, viewed at the smooth level and at the combinatorial level, actually correspond to each other under the reconstruction map (Theorem 1.8). After setting the foundations, we prove the following main result (stated informally here: see Sections 1 and 3 for precise definitions and statements).
###### Theorem 0.1
Let $`\xi `$ be an Euler structure with concave tangency circles. If $`P`$ is a branched standard spine which represents $`\xi `$, then $`P`$ allows to explicitly find a representative of $`\mathrm{\Psi }^1(\xi )\mathrm{Eul}^\mathrm{c}`$, and hence to compute the torsion of $`\xi `$ in terms of the finite combinatorial data which encode $`P`$.
After the foundations, we concentrate on pseudo-Legendrian knots $`(v,K)`$, assuming for simplicity the ambient manifold to be closed. The connection comes from the fact that the restriction of $`v`$ defines a concave Euler structure on the exterior $`E(K)`$ of $`K`$, with two parallel tangency lines on $`E(K)`$ determined by the framing defined by $`v`$ on $`K`$. We show that torsion, as an absolute invariant, contains (in a suitable sense) a lifting to pseudo-Legendrian knots of the classical Alexander invariant. Then we carefully analyze the relative information carried by torsion for pseudo-Legendrian pairs $`(v_0,K_0)`$ and $`(v_1,K_1)`$ such that $`(v_0,v_1)`$ are homotopic to each other and $`(K_0,K_1)`$ are framed-isotopic to each other. A relevant point which emerges from this analysis is that in general torsion does not provide a single-valued relative invariant, because the action of a certain mapping class group (which depends on the framed isotopy class only) must be taken into account. This leads us to the notion of โgoodโ framed knots, for which the action is trivial, and the study of torsion is simpler. We show that many knots are good (for instance, all knots in a homology sphere are good, and most knots with hyperbolic complement are good). Concentrating on good knots we then prove that in a homology sphere the relative torsion of two knots essentially coincides with the difference of their rotation numbers (Maslov indices), so torsion basically detects whether the knots are isotopic through pseudo-Legendrian immersions.
Moreover we analyze the effect on torsion of the framed first Reidemeister move (which does not change the framing but locally changes the winding number by $`\pm 2`$), and we show that for homology spheres the winding number is just the difference of Maslov indices, thus getting an alternative proof of the relation between torsion and rotation number. Using the fact (proved in ) that framed isotopy is generated by pseudo-Legendrian isotopy and the framed first Reidemeister move, we then obtain several interesting consequences, among which we state the following:
###### Theorem 0.2
Consider pseudo-Legendrian knots $`(v_0,K_0)`$ and $`(v_1,K_1)`$. Assume that $`K_0`$ is good and that the meridian of $`K_0`$ has infinite order in $`H_1(E(K_0);\text{})`$. Then the knots are pseudo-Legendrian isotopic if and only if they have trivial relative torsion invariants.
This paper is organized as follows. In Section 1 we provide the formal definitions of smooth and combinatorial Euler structure. In Section 2 we introduce torsion and state the equivariance property. In Section 3 we show how branched standard spines can be used for computing torsion. In Section 4 we specialize to pseudo-Legendrian knot exteriors and analyze torsion both as an absolute and as a relative invariant. In Section 5 we carry out a specific computation using the technology of Section 3. In Sections 1 to 4 proofs which are long and require the introduction of ideas and techniques not used elsewhere are omitted. Section 6 contains all these proofs.
We conclude this introduction by announcing related results which we have recently obtained and partially written down. In we extend to the case with boundary our combinatorial presentation of combed manifolds in terms of branched spines, and we provide similar presentations of framed and pseudo-Legendrian links, using $`\mathrm{C}^1`$ diagrams on branched spines. Some results from are actually used also in the present paper (see Sections 3 and 4). In we use the results of to develop an approach to torsion entirely based on combinatorial techniques, getting slightly different generalizations of Turaevโs theory. In we generalize the theory of Euler structures and (with some restrictions) of torsion to all dimensions and allowing any generic (Whitney-Morin-type) tangency to the boundary.
## 1 Euler structures
In this section we define smooth and combinatorial Euler structures and explain their correspondence. Fix once and for ever a compact oriented 3-manifold $`M`$, possibly with $`M=\mathrm{}`$. Using the Hauptvermutung, we will always freely intermingle the differentiable, piecewise linear and topological viewpoints. Homeomorphisms will always respect orientations. All vector fields mentioned in this paper will be non-singular, and they will be termed just fields for the sake of brevity.
#### Smooth and combinatorial Euler structures
We will call boundary pattern on $`M`$ a partition $`๐ซ=(W,B,V,C)`$ of $`M`$ where $`V`$ and $`C`$ are finite unions of disjoint circles, and $`W=B=VC`$. In particular, $`W`$ and $`B`$ are interiors of compact surfaces embedded in $`M`$. Even if $`๐ซ`$ can actually be determined by less data, e.g. the pair $`(W,V)`$, we will find it convenient to refer to $`๐ซ`$ as a quadruple. Points of $`W`$, $`B`$, $`V`$ and $`C`$ will be called white, black, convex and concave respectively. We define the set of smooth Euler structures on $`M`$ compatible with $`๐ซ`$, denoted by $`\mathrm{Eul}^\mathrm{s}(M,๐ซ)`$, as the set of equivalence classes of fields on $`M`$ which point inside on $`W`$, point outside on $`B`$ and have simple tangency to $`M`$ of convex type along $`V`$ and concave type along $`C`$, as shown in a cross-section in Fig. 1.
Two such fields are equivalent if they are obtained from each other by homotopy through fields of the same type and modifications supported into interior balls. The following variation on the Poincarรฉ-Hopf formula is established in Section 6:
###### Proposition 1.1
$`\mathrm{Eul}^\mathrm{s}(M,๐ซ)`$ is non-empty if and only if $`\chi (\overline{W})=\chi (M)`$.
We remark here that $`\chi (\overline{W})=\chi (W)`$, $`\chi (\overline{B})=\chi (B)`$, $`\chi (V)=\chi (C)=0`$ and $`\chi (W)+\chi (B)=\chi (M)=2\chi (M)`$, so there are various ways to rewrite the relation $`\chi (\overline{W})=\chi (M)`$, the most intrinsic of which is actually $`\chi (M)(\chi (\overline{W})\chi (C))=0`$ (see below for the reason).
Now, given $`\xi ,\xi ^{}\mathrm{Eul}^\mathrm{s}(M,๐ซ)`$ we can choose generic representatives $`v,v^{}`$, so that the set of points of $`M`$ where $`v^{}=v`$ is a union of loops contained in the interior of $`M`$. A standard procedure allows to give these loops a canonical orientation, thus getting an element $`\alpha ^\mathrm{s}(\xi ,\xi ^{})H_1(M;\text{})`$. The following result is easily obtained along the lines of the well-known analogue for closed manifolds.
###### Lemma 1.2
$`\alpha ^\mathrm{s}`$ is well-defined and turns $`\mathrm{Eul}^\mathrm{s}(M,๐ซ)`$ into an affine space over $`H_1(M;\text{})`$.
A (finite) cellularization $`๐`$ of $`M`$ is called suited to $`๐ซ`$ if $`VC`$ is a subcomplex, so $`W`$ and $`B`$ are unions of cells. Here and in the sequel by โcellโ we will always mean an open one. Let such a $`๐`$ be given. For $`\sigma ๐`$ define $`\mathrm{ind}(\sigma )=(1)^{\mathrm{dim}(\sigma )}`$. We define $`\mathrm{Eul}^\mathrm{c}(M,๐ซ)_๐`$ as the set of equivalence classes of integer singular 1-chains $`z`$ in $`M`$ such that
$$z=\underset{\sigma M(WV)}{}\mathrm{ind}(\sigma )p_\sigma $$
where $`p_\sigma \sigma `$ for all $`\sigma `$. Two chains $`z`$ and $`z^{}`$ with $`z=\mathrm{ind}(\sigma )p_\sigma `$ and $`z^{}=\mathrm{ind}(\sigma )p_\sigma ^{}`$ are defined to be equivalent if there exist $`\delta _\sigma :([0,1],0,1)(\sigma ,p_\sigma ,p_\sigma ^{})`$ such that
$$zz^{}+\underset{\sigma M(WV)}{}\mathrm{ind}(\sigma )\delta _\sigma $$
represents $`0`$ in $`H_1(M;\text{})`$. Elements of $`\mathrm{Eul}^\mathrm{c}(M,๐ซ)_๐`$ are called combinatorial Euler structures relative to $`๐ซ`$ and $`๐`$, and their representatives are called Euler chains. The definition implies that, for $`\xi ,\xi ^{}\mathrm{Eul}^\mathrm{c}(M,๐ซ)_๐`$, their difference $`\xi \xi ^{}`$ can be defined as an element $`\alpha ^\mathrm{c}(\xi ,\xi ^{})`$ of $`H_1(M;\text{})`$. The following is easy:
###### Lemma 1.3
$`\mathrm{Eul}^\mathrm{c}(M,๐ซ)_๐`$ is non-empty if and only if $`\chi (\overline{W})=\chi (M)`$, and in this case $`\alpha ^\mathrm{c}`$ turns it into an affine space over $`H_1(M;\text{})`$.
Since $`\overline{W}=WVC`$, the alternating sum of dimensions of cells in $`WV`$ is intrinsically interpreted as $`\chi (\overline{W})\chi (C)`$, which explains why the most meaningful way to write the relation $`\chi (\overline{W})=\chi (M)`$ is $`\chi (M)(\chi (\overline{W})\chi (C))=0`$. From now on we will always assume that this relation holds. Turaev only considers the case where $`V=C=\mathrm{}`$, so $`W=\overline{W}`$ and $`B=\overline{B}`$, and our relation takes the usual form $`\chi (M,W)=0`$. The following result was established by Turaev in in his setting, but the proof extends verbatim to our context, so we omit it. Only the first assertion is hard. We state the other two because we will use them.
###### Proposition 1.4
1. If $`๐^{}`$ is a subdivision of $`๐`$ then there exists a canonical $`H_1(M;\text{})`$-isomorphism $`\mathrm{Eul}^\mathrm{c}(M,๐ซ)_๐\mathrm{Eul}^\mathrm{c}(M,๐ซ)_๐^{}`$. In particular $`\mathrm{Eul}^\mathrm{c}(M;\text{})`$ is canonically defined up to $`H_1(M;\text{})`$-isomorphism independently of the cellularization.
2. If $`๐`$ is a cellularization of $`M`$ suited to $`๐ซ`$ and $`x_0M`$ is an assigned point, any element of $`\mathrm{Eul}^\mathrm{c}(M,๐ซ)`$ can be represented, with respect to $`๐`$, as a sum $`_{\sigma M(WV)}\mathrm{ind}(\sigma )\beta _\sigma `$ with $`\beta _\sigma :([0,1],0,1)(M,x_0,\sigma )`$.
3. If $`๐ฏ`$ is a triangulation of $`M`$ suited to $`๐ซ`$, any element of $`\mathrm{Eul}^\mathrm{c}(M,๐ซ)`$ can be represented, with respect to $`๐ฏ`$, as a simplicial $`1`$-chain in the first barycentric subdivision of $`๐ฏ`$.
Our first main result, proved in Section 6, is the extension to the case under consideration of Turaevโs correspondence between $`\mathrm{Eul}^\mathrm{c}`$ and $`\mathrm{Eul}^\mathrm{s}`$.
###### Theorem 1.5
There exists a canonical $`H_1(M;\text{})`$-equivariant isomorphism
$$\mathrm{\Psi }:\mathrm{Eul}^\mathrm{c}(M,๐ซ)\mathrm{Eul}^\mathrm{s}(M,๐ซ).$$
The definition of $`\mathrm{\Psi }`$ is based on an explicit geometric construction, but its bijectivity is only established through $`H_1(M;\text{})`$-equivariance. As already mentioned in the introduction, this makes in general a very difficult task to determine the inverse of $`\mathrm{\Psi }`$. One of the features of this paper is the description of $`\mathrm{\Psi }^1`$ in terms of the combinatorial encoding of fields by means of branched spines: Theorem 3.7 describes $`\mathrm{\Psi }^1`$ when $`๐ซ`$ is concave, and Theorem 1.8 shows that from a general $`๐ซ`$ we can effectively pass to a unique convex $`๐ซ`$, and hence to a unique concave $`๐ซ`$, and conversely.
In view of Theorem 1.5, when no confusion risks to arise, we shortly write $`\mathrm{Eul}(M,๐ซ)`$ for either $`\mathrm{Eul}^\mathrm{s}(M,๐ซ)`$ or $`\mathrm{Eul}^\mathrm{c}(M,๐ซ)`$, and $`\alpha `$ for the map giving the affine $`H_1(M;\text{})`$-structure on this space.
#### Convex Euler structure associated to an arbitrary one
Let $`M`$ and $`๐ซ=(W,B,V,C)`$ be as in the definition of $`\mathrm{Eul}(M,๐ซ)`$. The pattern $`\theta (๐ซ)=(W,B,VC,\mathrm{})`$ is a convex one canonically associated to $`๐ซ`$. We define a map
$$\mathrm{\Theta }^\mathrm{s}:\mathrm{Eul}^\mathrm{s}(M,๐ซ)\mathrm{Eul}^\mathrm{s}(M,\theta (๐ซ))$$
as geometrically described in Fig. 2. Concerning this figure, note that the loops in $`C`$ can be oriented as components of the boundary of $`B`$, which is oriented as a subset of the boundary of $`M`$.
###### Lemma 1.6
$`\mathrm{\Theta }^\mathrm{s}`$ is a well-defined $`H_1(M;\text{})`$-equivariant bijection.
Proof of1.6. The first two properties are easy and imply the third property. The inverse of $`\mathrm{\Theta }^\mathrm{s}`$ may actually be described geometrically by a figure similar to Fig. 2, but we leave this to the reader. 1.6
We define now a combinatorial version of $`\mathrm{\Theta }^\mathrm{s}`$. Consider a cellularization $`๐`$ suited to $`๐ซ`$, and denote by $`\gamma _1,\mathrm{},\gamma _n`$ the 1-cells contained in $`C`$. We choose the parameterizations $`\gamma _j:(0,1)C`$ so that they respect the natural orientation of $`C`$ already discussed above, and we extend the $`\gamma _j`$ to $`[0,1]`$, without changing notation. Now let $`z`$ be an Euler chain relative to $`๐ซ`$. It easily seen that $`z_{j=1}^n\gamma _j|_{[1/2,1]}`$ is an Euler chain relative to $`\theta (๐ซ)`$. Setting
$$\mathrm{\Theta }^\mathrm{c}([z])=\left[z\underset{j=1}{\overset{n}{}}\gamma _j|_{[1/2,1]}\right]$$
we get a map $`\mathrm{\Theta }^\mathrm{c}:\mathrm{Eul}^\mathrm{c}(M,๐ซ)\mathrm{Eul}^\mathrm{c}(M,\theta (๐ซ))`$.
###### Lemma 1.7
$`\mathrm{\Theta }^\mathrm{c}`$ is a well-defined $`H_1(M;\text{})`$-equivariant bijection.
Proof of1.7. Again, the first two properties are easy and imply the third one. 1.7
In Section 6 we will see the following:
###### Theorem 1.8
If $`\mathrm{\Psi }`$ is the reconstruction map of Theorem 1.5 then the following diagram is commutative:
$$\begin{array}{ccc}\mathrm{Eul}^\mathrm{c}(M,๐ซ)& \stackrel{\mathrm{\Theta }^\mathrm{c}}{}& \mathrm{Eul}^\mathrm{c}(M,\theta (๐ซ))\\ \mathrm{\Psi }& & \mathrm{\Psi }\\ \mathrm{Eul}^\mathrm{s}(M,๐ซ)& \stackrel{\mathrm{\Theta }^\mathrm{s}}{}& \mathrm{Eul}^\mathrm{s}(M,\theta (๐ซ)).\end{array}$$
Using this result we will sometimes just write $`\mathrm{\Theta }:\mathrm{Eul}(M,๐ซ)\mathrm{Eul}(M,\theta (๐ซ))`$.
## 2 Torsion of an Euler structure
In this section we define torsion. We set up the usual algebraic environment in which torsion can be defined, fixing a ring $`\mathrm{\Lambda }`$ with unit, with the property that if $`n`$ and $`m`$ are distinct positive integers then $`\mathrm{\Lambda }^n`$ and $`\mathrm{\Lambda }^m`$ are not isomorphic as $`\mathrm{\Lambda }`$-modules. The Whitehead group $`K_1(\mathrm{\Lambda })`$ is defined as the Abelianization of $`\mathrm{GL}_{\mathrm{}}(\mathrm{\Lambda })`$, and $`\overline{K}_1(\mathrm{\Lambda })`$ is the quotient of $`K_1(\mathrm{\Lambda })`$ under the action of $`1\mathrm{GL}_1(\mathrm{\Lambda })=\mathrm{\Lambda }_{}`$. (Later in this paper the symbol $`K_1`$ will also be used for a knot, but the meaning will always be clear from the context.)
We will directly define torsion only for a convex Euler structure, but the definition easily extends to any Euler structure $`\xi `$ with simple boundary tangency, taking the torsion of the convexified structure $`\mathrm{\Theta }(\xi )`$. So, we fix a manifold $`M`$, a convex boundary pattern $`๐ซ=(W,B,V,\mathrm{})`$ on $`M`$, a cellularization $`๐`$ suited to $`๐ซ`$ and a representation $`\phi :\pi _1(M)\mathrm{\Lambda }_{}`$. We will denote by $`\phi `$ again the extension $`\text{}[\pi _1(M)]\mathrm{\Lambda }`$ (a ring homomorphism).
We consider now the universal cover $`q:\stackrel{~}{M}M`$ and the twisted chain complex $`\mathrm{C}_{}^\phi (M,WV)`$, where $`\mathrm{C}_i^\phi (M,WV)`$ is defined as $`\mathrm{\Lambda }_\phi \mathrm{C}_i^{\mathrm{cell}}(\stackrel{~}{M},q^1(WV);\text{})`$, and the boundary operator is induced from the ordinary boundary. The homology of this complex is denoted by $`H_{}^\phi (M,WV)`$ and called the $`\phi `$-twisted homology. We assume that each $`H_i^\phi (M,WV)`$ is a free $`\mathrm{\Lambda }`$-module and fix a basis $`\text{h}_i`$.
###### Remark 2.1
1. To have a formal completely intrinsic definition of $`H_{}^\phi (M,WV)`$, one should fix from the beginning a basepoint $`x_0M`$ for $`\pi _1(M)`$, and consider pointed universal covers $`q:(\stackrel{~}{M},\stackrel{~}{x}_0)(M,x_0)`$, because any two such covers are canonically isomorphic, and the action of $`\pi _1(M)`$ on $`\stackrel{~}{M}`$ is canonically defined on them.
2. To define $`H_{}^\phi (M,WV)`$ we have used in an essential way the fact that $`WV=\overline{W}`$ is closed, because otherwise $`\mathrm{C}_{}^\phi (M,WV)`$ cannot be defined.
3. $`\mathrm{C}_i^\phi (M,WV)`$ is a free $`\mathrm{\Lambda }`$-module, and each $`\text{}[\pi _1(M)]`$-basis of $`\mathrm{C}_i^{\mathrm{cell}}(\stackrel{~}{M},q^1(WV);\text{})`$ determines a $`\mathrm{\Lambda }`$-basis of $`\mathrm{C}_i^\phi (M,WV)`$.
4. If we compose $`\phi `$ with the projection $`\mathrm{\Lambda }_{}\overline{K}_1(\mathrm{\Lambda })`$ we get a homomorphism of $`\pi _1(M)`$ into an Abelian group, so we get a homomorphism $`\overline{\phi }:H_1(M;\text{})\overline{K}_1(\mathrm{\Lambda }).`$
Now let $`\xi \mathrm{Eul}^\mathrm{c}(M,๐ซ)`$ and choose a representative of $`\xi `$ as in point 2 of Proposition 1.4, namely
$$\underset{\sigma ๐,\sigma M(WV)}{}\mathrm{ind}(\sigma )\beta _\sigma $$
with $`\beta _\sigma (0)=x_0`$ for all $`\sigma `$, $`x_0`$ being a fixed point of $`M`$. We choose $`\stackrel{~}{x}_0q^1(x_0)`$ and consider the liftings $`\stackrel{~}{\beta }_\sigma `$ which start at $`\stackrel{~}{x}_0`$. For $`\sigma M(WV)`$ we select its preimage $`\stackrel{~}{\sigma }`$ which contains $`\stackrel{~}{\beta }_\sigma (1)`$, and define $`\text{g}(\xi )`$ as the collection of all these $`\stackrel{~}{\sigma }`$. Arranging the $`i`$-dimensional elements of $`\text{g}(\xi )`$ in any order, by Remark 2.1(3) we get a $`\mathrm{\Lambda }`$-basis $`\text{g}_i(\xi )`$ of $`\mathrm{C}_i^\phi (M,WV)`$. We consider a set $`\stackrel{~}{\text{h}}_i`$ of elements of $`\mathrm{C}_i^\phi (M,WV)`$ which project to the fixed basis $`\text{h}_i`$ of $`H_i^\phi (M,WV)`$.
Now note that, given a free $`\mathrm{\Lambda }`$-module $`L`$ and two finite bases $`\text{b}=(b_k)`$, $`\text{b}^{}=(b_k^{})`$ of $`M`$, the assumption made on $`\mathrm{\Lambda }`$ guarantees that b and $`\text{b}^{}`$ have the same number of elements, so there exists an invertible square matrix $`(\lambda _k^h)`$ such that $`b_k^{}=_h\lambda _k^hb_h`$. We will denote by $`[\text{b}^{}/\text{b}]`$ the image of $`(\lambda _k^h)`$ in $`K_1(\mathrm{\Lambda })`$.
###### Proposition 2.2
If $`\text{b}_i\mathrm{C}_i^\phi (M,WV)`$ is such that $`\text{b}_i`$ is a $`\mathrm{\Lambda }`$-basis of $`(\mathrm{C}_i^\phi (M,WV))`$, then $`(\text{b}_{i+1})\stackrel{~}{\text{h}}_i\text{b}_i`$ is a $`\mathrm{\Lambda }`$-basis of $`\mathrm{C}_i^\phi (M,WV)`$, and
$$\tau ^\phi (M,๐ซ,\xi ,\text{h})=\pm \underset{i=0}{\overset{3}{}}\left[\left((\text{b}_{i+1})\stackrel{~}{\text{h}}_i\text{b}_i\right)/\text{g}_i(\xi )\right]^{(1)^i}\overline{K}_1(\mathrm{\Lambda })$$
is independent of all choices made. Moreover
$$\tau ^\phi (M,๐ซ,\xi ^{},\text{h})=\tau ^\phi (M,๐ซ,\xi ,\text{h})\overline{\phi }(\alpha ^\mathrm{c}(\xi ^{},\xi )).$$
(2)
Proof of2.2. The first assertion and independence of the $`\text{b}_i`$โs is purely algebraic and classical, see . Now note that $`\xi \mathrm{Eul}^\mathrm{c}(M,๐ซ)`$ was used to select the bases $`\text{g}_i(\xi )`$. The $`\text{g}_i(\xi )`$ are of course not uniquely determined themselves, but we can show that different choices lead to the same value of $`\tau ^\phi `$.
First of all, the arbitrary ordering in the $`\text{g}_i(\xi )`$ is inessential because torsion is only regarded up to sign. Second, consider the effect of choosing a different representative of $`\xi `$. This leads to a new family $`\stackrel{~}{\sigma }^{}`$ of cells. If $`\stackrel{~}{\sigma }^{}=a(\sigma )\stackrel{~}{\sigma }`$, with $`a(\sigma )\pi _1(M)`$, and $`\overline{a}(\sigma )`$ is the image in $`H_1(M;\text{})`$, we automatically have
$$\underset{\sigma MWV}{}\mathrm{ind}(\sigma )\overline{a}(\sigma )=0H_1(M;\text{}),$$
which allows to conclude that also the representative chosen is inessential. The choice of the lifting $`\stackrel{~}{x}_0`$ can be shown to be inessential either in the spirit of Remark 2.1(1), or by showing that a simultaneous $`a`$-translation of all $`\stackrel{~}{\sigma }`$, for $`a\pi _1(M)`$, multiplies the torsion by $`\overline{\phi }(a)^{\chi (M)\chi (WV)}=1`$.
Formula (2) is readily established by choosing representatives $`\mathrm{ind}(\sigma )\beta _\sigma `$ and $`\mathrm{ind}(\sigma )\beta _\sigma ^{}`$ of $`\xi `$ and $`\xi ^{}`$ such that $`\beta _\sigma ^{}=\beta _\sigma `$ for all $`\sigma `$ but one. 2.2
Since the above construction uses the cellularization $`๐`$ in a way which may appear to be essential, we add a subscript $`๐`$ to the torsion we have defined. The next result, which can be established following Turaev , shows that dependence on $`๐`$ is actually inessential.
###### Proposition 2.3
Let $`๐`$ and $`๐^{}`$ be cellularizations suited to $`๐ซ`$. Assume that $`๐^{}`$ subdivides $`๐`$, and consider the bijection $`๐ฎ_{(๐^{},๐)}:\mathrm{Eul}^\mathrm{c}(M,๐ซ)_๐\mathrm{Eul}^\mathrm{c}(M,๐ซ)_๐^{}`$ of Proposition 1.4, and the canonical isomorphism $`j_{(๐^{},๐)}:H_{}^\phi (M,WV)_๐H_{}^\phi (M,WV)_๐^{}`$. Then, with obvious meaning of symbols we have:
$$\tau _๐^\phi (M,๐ซ,\xi ,\text{h})=\tau _๐^{}^\phi (M,๐ซ,๐ฎ_{(๐^{},๐)}(\xi ),j_{(๐^{},๐)}(\text{h})).$$
It is maybe appropriate here to remark that the choice of a basis h of $`H_{}^\phi (M,WV)`$ and the definition of $`\tau ^\phi (M,๐ซ,\xi ,\text{h})`$ implicitly assume a description of the universal cover of $`M`$, which is typically undoable in practical cases. However, if one starts from a representation of $`\pi _1(M)`$ into the units of a commutative ring $`\mathrm{\Lambda }`$, i.e. a representation which factors through one of $`H_1(M;\text{})`$, one can use from the very beginning the maximal Abelian rather than the universal cover, which makes computations more feasible.
###### Remark 2.4
Turaev has shown that a homological orientation yields a sign-refinement of torsion, i.e. a lifting from $`\overline{K}_1(\mathrm{\Lambda })`$ to $`K_1(\mathrm{\Lambda })`$. This refinement extends with minor modifications to our setting of boundary tangency. This sign-refinement, in the closed and monochromatic case, is often an essential component of the theory (for instance, it is crucial for the relation with the 3-dimensional Seiberg-Witten invariants and for the definition of the Casson invariant ), but we will not address it in the present paper.
#### Computation of torsion via disconnected spiders
In this paragraph we show that to determine the family of lifted cells necessary to define torsion one can use representatives of Euler structures more general than those used above. This is a technical point which we will use below to compute torsions using branched spines (Section 3).
We fix $`M`$, $`๐ซ`$, $`๐`$ and $`\phi `$ as above, and $`\xi \mathrm{Eul}^\mathrm{c}(M,๐ซ)`$. Let $`\text{g}(\xi )=\{\stackrel{~}{\sigma }\}`$ be the family of liftings of the cells lying in $`M(WV)`$ determined by a connected spider as explained above. Note that if $`\text{g}^{}=\{\stackrel{~}{\sigma }^{}\}`$ is any other family of liftings we have $`\stackrel{~}{\sigma }^{}=a(\sigma )\stackrel{~}{\sigma }`$ for some $`a\pi _1(M)`$, and we can define
$$h(\text{g}^{},\text{g}(\xi ))=\underset{\sigma M(WV)}{}\mathrm{ind}(\sigma )\overline{a}(\sigma )H_1(M;\text{}).$$
###### Proposition 2.5
Assume there exists a partition $`๐_1\mathrm{}๐_k`$ of the set of cells lying in $`M(WV)`$, and let $`\xi \mathrm{Eul}^\mathrm{c}(M,๐ซ)`$ have a representative of the form
$$z=\underset{j=1}{\overset{k}{}}\left(\underset{\sigma ๐_j\{\sigma _j\}}{}\mathrm{ind}(\sigma )\gamma _\sigma ^{(j)}\right)$$
where $`\sigma _j๐_j`$ and $`\gamma _\sigma ^{(j)}:([0,1],0,1)(M,p_{\sigma _j},p_\sigma )`$. Choose any lifting $`\stackrel{~}{p}_{\sigma _j}`$ of $`p_{\sigma _j}`$, lift $`\gamma _\sigma ^{(j)}`$ to $`\stackrel{~}{\gamma }_\sigma ^{(j)}`$ starting from $`\stackrel{~}{p}_{\sigma _j}`$, let $`\stackrel{~}{\sigma }^{}`$ be the lifting of $`\sigma `$ containing $`\stackrel{~}{\gamma }_\sigma ^{(j)}(1)`$, and let $`\text{g}^{}`$ be the family of all these liftings. Then $`h(\text{g}^{},\text{g}(\xi ))=0H_1(M;\text{})`$. In particular $`\text{g}^{}`$ can be used to compute $`\tau ^\phi (M,๐ซ,\xi ,\text{h})`$.
Proof of2.5. Note first that the coefficient of $`p_{\sigma _j}`$ in $`z`$ is exactly
$$\underset{\sigma ๐_j\{\sigma _j\}}{}\mathrm{ind}(\sigma ).$$
On the other hand this coefficient must be equal to $`\mathrm{ind}(\sigma _j)`$. Summing up we deduce that $`_{\sigma ๐_j}\mathrm{ind}(\sigma )=0`$.
Now choose $`x_0M`$ and $`\delta ^{(j)}:([0,1],0,1)(M,x_0,p_{\sigma _j})`$. For $`\sigma ๐_j`$ define
$$\beta _\sigma =\{\begin{array}{cc}\delta ^{(j)}\hfill & \text{if }\sigma =\sigma _j\hfill \\ \delta ^{(j)}\gamma _\sigma ^{(j)}\hfill & \text{otherwise,}\hfill \end{array}$$
so that $`\beta _\sigma :([0,1],0,1)(M,x_0,p_\sigma )`$, whence $`w=_{\sigma M(WV)}\beta _\sigma `$ is an Euler chain. Moreover:
$$wz=\underset{j=1}{\overset{k}{}}\left(\underset{\sigma ๐_j}{}\mathrm{ind}(\sigma )\right)\delta ^{(j)}=0H_1(M;\text{}),$$
so $`[w]=\xi `$. Now choose $`\stackrel{~}{x}_0`$ over $`x_0`$, lift the $`\delta ^{(j)}`$ and $`\beta _\sigma `$ starting from $`\stackrel{~}{x}_0`$, and let $`a^{(j)}\pi _1(M)`$ be such that $`\stackrel{~}{p}_{\sigma _j}=a^{(j)}\stackrel{~}{\delta }^{(j)}(1)`$. Then
$$h(\text{g}^{},\text{g}(\xi ))=\underset{j=1}{\overset{k}{}}\left(\underset{\sigma ๐_j}{}\mathrm{ind}(\sigma )\right)\overline{a}^{(j)}=0H_1(M;\text{}),$$
and the proof is complete. 2.5
The next result follows directly from the definition, but it is worth stating because it shows how torsions may be used to distinguish triples $`(M,๐ซ,\xi )`$ from each other.
###### Proposition 2.6
Let $`f:MM^{}`$ be a homeomorphism, consider $`\xi \mathrm{Eul}(M,๐ซ)`$, $`\phi :\pi _1(M)\mathrm{\Lambda }_{}`$ and a $`\mathrm{\Lambda }`$-basis h of $`H_{}^\phi (M,\overline{W})`$. Then
$$\tau ^{\phi f_{}^1}(M^{},f_{}(๐ซ),f_{}(\xi ),f_{}(\text{h}))=\tau ^\phi (M,๐ซ,\xi ,\text{h}).$$
## 3 Spines and computation of torsion
In this section we show how to geometrically invert the reconstruction map $`\mathrm{\Psi }`$, and how to compute torsions starting from a combinatorial encoding of vector fields. We first review the theory developed in . See the beginning of Section 1 for our conventions on manifolds, maps, and fields. In addition to the terminology introduced there, we will need the notion of traversing field on a manifold $`M`$, defined as a field whose orbits eventually intersect $`M`$ transversely in both directions (in other words, orbits are compact intervals).
#### Branched spines
A simple polyhedron $`P`$ is a finite connected 2-dimensional polyhedron with singularity of stable nature (triple lines and points where six non-singular components meet). Such a $`P`$ is called standard if all the components of the natural stratification given by singularity are open cells. Depending on dimension, we will call the components vertices, edges and regions.
A standard spine of a $`3`$-manifold $`M`$ with $`M\mathrm{}`$ is a standard polyhedron $`P`$ embedded in $`\mathrm{Int}(M)`$ so that $`M`$ collapses onto $`P`$. Standard spines of oriented $`3`$-manifolds are characterized among standard polyhedra by the property of carrying an orientation, defined (see Definition 2.1.1 in ) as a โscrew-orientationโ along the edges (as in the left-hand-side of Fig. 3), with an obvious compatibility at vertices (as in the centre of Fig. 3).
It is the starting point of the theory of standard spines that every oriented $`3`$-manifold $`M`$ with $`M\mathrm{}`$ has an oriented standard spine, and can be reconstructed (uniquely up to homeomorphism) from any of its oriented standard spines. See for the non-oriented version of this result and or Proposition 2.1.2 in for the (slight) oriented refinement.
A branching on a standard polyhedron $`P`$ is an orientation for each region of $`P`$, such that no edge is induced the same orientation three times. See the right-hand side of Fig. 3 and Definition 3.1.1 in for the geometric meaning of this notion. An oriented standard spine $`P`$ endowed with a branching is shortly named branched spine. We will never use specific notations for the extra structures: they will be considered to be part of $`P`$. The following result, proved as Theorem 4.1.9 in , is the starting point of our constructions.
###### Proposition 3.1
To every branched spine $`P`$ there corresponds a manifold $`M(P)`$ with non-empty boundary and a concave traversing field $`v(P)`$ on $`M(P)`$. The pair $`(M(P),v(P))`$ is well-defined up to diffeomorphism. Moreover an embedding $`i:P\mathrm{Int}(M(P))`$ is defined, and has the property that $`v(P)`$ is positively transversal to $`i(P)`$.
The topological construction which underlies this proposition is actually quite simple, and it is illustrated in Fig. 4. Concerning the last
assertion of the proposition, note that the branching allows to define an oriented tangent plane at each point of $`P`$.
#### Combinatorial encoding of combings
Let $`P`$ be a branched spine, and define $`v(P)`$ on $`M(P)`$ as just explained. Assume that in $`M(P)`$ there is only one component which is homeomorphic to $`S^2`$ and is split by the tangency line of $`v(P)`$ to $`M(P)`$ into two discs. (Such a component will be denoted by $`S_{\mathrm{triv}}^2`$.) Now, notice that $`S_{\mathrm{triv}}^2`$ is also the boundary of the closed $`3`$-ball with constant vertical field, denoted by $`B_{\mathrm{triv}}^3`$. This shows that we can cap off $`S_{\mathrm{triv}}^2`$ by attaching a copy of $`B_{\mathrm{triv}}^3`$, getting a compact manifold $`\widehat{M}(P)`$ and a field $`\widehat{v}(P)`$ on $`\widehat{M}(P)`$. If we denote by $`\widehat{๐ซ}(P)`$ the boundary pattern of $`\widehat{v}(P)`$ on $`\widehat{M}(P)`$, we easily see that the pair $`(\widehat{M}(P),\widehat{v}(P))`$ is only well-defined up to homeomorphism of $`\widehat{M}(P)`$ and homotopy of $`\widehat{v}(P)`$ through fields compatible with $`\widehat{๐ซ}(P)`$. Note also that $`\widehat{๐ซ}(P)`$ is automatically concave.
If $`๐ซ`$ is a boundary pattern on $`M`$, we define $`\mathrm{Comb}(M,๐ซ)`$ as the set of fields compatible with $`๐ซ`$ under homotopy through fields also compatible with $`๐ซ`$. An element of $`\mathrm{Comb}(M,๐ซ)`$ is called a combing on $`(M,๐ซ)`$. Note that we have a projection $`\mathrm{Comb}(M,๐ซ)\mathrm{Eul}(M,๐ซ)`$.
The above construction shows that a branched spine $`P`$ with only one $`S_{\mathrm{triv}}^2`$ on $`M(P)`$ defines an element $`\mathrm{\Phi }(P)`$ of $`\mathrm{Comb}(\widehat{M}(P),\widehat{๐ซ}(P))`$. In we will establish the following:
###### Theorem 3.2
If $`M`$ is any compact oriented $`3`$-manifold and $`๐ซ`$ is a concave boundary pattern on $`M`$ not containing $`S_{\mathrm{triv}}^2`$ components, then $`\mathrm{\Phi }`$ maps surjectively $`\{P:\widehat{M}(P)M,\widehat{๐ซ}(P)๐ซ\}`$ onto $`\mathrm{Comb}(M,๐ซ)`$.
This theorem generalizes the main achievement of (Theorems 1.4.1 and 5.2.1), where it is proved in the special case of closed $`M`$. The complete statement includes also the description of a finite set of local moves on branched spines generating the equivalence relation induced by $`\mathrm{\Phi }`$. We will not need the moves in this paper. The following geometric interpretation the theorem may however be of some interest.
###### Remark 3.3
In general, the dynamics of a field, even a concave one, can be very complicated, whereas the dynamics of a traversing field (in particular, $`B_{\mathrm{triv}}^3`$) is simple. Theorem 3.2 means that for any (complicated) concave field there exists a sphere $`S^2`$ which splits the field into two (simple) pieces: a standard $`B_{\mathrm{triv}}^3`$ and a concave traversing field.
We can give here an easy special proof of Theorem 3.2 for the case we are most interested in, namely link exteriors. Note that our argument relies on the results of .
Proof of 3.2 for link exteriors. We have to show that if $`M`$ is closed, $`v`$ is a field on $`M`$ and $`L`$ is transversal to $`v`$, then the exterior $`E(L)`$ of $`L`$ with the restricted field is represented by some branched spine in the sense explained above.
The construction explained in Section 5.1 of shows that there exists a branched standard spine $`P`$ such that $`v`$ is positively transversal to $`P`$, and the complement of $`P`$, with the restriction of $`v`$, is isomorphic to the open 3-ball with the constant vertical field. The last condition easily implies that $`L`$ can be isotoped through links transversal to $`v`$ to a link lying in an arbitrarily small neighbourhood of $`P`$, with the further property that its natural projection on $`P`$ is $`\mathrm{C}^1`$, possibly with crossings.
Once $`L`$ has been isotoped to a $`\mathrm{C}^1`$ link on $`P`$, a branched spine of $`(E(L),v|_{E(L)})`$ is obtained by digging a tunnel in $`P`$ along the projection of $`L`$, as shown in Fig. 5.
A crossing in the projection will of course give rise to 4 vertices in the spine. Note that the spine which results from the digging may occasionally be non-standard, but it is standard as soon as the projection is complicated enough (e.g. if on each component there are both a crossing and an intersection with $`S(P)`$). 3.2$`E(L)`$
#### Spines and ideal triangulations
We remind the reader that an ideal triangulation of a manifold $`M`$ with non-empty boundary is a partition $`๐ฏ`$ of $`\mathrm{Int}(M)`$ into open cells of dimensions 1, 2 and 3, induced by a triangulation $`๐ฏ^{}`$ of the space $`Q(M)`$, where:
1. $`Q(M)`$ is obtained from $`M`$ by collapsing each component of $`M`$ to a point;
2. $`๐ฏ^{}`$ is a triangulation only in a loose sense, namely self-adjacencies and multiple adjacencies of tetrahedra are allowed;
3. The vertices of $`๐ฏ^{}`$ are precisely the points of $`Q(M)`$ which correspond to components of $`M`$.
It turns out (see for instance ) that there exists a natural bijection between standard spines and ideal triangulations of a 3-manifold. Given an ideal triangulation, the corresponding standard spine is just the 2-skeleton of the dual cellularization, as illustrated in Figure 6.
The inverse of this correspondence will be denoted by $`P๐ฏ(P)`$.
Now let $`P`$ be a branched spine. First of all, we can realize $`๐ฏ(P)`$ in such a way that its edges are orbits of the restriction of $`v(P)`$ to $`\mathrm{Int}(M(P))`$, and the 2-faces are unions of such orbits. Being orbits, the edges of $`๐ฏ(P)`$ have a natural orientation, and the branching condition, as remarked in , is equivalent to the fact that on each tetrahedron of $`๐ฏ(P)`$ exactly one of the vertices is a sink and one is a source.
###### Remark 3.4
It turns out that if $`P`$ is a branched spine, not only the edges, but also the faces and the tetrahedra of $`๐ฏ(P)`$ have natural orientations. For tetrahedra, we just restrict the orientation of $`M(P)`$. For faces, we first note that the edges of $`P`$ have a natural orientation (the prevailing orientation induced by the incident regions). Now, we orient a face of $`๐ฏ(P)`$ so that the algebraic intersection in $`M(P)`$ with the dual edge is positive.
#### Euler chain defined by a branched spine
We fix in this paragraph a standard spine $`P`$ and consider its manifold $`M=M(P)`$. We start by noting that the ideal triangulation $`๐ฏ=๐ฏ(P)`$ defined by $`P`$ can be interpreted as a realization of $`\mathrm{Int}(M)`$ by face-pairings on a finite set of tetrahedra with vertices removed. If, instead of removing vertices, we remove open conic neighbourhoods of the vertices, thus getting truncated tetrahedra, after the face-pairings we obtain $`M`$ itself. This shows that $`P`$ determines a cellularization $`\overline{๐ฏ}=\overline{๐ฏ}(P)`$ of $`M`$ with vertices only on $`M`$ and 2-faces which are either triangles contained in $`M`$ or hexagons contained in $`\mathrm{Int}(M)`$, with edges contained alternatingly in $`M`$ and in $`\mathrm{Int}(M)`$.
Now assume that $`P`$ is branched and that $`M`$ contains only one $`S_{\mathrm{triv}}^2`$ component, so $`\widehat{M}=\widehat{M}(P)`$ is defined. Note that $`\widehat{M}`$ can be thought of as the space obtained from $`M`$ by contracting $`S_{\mathrm{triv}}^2`$ to a point, so a projection $`\pi :M\widehat{M}`$ is defined, and $`\pi (\overline{๐ฏ})`$ is a cellularization of $`\widehat{M}`$. Next, we modify $`\pi (\overline{๐ฏ})`$ by subdividing the triangles on $`\widehat{M}`$ as shown in Fig. 7.
The result is a cellularization $`\widehat{๐ฏ}=\widehat{๐ฏ}(P)`$ of $`\widehat{M}`$. Note that $`\widehat{๐ฏ}`$ on $`\widehat{M}`$ consists of โkitesโ, with long edges coming from tetrahedra and short edges coming from subdivision. Note also that $`\widehat{๐ฏ}`$ has exactly one vertex $`x_0`$ in $`\mathrm{Int}(\widehat{M})`$, and that the cells contained in $`\mathrm{Int}(\widehat{M})`$, except $`x_0`$, are the duals to the cells of the natural cellularization $`๐ฐ=๐ฐ(P)`$ of $`P`$. For $`u๐ฐ`$ we denote by $`\widehat{u}`$ its dual and by $`p_u=p_{\widehat{u}}`$ the point where $`u`$ and $`\widehat{u}`$ intersect, called the centre of both.
We will now use the field $`\widehat{v}=\widehat{v}(P)`$ to construct a combinatorial Euler chain on $`\widehat{M}`$ with respect to $`\widehat{๐ฏ}`$. It is actually convenient to consider, instead of $`\widehat{v}`$, the field $`\overline{v}=\pi (v)`$, which coincides with $`\widehat{v}`$ except near $`x_0`$, where it has a (removable) singularity. For $`u๐ฐ`$ we denote by $`\beta _u`$ the arc obtained by integrating $`\overline{v}(P)`$ in the positive direction, starting from $`p_u`$, until the boundary or the singularity is reached. We define:
$$s(P)=\underset{u๐ฐ}{}\mathrm{ind}(u)\beta _u.$$
Let us consider now the pattern $`\widehat{๐ซ}=\widehat{๐ซ}(P)=(W,B,\mathrm{},C)`$ defined by $`P`$. If $`p`$ is a vertex of $`\pi (\overline{๐ฏ})`$ contained in $`B`$, we define its star $`\mathrm{St}(p)`$ as the sum of the straight segments going from $`p`$ to the centres of all the kites containing $`p`$, minus the sum of the straight segments going from $`p`$ to the centres of all the long edges containing $`p`$. If $`\sigma `$ is an edge of $`\pi (\overline{๐ฏ})`$ contained in $`B`$ we define its bi-arrow $`\mathrm{Ba}(\sigma )`$ as the sum of the two straight segments going from the centre $`p_\sigma `$ of $`\sigma `$ to the centres of the two short kite-edges containing $`p_\sigma `$. A star and a bi-arrow are shown in Fig. 8.
We define:
$$s^{}(P)=s(P)+\underset{pB\overline{๐ฏ}(P)^{(0)}}{}\mathrm{St}(p)+\underset{\sigma \widehat{๐ฏ}(P)^{(1)},\sigma B}{}\mathrm{Ba}(\sigma ).$$
###### Lemma 3.5
$`s^{}(P)`$ defines an element of $`\mathrm{Eul}^\mathrm{c}(\widehat{M},\theta (P))`$.
Proof of3.5. Recall that $`\theta (\widehat{๐ซ})=(W,B,C,\mathrm{})`$, i.e. the concave line $`C`$ is turned into a convex one. So by definition we have to show that $`s^{}(P)`$ contains, with the right sign, the centres of all cells of $`\widehat{๐ฏ}`$ except those of $`WC`$.
It will be convenient to analyze first the natural lifting of $`s(P)`$ to $`M`$, denoted by $`\stackrel{~}{s}(P)=_{u๐ฐ}\mathrm{ind}(u)\stackrel{~}{\beta }_u`$ with obvious meaning of symbols. So
$$\stackrel{~}{s}(P)=\underset{u๐ฐ}{}\mathrm{ind}(u)\stackrel{~}{\beta }_u(0)+\underset{u๐ฐ}{}\mathrm{ind}(u)\stackrel{~}{\beta }_u(1).$$
(3)
Since the cellularization $`\overline{๐ฏ}`$ of $`M`$ is dual to $`๐ฐ`$, the first half of (3) gives the centres of the cells contained in $`\mathrm{Int}(M)`$, with right sign. One easily sees that the second half gives exactly the centres of the cells (of $`\overline{๐ฏ}`$) contained in $`B`$, also with right sign.
When we project to $`\widehat{M}`$ and consider $`s(P)`$, the first half of (3) again provides (with right sign) the centres of the all cells contained in $`\mathrm{Int}(\widehat{M})`$, except the special vertex $`x_0`$ obtained by collapsing $`S_{\mathrm{triv}}^2`$. We can further split the points of the second half of (3) into those which lie on $`S_{\mathrm{triv}}^2`$ and those which do not. The points of the first type project to $`x_0`$, and the resulting coefficient of $`x_0`$ is $`\chi (BS_{\mathrm{triv}}^2)`$, but $`BS_{\mathrm{triv}}^2`$ is an open 2-disc, so the coefficient is 1. (We are here using the very special property of dimension 2 that $`\chi `$ can be computed using a finite cellularization of an open manifold, because the boundary of the closure has $`\chi =0`$.) The points of the second type faithfully project to $`\widehat{M}`$, giving the centres of the simplices contained in $`B`$ of the triangulation $`\pi (\overline{๐ฏ})|_{\widehat{M}}`$. However $`\widehat{๐ฏ}`$ on $`\widehat{M}`$ is a subdivision of $`\pi (\overline{๐ฏ})`$, and this is the reason why we have added the stars and the bi-arrows to $`s(P)`$ getting $`s^{}(P)`$. The following computation of the coefficients in $`s^{}(P)`$ of the centres of the cells of $`\widehat{๐ฏ}`$ contained in $`B`$ concludes the proof.
1. Cells of dimension 0 are listed as follows:
1. Centres of triangles of $`\pi (\overline{๐ฏ})`$, which receive coefficient $`+1`$ from $`s(P)`$;
2. Midpoints of edges of $`\pi (\overline{๐ฏ})`$, which receive coefficient $`1`$ from $`s(P)`$ and $`+2`$ from the bi-arrows they determine;
3. Vertices of $`\pi (\overline{๐ฏ})`$ receive $`+1`$ from $`s(P)`$ and (algebraically) $`0`$ from the star they determine;
2. Cells of dimension 1 are:
1. Short edges of kites, whose midpoints receive $`1`$ from the bi-arrows;
2. Long edges of kites, whose midpoints receive $`1`$ from the stars;
3. Cells of dimension 2 are kites, and their centres receive $`+1`$ from the stars.
3.5
Now we denote by $`\gamma _j:(0,1)C`$, for $`j=1,\mathrm{},n`$, orientation-preserving parameterizations of the 1-cells of $`\widehat{๐ฏ}`$ contained in $`C`$, and we extend the $`\gamma _j`$ to $`[0,1]`$, without changing notation. We define
$$s^{\prime \prime }(P)=s^{}(P)+\underset{j=1}{\overset{n}{}}\gamma _j|_{[1/2,1]}.$$
###### Lemma 3.6
$`s^{\prime \prime }(P)`$ defines an element of $`\mathrm{Eul}^\mathrm{c}(\widehat{M},\widehat{๐ซ})`$, and
$$[s^{}(P)]=\mathrm{\Theta }^\mathrm{c}([s^{\prime \prime }(P)])\mathrm{Eul}^\mathrm{c}(\widehat{M},\theta (๐ซ)).$$
Proof of3.6. At the level of representatives, the second assertion is obvious, and it implies the first assertion. 3.6
We defer to Section 6 the proof of the next result, which shows that the map $`P[s^{\prime \prime }(P)]\mathrm{Eul}^\mathrm{c}(\widehat{M},\widehat{๐ซ})`$ allows, using branched spines, to explicitly find the inverse of the reconstruction map $`\mathrm{\Psi }`$ of Theorem 1.5. This result was informally announced as Theorem 0.1 in the Introduction.
###### Theorem 3.7
$`\mathrm{\Psi }([s^{\prime \prime }(P)])=[\widehat{v}(P)]\mathrm{Eul}^\mathrm{s}(\widehat{M},\widehat{๐ซ})`$.
Recall now that we have defined torsions directly only for convex patterns, and we have extended the definition to concave patterns via the map $`\mathrm{\Theta }`$. As a consequence of Lemma 3.6 and Theorem 3.7, and by direct inspection of $`s^{}(P)`$, we have the following result which summarizes our investigations on the relation between spines, Euler structures, and torsion:
###### Theorem 3.8
If $`P`$ is a branched spine which represents a manifold $`\widehat{M}`$ with concave boundary pattern $`\widehat{๐ซ}=(W,B,\mathrm{},C)`$ in the sense of Theorem 3.2, then for any representation $`\phi :\pi _1(M)\mathrm{\Lambda }_{}`$ and any $`\mathrm{\Lambda }`$-basis h of $`H_{}^\phi (\widehat{M},WC)`$, the torsion $`\tau ^\phi (\widehat{M},\widehat{๐ซ},[\widehat{v}(M)],\text{h})`$ can be computed using (in the sense of Proposition 2.5) the lifting to the universal cover of $`\widehat{M}`$ of the chain $`s^{}(P)`$ defined above. In particular, $`s^{}(P)`$ can be used directly, without replacing it by a connected spider.
#### Computational hints
To actually compute torsion starting from a branched spine $`P`$, besides describing the universal (or maximal Abelian) cover of $`\widehat{M}=\widehat{M}(P)`$ and determining the preferred liftings of the cells in $`\widehat{M}(WC)`$, one needs to compute the boundary operators in the twisted chain complex $`\mathrm{C}_{}^\phi (M,WC)`$. These operators are of course twisted liftings of the corresponding operators in the cellular chain complex of $`(\widehat{M},WC)`$, with respect to $`\widehat{๐ฏ}`$. We briefly describe here the form of the latter operators. Recall first that $`\widehat{๐ฏ}`$ consists of a special vertex $`x_0`$, the kites (with their vertices and edges) on $`\widehat{M}`$, and the duals of the cells of $`P`$. On $`\widehat{M}`$ the situation is easily described, so we consider the internal cells.
1. If $`R`$ is a region of $`P`$, the ends of its dual edge $`\widehat{R}`$ are either $`x_0`$ or vertices of $`\widehat{M}`$ contained only in long edges of kites.
2. If $`e`$ is an edge of $`P`$ then $`\widehat{e}`$ is given by $`\widehat{R}_1+\widehat{R}_2\widehat{R}_0`$ plus 3 long edges of kites, where $`R_0,R_1,R_2`$ are the regions incident to $`e`$, numbered so that $`R_1`$ and $`R_2`$ induce on $`e`$ the same orientation. Here $`R_0,R_1,R_2`$ need not be different from each other, so the formula may actually have some cancellation. The 3 long edges of kites must be given an appropriate sign, and some of them may actually be collapsed to the point $`x_0`$. Note that we have only 3 kite-edges, out of the 6 which geometrically appear on $`\widehat{e}`$, because the other 3 are white.
3. If $`v`$ is a vertex of $`P`$ then $`\widehat{v}`$ is given by $`\widehat{e}_1+\widehat{e}_2\widehat{e}_3\widehat{e}_4`$ plus 6 kites, where $`e_1,e_2`$ are the edges which (with respect to the natural orientation) are leaving $`v`$, and $`e_3,e_4`$ are those which are reaching it. Again, there could be repetitions in the $`e_i`$โs. The kites all have coefficient $`+1`$, and again some of them may actually be collapsed to $`x_0`$. As above, we have only 6 kites because the other 6 are white.
###### Remark 3.9
To define the cellularization $`\widehat{๐ฏ}(P)`$ associated to a spine we have decided to subdivide all the triangles on $`\widehat{M}`$ into 3 kites, but when doing actual computations this is not necessary and impractical. The only triangles which we really need to subdivide are those intersected by $`C`$, because we need the cellularization to be suited to the pattern. Let us consider the 4 triangles corresponding to the ends of a certain tetrahedron. If in each of them we count the number of black kites and the number of white kites, we get respectively $`(3,0)`$, $`(2,1)`$, $`(1,2)`$, $`(0,3)`$. So, the first and last triangles do not have to be subdivided, and the other two can be subdivided using a segment only. Summing up, for each vertex of $`P`$ we only need to add two segments on the boundary. Before projecting $`M(P)`$ to $`\widehat{M}(P)`$ one sees that the number of cells, with respect to $`\overline{๐ฏ}(P)`$, is increased in all dimensions 0, 1 and 2 by twice the number of vertices of $`P`$. When projecting to $`\widehat{M}(P)`$ the cells lying in $`S_{\mathrm{triv}}^2`$ get collapsed to points.
## 4 Pseudo-Legendrian knots
We fix in this section a compact oriented manifold $`M`$ and a boundary pattern $`๐ซ`$ on $`M`$. The boundary of $`M`$ may be empty or not. Recall that if $`v`$ is a vector field on $`M`$ and $`K`$ is a knot in $`\mathrm{Int}(M)`$, we have defined $`K`$ to be pseudo-Legendrian in $`(M,v)`$ if $`v`$ is transversal to $`K`$. We will also call $`(v,K)`$ a pseudo-Legendrian pair. Having fixed $`๐ซ`$, we will only consider fields $`v`$ compatible with $`๐ซ`$. Some of the results we will establish hold also for links, but we will stick to knots for the sake of simplicity. First, we need to spell out the equivalence relation which we consider.
Let $`v_0,v_1`$ be compatible with $`๐ซ`$ and let $`K_0,K_1`$ be pseudo-Legendrian in $`(M,v_0)`$ and $`(M,v_1)`$ respectively. We define $`(v_0,K_0)`$ to be weakly equivalent to $`(v_1,K_1)`$ if there exist a homotopy $`(v_t)_{t[0,1]}`$ through fields compatible with $`๐ซ`$ and an isotopy $`(K_t)_{t[0,1]}`$ such that $`K_t`$ is transversal to $`v_t`$ for all $`t`$. If $`v_0=v_1`$ then $`K_0`$ and $`K_1`$ are called strongly equivalent if the homotopy $`(v_t)`$ can be chosen to be constant.
###### Remark 4.1
1. Of course strong equivalence implies weak equivalence. Weak equivalence is the natural relation to consider on pseudo-Legendrian pairs $`(v,K)`$, while strong equivalence is natural for pseudo-Legendrian knots in a fixed $`(M,v)`$. See for a further discussion on this point.
2. The term โpseudo-Legendrian isotopyโ, used in the introduction and in , corresponds to โweak equivalence.โ For the sake of brevity, and to emphasize the difference with strong equivalence, we will only use the term โweak equivalenceโ in the rest of this paper.
Before proceeding note that if $`K`$ is pseudo-Legendrian in $`(M,v)`$ then $`v`$ turns $`K`$ into a framed knot, which we will denote by $`K^{(v)}`$. The framed-isotopy class of $`K^{(v)}`$ is of course invariant under weak equivalence.
#### Euler structures on knot exteriors
For a knot $`K`$ in $`M`$ we consider a (closed) tubular neighbourhood $`U(K)`$ of $`K`$ in $`M`$ and we define $`E(K)`$ as the closure of the complement of $`U(K)`$. If $`F`$ is a framing on $`K`$ we extend the boundary pattern $`๐ซ`$ previously fixed on $`M`$ to a boundary pattern $`๐ซ(K^F)`$ on $`E(K)`$, by splitting $`U(K)`$ into a white and a black longitudinal annuli, the longitude being the one defined by the framing $`F`$. As a direct application of Proposition 1.1 one sees that $`\mathrm{Eul}(E(K),๐ซ(K^F))`$ is non-empty (assuming $`\mathrm{Eul}(M,๐ซ)`$ to be non-empty).
A convenient way to think of $`๐ซ(K^F)`$ is as follows. The framing $`F`$ determines a transversal vector field along $`K`$. If we extend this field near $`K`$ and choose $`U(K)`$ small enough then the pattern we see on $`U(K)`$ is exactly as required. With this picture in mind, it is clear that if $`K`$ is pseudo-Legendrian in $`(M,v)`$, where $`v`$ is compatible with $`๐ซ`$, then the restriction of $`v`$ to $`E(K)`$ defines an element
$$\xi (v,K)\mathrm{Eul}(E(K),๐ซ(K^{(v)})$$
(this structure was already considered in the partial proof of Theorem 3.2). So the theory of torsion applies. We will discuss in this section torsion both as an absolute and as a relative invariant, splitting the section into two subsections.
### 4.1 Absolute torsion and the Alexander invariant
We will show in this section that torsion as an absolute invariant of pseudo-Legendrian knots lifts the classical Alexander invariant, but the relation between the two objects is more complicated than in Turaevโs situation ( and ), because here two different algebraic complexes are involved at the same time.
#### Turaevโs lifting of Milnor torsion
Let us first recall again in what sense Turaevโs torsion lifts the classical one. Let $`M`$ be a manifold which is closed or bounded by tori, and take a representation $`\phi :H_1(M;\text{})\mathrm{\Lambda }_{}`$, where $`\mathrm{\Lambda }`$ is as usual. The classical theory allows to define an invariant
$$\tau ^\phi (M)K_1(\mathrm{\Lambda })/(\pm \phi (H_1(M;\text{}))),$$
usually stipulated to be $`1`$ if the $`\phi `$-twisted homology of $`M`$ does not vanish, i.e., using the above notation, if the complex $`C_{}^\phi (M)=\mathrm{\Lambda }_\phi C_{}^{\mathrm{cell}}(\stackrel{~}{M};\text{})`$ is not acyclic, where $`q:\stackrel{~}{M}M`$ is the maximal Abelian cover. When $`\xi `$ is an Euler structure on $`M`$ with monochromatic boundary components, Turaev shows that his torsion $`\tau ^\phi (M,\xi )\overline{K}_1(\mathrm{\Lambda })`$ is a lifting of $`\tau ^\phi (M)`$ with respect to the obvious projection of $`K_1(\mathrm{\Lambda })/\pm \phi (H_1(M;\text{})`$ onto $`\overline{K}_1(\mathrm{\Lambda })`$.
In the special case where $`\mathrm{\Lambda }`$ is the field of fractions obtained from the group ring of $`H_1(M;\text{})`$ modulo torsion, and $`\phi :H_1(M;\text{})\mathrm{\Lambda }`$ is the natural projection, the invariant $`\tau ^\phi (M)`$ is called Milnor torsion, and its sign-refinement provided by Turaev in has been shown to be equivalent to the classical Alexander invariant. So Turaevโs torsion for Euler structures contains a lifting of the Alexander invariant. We will discuss in the rest of this subsection the extent to what the same holds when the Euler structure arises from a pseudo-Legendrian knot. What we will say applies to any allowed representation $`\phi :H_1(M;\text{})\mathrm{\Lambda }`$, but we keep in mind that the relation with the Alexander invariant emerges for a special choice of $`\phi `$ and $`\mathrm{\Lambda }`$. Since we will also drop the condition that the involved complexes be acyclic, we note that torsion is only defined when the resulting homology is free. This is not true in general, but it is for instance when $`\mathrm{\Lambda }`$ is a field.
#### Torsion of a knot complement
Let us restrict to the case of a closed manifold $`M`$, and let us consider a pseudo-Legendrian pair $`(v,K)`$ in $`M`$ and a representation $`\phi :H_1(E(K);\text{})\mathrm{\Lambda }`$ as usual. We would like to interpret the torsion of the Euler structure $`\xi (v,K)`$ on $`E(K)`$ with respect to $`\phi `$ as a lifting of $`\tau ^\phi (E(K))`$, but a difficulty immediately emerges, because the algebraic complexes used to compute these torsions do not coincide.
To be more specific, let us first spell out how the torsion of $`\xi (v,K)`$ is defined. Let $`๐ซ(K^{(v)})=(B,W,\mathrm{},C)`$ be the boundary pattern defined on $`E(K)`$. Then we define $`\tau ^\phi (M,v,K,\text{h})`$ as $`\tau ^\phi (M,๐ซ(K^{(v)}),\mathrm{\Theta }(\xi (v,K)),\text{h})`$. More specifically, $`\tau ^\phi (M,v,K,\text{h})`$ is the torsion of the complex $`C_{}^\phi (E(K),\overline{W})`$, where $`W`$ is the (open) white annulus on $`E(K)`$, as above the maximal Abelian cover of $`E(K)`$ is used to define the complex, the preferred cell lifting is obtained using an Euler chain for the convexified structure $`\mathrm{\Theta }(\xi (v,K))`$, and h is a basis of the twisted homology of $`E(K)`$ relative to $`\overline{W}`$.
Now, $`\tau ^\phi (E(K))`$ is the torsion of $`C_{}^\phi (E(K))`$, and this complex can be radically different from the previous one. For instance, when $`M`$ is a homology sphere, the absolute complex is always acyclic, while the complex relative to $`\overline{W}`$, which depends only on the framed knot $`K^{(v)}`$, very often is not. We will see how to overcome this difficulty using the fundamental multiplicativity properties of torsion.
#### How to turn a torus into black
We will describe in this paragraph two explicit methods for modifying $`\xi (v,K)`$ to an Euler structure $`\beta (v,K)`$ such that $`E(K)`$ becomes monochromatic black. These methods are respectively a geometric and an algebraic one. The fact that they actually lead to the same result is true but not very important, so we will omit the proof. Both methods involve the choice of an orientation of $`K`$. The first method is explained in a cross-section in Fig. 9
The cross-section is transversal to $`K`$, and the apparent singularity of the modified field is removed by summing a field parallel to $`K`$ and supported near the singularity (cf. Fig. 2, where a similar method was used).
To describe the algebraic construction of $`\beta `$, recall that if $`z`$ is a 1-chain representing $`\xi (v,K)`$ then $`z`$ contains, with the appropriate sign, the centres of all cells in $`E(K)W`$. Knowing the subdivision rule for Euler chains (Proposition 1.4) we can also assume that the cellularization on $`W`$ has a particularly simple shape. We assume it consists of rectangles as in Fig. 10 (left), where we also show a 1-chain $`z_W`$ having the property that $`z_W`$ contains the centres of all cells in $`W`$. We can now define $`\beta (v,K)`$ as the Euler structure carried by $`z+z_W`$. The boundary of $`E(K)`$ is completely black with respect to this structure, because $`(z+z_W)`$ contains the centres of all cells of $`E(K)`$.
One easily sees from both our descriptions of $`\beta `$ that it is canonically defined and $`H_1`$-equivariant. Since we will need these properties, we spell out their meaning, starting from an oriented framed knot $`K^F`$ rather than a pseudo-Legendrian knot. Let $`(W,B,\mathrm{},C)`$ be the concave boundary pattern determined by $`F`$ on $`E(K)`$: then $`\beta :\mathrm{Eul}(E(K),(W,B,\mathrm{},C))\mathrm{Eul}(E(K);(\mathrm{},E(K),\mathrm{},\mathrm{})`$ is well-defined (depending on $`K^F`$ only) and $`H_1(E(K);\text{})`$-equivariant.
###### Remark 4.2
If $`K`$ denotes the same knot $`K`$ with reversed orientation then
$$\alpha (\beta (v,K),\beta (v,K))=[\lambda ]H_1(E(K);\text{})$$
where $`\lambda `$ is the longitude on $`E(K)`$ determined by the framing $`K^{(v)}`$.
A geometric interpretation of the second description of $`\beta `$ is possible and used below. We have mentioned that a theory of Euler structures exists in all dimensions. While the case $`n4`$ requires some technicalities , the reader can easily work out the case $`n=2`$ using the case $`n=3`$ treated in the present paper. And one easily sees that $`z_W`$ is just an Euler chain of the inward-pointing Reeb field $`r_W`$ on $`W`$ shown in Fig. 10 (right). Moreover $`r_W`$ can be canonically turned into an outward-pointing field $`\mathrm{\Theta }(r_W)`$, which of course is the outward-pointing Reeb field (but the core spins in the opposite direction). So a torsion $`\tau ^\psi (W,\mathrm{\Theta }(r_W))`$ can be computed (possibly with a basis of the twisted homology added to the data).
#### Knot torsion as a lifting of Milnorโs torsion
Let as above $`(v,K)`$ be pseudo-Legendrian and let $`\phi :H_1(E(K);\text{})\mathrm{\Lambda }`$ be a representation. If $`i:WE(K)`$ is the inclusion, we set $`\phi _W=\phi i_{}`$. Considering the twisted homology of the pair $`(E(K),W)`$ we get an exact sequence
$$=\left(\mathrm{}H_i^{\phi _W}(W)H_i^\phi (E(K))H^\phi (E(K),W)H_{i1}^{\phi _W}(W)\mathrm{}\right).$$
We choose bases h, $`\text{h}^{}`$, and $`\text{h}^{\prime \prime }`$ respectively for $`H_{}^\phi (E(K);\overline{W})`$, $`H_{}^\phi (E(K))`$ and $`H_{}^{\phi _W}(W)`$, so we can compute $`\tau ^\phi (M,v,K,\text{h})`$, $`\tau ^\phi (E(K),\beta (v,K),\text{h}^{})`$ and $`\tau ^{\phi _W}(W,\mathrm{\Theta }(r_W),\text{h}^{\prime \prime })`$. Moreover we can compute $`\tau (,\text{h},\text{h}^{},\text{h}^{\prime \prime })`$. The following result is a refinement of Theorem 3.2 in , and a proof can be given imitating the argument given in (where a special case of the result is established).
###### Proposition 4.3
The following equality holds:
$$\tau ^\phi (E(K),\beta (v,K),\text{h}^{})=\tau ^\phi (M,v,K,\text{h})\tau ^{\phi _W}(W,\mathrm{\Theta }(r_W),\text{h}^{\prime \prime })\tau (,\text{h},\text{h}^{},\text{h}^{\prime \prime }).$$
(4)
The following comments on the previous proposition eventually explain in what sense our torsion can be viewed as a lifting of the classical torsion (in particular, Milnor torsion and the Alexander invariant).
###### Remark 4.4
In equation (4) the term $`\tau ^\phi (E(K),\beta (v,K),\text{h}^{})`$ is one of Turaevโs torsion, so it is indeed a lifting of the classical torsion. The term $`\tau ^\phi (M,v,K,\text{h})`$ is the torsion for pseudo-Legendrian knots introduced in this paper, while $`\tau ^{\phi _W}(W,\mathrm{\Theta }(r_W),\text{h}^{\prime \prime })`$ and $`\tau (,\text{h},\text{h}^{},\text{h}^{\prime \prime })`$ can be viewed as normalizing terms. One can for instance choose homology bases so that $`\tau (,\text{h},\text{h}^{},\text{h}^{\prime \prime })=1`$, and note that $`\tau ^{\phi _W}(W,\mathrm{\Theta }(r_W),\text{h}^{\prime \prime })`$ depends only on the framed knot $`K^{(v)}`$, not on the Euler structure.
###### Remark 4.5
The factor $`\tau ^{\phi _W}(W,\mathrm{\Theta }(r_W),\text{h}^{\prime \prime })`$ may be understood quite easily. Denoting by $`1`$ the generator of $`H_1(W;\text{})`$, the result depends on $`\phi _W(1)`$. If $`\phi _W(1)=1`$ then the $`\phi _W`$-twisted homology of $`W`$ is not twisted at all, so it is non-zero and free, and we can choose $`\text{h}^{\prime \prime }`$ so that $`\tau ^{\phi _W}(W,\mathrm{\Theta }(r_W),\text{h}^{\prime \prime })=1`$. On the contrary, if $`\phi _W(1)1`$ is invertible, then the $`\phi _W`$-twisted homology is zero, and $`\tau ^{\phi _W}(W,\mathrm{\Theta }(r_W))`$ is computed to be $`\phi _W(1)1`$. In the intermediate cases where $`\phi _W(1)1`$ is neither zero nor a unit, which can only occur when $`\mathrm{\Lambda }`$ is not a field, $`\tau ^{\phi _W}(W,\mathrm{\Theta }(r_W))`$ is likely not to be defined.
###### Remark 4.6
We can further specialize the understanding of $`\tau ^{\phi _W}(W,\mathrm{\Theta }(r_W),\text{h}^{\prime \prime })`$ when $`M`$ is a homology sphere and $`\phi :H_1(E(K);\text{})\mathrm{\Lambda }`$ is the representation which gives the Milnor torsion. In this case $`\phi _W(1)=\phi (\mu )^n`$, where $`\mu `$ is the meridian of $`K`$ and $`n`$ is the framing. So $`\tau ^{\phi _W}(W,\mathrm{\Theta }(r_W),\text{h}^{\prime \prime })`$ can be normalized to be $`1`$ for $`\phi (\mu )^n=1`$, and it equals $`\phi (\mu )^n1`$ when $`\phi (\mu )^n1`$
### 4.2 Torsion as a relative invariant
We study in this section how torsion can be employed to distinguish pseudo-Legendrian knots from each other. We first show that as a relative invariant torsion is only well-defined as a multi-valued function, the ambiguity being given by the action of a group. Then we concentrate on the knots (called โgoodโ below) for which this action is trivial, and we interpret the relative information carried by torsion as a difference of winding numbers (or Maslov indices).
#### Group action on Euler structures
Consider a knot $`K`$ and a self-diffeomorphism $`f`$ of $`E(K)`$ which is the identity near $`E(K)`$. Then $`f`$ extends to a self-diffeomorphism $`\widehat{f}`$ of $`M`$, where $`\widehat{f}|_{U(K)}=\mathrm{id}_{U(K)}`$. We define $`G(K)`$ as the group of all such $`f`$โs with the property that $`\widehat{f}`$ is isotopic to the identity on $`M`$. Elements of $`G(K)`$ are regarded up to isotopy relative to $`E(K)`$. If $`F`$ is a framing on $`K`$ then the pull-forward of vector fields induces an action of $`G(K)`$ on $`\mathrm{Eul}(E(K),๐ซ(K^{(v)})`$. We will now see that an obstruction to weak equivalence can be expressed in terms this group action.
Let $`(v_0,K_0)`$ and $`(v_1,K_1)`$ be pseudo-Legendrian pairs in $`M`$, and assume that $`K_0^{(v_0)}`$ is framed-isotopic to $`K_1^{(v_1)}`$ under a diffeomorphism $`f`$ relative to $`M`$. Using the restriction of $`f`$ and the pull-back of vector fields we get a bijection
$$f^{}:\mathrm{Eul}(E(K_1),๐ซ(K_1^{(v_1)}))\mathrm{Eul}(E(K_0),๐ซ(K_0^{(v_0)})).$$
###### Proposition 4.7
Under the current assumptions, if $`(v_0,K_0)`$ and $`(v_1,K_1)`$ are weakly equivalent to each other then $`f^{}(\xi (v_1,K_1))`$ belongs to the $`G(K_0)`$-orbit of $`\xi (v_0,K_0)`$ in $`\mathrm{Eul}(E(K_0),๐ซ(K_0^{(v_0)})`$.
Proof of4.7. By assumption $`K_0,K_1`$ and $`v_0,v_1`$ embed in continuous families $`(K_t)_{t[0,1]}`$ and $`(v_t)_{t[0,1]}`$, where $`v_t`$ is transversal to $`K_t`$ for all $`t`$. Now $`(K_t^{(v_t)})_{t[0,1]}`$ is a framed-isotopy, so there exists a continuous family $`(g_t)_{t[0,1]}`$ of diffeomorphisms of $`M`$ fixed on $`M`$ and such that $`g_0=\mathrm{id}_M`$ and $`g_t(K_0^{(v_0)})=K_t^{(v_t)}`$. So we get a map
$$[0,1]t\alpha (\xi (v_0,K_0),g_t^{}(\xi (v_t,K_t)))H_1(E(K_0);\text{}).$$
Since $`H_1(E(K_0);\text{})`$ is discrete and the map is continuous, we deduce that the map is identically 0. So $`g_1^{}(\xi (v_1,K_1))=\xi (v_0,K_0)`$. Now
$$f^{}(\xi (v_1,K_1))=(f^{}(g_1)_{}g_1^{})(\xi (v_1,K_1))=(f^1g_1)_{}(\xi (v_0,K_0))$$
and the conclusion follows because $`f^1g_1`$ defines an element of $`G(K_0)`$. 4.7
The group $`G(K)`$ is in general rather difficult to understand (see ), so we introduce a special terminology for the case where its action can be neglected. We will say that a framed knot $`K^F`$ is good if $`G(K)`$ acts trivially on $`\mathrm{Eul}(E(K),๐ซ(K^F))`$. If $`K^F`$ is good for all framings $`F`$, we will say that $`K`$ itself is good. The following are easy examples of good knots:
* $`M`$ is $`S^3`$ and $`K`$ is the trivial knot;
* $`M`$ is a lens space $`L(p,q)`$ and $`K`$ is the core of one of the handlebodies of a genus-one Heegaard splitting of $`M`$.
The reason is that in both cases $`E(K)`$ is a solid torus, and we know that an automorphism of the solid torus which is the identity on the boundary is isotopic to the identity relatively to the boundary, so $`G(K)`$ is trivial. The next three results show that on one hand $`G(K)`$ is very seldom trivial, but on the other hand many knots are good. We will give proofs in the sequel, after introducing some extra notation. In the statements, by โ$`E(K)`$ is hyperbolicโ we mean โ$`\mathrm{Int}(E(K))`$ is complete, finite-volume hyperbolic.โ
###### Proposition 4.8
If $`M`$ is closed and $`E(K)`$ is hyperbolic then $`G(K)`$ is non-trivial.
###### Theorem 4.9
If $`M`$ is closed, $`E(K)`$ is hyperbolic and either $`\mathrm{Out}(\pi _1(E(K)))`$ is trivial or $`H_1(E(K);\text{})`$ is torsion-free then $`K`$ is good.
###### Theorem 4.10
If $`M`$ is a homology sphere then every knot in $`M`$ is good.
The next result, which follows directly from Proposition 4.7, the definition of goodness, and Proposition 2.6, shows that for good knots torsion can be used as an obstruction to weak (and hence strong) equivalence.
###### Proposition 4.11
Let $`(v_0,K_0)`$ and $`(v_1,K_1)`$ be pseudo-Legendrian pairs in $`M`$, and assume that $`K_0^{(v_0)}`$ is framed-isotopic to $`K_1^{(v_1)}`$ under a diffeomorphism $`f`$ relative to $`M`$. Suppose that $`K_0^{(v_0)}`$ is good, and that for some representation $`\phi :\pi _1(E(K_0))\mathrm{\Lambda }`$ and some $`\mathrm{\Lambda }`$-basis h of $`H_{}^\phi (E(K_0),\overline{W(๐ซ(K_0^{(v_0)}))})`$ we have
$$\tau ^\phi (E(K_0),๐ซ(K_0^{(v_0)}),\xi (v_0,K_0),\text{h})\tau ^{\phi f_{}^1}(E(K_1),๐ซ(K_1^{(v_1)}),\xi (v_1,K_1),f_{}(\text{h})).$$
(5)
Then $`(v_0,K_0)`$ and $`(v_1,K_1)`$ are not weakly equivalent.
###### Remark 4.12
1. The right-hand side of equation (5) actually equals
$`\tau ^\phi (E(K_0),๐ซ(K_0^{(v_0)}),f^{}(\xi (v_1,K_1)),\text{h})`$
$`=`$ $`\overline{\phi }(\alpha (v_0,f^{}(v_1))\tau ^\phi (E(K_0),๐ซ(K_0^{(v_0)}),\xi (v_0,K_0),\text{h}).`$
This shows that the most torsion can capture as a relative invariant of $`(v_0,K_0)`$ and $`(v_1,K_1)`$ is $`\alpha (v_0,f^{}(v_1))`$. We will show below that in some cases torsion indeed allows to determine $`\alpha (v_0,f^{}(v_1)`$ completely.
2. By definition of goodness the homology class $`\alpha (v_0,f^{}(v_1))`$ just considered is actually independent of $`f`$. We will denote it by $`\alpha ((K_0,v_0),(K_1,v_1))`$.
3. For non-good knots the relative invariant is an orbit of the action of $`G(K_0)`$. So an obstruction in terms of torsion could be given also for non-good knots, but the statement would become awkward, and we have refrained from giving it.
4. If equation (5) holds for some basis h then it holds for any basis.
To conclude this paragraph we note that using the technology of Turaev , one can actually see that the action on Euler structures of an automorphism is invariant under homotopy (not only isotopy) relative to the boundary. We will not use this fact.
#### Good knots
We introduce now some notation needed for the proofs of Proposition 4.8 and Theorem 4.9 (for Theorem 4.10 we will use a different approach, see below). Recall that $`(M,๐ซ)`$ is fixed for the whole section. We temporarily fix also a framed knot $`K^F`$ in $`M`$, a regular neighbourhood $`U`$ of $`K`$, and we denote by $`T`$ the boundary torus of $`U`$. On $`T`$ we consider 1-periodic coordinates $`(x,y)`$ such that $`x(x,0)`$ is a meridian of $`U`$ and $`y(0,y)`$ is a longitude compatible with $`F`$. We denote a collar of $`T`$ in $`E(K)`$ by $`V`$ and parameterize $`V`$ as $`T\times [0,1]`$, where $`T=T\times \{0\}`$. We consider on $`[0,1]`$ a coordinate $`s`$. For $`p,q\text{}`$ we define automorphisms $`๐_{(p,q)}`$ of $`E(K)`$ as follows. Each $`๐_{(p,q)}`$ is supported in $`V`$, and on $`V`$, using the coordinates just described, it is given by
$$๐_{(p,q)}(x,y,s)=(x+ps,y+qs,s).$$
We will call such a map a Dehn twist. It is easy to verify that the extension of $`๐_{(p,q)}`$ to $`M`$ is isotopic to the identity of $`M`$. Note that $`๐_{(p,q)}`$ is actually not smooth on $`T\times \{1\}`$, but we can consider some smoothing and identify $`๐_{(p,q)}`$ to an element of $`G(K)`$, because the equivalence class is independent of the smoothing.
Proof of4.8. We show that $`๐_{(p,q)}`$ is non-trivial in $`G(K)`$ for all $`(p,q)(0,0)`$. Fix the basepoint $`a_0=(0,0)T`$ for the fundamental groups of $`T`$ and $`E(K)`$. Then $`๐_{(p,q)}`$ acts on $`\pi _1(E(K),a_0)`$ as the conjugation by $`i_{}(p,q)`$, where $`i:TE(K)`$ is the inclusion and $`(p,q)\text{}\times \text{}=\pi _1(T,a_0)`$. If $`๐_{(p,q)}`$ is trivial in $`G(K)`$, i.e. it is isotopic to the identity relatively to $`E(K)`$, in particular it acts trivially on $`\pi _1(E(K),a_0)`$. This implies that $`i_{}(p,q)`$ is in the centre of $`\pi _1(E(K),a_0)`$. Now it follows from hyperbolicity that this centre is trivial and $`i_{}`$ is injective, whence the conclusion. 4.8
The proof of Theorem 4.9 will rely on properties of hyperbolic manifolds and on the following fact, which we consider to be quite remarkable (note that the 2-dimensional analogue, which may be stated quite easily, is false). Remark that the result applies in particular to Dehn twists.
###### Proposition 4.13
If $`[f]G(K)`$ and $`f`$ is supported in the collar $`V`$ of $`U`$ then $`[f]`$ acts trivially on $`\mathrm{Eul}(E(K),๐ซ(K^F))`$.
Proof of4.13. Consider a vector field $`v`$ on $`E(K)`$ compatible with $`๐ซ(K^F)`$. Since $`v`$ and $`f_{}(v)`$ differ only on $`V`$, their difference belongs to the image of $`H_1(V;\text{})`$ in $`H_1(E(K);\text{})`$. So we may as well assume that $`E(K)=V`$, i.e. $`M`$ is the solid torus $`UV`$.
By contradiction, let $`\xi \mathrm{Eul}(V,๐ซ(K^F))`$ be such that $`\alpha (\xi ,(๐_{(p,q)})_{}(\xi ))`$ is non-zero in $`H_1(V;\text{})`$, so it is given by $`k[\gamma ]`$ for some $`k\text{}\{0\}`$ and some simple closed curve $`\gamma `$ on $`T\times \{1\}V`$. Let us now take another simple closed curve $`\delta `$ on $`T\times \{1\}`$ which intersects $`\gamma `$ transversely at one point. Let us define $`N`$ as the manifold obtained by attaching the solid torus to $`V`$ along $`T\times \{1\}`$, in such a way that the meridian of the solid torus gets identified with $`\delta `$. Note that $`N`$ is again a solid torus and that the homology class of $`\gamma `$ in $`H_1(N;\text{})\text{}`$ is a generator. Now we can apply Proposition 1.1 to extend $`\xi `$ to an Euler structure $`\xi _N`$ on $`N`$. Moreover we can extend $`f`$ to an automorphism $`g`$ of $`N`$ which is the identity on $`N=T\times \{0\}`$. Now by construction $`\alpha (\xi _N,g_{}(\xi _N))`$ equals $`k[\gamma ]`$ in $`H_1(N;\text{})\text{}`$, so it is non-zero. But $`g`$ is isotopic to the identity of $`N`$ relatively to the boundary of $`N`$, so we have a contradiction. 4.13
For the proof of Theorem 4.9 we will also need the following easy fact.
###### Lemma 4.14
Let $`f`$ be an automorphism of $`M`$ relative to $`M`$, and consider the induced automorphisms of $`H_1(M;\text{})`$ and $`\mathrm{Eul}(M,๐ซ)`$, both denoted by $`f_{}`$. Then:
$$\alpha (f_{}(\xi _0),f_{}(\xi _1))=f_{}(\alpha (\xi _0,\xi _1)),\xi _0,\xi _1\mathrm{Eul}(M,๐ซ).$$
Proof of4.14. Take representatives of $`\xi _0`$ and $`\xi _1`$ such that $`\alpha (\xi _0,\xi _1)`$ can be viewed as the anti-parallelism locus. The formula is then obvious. 4.14
Proof of4.9. Consider $`[f]G(K)`$. It follows from the work of Johansson (see ) that, under the assumption that $`E(K)`$ is hyperbolic, the group generated by Dehn twists has finite index in the mapping class group of $`E(K)`$ relative to the boundary. More precisely, the quotient group can be identified to a subgroup of $`\mathrm{Out}(\pi _1(E(K))`$, which is finite as a consequence of Mostowโs rigidity. If $`\mathrm{Out}(\pi _1(E(K))`$ is trivial then $`[f]`$ is equivalent to a Dehn twist, so $`f`$ acts trivially on $`\mathrm{Eul}(E(K),๐ซ(K^F))`$ by Proposition 4.13.
We are left to deal with the case where $`H_1(E(K);\text{})`$ is torsion-free. By Johanssonโs result, there exists an integer $`n`$ such that $`f^n`$ acts trivially on $`\mathrm{Eul}(E(K),๐ซ(K^F))`$. Consider now $`\xi \mathrm{Eul}(E(K),๐ซ(K^F))`$, and set $`\alpha =\alpha (\xi ,f_{}(\xi ))`$. We must show that $`\alpha =0`$. We denote by $`\widehat{\alpha }`$ the image of $`\alpha `$ in $`H_1(M;\text{})`$, and by $`\widehat{f}`$ the extension of $`f`$ to $`M`$. Since $`\widehat{f}`$ is isotopic to the identity, we have $`\widehat{f}_{}(\widehat{\alpha })=\widehat{\alpha }`$. If we take an oriented 1-manifold $`a`$ representing $`\alpha `$ and disjoint from $`U(K)`$, this means that there exists an oriented surface $`\mathrm{\Sigma }`$ in $`M`$ such that $`\mathrm{\Sigma }=a(f(a))`$. Up to isotopy we can assume that $`\mathrm{\Sigma }`$ intersects $`U(K)`$ transversely in a union of circles. This shows that $`f_{}(\alpha )=\alpha +k\mu `$, where $`\mu `$ is the meridian of $`T`$. Note that $`f_{}(\mu )=\mu `$, so for all integers $`m`$ we have $`f_{}^m(\alpha )=\alpha +mk\mu `$. Now, using Lemma 4.14, we have:
$`0`$ $`=`$ $`\alpha (\xi ,f_{}^n(\xi )={\displaystyle \underset{m=0}{\overset{n1}{}}}\alpha (f_{}^m(\xi ),f_{}^{m+1}(\xi ))`$
$`=`$ $`{\displaystyle \underset{m=0}{\overset{n1}{}}}f_{}^m(\alpha (\xi ,f_{}(\xi )))={\displaystyle \underset{m=0}{\overset{n1}{}}}f_{}^m(\alpha )={\displaystyle \underset{m=0}{\overset{n1}{}}}(\alpha +mk\mu )`$
$`=`$ $`n\alpha +{\displaystyle \frac{n(n1)}{2}}k\mu .`$
This shows that $`2\alpha +(n1)k\mu `$ is a torsion element of $`H_1(E(K);\text{})`$, so it is null by assumption. So $`(1n)k\mu =2\alpha `$. If we apply $`f_{}`$ to both sides of this equality we get $`(1n)kf_{}(\mu )=2f_{}(\alpha )`$. Using the equality again and the relations $`f_{}(\mu )=\mu `$ and $`f_{}(\alpha )=\alpha +k\mu `$ we get
$$(1n)k\mu =2\alpha +2k\mu =(1n)k\mu +2k\mu .$$
Therefore $`k\mu `$ is a torsion element, and hence null. But $`2\alpha =(1n)k\mu `$, so also $`\alpha `$ is null. 4.9
#### Rotation number, and more good knots
We will show in this section that in a homology sphere the rotation number of a pseudo-Legendrian knot can be (defined and) expressed in terms of Euler structures on the exterior. This will lead us to a simple interpretation of torsion as a relative invariant of knots, and it will allow us to show that in a homology sphere all knots are good (Theorem 4.10).
To begin, we note that the definition of the rotation number, classically defined in the contact case, actually extends to the situation we are considering. Since we will need this definition, we recall it. Let $`M`$ be a homology sphere, let $`v`$ be a field on $`M`$ and let $`K`$ be an oriented pseudo-Legendrian knot in $`(M,v)`$. Take a plane field $`\eta `$ transversal to $`v`$ and tangent to $`K`$, and a Seifert surface $`S`$ for $`K`$. Up to isotopy of $`S`$ we can assume that $`\eta `$ is tangent to $`S`$ only at isolated points. Then $`\mathrm{rot}_v(K)`$ is the sum of a contribution for each of these tangency points $`p`$. Define $`\mathrm{o}(p)`$ to be $`+1`$ if $`\eta _p=T_pS`$ and $`1`$ if $`\eta _p=T_pS`$. If $`pS=K`$ then $`p`$ contributes just with $`\mathrm{o}(p)`$. If $`p\mathrm{Int}(S)`$ we can consider near $`p`$ a section of $`\eta TS`$ which vanishes at $`p`$ only, and denote by $`\mathrm{i}(p)`$ its index. Then $`p`$ contributes to $`\mathrm{rot}_v(K)`$ with $`\mathrm{o}(p)\mathrm{i}(p)`$.
It is quite easy to see that the resulting number is indeed independent from $`\eta `$ and $`S`$. Moreover $`\mathrm{rot}_v(K)`$ is unchanged under homotopies of $`v`$ relative to $`K`$, and local modifications away from $`K`$, so we can actually define $`\mathrm{rot}_\xi (K)`$ where $`\xi =\xi (v,K)\mathrm{Eul}(E(K),๐ซ(K^{(v)})`$.
###### Proposition 4.15
Let $`M`$ be a homology sphere, let $`v`$ be a field on $`M`$ and let $`K_0`$ and $`K_1`$ be oriented pseudo-Legendrian knots in $`(M,v)`$. Assume that there exists a framed-isotopy $`f`$ which maps $`K_1^{(v)}`$ to $`K_0^{(v)}`$. Identify $`H_1(E(K_0);\text{})`$ to using a meridian. Then:
$$\mathrm{rot}_v(K_1)=\mathrm{rot}_v(K_0)+2\alpha (f_{}(\xi (v,K_1)),\xi (v,K_0)).$$
Proof of4.15. Let $`K:=K_0`$, $`v_0:=v`$ and $`v_1:=f_{}(v)`$. Note that $`v_0`$ and $`v_1`$ coincide along $`K`$. Of course $`\mathrm{rot}_v(K_1)=\mathrm{rot}_{v_1}(K)`$. We are left to show that
$$\mathrm{rot}_{v_1}(K)=\mathrm{rot}_{v_0}(K)+2\alpha (\xi (v_1,K)),\xi (v_0,K)).$$
We can now homotope $`v_0`$ and $`v_1`$ away from $`K`$ until they differ only in the neighbourhood $`W(L)`$ of an oriented link $`L`$, and within this neighbourhood they differ exactly by a โPontrjagin moveโ, as defined for instance in . Namely, $`v_0`$ runs parallel to $`L`$ in $`W(L)`$, while $`v_1`$ runs opposite to $`L`$ on $`L`$ and has non-positive radial component on $`W(L)`$ (see below for a picture). Note that $`L`$ represents $`\alpha (\xi (v_1,K)),\xi (v_0,K))`$.
Let us choose now a Seifert surface $`S`$ for $`K`$ and a Riemannian metric on $`M`$, and define $`\eta _i=v_i^{}`$, for $`i=0,1`$. Since $`\eta _0|_K=\eta _1|_K`$, the contributions along $`K`$ to $`\mathrm{rot}_{v_0}(K)`$ and $`\mathrm{rot}_{v_1}(K)`$ are the same. Up to isotoping $`S`$ we may assume that $`L`$ is transversal but never orthogonal to $`S`$. At the points where $`\eta _0`$ is tangent to $`S`$ also $`\eta _1`$ is tangent to $`S`$, and the contributions are the same. So $`\mathrm{rot}_{v_1}(K)\mathrm{rot}_{v_0}(K)`$ is given by the sum of the contributions of the tangency points of $`\eta _1`$ to $`S`$ within $`W(L)`$. We will show that each point of $`LS`$ gives rise to exactly two tangency points, which both contribute with $`+1`$ or $`1`$ according to the sign of the intersection of $`L`$ and $`S`$ at the point. This will show that $`\mathrm{rot}_{v_1}(K)\mathrm{rot}_{v_0}(K)`$ is twice the algebraic intersection of $`L`$ and $`S`$. This algebraic intersection is exactly the value of $`[L]=\alpha (\xi (v_1,K)),\xi (v_0,K))`$ as a multiple of $`[m]`$, so the local analysis at $`LS`$ will imply the desired conclusion.
For the sake of simplicity we only examine a positive intersection point of $`L`$ and $`S`$. This is done in a cross-section in Fig. 11, which shows the local effect
of the move. The fields pictured both have a rotational symmetry, suggested in the figure. The two tangency points which arise are a positive focus (on the right) and a negative saddle (on the left), so the local contribution is indeed $`+2`$, and the proof is complete. 4.15
###### Remark 4.16
The definition of rotation number and Proposition 4.15 easily extend to the case of manifolds which are not homology spheres, by restricting to homologically trivial knots and choosing a relative homology class in the exterior.
We can now prove that in a homology sphere all knots are good.
Proof of4.10. Consider $`[f]G(K)`$, a framing $`F`$ on $`K`$ and $`\xi \mathrm{Eul}(E(K),๐ซ(K^F))`$. We must show that $`f_{}(\xi )=\xi `$. Let $`\xi =[v]`$ and denote by $`\widehat{v}`$ the obvious extension of $`v`$ to $`M`$. As above, let $`\widehat{f}`$ be the extension of $`f`$ to $`M`$. During the proof of Proposition 4.15 we have shown that
$$\mathrm{rot}_{\widehat{f}_{}(\widehat{v})}(K)\mathrm{rot}_{\widehat{v}}(K)=2\alpha (f_{}(v),v).$$
But $`\mathrm{rot}_{\widehat{f}_{}(\widehat{v})}(K)`$ is actually equal to $`\mathrm{rot}_{\widehat{v}}(K)`$, because $`\widehat{f}`$ is the identity near $`K`$. Therefore $`f_{}(v)`$ and $`v`$ differ by a torsion element of $`H_1(E(K);\text{})\text{}`$, so they are equal. By definition $`f_{}(\xi )=[f_{}(v)]`$ and $`\xi =[v]`$, and the proof is complete. 4.10
Theorems 4.9 and 4.10 provide a partial answer to the problem of determining which knots are good. The general problem does not appear to be straight-forward, and we leave it for further investigation. We will only show below an example of knot which is not good.
#### Curls and the winding number
We show in this paragraph the relation between the relative invariant $`\alpha ((v_0,K_0),(v_1,K_1))`$ of two pseudo-Legendrian knots (when this invariant is well-defined) and an analogue of the winding number (the invariant which allows to distinguish framed-isotopic planar link diagrams which are not equivalent under the second and third of Reidemeisterโs moves, see ). Moreover we will give an example of knot which is not good. The proof of the next result uses the example of Section 5, so it is deferred to Section 6.
###### Proposition 4.17
Consider a field $`v`$ on $`M`$ and a portion of $`M`$ on which $`v`$ can be identified to the vertical field in $`\text{}^3`$. Consider oriented knots $`K_0`$ and $`K_{\pm 1}`$ which are transversal to $`v`$ and differ only within the chosen portion of $`M`$, as shown in Fig. 12.
Choose the positive meridian $`m`$ of $`K_0`$, as also shown in the figure. Let $`f`$ be an isotopy which maps $`K_{\pm 1}^{(v)}`$ to $`K_0^{(v)}`$ and is supported in a tubular neighbourhood of $`K_0`$. Then:
$$\alpha (\xi (v,K_0),f_{}(\xi (v,K_{\pm 1})))=\pm [m]H_1(E(K_0);\text{}).$$
###### Proposition 4.18
Let $`(v,K_0)`$ be a pseudo-Legendrian pair in $`M`$, and denote by $`[m]H_1(E(K_0);\text{})`$ the homology class of the meridian of $`U(K_0)`$. Assume either that $`K_0^{(v)}`$ is good and $`[m]0`$ or that $`E(K_0)`$ is hyperbolic and $`[m]`$ has infinite order. Let $`K_1`$ be a knot obtained from $`K_0`$ as in Fig. 12. Then $`(v,K_0)`$ and $`(v,K_1)`$ are not weakly equivalent.
Proof of4.18. By contradiction, using Propositions 4.7 and 4.17, we would get elements $`\xi _0,\xi _1`$ of $`\mathrm{Eul}(E(K_0),๐ซ(K_0^{(v)})`$ such that $`\alpha (\xi _0,\xi _1)=[m]`$ and $`\xi _1=f_{}(\xi _0)`$ for some $`[f]G(K_0)`$. If $`K_0^{(v)}`$ is good and $`[m]0`$ this is a contradiction. Assume now that $`E(K_0)`$ is hyperbolic and $`[m]`$ has infinite order. Since $`f_{}([m])=[m]`$, using Lemma 4.14 we easily see that $`\alpha (\xi _0,f_{}^k(\xi _0))=k[m]`$ for all $`k`$. Proposition 4.13 and the result of Johansson already used in the proof of Theorem 4.9 now imply that $`f^k`$ acts trivially on $`\mathrm{Eul}(E(K_0),๐ซ(K_0^{(v)})`$ for some $`k`$, whence the contradiction. 4.18
As an application of Proposition 4.17, we can show that there exist knots which are not good. Consider $`S^2\times [0,1]`$ with vector field parallel to the $`[0,1]`$ factor. Let $`K_0`$ be the equator of $`S^2\times \{1/2\}`$, and let $`K_1`$ be obtained from $`K_0`$ by the modification described in Fig. 12. Using Proposition 4.17, if we choose a framed-isotopy $`g`$ of $`K_1^{(v)}`$ onto $`K_0^{(v)}`$ supported in $`U(K_0)`$, we have
$$\alpha (\xi (v,K_0),(g|_{E(K_1)})_{}(\xi (v,K_1)))=[m],$$
where $`[m]`$ is a generator of $`H_1(E(K_0);\text{})\text{}`$. On the other hand, $`K_1`$ is strongly equivalent to $`K_0`$ in $`(M,v)`$ (the winding number only exists on $`\text{}^2`$, not on $`S^2`$). So there exists an isotopy $`h`$ of $`K_1^{(v)}`$ onto $`K_0^{(v)}`$ through links transversal to $`v`$, and we have
$$\alpha (\xi (v,K_0),(h|_{E(K_1)})_{}(\xi (v,K_1)))=0.$$
This implies that $`(hg^1)|_{E(K_0)}`$ acts non-trivially on $`\xi (v,K_0)\mathrm{Eul}(E(K_0),๐ซ(K_0^{(v)}))`$.
To conclude our discussion on the relative invariant $`\alpha `$ between two pseudo-Legendrian knots, we state now a result proved in . The consequences we deduce easily follow from Proposition 4.17.
###### Proposition 4.19
Let $`(v_0,K_0)`$ and $`(v_1,K_1)`$ be pseudo-Legendrian in $`M`$, assume that $`v_0`$ and $`v_1`$ are homotopic fields, and that $`K_0^{(v_0)}`$ and $`K_1^{(v_1)}`$ are isotopic as framed knots. Then $`(v_0,K_0)`$ and $`(v_1,K_1)`$ become weakly equivalent up to a finite number of local moves $`K_0K_{\pm 1}`$ as in Fig. 12.
###### Corollary 4.20
Under the same assumptions, assume also that $`K_0^{(v_0)}`$ is good, so $`\alpha ((K_0,v_0),(K_1,v_1))`$ is defined. Then
$$\alpha ((K_0,v_0),(K_1,v_1))=\mathrm{w}((K_0,v_0),(K_1,v_1))[m_0]H_1(E(K_0);\text{})$$
where $`m_0`$ is the meridian of $`K_0`$ and $`\mathrm{w}((K_0,v_0),(K_1,v_1))`$ is the (non-well-defined) algebraic number of moves $`K_0K_{\pm 1}`$ needed to make $`K_0`$ and $`K_1`$ weakly equivalent.
Concerning the statement of the previous corollary, note that both $`\mathrm{w}((K_0,v_0),(K_1,v_1))`$ and $`[m_0]`$ depend on the choice of an orientation on $`K_0`$, but their product does not.
###### Corollary 4.21
Assume furthermore that $`[m_0]`$ has infinite order in $`H_1(E(K_0);\text{})`$. Then $`\mathrm{w}((K_0,v_0),(K_1,v_1))\text{}`$ is a well-defined integer relative invariant, which we call the relative winding number.
###### Remark 4.22
If $`M`$ is a homology sphere then the local moves of Fig. 12 which modify the winding number also change the rotation number, and Corollary 4.20 is consistent with Proposition 4.15.
The next proposition contains in particular Theorem 0.2 stated in the introduction.
###### Proposition 4.23
Under the assumptions of Proposition 4.19, assume that $`K_0^{(v_0)}`$ is good and that $`[m_0]`$ has infinite order in $`H_1(E(K_0);\text{})`$. The following facts are pairwise equivalent:
1. the relative winding number of $`(K_0,v_0)`$ and $`(K_1,v_1)`$ vanishes;
2. all relative torsion invariants of $`(K_0,v_0)`$ and $`(K_1,v_1)`$ are trivial;
3. $`(K_0,v_0)`$ and $`(K_1,v_1)`$ are weakly equivalent.
Proof of4.23. Equivalence of (1) and (3) comes from the previous discussion and from the fact that a positive double curl and a negative double curl cancel via weak equivalence. To show that (1) and (2) are equivalent we only need to consider torsion with respect to a representation $`\phi :H_1(E(K_0);\text{})\mathrm{\Lambda }`$ such that $`\phi ([m_0])`$ has infinite order. 4.23
###### Corollary 4.24
Under the assumptions of Proposition 4.19, assume that $`M`$ is a homology sphere. Then the facts (1), (2), and (3) of Proposition 4.23 are also equivalent to the following:
1. $`(K_0,v_0)`$ and $`(K_1,v_1)`$ have the same rotation number.
Proof of4.24. Equivalence of (1) and ((4)) comes from the previous discussion and Proposition 4.17. 4.24
Since in a homology sphere two pseudo-Legendrian knots which are homotopic through pseudo-Legendrian immersions certainly have the same Maslov index, the previous corollary seems to suggest that all torsion can capture in a homology sphere is the homotopy class through immersions. We believe that it would be interesting to check if also for a general manifold $`M`$, under the assumptions of Corollary 4.20, homotopy through pseudo-Legendrian immersions implies $`\mathrm{w}((K_0,v_0),(K_1,v_1))[m_0]=0`$. We conclude by informing the reader that in we have discussed the extent to which the category of pseudo-Legendrian knots can be represented by the category of genuine Legendrian knots in overtwisted contact structures.
## 5 An example
Figure 13 shows a neighbourhood of the singular set of the so-called
abalone, a branched standard spine of $`S^3`$, which we denote by $`A`$. Note that $`A`$ has one vertex, two edges and two regions. The figure on the left is easier to understand, but it does not represent the genuine embedding of $`A`$ in $`S^3`$, which is instead shown in the centre (hint: compute linking numbers). On the right we show (using the easier picture) a $`\mathrm{C}^1`$ knot $`K`$ on $`A`$. Using the genuine picture one sees that $`K`$ is actually trivial in $`S^3`$, and its framing is $`1`$. So the knot exterior $`E(K)`$ is actually a solid torus, with an induced Euler structure $`\xi `$, and the white annulus $`WE(K)`$ is a longitudinal one. Let us now take the representation $`\phi :\pi _1(E(K))\text{}[t^{\pm 1}]`$ which maps the generator to $`t`$. It is not hard to see that $`H_{}^\phi (E(K),\overline{W})=0`$, so we can compute $`\tau ^\phi (E(K),\xi )`$. We describe the method to be followed, skipping several details and all explicit formulae.
We can apply directly the method described in the (partial) proof of Theorem 3.2, to get a branched standard spine $`P`$ (in the sense of Theorem 3.2) of $`E(K)`$. This $`P`$ is easily recognized to have 5 vertices (denoted $`v_1,\mathrm{},v_5`$), 10 edges (denoted $`e_0,\mathrm{},e_9`$) and 6 regions (denoted $`r_1,\mathrm{}r_6`$). Figure 14 shows the truncated ideal triangulation dual to $`P`$.
In the figure the hat denotes duality as usual. We have written $`\widehat{e}_i`$ instead of $`\widehat{e}_i`$ when $`\widehat{e}_i`$ lies on $`\widehat{v}_j`$ but the natural orientation of $`\widehat{e}_i`$ is not induced by the orientation of $`\widehat{v}_j`$. The letters $`S`$ and $`T`$ refer to the boundary sphere and torus respectively ($`S`$ should actually be collapsed to one point $`x_0`$, but the picture is easier to understand before collapse).
Recall that the algebraic complex of which we must compute the torsion has one generator for each cell in the cellularization of $`E(K)`$ arising from $`P`$, excluding the white cells and the tangency circles on the boundary. From Fig. 14 we can see how many such cells there will be in each dimension, namely 3 in dimension 0 ($`x_0`$ and two vertices on $`T`$), 14 in dimension 1 (the $`\widehat{r}_i`$โs and 8 edges on $`T`$), 16 in dimension 2 (the $`\widehat{e}_i`$โs and the 6 black kites on $`T`$) and 5 in dimension 3 (the $`\widehat{v}_i`$โs). We can also easily describe the combinatorial Euler chain $`s^{}(P)`$ which will be used to find the preferred cell liftings: besides the orbits of the field there are only one star and one bi-arrow; the support of $`s^{}(P)`$ has 3 connected components (one spider with 19 legs and head at $`x_0`$, the star union the second half of $`\widehat{r}_2`$, and the bi-arrow union a segment contained in $`\widehat{e}_3`$).
To actually determine the preferred liftings we need an effective description of the lifting of the cellularization to the universal cover $`\stackrel{~}{E}(K)E(K)`$. Since $`\pi _1(E(K))=\text{}`$, each cell $`c`$ will have liftings $`\stackrel{~}{c}^{(n)}`$ for $`n\text{}`$, where $`\stackrel{~}{c}^{(n)}`$ is the $`n`$-th translate of $`\stackrel{~}{c}^{(0)}`$. The choice of $`\stackrel{~}{c}^{(0)}`$ itself is of course arbitrary, but to understand the cover we must express the $`\stackrel{~}{c}^{(0)}`$โs in terms of the other $`\stackrel{~}{d}^{(n)}`$โs. To do this we start with a lifting $`\stackrel{~}{x}_0`$ of the basepoint $`x_0`$ and we lift the other cells one after each other, taking into account the relations in $`\pi _1(E(K)`$ and making sure that the union of cells already lifted is always connected. When a cell $`c`$ is reached for the first time, its lifting is chosen arbitrarily and declared to be $`\stackrel{~}{c}^{(0)}`$, but its boundary will involve in general $`\stackrel{~}{d}^{(n)}`$โs with $`n0`$. Once the lifted cellularization is known, it is a simple matter to determine preferred cell liftings: since the support of $`s^{}(P)`$ consists of 3 spiders, we only need to choose liftings of the 3 heads and then lift the legs.
Carrying out the computations we have explicitly found the algebraic complex with coefficients in $`\text{}[t^{\pm 1}]`$, and the preferred generators of the 4 moduli appearing. Then, using Maple, we have checked that indeed the complex is acyclic, and we have computed its torsion as follows:
$$\tau ^\phi (E(K),\xi )=\pm t^1.$$
Note that as an application of Proposition 4.17, by adding curls, we can easily construct a family $`\{K_n\}`$ of pseudo-Legendrian knots such that $`\tau ^\phi (E(K_n),\xi _n)=\pm t^n`$.
## 6 Main proofs
In this section we provide all the proofs which we have omitted in the rest of the paper. We will always refer to the statements for the notation.
Proof of1.1. Let us first recall the classical Poincarรฉ-Hopf formula, according to which if $`v`$ is a vector field with isolated singularities on a manifold $`M`$, and $`v`$ points outwards on $`M`$ (i.e. $`M`$ is black), then the sum of the indices of all singularities is $`\chi (M)`$. Assume now that $`v`$ has isolated singularities and on $`M`$ it is compatible with a pattern $`๐ซ=(W,B,V,C)`$. We claim that if $`๐`$ is a cellularization of $`M`$ suited to $`๐ซ`$ we have:
$$\underset{x\mathrm{Sing}(v)}{}\mathrm{ind}_x(v)=\chi (M)\underset{\sigma ๐,\sigma WV}{}\mathrm{ind}(\sigma ).$$
(6)
This formula is enough to prove the statement: if a non-singular field $`v`$ compatible with $`๐ซ`$ exists then the left-hand side of 6 vanishes, and the right-hand side of 6 equals the obstruction of the statement. On the other hand, if the obstruction vanishes, then one can first consider a singular field compatible with $`๐ซ`$, then group up the singularities in a ball, and remove them.
To prove 6 we consider the manifold $`M^{}`$ obtained by attaching a collar $`M\times [0,1]`$ to $`M`$ along $`M=M\times \{0\}`$. Of course $`M^{}M`$. We will now extend $`v`$ to a field $`v^{}`$ on $`M^{}`$ with the property that $`v^{}`$ points outwards on $`M^{}`$, and in $`M\times (0,1)`$ the field $`v^{}`$ has exactly one singularity for each cell $`\sigma WV`$, with index $`\mathrm{ind}(\sigma )`$. An application of the classical Poincarรฉ-Hopf formula then implies the conclusion. The construction of $`v^{}`$ is done cell by cell. We first show how the construction goes in dimension 2, see Fig. 15.
For the 3-dimensional case, it is actually convenient to choose a cellularization $`๐`$ of special type. Namely, we assume that $`๐|_M`$ consists of rectangles and triangles, each rectangle having exactly one edge on $`VC`$, and the union of rectangles covering a neighbourhood of $`VC`$. We suggest in Fig. 16 how to
define $`v^{}`$ on $`\sigma \times [0,1]`$ for $`\sigma W`$ of dimension 0, 1 and 2 respectively. By the choices we have made the situation near $`W`$ contains the 2-dimensional situation as a transversal cross-section, and it is not too difficult to extend $`v^{}`$ further and check that indices of the singularities are as required. As an example, we suggest in Fig. 17
how to do this near a convex edge. 1.1
Proof of1.5. Our proof follows the scheme given by Turaev in , with some technical simplifications and some extra difficulties due to the tangency circles. We first recall that if $`๐ฎ`$ is a (smooth) triangulation of a manifold $`N`$, a (singular) vector field $`w_๐ฎ`$ on $`N`$ can be defined by the requirements that: (1) each simplex is a union of orbits; (2) the singularities are exactly the barycentres of the simplices; (3) barycentres of higher dimensional simplices are more attractive that those of lower dimensional simplices. More precisely, each orbit (asymptotically) goes from a barycentre $`p_\sigma `$ to a barycentre $`p_\sigma ^{}`$, where $`\sigma \sigma ^{}`$. It is automatic that $`\mathrm{ind}_{p_\sigma }(w_๐ฎ)=\mathrm{ind}(\sigma )`$. See Fig. 18 for a description of $`w_๐ฎ`$ on a 2-simplex of $`๐ฎ`$.
Let us consider now a triangulation $`๐ฏ`$ of $`M`$, and let us choose a representative $`z`$ of the given $`\xi \mathrm{Eul}^\mathrm{c}(M,๐ซ)`$ as in Proposition 1.4(3). We consider now the manifold $`M`$ obtained by attaching $`M\times [0,\mathrm{})`$ to $`M`$ along $`M=M\times \{0\}`$. Note that $`M^{}\mathrm{Int}(M)`$. Moreover $`๐ฏ`$ extends to a โtriangulationโ $`๐ฏ^{}`$ of $`M^{}`$, where on $`M\times (0,\mathrm{})`$ we have simplices with exactly one ideal vertex, obtained by taking cones over the simplices in $`M`$ and then removing the vertex. Even if $`๐ฏ^{}`$ is not strictly speaking a triangulation, the construction of $`w_๐ฏ^{}`$ makes sense, because the missing vertex at infinity would be a repulsive singularity anyway. We arrange things in such a way that if $`\sigma M`$ then the singularity in $`\sigma \times (0,\mathrm{})`$ is at height $`1`$, so it is $`p_\sigma \times \{1\}`$.
We will define now a smooth function $`h:M(0,\mathrm{})`$ and set $`M_h=M\{(x,t)M\times [0,\mathrm{}):th(x)\}`$, in such a way that $`w_๐ฏ^{}`$ is non-singular on $`M_h`$, and, modulo the natural homeomorphism $`MM_h`$, it induces on $`M_h`$ the desired boundary pattern $`๐ซ`$. Later we will show how to use $`z`$ to remove the singularities of $`w_๐ฏ^{}`$ on $`M_h`$.
To define the function $`h`$ we consider a (very thin) left half-collar $`L`$ of $`V`$ on $`M`$ and a right half-collar $`R`$ of $`C`$. Here โleftโ and โrightโ refer to the natural orientations of $`M`$ and of $`VC`$. Note that $`LB`$ and $`RW`$. Now we set $`h|_{BL}1/2`$, and $`h|_{WR}2`$. Figures 19 and 20
respectively show that away from $`VC`$ indeed the pattern of $`w_๐ฏ^{}`$ on
$`M_h`$ is as required. Now we identify $`L`$ to $`V\times [1,0]`$ and $`R`$ to $`C\times [0,1]`$, and we define $`h(x,s)=f(s)`$ for $`(x,s)V\times [1,0]`$ and $`h(x,s)=f(s1)`$ for $`(x,s)C\times [0,1]`$, where $`f:[1,0][1/2,2]`$ is a smooth monotonic function with all the derivatives vanishing at $`1`$ and $`0`$. Instead of describing $`f`$ explicitly we picture it and show that also near $`VC`$ the pattern is as required. This is done near $`V`$ and $`C`$ respectively in Figg. 21
and 22. In both pictures we have only considered a special configuration for the triangulation on $`M`$, and we have refrained from picturing the orbits of the field in the 3-dimensional figure. Instead, we have separately shown the orbits on the vertical simplices on which the value of $`h`$ changes.
The conclusion is now exactly as in Turaevโs argument, so we only give a sketch. The chosen representative $`z`$ of $`\xi \mathrm{Eul}^\mathrm{c}(M,๐ซ)`$ can be described as an integer linear combination of orbits of $`w_๐ฏ^{}`$, which we can describe as segments $`[p_\sigma ,p_\sigma ^{}]`$ with $`\sigma \sigma ^{}`$. Now we consider the chain
$$z^{}=z\underset{\sigma WV}{}\mathrm{ind}(\sigma )p_\sigma \times [0,1].$$
(7)
By definition of $`h`$ we have that $`z^{}`$ is a 1-chain in $`M_h`$, and $`z^{}`$ consists exactly of the singularities of $`w_๐ฏ^{}`$ contained in $`M_h`$, each with its index. For each segment $`s`$ which appear in $`z^{}`$ we first modify $`w_๐ฏ^{}`$ to a field which is โconstantโ on a tube $`T`$ around $`s`$, and then we modify the field again within $`T`$, in a way which depends on the coefficient of $`s`$ in $`z^{}`$. The resulting field has the same singularities as $`w_๐ฏ^{}`$, but one checks that these singularities can be removed by a further modification supported within small balls centred at the singular points. We define $`\mathrm{\Psi }(\xi )`$ to be the class in $`\mathrm{Eul}^\mathrm{s}(M,๐ซ)`$ of this final field. Turaevโs proof that $`\mathrm{\Psi }`$ is indeed well-defined and $`H_1(M;\text{})`$-equivariant
applies without essential modifications. 1.5
###### Remark 6.1
In the previous proof we have defined $`\mathrm{\Psi }`$ using triangulations, in order to apply directly Turaevโs technical results (in particular, invariance under subdivision). However the geometric construction makes sense also for cellularizations $`๐`$ more general than triangulations, the key point being the possibility of defining a field $`w_๐`$ satisfying the same properties as the field defined for triangulations. This is certainly true, for instance, for cellularizations $`๐`$ of $`M`$ induced by realizations of $`M`$ by face-pairings on a finite number of polyhedra, assuming that the projection of each polyhedron to $`M`$ is smooth.
Proof of1.8. To help the reader follow the details, we first outline the scheme of the proof:
1. By identifying $`M`$ to a collared copy of itself, we choose a representative $`z`$ of the given $`\xi \mathrm{Eul}^\mathrm{c}(M,๐ซ)`$ such that the extra terms added to define $`\mathrm{\Theta }^\mathrm{c}(\xi )`$ cancel with terms already appearing in $`z`$. (We know a priori that this happens at the level of boundaries, but it may well not happen at the level of $`1`$-chains.)
2. We apply Remark 6.1 and choose a cellularization of $`M`$ in which it is particularly easy to analyze $`\mathrm{\Psi }(\xi )`$ and $`\mathrm{\Psi }(\mathrm{\Theta }^\mathrm{c}(\xi ))`$, both constructed using the representative $`z`$ already obtained.
We consider a cellularization $`๐`$ of $`M`$ satisfying the same assumptions on $`M`$ as those considered in the proof of Proposition 1.1, namely $`CV`$ is surrounded on both sides by a row of rectangular tiles, and the other tiles are triangular. We denote by $`\gamma _1,\mathrm{},\gamma _n`$ the segments in $`C`$, oriented as $`C`$.
Let us consider a representative $`z`$ relative to $`๐`$ of the given $`\xi \mathrm{Eul}^\mathrm{c}(M,๐ซ)`$. We construct a new copy $`M_1`$ of $`M`$ by attaching $`M\times [1,0]`$ to $`M`$ along $`M=M\times \{1\}`$, and we extend $`๐`$ to $`๐_1`$ by taking the product cellularization on $`M\times [1,0]`$. We define a new chain as
$`z_1`$ $`=`$ $`z+{\displaystyle \underset{\sigma B}{}}\mathrm{ind}(\sigma )p_\sigma \times [1/2,0]{\displaystyle \underset{\sigma WV}{}}\mathrm{ind}(\sigma )p_\sigma \times [1,1/2]`$
$`+{\displaystyle \underset{j=1}{\overset{n}{}}}\left(\gamma _j|_{[1/2,1]}\times \{1/2\}\gamma _j|_{[1/2,1]}\times \{0\}\right).`$
Note that $`z_1`$ is an Euler chain in $`M_1`$ with respect to $`๐_1`$. Consider the natural homeomorphism $`f:MM_1`$ and the class
$$a=\alpha ^\mathrm{c}(f_{}(\xi ),[z_1])H_1(M_1;\text{})$$
which may or not be zero. Since the inclusion of $`M`$ into $`M_1`$ is an isomorphism at the $`H_1`$-level, $`a`$ can be represented by a $`1`$-chain in $`M`$, so $`z_1`$ can be replaced by a new Euler chain $`z_2`$ such that $`[z_2]=f_{}(\xi )`$ and $`z_2`$ differs from $`z_1`$ only on $`M`$.
Renaming $`M_1`$ by $`M`$ and $`z_2`$ by $`z`$ we have found a representative $`z`$ of $`\xi `$ such that $`z=z_\theta +_{j=1}^n\gamma _j|_{[1/2,1]}`$, where $`z_\theta `$ is a sum of simplices contained in $`B\mathrm{Int}M`$. Note that of course $`\mathrm{\Theta }^\mathrm{c}(\xi )=[z_\theta ]`$. To conclude the proof we need now to analyze $`\mathrm{\Psi }(\xi )`$, constructed using $`z`$, and $`\mathrm{\Psi }(\mathrm{\Theta }^\mathrm{c}(\xi ))`$, constructed using $`[z_\theta ]`$, and show that $`\mathrm{\Psi }(\mathrm{\Theta }^\mathrm{c}(\xi ))=\mathrm{\Theta }^\mathrm{s}(\mathrm{\Psi }(\xi ))`$. By construction $`\mathrm{\Psi }(\xi )`$ and $`\mathrm{\Psi }(\mathrm{\Theta }^\mathrm{c}(\xi ))`$ will only differ near $`C`$, and we concentrate on one component of $`C`$ to show that the difference is exactly (up to homotopy) as in the definition of $`\mathrm{\Theta }^\mathrm{s}`$, i.e. as in Fig. 2.
The difference between $`\mathrm{\Psi }(\xi )`$ and $`\mathrm{\Psi }(\mathrm{\Theta }^\mathrm{c}(\xi ))`$ is best visualized on a cross-section of the form $`C\times [0,\mathrm{})`$. We leave to the reader to analyze the complete 3-dimensional pictures. To understand the cross-section, we follow the various steps in the proof of Theorem 1.5.
The first step in the definition of $`\mathrm{\Psi }(\xi )`$ (respectively, $`\mathrm{\Psi }(\mathrm{\Theta }^\mathrm{c}(\xi ))`$) consists in choosing the height function $`h`$ (respectively, $`h_\theta `$) and replacing the chains $`z`$ (respectively, $`z_\theta `$) by a chain $`z^{}`$ (respectively, $`z_\theta ^{}`$) as in formula (7). This is done in Fig. 23
where only the difference between the chains is shown.
To conclude we must modify the field $`w_๐`$ within a small neighbourhood of the support of $`z^{}`$ and $`z_\theta ^{}`$. This is done in Figg. 24
and 25 respectively. The rightmost picture in Fig. 25
is obtained by homotopy on the previous one. The representatives of $`\mathrm{\Psi }(\xi )`$ and $`\mathrm{\Psi }(\mathrm{\Theta }^\mathrm{c}(\xi ))`$ can be compared directly, and indeed they differ by a curve parallel to $`C`$ and directed consistently with $`C`$, so $`\mathrm{\Psi }(\mathrm{\Theta }^\mathrm{c}(\xi ))=\mathrm{\Theta }^\mathrm{s}(\mathrm{\Psi }(\xi ))`$. 1.8
We give now the proof omitted in Section 4.
Proof of4.17. Let us first note that the comparison class which we must show to be $`[m]`$ is independent of $`f`$ by Proposition 4.13. We will give two completely independent (but somewhat sketchy) proofs that this class is indeed $`[m]`$.
For a first proof, instead of comparing a โstraightโ knot with one with two curls, we compare two knots with one curl, chosen so that the framing is the same but the winding number is different. This is of course equivalent. The two knots are shown in Fig. 26 as thick tubes, together with one specific orbit of the field they are immersed in. The resulting bicoloration on the boundary of the tubes is also outlined.
To compare the curls we isotope the bicolorated tubes to the same straight tube, and we show how the orbit of the field is transformed under this isotopy. This is done in Fig. 27.
Also from this very partial picture it is quite evident that the resulting fields wind in opposite directions around the tube. A more accurate picture would show that the difference is actually a meridian of the tube.
Another (indirect) proof goes as follows. Note first that the comparison class which we must compute certainly is a multiple of $`[m]`$, say $`k[m]`$. Note also that $`k`$ is independent of the ambient manifold $`(M,v)`$. Moreover, by symmetry, we will have $`\alpha (\xi (v,K_0),f_{}(\xi (v,K_1)))=k[m]`$ if $`K_1`$ is obtained by locally adding a double curl with opposite winding number.
We take now $`M`$ to be $`S^3`$, with the field $`v`$ carried by the abalone $`P`$ as in Section 5, and $`K`$ to be a trivial knot contained in the โsmallerโ disc of the $`P`$. Using either the classical machinery of obstruction theory or the techniques developed in , one can see that there exists another pseudo-Legendrian knot $`K^{}`$ in $`(S^3,v)`$ such that $`\alpha (\xi (v,K),\xi (v,K^{}))=[m]`$, where by simplicity we are omitting the framed-isotopies necessary to compare these classes. As already remarked in Section 3, we can assume that $`K^{}`$ has a $`\mathrm{C}^1`$-projection on $`P`$. If one examines $`P`$ carefully one easily sees that $`K^{}`$ can actually be slid over $`P`$ to lie again in the small disc of $`P`$. Now $`K^{}`$ is a planar projection of the trivial knot, so through Reidemeister moves of types II and III, which correspond to isotopies through knots transversal to $`v`$, it can be transformed into a projection which differs from the trivial one only for a finite (even) number of curls. Summing up, we have a knot $`K^{}`$ such that $`\alpha (\xi (v,K),\xi (v,K^{}))=[m]`$ and $`K^{}`$ differs from $`K`$ only for a finite number of transformations of the form $`KK_1`$ or $`KK_1`$. This shows that $`[m]`$ is a multiple of $`k[m]`$, so $`k=\pm 1`$. 4.17
We conclude the paper by establishing the only statement given in Section 3 and not proved there. As above, we do not recall all the notation.
Proof of3.7. We fix $`P`$ and set $`s^{\prime \prime }=s^{\prime \prime }(P)`$, $`\widehat{v}=\widehat{v}(P)`$. Using Remark 6.1 we see that the construction of $`\mathrm{\Psi }([s^{\prime \prime }])`$ explained in the proof of Theorem 1.5 can be directly applied to the cellularization $`\widehat{๐ฏ}=\widehat{๐ฏ}(P)`$ of $`\widehat{M}`$. Recall that this construction requires identifying $`\widehat{M}`$ to a collared copy of itself, and extending $`s^{\prime \prime }`$ to a chain $`s^{\prime \prime \prime }`$ whose boundary consists precisely of the singularities of a field $`w`$. A representative of $`\mathrm{\Psi }([s^{\prime \prime }])`$ is then obtained by applying to $`w`$ a certain desingularization procedure. This desingularization is supported in a neighborhood of $`s^{\prime \prime \prime }`$, and one can easily check that the connected components of the support of $`s^{\prime \prime \prime }`$ (denoted henceforth by $`S`$) are actually contractible. Therefore, any desingularization of $`w`$ supported in a neighbourhood of $`s^{\prime \prime \prime }`$ will give a representative of $`\widehat{v}`$. We will prove the desired formula $`\mathrm{\Psi }([s^{\prime \prime }])=[\widehat{v}]`$ by exhibiting one such desingularization which is nowhere antipodal to $`\widehat{v}`$. In our argument we will always neglect the contraction of $`S_{\mathrm{triv}}^2`$ which maps $`M`$ onto $`\widehat{M}`$. (The desired formula actually holds at the level of $`M`$, and it easily implies the formula for $`\widehat{M}`$.)
By the above observations, the following claims easily imply the conclusion of the proof:
1. The set of points where $`w`$ is antipodal to $`\widehat{v}`$ is contained in $`S`$.
2. If $`S_0`$ is a component of $`S`$ then $`w`$ can be desingularized within a neighbourhood of $`S_0`$ to a field which is not antipodal to $`\widehat{v}`$ in the neighbourhood.
We prove claim 1 by first noting that the cells dual to those of $`P`$ are unions of orbits of both $`w`$ and $`\widehat{v}`$. So we can analyze cells separately. We do this explicitly only for 2-dimensional cells, leaving to the reader the other cases. In Fig. 28
we describe $`\widehat{v}`$. In the left-hand side of Fig. 29
we describe $`w`$ on the collared hexagon. In the right-hand side of the same figure we only show the singularities of $`w`$ on the renormalized hexagon, and the intersection of $`S`$ with the hexagon. In this figure the 7 short segments come from $`s^{\prime \prime \prime }s^{\prime \prime }`$; the other bits of $`S`$ have been labeled by โOrโ, โStโ, โBaโ or โHeโ to indicate that they come from orbits of $`\widehat{v}`$, stars, bi-arrows or half-edges.
This proves claim 1. Comparing Fig. 29 with Fig. 28, and carrying out the same analysis for 3-cells, one actually shows also claim 2 for components $`S_0`$ coming from $`s^{\prime \prime \prime }s^{\prime \prime }`$. Components of $`S`$ other than these can be described in one of the following ways:
* an orbit of $`\widehat{v}`$ emanating from a vertex of $`P`$;
* a half-edge of $`C`$;
* a bi-arrow together with an orbit of $`\widehat{v}`$ emanating from the centre of an edge of $`P`$ and reaching the centre of the bi-arrow;
* a star together with an orbit of $`\widehat{v}`$ emanating from the centre of a disc of $`P`$ and reaching the centre of the star.
All cases can be treated with the same method, we only do case (c). Figure 30
shows the component placed so that $`\widehat{v}`$ can be thought of as the constant vertical field pointing upwards, and the field $`w`$ near the component. The conclusion easily follows. 3.7
benedett@dm.unipi.it
petronio@dm.unipi.it
Dipartimento di Matematica
Via F. Buonarroti, 2
I-56127, PISA (Italy) |
warning/0002/nlin0002009.html | ar5iv | text | # Contents
## 1 Introduction: The tensorial approach and the birth of the method of Poisson pairs
This lecture is an introduction to the Hamiltonian analysis of PDEs form an โexperimentalโ point of view. This means that we are more concerned in unveiling the spirit of the method than in working out the theoretical details. Therefore the style of the exposition will be informal, and proofs will be mainly omitted. We shall follow, step by step, the birth and the evolution of the Hamiltonian analysis of the KdV equation
$$u_t=\frac{1}{4}(u_{xxx}6uu_x),$$
(1.1)
from its โinfancyโ to the final representation of the KdV flow as a linear flow on an infinite-dimensional Grassmannian due to Sato . The route is long and demanding. Therefore the exposition is divided in two parts, to be carried out in this and in the fourth lecture (see Section 4). Here our primary aim is to show the birth of the method of Poisson pairs. It is reached by means of a suitable use of the well-known methods of tensor analysis. We proceed in three steps. First, by using the transformation laws of vector fields, we construct the Miura map and the so called modified KdV equation (mKdV). This result leads quite simply to the theory of (elementary) Darboux transformations and to the concept of Poisson pair. Indeed, a peculiarity of mKdV is to possess an elementary Hamiltonian structure. By means of the transformation law of Poisson bivectors, we are then able to transplant this structure to the KdV equation, unraveling its โbi-Hamiltonian structureโ. This structure can be used in turn to define the concept of Lenard chain and to plunge the KdV equation into the โKdV hierarchyโ. This step is rather important from the point of view of finding classes of solutions to the KdV equation. Indeed the hierarchy is a powerful instrument to construct finite-dimensional invariant submanifolds of the equation and, therefore, finite-dimensional reductions of the KdV equation. The study of this process of restriction and of its use to construct solutions will be one of the two leading themes of these lectures. It is intimately related to the theory of separation of variables dealt with in the last two lectures. The second theme is that of the linearization of the full KdV flow on the infinite-dimensional Sato Grassmannian. The starting point of this process is surprisingly simple, and once again based on a simple procedure of tensor calculus. By means of the transformation laws of oneโforms, we pull back the KdV hierarchy from its phase space onto the phase space of the mKdV equation. In this way we obtain the โmKdV hierarchyโ. In the fourth lecture we shall show that this hierarchy can be written as a flow on an infiniteโdimensional Grassmann manifold, and that this flow can be linearized by means of a (generalized) Darboux transformation.
### 1.1 The Miura map and the KdV equation
As an effective way of probing the properties of equation (1.1) we follow the tensorial approach. Accordingly, we regard equation (1.1) as the definition of a vector field
$$u_t=X(u,u_x,u_{xx},u_{xxx})$$
(1.2)
on a suitable function space, and we investigate how it transforms under a point transformation in this space. Since our โcoordinateโ $`u`$ is a function and not simply a number, we are allowed to consider transformations of coordinates depending also on the derivatives of the new coordinate function of the type<sup>2</sup><sup>2</sup>2For further details on these kind of transformations, see .
$$u=\mathrm{\Phi }(h,h_x).$$
(1.3)
We ask whether there exists a transformed vector field
$$h_t=Y(h,h_x,h_{xx},h_{xxx})$$
(1.4)
related to the KdV equation according to the transformation law for vector fields,
$$X(\mathrm{\Phi }(h))=\mathrm{\Phi }_h^{}Y(h),$$
(1.5)
where $`\mathrm{\Phi }_h^{}`$ is the (Frรฉchet) derivative of the operator $`\mathrm{\Phi }`$ defining the transformation. This condition gives rise to a (generally speaking) over-determined system of partial differential equations on the unknown functions $`\mathrm{\Phi }(h,h_x)`$ and $`Y(h,h_x,h_{xx},h_{xxx}).`$ In the specific example the over-determined system can be solved. Apart form the trivial solution $`u=h_x`$, we find the Miura transformation ,
$$u=h_x+h^2\lambda ,$$
(1.6)
depending on an arbitrary parameter $`\lambda `$. The transformed equation is the modified KdV equation:
$$h_t=\frac{1}{4}(h_{xxx}6h^2h_x+6\lambda h_x).$$
(1.7)
###### Exercise 1.1
Work out in detail the transformation law (1.5), checking that $`X`$, $`\mathrm{\Phi }`$, and $`Y`$, defined respectively by equations (1.1), (1.6) and (1.7) do satisfy equation (1.5). $`\mathrm{}`$
The above result is plenty of consequences. The first one is a simple method for constructing solutions of the KdV equation. It is called the method of (elementary) Darboux transformations . It rests on the remark that the mKdV equation (1.7) admits the discrete symmetry
$$hh^{}=h.$$
(1.8)
Let us exploit this property to construct the well-known oneโsoliton solution of the KdV equation. We notice that the point $`u=0`$ is a very simple (singular) invariant submanifold of the KdV equation. Its inverse image under the Miura transformation is the 1โdimensional submanifold $`S_1`$ formed by the solutions of the special Riccati equation
$$h_x+h^2=\lambda .$$
(1.9)
This submanifold, in its turn, is invariant with respect to equation (1.7). A straightforward computation shows that, on this submanifold,
$$\frac{1}{4}(h_{xxx}6h^2h_x+6\lambda h_x)=\lambda h_x.$$
(1.10)
Therefore, on $`S_1`$ the mKdV equation takes the simple form $`h_t=\lambda h_x.`$ Solving the first order system formed by this equation and the the Riccati equation $`h_x+h^2=\lambda `$, and setting $`\lambda =z^2`$, we find the general solution
$$h(x,t)=z\mathrm{tanh}(zx+z^3t+c)$$
(1.11)
of the mKdV equation on the invariant submanifold $`S_1.`$ At this point we use the symmetry property and the Miura map. By the symmetry property (1.8) the function $`h(x,t)`$ is a new solution of the modified equation, and by the Miura map the function
$$u^{}(x,t)=h_x+h^2z^2=2z^2\text{sech}^2(zx+z^3t+c)$$
(1.12)
is a new solution of the KdV equation. It is called the one soliton solution<sup>3</sup><sup>3</sup>3For a very nice account of the origin and of the properties of the KdV equation and of other soliton equations and their solutions, see, e.g., .. It can also be interpreted in terms of invariant submanifolds. To this end, we have to notice that the Miura map (1.6) transforms the invariant submanifold $`S_1`$ of the modified equation into the submanifold formed by the solutions of the first order differential equation
$$\frac{1}{2}\left(\frac{1}{2}u_x^2+u^3\right)+\lambda u^2=0.$$
(1.13)
As one can easily check, this set is preserved by the KdV equation, and therefore it is an invariant oneโdimensional submanifold of the KdV equation, built up from the singular manifold $`u=0`$. On this submanifold, the KdV equation takes the simple form $`u_t=\lambda u_x`$, and the flow can be integrated to recover the solution (1.12).
This example clearly shows that the Darboux transformations are a mechanism to build invariant submanifolds of the KdV equation. Some of these submanifolds will be examined in great detail in the present lectures. The purpose is to show that the reduced equations on these submanifolds are classical Hamiltonian vector fields whose associated HamiltonโJacobi equations can be solved by separation of variables. In this way, we hope, the interest of the Hamiltonian analysis of the KdV equations can better emerge.
### 1.2 Poisson pairs and the KdV hierarchy
We shall now examine a more deep and far reaching property of the Miura map. It is connected with the concept of Hamiltonian vector field. From Analytical Mechanics, we know that the Hamiltonian vector fields are the images of exact oneโforms through a suitable linear map, associated with a soโcalled Poisson bivector. We shall formally define these notions in the next lecture. These definition can be easily extended to vector fields on infiniteโdimensional manifolds. Let us give an example, by showing that the mKdV equation is a Hamiltonian vector field. This requires a series of three consecutive remarks. First we notice that equation (1.7) can be factorized as
$$h_t=\left[\frac{1}{2}_x\right]\left[\frac{1}{2}h_xx\frac{3}{2}h^3+3\lambda h\right].$$
(1.14)
Then, we notice that the linear operator in the first bracket, $`\frac{1}{2}_x`$, is a constant skewsymmetric operator which we can recognize as a Poisson bivector. Finally, we notice that in the differential polynomial appearing in the second bracket in the right hand side of equation (1.14), we can easily recognize an exact oneโform. Indeed,
$$\left(\frac{1}{2}h_{xx}\frac{3}{2}h^3+3\lambda h\right)\dot{h}๐x=\frac{d}{dt}\left(\frac{1}{2}h_x^2\frac{3}{8}h^4+\frac{3}{2}\lambda h^2\right)๐x$$
(1.15)
for any tangent vector $`\dot{h}`$. These statements are true under suitable boundary conditions, as explained in, e.g., . Here and in the rest of these lectures we will tacitly use periodic boundary conditions.
The Hamiltonian character of the mKdV equation is obviously independent of the existence of the Miura map. However, this map finely combines this property from the point of view of tensor analysis. Let us recall that a Poisson bivector is a skewsymmetric linear map from the cotangent to the tangent spaces satisfying a suitable differential condition (see Lecture 2). It obeys the transformation law
$$Q_{\mathrm{\Phi }(h)}=\mathrm{\Phi }_h^{}P_h\mathrm{\Phi }_h^{}$$
(1.16)
under a change of coordinates (or a map between two different manifolds). In this formula the point transformation is denoted (in operator form) by $`u=\mathrm{\Phi }(h)`$. The operators $`\mathrm{\Phi }_h^{}`$ and $`\mathrm{\Phi }_h^{}`$ are the Frรฉchet derivative of $`\mathrm{\Phi }`$ and its adjoint operator. The symbols $`P_h`$ and $`Q_u`$ denote the Poisson bivectors in the space of the functions $`h`$ and $`u`$, respectively. Since the Miura map $`u=h_x+h^2\lambda `$ is not invertible, it is rather nontrivial that there exists a Poisson bivector $`Q_u`$, on the phase space of the KdV equation, which is $`\mathrm{\Phi }`$โrelated (in the sense of equation (1.16)) to the Poisson bivector $`P_h=\frac{1}{2}_x`$ associated with the modified equation. Surprisingly, this is the case. One can check that the operator $`Q_u`$ is defined by
$$Q_u=\frac{1}{2}_{xxx}+2(u+\lambda )_x+u_x.$$
(1.17)
###### Exercise 1.2
Verify the above claim by computing the product (in the appropriate order) of the operators $`\mathrm{\Phi }_h^{}=_x+h`$, $`P_h=\frac{1}{2}_x`$, and $`\mathrm{\Phi }_h^{}=_x+h`$, and by expressing the results in term of $`u=h_x+h^2\lambda `$. $`\mathrm{}`$
This exercise shows that the Miura map is a peculiar Poisson map. Since it depends on the parameter $`\lambda `$, the final result is that the phase space of the KdV equation is endowed with a oneโparameter family of Poisson bivectors,
$$Q_\lambda =Q_1\lambda Q_0,$$
(1.18)
which we call a Poisson pencil. The operators $`(Q_1,Q_0)`$ defining the pencil are said to form a Poisson pair, a concept to be systematically investigated in the next lecture.
These operators enjoy a number of interesting properties, and define new geometrical structures associated with the equation. One of the simplest but farโreaching is the concept of Lenard chain. The idea is to use the pair of bivectors to define a recursion relation on oneโforms:
$$Q_0\alpha _{j+1}=Q_1\alpha _j.$$
(1.19)
In the applications a certain care must be taken in dealing with this recursion relation, since it does not define uniquely the forms $`\alpha _j`$ (the operator $`Q_0`$ is seldom invertible). Furthermore, it is still less apparent that it can be solved in the class of exact oneโforms. However, in the KdV case we bonafide proceed and we find
$$\begin{array}{cc}\hfill \alpha _0& =1\hfill \\ \hfill \alpha _1& =\frac{1}{2}u\hfill \\ \hfill \alpha _2& =\frac{1}{8}(u_{xx}+3u^2)\hfill \\ \hfill \alpha _3& =\frac{1}{32}(10u^3+10uu_{xx}u_{xxxx}+5u_x^2)\hfill \end{array}$$
(1.20)
as first terms of the recurrence. The next step is to consider the associated vector fields (the meaning of numbering them with odd integers will be explained in Lecture 4):
$$\begin{array}{cc}\hfill \frac{u}{t_1}& =Q_1\alpha _0=Q_0\alpha _1=u_x\hfill \\ \hfill \frac{u}{t_3}& =Q_1\alpha _1=Q_0\alpha _2=\frac{1}{4}(u_{xxx}6uu_x)\hfill \\ \hfill \frac{u}{t_5}& =Q_1\alpha _2=Q_0\alpha _3=\frac{1}{16}(u_{xxxxx}10uu_{xxx}20u_xu_{xx}+30u^2u_x).\hfill \end{array}$$
(1.21)
They are the first members of the KdV hierarchy. In the fourth lecture, we shall show that it is a special instance of a general concept, the GelโfandโZakharevich hierarchy associated with any Poisson pencil of a suitable class.
### 1.3 Invariant submanifolds and reduced equations
The introduction of the KdV hierarchy has important consequences on the problem of constructing solutions of the KdV equation. The hierarchy is indeed a basic supply of invariant submanifolds of the KdV equation. This is due to the property of the vector fields of the hierarchy to commute among themselves. From this property it follows that the set of singular point of any linear combination (with constant coefficients) of the vector fields of the hierarchy is a finiteโdimensional invariant submanifold of the KdV flow. These submanifold can be usefully exploited to construct classes of solutions of this equation.
As a first example of this technique we consider the submanifold defined by the condition
$$\frac{u}{t_3}=\lambda \frac{u}{t_1},$$
(1.22)
that is, the submanifold where the second vector field of the hierarchy is a constant multiple of the first one. It is formed by the solutions of the third order differential equation
$$\frac{1}{4}(u_{xxx}6uu_x)\lambda u_x=0.$$
(1.23)
Therefore it is a three dimensional manifold, which we denote by $`M_3`$. We can use as coordinates on $`M_3`$ the values of the function $`u`$ and its derivatives $`u_x`$ and $`u_{xx}`$ at any point $`x_0`$. To avoid cumbersome notations, we will continue to denote these three numbers with the same symbols, $`u,u_x,u_{xx}`$, but the reader should be aware of this subtlety. To perform the reduction of the first equation of the hierarchy (1.21) on $`M_3`$, we consider the first two differential consequences of the equation $`{\displaystyle \frac{u}{t_1}}=u_x`$ and we use the constraint (1.22) to eliminate the third order derivative. We obtain the system
$$\frac{u}{t_1}=u_x,\frac{u_x}{t_1}=u_{xx},\frac{u_{xx}}{t_1}=6uu_x+4\lambda u_x.$$
(1.24)
We call $`X_1`$ the vector field defined by these equations on $`M_3`$. It shares many of the properties of the KdV equation. For instance, it is related to a Poisson pair. The simplest way to display this property is to remark that $`X_1`$ possesses two integrals of motion,
$$\begin{array}{cc}\hfill H_0& =u_{xx}3u^24\lambda u\hfill \\ \hfill H_1& =\frac{1}{2}u_x^2+u^3+2\lambda u^2+uH_0.\hfill \end{array}$$
(1.25)
Then we notice that on $`M_3`$ there exists a unique Poisson bracket $`\{,\}_0`$ with the following two properties:
1. The function $`H_0`$ is a Casimir function, that is, $`\{F,H_0\}_0=0`$ for every smooth function $`F`$ on $`M_3`$.
2. $`X_1`$ is the Hamiltonian vector field associated with the function $`H_1`$.
Such a Poisson bracket $`\{,\}_0`$ is defined by the relations
$$\{u,u_x\}_0=1,\{u,u_{xx}\}_0=0,\{u_x,u_{xx}\}_0=6u+4\lambda .$$
(1.26)
Similarly, one can notice that on $`M_3`$ there exists a unique Poisson bracket $`\{,\}_1`$ with the following โdualโ properties:
1. The function $`H_1`$ is a Casimir function, that is, $`\{F,H_1\}_1=0`$ for every smooth function $`F`$ on $`M_3`$.
2. $`X_1`$ is the Hamiltonian vector field associated with the function $`H_0`$.
This second Poisson bracket $`\{,\}_1`$ is defined by the relations
$$\{u,u_x\}_1=u,\{u,u_{xx}\}_1=u_x,\{u_x,u_{xx}\}_1=u_{xx}u(6u+4\lambda ).$$
(1.27)
###### Exercise 1.3
Verify the stated properties of the Poisson pair (1.26) and (1.27).
$`\mathrm{}`$
We now exploit the previous remarks to understand the geometry of the flow associated with $`X_1`$. First we use the Hamiltonian representation
$$\frac{dF}{dt}=X_1(F)=\{F,H_1\}_0.$$
(1.28)
It entails that the level surfaces of the Casimir function $`H_0`$ are twoโdimensional (symplectic) manifolds to which $`X_1`$ is tangent. Let us pick up any of these symplectic leaves, for instance the one passing through the origin $`u=0,u_x=0,u_{xx}=0`$. Let us call $`S_2`$ this leaf. The coordinates $`(u,u_x)`$ are canonical coordinates on $`S_2`$. The level curves of the Hamiltonian $`H_1`$ define a Lagrangian foliation of $`S_2`$. Our problem is to find the flow of the vector field $`X_1`$ along these Lagrangian submanifold. We have already given the solution of this problem in the particular case of the Lagrangian submanifold passing through the origin. This submanifold is the oneโdimensional invariant submanifold (1.13) previously discussed in connection with Darboux transformations. The relative flow is the oneโsoliton solution to KdV. To deal with the generic Lagrangian submanifolds on an equal footing, it is useful to change strategy and to use the HamiltonโJacobi equation
$$H_1(u,\frac{dW}{du})=E.$$
(1.29)
In this rather simple example, there is almost nothing to say about this equation (which is obviously solvable), and the second Poisson bracket (1.26) seems not to play any role in the theory.
This (wrong!) impression is promptly corrected by the study of a more elaborated example. Let us consider the fiveโdimensional submanifold $`M_5`$ of the singular points of the third vector field of the KdV hierarchy. It is defined by the equation
$$u_{xxxxx}10uu_{xxx}20u_xu_{xx}+30u^2u_x=0.$$
(1.30)
On this submanifold we can consider the restrictions of the first two vector fields of the same hierarchy. To compute the reduced equation we proceed as before. We regard the Cauchy data $`(u,u_x,u_{xx},u_{xxx},u_{xxxx})`$ as coordinates on $`M_5.`$ then we compute the time derivatives $`{\displaystyle \frac{u}{t_1}},{\displaystyle \frac{u_x}{t_1}},{\displaystyle \frac{u_{xx}}{t_1}},{\displaystyle \frac{u_{xxx}}{t_1}},{\displaystyle \frac{u_{xxxx}}{t_1}}`$ by taking the differential consequences of $`{\displaystyle \frac{u}{t_1}}=u_x`$, and by using (1.30) and its differential consequences as a constraint to eliminate all the derivatives of degree higher than four. We obtain the equations
$$\begin{array}{cc}\hfill \frac{u}{t_1}& =u_x\hfill \\ \hfill \frac{u_x}{t_1}& =u_{xx}\hfill \\ \hfill \frac{u_{xx}}{t_1}& =u_{xxx}\hfill \\ \hfill \frac{u_{xxx}}{t_1}& =u_{xxxx}\hfill \\ \hfill \frac{u_{xxxx}}{t_1}& =10uu_{xxx}+20u_xu_{xx}30u^2u_x\hfill \end{array}$$
(1.31)
In the same way, for the reduction of the KdV equation, we get
$$\begin{array}{cc}\hfill \frac{u}{t_3}& =\frac{1}{4}(u_{xxx}6uu_x)\hfill \\ \hfill \frac{u_x}{t_3}& =\frac{1}{4}(u_{xxxx}6uu_{xx}6u_x^2\hfill \\ \hfill \frac{u_{xx}}{t_3}& =\frac{1}{4}(4uu_{xxxx}+2u_xu_{xx}30u^2u_x)\hfill \\ \hfill \frac{u_{xxx}}{t_3}& =\frac{1}{4}(2u_{xx}^2+6u_xu_{xxx}+4uu_{xxxx}60uu_x^230u^2u_{xx})\hfill \\ \hfill \frac{u_{xxxx}}{t_3}& =\frac{1}{4}(10u_{xx}u_{xxx}+10u_xu_{xxxx}120u^3u_x100uu_xu_{xx}\hfill \\ & +10u^2u_{xxx}60u_x^3).\hfill \end{array}$$
(1.32)
###### Exercise 1.4
Verify the previous computations. $`\mathrm{}`$
To find the corresponding solutions of the KdV equation, regarded as a partial differential equation in $`x`$ and $`t`$, we have to:
1. Construct a common solution to the ordinary differential equations (1.31) and (1.32);
2. Consider the first component $`u(t_1,t_3)`$ of such a solution;
3. Set $`t_1=x`$ and $`t_3=t`$.
The function $`u(x,t)`$ obtained in this way is the solution we were looking for. What makes this procedure worth of interest is that the ODEs (1.31)โ(1.32) can be solved by means of a variety of methods. In particular, they can be solved by means of the method of separation of variables<sup>4</sup><sup>4</sup>4The fact that the stationary reductions of KdV can be solved by separation of variables is well-known (see, e.g., ). This classical method has recently found a lot of interesting new applications, as shown in the survey .. It can be shown that they are rather special equations: They are Hamiltonian with respect to a Poisson pair; this Poisson pair allows to foliate the manifold $`M_5`$ into fourโdimensional symplectic leaves with special properties; each symplectic leaf $`S_4`$ carries a Lagrangian foliation to which the vector fields (1.31) and (1.32) are tangent; the Poisson pair defines a special set of coordinates on each $`S_4`$; in these coordinates the HamiltonโJacobi equations associated with the Hamiltonian equations (1.31) and (1.32) can be simultaneously solved by additive separation of variables. Most of these properties will be proved in the next lecture.
### 1.4 The modified KdV hierarchy
We leave for the moment the theme of the reduction, and come back to the KdV hierarchy in its general form. We notice that the first equations (1.21) can also be written in the form
$$\begin{array}{cc}\hfill \frac{u}{t_1}& =(Q_1\lambda Q_0)\alpha _0\hfill \\ \hfill \frac{u}{t_3}& =(Q_1\lambda Q_0)(\lambda \alpha _0+\alpha _1)\hfill \\ \hfill \frac{u}{t_5}& =(Q_1\lambda Q_0)(\lambda ^2\alpha _0+\lambda \alpha _1+\alpha _2)\hfill \end{array}$$
(1.33)
This representation shows that these equations are Hamiltonian with respect to the whole Poisson pencil. This elementary property can be exploited to simply construct the modified KdV hierarchy. Let us write in general
$$\frac{u}{t_{2j+1}}=(Q_1+\lambda Q_0)\alpha ^{(j)}(\lambda ),$$
(1.34)
where
$$\alpha ^{(j)}(\lambda )=\lambda ^j\alpha _0+\lambda ^{j1}\alpha _1+\mathrm{}+\alpha _j.$$
(1.35)
By means of the Miura map (1.6) we pullโback the oneโforms $`\alpha ^{(j)}`$ to oneโform $`\beta ^{(j)}`$ defined on the phase space of the modified equation according to the transformation law of one-forms,
$$\beta (h)=\mathrm{\Phi }_{h}^{}{}_{}{}^{}\alpha (\mathrm{\Phi }(h)).$$
(1.36)
We then define the equations
$$\frac{h}{t_{2j+1}}=P_h\beta ^{(j)}(\lambda ).$$
(1.37)
They are $`\mathrm{\Phi }`$โrelated to the corresponding equations of the KdV hierarchy, exactly as the mKdV equation is $`\mathrm{\Phi }`$โrelated to the KdV equation. Indeed,
$$\frac{u}{t_{2j+1}}=\mathrm{\Phi }_h^{}\frac{h}{t_{2j+1}}=\mathrm{\Phi }_h^{}\mathrm{\Phi }_{h}^{}{}_{}{}^{}\alpha ^{(j)}(\mathrm{\Phi }(h);\lambda )=Q_u\alpha ^{(j)}(u;\lambda ).$$
(1.38)
It is therefore natural to call equations (1.37) the modified KdV hierarchy. By using the explicit form of the operators $`P_h`$ and $`\mathrm{\Phi }_{h}^{}{}_{}{}^{}`$, it is easy to check that the modified hierarchy is defined by the conservation laws
$$\frac{h}{t_{2j+1}}=_xH^{(2j+1)},$$
(1.39)
where
$$H^{(2j+1)}=\frac{1}{2}\alpha _x^{(j)}+\alpha ^{(j)}h.$$
(1.40)
###### Exercise 1.5
Write down explicitly the first three equations of the modified hierarchy. $`\mathrm{}`$
The above formulas are basic in the Sato approach. In the fourth lecture, after a more accurate analysis of the Hamiltonian structure of the KdV equations, we shall be led to consider the currents $`H^{(2j+1)}`$ as defining a point of an infiniteโdimensional Grassmannian. This point evolves in time as the point $`u`$ moves according to the KdV equation. We shall determine the equation of motion of the currents $`H^{(2j+1)}`$. They define a โbiggerโ hierarchy called the Central System. This system contains the KdV hierarchy as a particular reduction. It enjoys the property of being linearizable. In this way, by a continuous process of extension motivated by the Hamiltonian structure of the equations (from the single KdV equation to the KdV hierarchy and to the Central System), we arrive at the result that the KdV flow can be linearized. At this point the following picture of the possible strategies for solving the KdV equations emerges: either we pass to the Sato infiniteโdimensional Grassmannian and we use a linearization technique, or we restrict the equation to a finiteโdimensional invariant submanifold and we use a technique of separation of variables. The two strategies complement themselves rather well. Our attitude is to see the Grassmannian picture as a compact way of defining the equations, and the โreductionistโ picture as an effective way for finding interesting classes of solutions.
### The plan of the lectures
This is the web of ideas which we would like to make more precise in the following lectures. As cornerstone of our presentation we choose the concept of Poisson pairs. In the second lecture, we develop the theory of these pairs up to the point of giving a sound basis to the concept of Lenard chain. In the third lecture we exhibit a first class of examples, and we explain a reduction technique allowing to construct the Poisson pairs of the reduced flows. In the fourth lecture we give a second look at the KdV theory, and we explain the reasons which, according to the Hamiltonian standpoint, suggest to pass on the infiniteโdimensional Sato Grassmannian. In the fifth lecture we better explore the relation between the two strategies, and we touch the concept of Lax representation. Finally, the last lecture is devoted to the method of separation of variables. The purpose is to show how the geometry of the reduced Poisson pairs can be used to define the separation coordinates.
## 2 The method of Poisson pairs
In the previous lecture we have outlined the birth of the method of Poisson pairs and its main purpose: To define integrable hierarchies of vector fields. In this lecture we dwell on the theoretical basis of this construction presenting the concept of GelโfandโZakharevich system.
The starting point is the notion of Poisson manifold. A manifold is said to be a Poisson manifold<sup>5</sup><sup>5</sup>5The books and are very good references for this topic. if a composition law on scalar functions has been defined obeying the usual properties of a Poisson bracket: bilinearity, skewsymmetry, Jacobi identity and Leibnitz rule. The last condition means that that the Poisson bracket is a derivation in each entry:
$$\{fg,h\}=\{f,h\}g+f\{g,h\}.$$
(2.1)
Therefore, by fixing the argument of one of the two entries and keeping free the remaining one, we obtain a vector field,
$$X_h=\{,h\}.$$
(2.2)
It is called the Hamiltonian vector field associated with the function $`h`$ with respect to the given Poisson bracket. Due to the remaining conditions on the Poisson bracket, these vector fields are closed with respect to the commutator. They form a Lie algebra homeomorphic to the algebra of functions defined by the Poisson bracket:
$$[X_f,X_g]=X_{\{f,g\}}.$$
(2.3)
Therefore a Poisson bracket on a manifold has a twofold role: it defines a Lie algebra structure on the ring of $`C^{\mathrm{}}`$โfunctions, and provides a representation of this algebra on the manifold by means of the Hamiltonian vector fields.
Instead of working with the Poisson bracket, it is often suitable to work (especially in the infiniteโdimensional case) with the associated Poisson tensor. It is the bivector field $`P`$ on $`M`$ defined by
$$\{f,g\}=df,Pdg.$$
(2.4)
In local coordinates, its components $`P^{jk}(x^1,\mathrm{},x^n)`$ are the Poisson brackets of the coordinate functions,
$$P^{jk}(x^1,\mathrm{},x^n)=\{x^j,x^k\}.$$
(2.5)
By looking at this bivector field as a linear skewsymmetric map $`P:T^{}MTM`$, we can define the Hamiltonian vector fields $`X_f`$ as the images through $`P`$ of the exact one-forms,
$$X_f=Pdf.$$
(2.6)
In local coordinates this means
$$X_f^j=P^{jk}\frac{f}{x^k}.$$
(2.7)
###### Exercise 2.1
Show that the components of the Poisson tensor satisfy the cyclic condition
$$\underset{l}{}\left(P^{jl}\frac{P^{km}}{x^l}+P^{kl}\frac{P^{mj}}{x^l}+P^{ml}\frac{P^{jk}}{x^l}\right)=0.$$
(2.8)
$`\mathrm{}`$
###### Exercise 2.2
Suppose that $`M`$ is an affine space $`A`$. Call $`V`$ the vector space associated with $`A`$. Define a bivector field on $`A`$ as a mapping $`P:A\times V^{}V`$ which satisfies the skewsymmetry condition
$$\alpha ,P_u\beta =\beta ,P_u\alpha $$
(2.9)
for every pair of covector $`(\alpha ,\beta )`$ in $`V^{}`$ and at each point $`uA`$. Denote the directional derivative at the point $`u`$ of the mapping $`uP_u\alpha `$ along the vector $`v`$ by
$$P_u^{}(\alpha ;v)=\frac{d}{dt}P_{u+tv}\alpha |_{t=0}.$$
(2.10)
Show that the bivector $`P`$ is a Poisson bivector if and only if it satisfies the cyclic condition
$$\alpha ,P_u^{}(\beta ;P_u\gamma )+\beta ,P_u^{}(\gamma ;P_u\alpha )+\gamma ,P_u^{}(\alpha ;P_u\beta )=0.$$
(2.11)
$`\mathrm{}`$
###### Exercise 2.3
Check that the bivector $`Q_\lambda `$ of equation (1.17), associated with the KdV equation, fulfills the conditions (2.9) and (2.11). $`\mathrm{}`$
No condition is usually imposed on the rank of the Poisson bracket, that is, on the dimension of the vector space spanned by the Hamiltonian vector fields at each point of the manifold. If these vector fields span the whole tangent space the bracket is said to be regular, and the manifold $`M`$ turns out to be a symplectic manifold. Indeed there exists, in this case, a unique symplectic 2-form $`\omega `$ such that
$$\{f,g\}=\omega (X_f,X_g).$$
(2.12)
More interesting is the case where the bracket is singular. In this case, the Hamiltonian vector fields span a proper distribution $`D`$ on $`M`$. It is involutive but, generically, not of constant rank. Nonetheless, this distribution is completely integrable: at each point there exists an integral submanifold of maximal dimension which is tangent to the distribution. These submanifolds are symplectic manifolds, and are called the symplectic leaves of the Poisson structure. The symplectic form is still defined by equation (2.12). Indeed, even if there is not a 1โ1 correspondence between (differentials of) functions and Hamiltonian vector fields, this formula keeps its meaning, since the value of the Poisson bracket does not depend on the particular choice of the Hamiltonian function associated with a given Hamiltonian vector field. We arrive thus at the following conclusion: a Poisson manifold is either a symplectic manifold or a stratification of symplectic manifolds possibly of different dimensions. It can be proven that, in a sufficiently small open set where the rank of the Poisson tensor is constant, these symplectic manifolds are the level sets of some smooth functions $`F_1,\mathrm{},F_k`$, whose differentials span the kernel of the Poisson tensor. They are called Casimir functions of $`P`$ (see below).
###### Exercise 2.4
Let $`\{x_1,x_2,x_3\}`$ be Cartesian coordinates in $`M=^3`$. Prove that the assignment
$$\{x_1,x_2\}=x_3,\{x_1,x_3\}=x_2,\{x_2,x_3\}=x_1$$
(2.13)
defines a Poisson tensor on $`M`$. Find its Casimir function, and describe the symplectic foliation associated with it. $`\mathrm{}`$
After these brief preliminaries on Poisson manifolds as natural settings for the theory of Hamiltonian vector fields, we pass to the theory of bi-Hamiltonian manifolds. Our purpose is to provide evidence that they are a suitable setting for the theory of integrable Hamiltonian vector fields. The simplest connection between the theory of integrable Hamiltonian vector fields and the theory of bi-Hamiltonian manifold is given by the GelโfandโZakharevich (GZ) theorem we shall discuss in this lecture.
A bi-Hamiltonian manifold is a Poisson manifold endowed with a pair of compatible Poisson brackets. We shall denote these brackets with $`\{f,g\}_0`$ and $`\{f,g\}_1`$. They are compatible if the Poisson pencil
$$\{f,g\}_\lambda :=\{f,g\}_1\lambda \{f,g\}_0$$
(2.14)
verifies the Jacobi identity for any value of the continuous (say, real) parameter $`\lambda `$. By means of this concept we catch the main features of the situation first met in the KdV example of Lecture 1. The new feature deserving attention is the dependence of the Poisson bracket (2.14) on the parameter $`\lambda `$. It influences all the objects so far introduced on a Poisson manifold: Hamiltonian fields and symplectic foliation. In particular, this foliation changes with $`\lambda `$. The useful idea is to extract from this moving foliation the invariant part. It is defined as the intersection of the symplectic leaves of the pencil when $`\lambda `$ ranges over $`\{\mathrm{}\}`$. The GZ theorem describes the structure of these intersections in particular cases.
Let us suppose that the dimension of $`M`$ is odd, $`\text{dim}M=2n+1`$, and that the rank of the Poisson pencil is maximal. This means that the dimension of the characteristic distribution spanned by the Hamiltonian vector field is $`2n`$ for almost all the values of the parameter $`\lambda `$, and almost everywhere on the manifold $`M`$. In this situation the generic symplectic leaf of the pencil has accordingly dimension $`2n`$ and the intersection of all the symplectic leaves are submanifolds of dimension $`n`$. For brevity, we shall call this intersection the support of the pencil. The GZ theorem displays an important property of the leaves of the support of the pencil.
###### Theorem 2.5
On a $`(2n+1)`$โdimensional bi-Hamiltonian manifold, whose Poisson pencil has maximal rank, the leaves of the support are generically Lagrangian submanifolds of dimension $`n`$ contained on each symplectic leaf of dimension $`2n`$.
This theorem contains two different statements. First of all it states that the dimension of the support is exactly half of the dimension of the generic symplectic leaf. It is the โhardโ part of the theorem. Then it claims that the leaves of the support are Lagrangian submanifolds. Contrary to the appearances, this is the easiest part of the theorem, as we shall see. To better understand the content of the GZ theorem, we deem suitable to look at it from a different and, so to say, more constructive, point of view. It requires the use of the concept of Casimir function, defined as a function which commutes with all the other functions with respect to the Poisson bracket. Equivalently, it can be defined as a function whose differential belongs to the kernel of the Poisson tensor, i.e., a function generating the null vector field. In the case of a Poisson pencil, the Casimir functions depend on the parameter $`\lambda `$. If the rank of the Poisson pencil is maximal, the Casimir function is essentially unique (two Casimir functions are functionally dependent). The main content of the GZ theorem is that there exists a Casimir function depending polynomially on the parameter $`\lambda `$, and that the degree of this polynomial is exactly $`n`$ if $`\text{dim}M=2n+1`$. Thus we can write the Casimir function in the form
$$C(\lambda )=C_0\lambda ^n+C_1\lambda ^{n1}+\mathrm{}+C_n.$$
(2.15)
This result means that the Poisson pencil selects $`n+1`$ distinguished functions $`(C_0,C_1,\mathrm{}C_n)`$. Generically these functions are independent. Their common level surfaces are the leaves of the support of the pencil. Indeed, on the support the function $`C(\lambda )`$ must be constant independently of the particular value of $`\lambda `$. Thus all the coefficients $`(C_0,C_1,\mathrm{}C_n)`$ must be separately constant. Furthermore, as a consequence of the fact that $`C(\lambda )`$ is a Casimir function, it is easily seen that the coefficients $`C_k`$ verify the Lenard recursion relations,
$$\{,C_k\}_1=\{,C_{k+1}\}_0,$$
(2.16)
together with the vanishing conditions
$$\{,C_0\}_0=\{,C_n\}_1=0.$$
(2.17)
In the language of the previous lecture, the functions $`(C_0,C_1,\mathrm{}C_n)`$ form a Lenard chain. A typical property of these functions is to be mutually in involution:
$$\{C_j,C_k\}_0=\{C_j,C_k\}_1=0.$$
(2.18)
This is proved by repeatedly using the recursion relation (2.16) to go back and forth along the chain. It follows that the leaves of the support are isotropic submanifolds, but since they are of maximal dimension $`n`$ they are actually Lagrangian submanifolds. These short remarks should give a sufficiently detailed idea of the meaning of the GZ theorem.
###### Exercise 2.6
Check that that the integrals of motion $`H_0`$ and $`H_1`$ of the reduced flow $`X_1`$ on the invariant submanifold $`M_3`$ considered in Section 1.3 are the coefficients of the Casimir function $`C(\lambda )=\lambda H_0+H_1`$ of the Poisson pencil defined on $`M_3`$. $`\mathrm{}`$
###### Exercise 2.7
Prove the claim (2.18) about the involutivity of the coefficients of a Casimir polynomial. $`\mathrm{}`$
From our standpoint, the above results are worthwhile of interest for two different reasons: First of all they show how the Lenard recursion relations, characteristic of the theory of โsoliton equationsโ, arise in a theoretically sound way in the framework of bi-Hamiltonian manifold. Secondly, they highlight the existence of a distinguished set of Hamiltonian $`(C_0,C_1,\mathrm{}C_n)`$ on the manifold $`M`$. Let us now choose one of the brackets of the pencil, say the bracket $`\{,\}_0`$. The function $`C_0`$ is a Casimir function for this bracket, and therefore its level surfaces are the synplectic leaves of the bracket $`\{,\}_0`$. Let us call $`\omega _0`$ the symplectic 2โform defined on these submanifolds. As a consequence of the involution relation (2.18), the restrictions of the $`n`$ functions $`(C_1,\mathrm{}C_n)`$ to the symplectic leaf are in involution with respect to $`\omega _0`$. According to the ArnolโdโLiouville theorem, they define a family (or โhierarchyโ) of $`n`$ completely integrable Hamiltonian vector fields on the symplectic leaf.
###### Definition 2.8
The family of completely integrable Hamiltonian systems defined by the functions $`(C_1,\mathrm{}C_n)`$ on each symplectic leaf of the Poisson bracket $`\{,\}_0`$ will be called the GZ hierarchy associated with the Poisson pencil $`\{,\}_\lambda `$ defined on the bi-Hamiltonian manifold $`M`$.
We shall be particularly interested in the study of this hierarchy for two reasons. First we want to show that the previous simple concepts allow to reconstruct a great deal of the KdV hierarchy, up to the linearization process on the infiniteโdimensional Sato Grassmannian. In other words, we want to show that the theory of Poisson pairs is a natural gate to the theory of infinite-dimensional Hamiltonian systems described by partial differential equations of evolutionary type. Secondly, in a finite-dimensional setting, we want to show that the GZ vector fields are often more than integrable in the Liouville sense. Indeed, under some mild additional assumptions on the Poisson pencil, they are separable, and the separation coordinates are dictated by the geometry of the bi-Hamiltonian manifold. This result strenghtens the connection between Poisson pairs and integrability.
## 3 A first class of examples and the reduction technique
The aim of this lecture is to present a first class of nontrivial examples of GZ hierarchies. The examples are constructed to reproduce the reduced KdV flows discussed in the first lecture. The relation, however, will not be immediately manifest, and the reader has to wait until the fifth lecture for a full understanding of the motivations for some particular choice herewith made.
This lecture is split into three parts. In the first one we introduce a simple class of bi-Hamiltonian manifolds called LieโPoisson manifolds. They are duals of Lie algebras endowed with a special Poisson pencil of Lie-theoretical origin. The Hamiltonian vector fields defined on these manifolds admit a Lax representation with a Lax matrix depending linearly on the parameter $`\lambda `$. In the second part we show how to combine several copies of these manifolds, in such a way to obtain Hamiltonian vector fields admitting a Lax representation depending polynomially on the parameter $`\lambda `$. Finally, in the third part, we introduce the geometrical technique of reduction of Marsden and Ratiu. It will allow us to specialize the form of the Lax matrix. The contact with the KdV theory, to be done in the fifth lecture, will then consist in showing that the reduced KdV flows admit exactly the Lax representation of the Hamiltonian vector fields considered in this lecture. This will ascertain the bi-Hamiltonian character of the reduced KdV flows. The lecture ends with an example worked out in detail.
### 3.1 LieโPoisson manifolds
In this section $`M=๐ค^{}`$ is the dual of a Lie algebra $`๐ค`$. We denote by $`S`$ a point in $`M`$, and by $`{\displaystyle \frac{F}{S}}`$ the differential of a function $`F:M`$. This differential is the unique element of the algebra $`๐ค`$ such that
$$\frac{dF}{dt}=\frac{F}{S},\dot{S}$$
(3.1)
along any curve passing through the point $`S`$. The Poisson pencil on $`M`$ is defined by
$$\{F,G\}_\lambda =S+\lambda A,[\frac{F}{S},\frac{G}{S}],$$
(3.2)
where $`A`$ is any fixed element in $`๐ค^{}`$. In all the examples related to the KdV theory, $`๐ค=๐ฐ๐ฉ(2)`$, $`S`$ and $`{\displaystyle \frac{F}{S}}`$ are traceless $`2\times 2`$ matrices, and
$$A=\left(\begin{array}{cc}0& 0\\ 1& 0\end{array}\right).$$
(3.3)
The Hamiltonian vector field $`X_F`$ has the form
$$\dot{S}=[S+\lambda A,\frac{F}{S}].$$
(3.4)
It is already in Lax form, with a Lax matrix given by
$$L(\lambda )=\lambda A+S.$$
(3.5)
###### Exercise 3.1
Compute the Poisson tensor and the Hamiltonian vector fields associated with the pencil (3.2). $`\mathrm{}`$
### 3.2 Polynomial extensions
We consider two copies of the algebra $`๐ค`$. Accordingly, we denote by $`(S_0,S_1)`$ a point in $`M`$ and by $`({\displaystyle \frac{F}{S_0}},{\displaystyle \frac{F}{S_1}})`$ the differential of the function $`F:M`$. By definition, along any curve $`t(S_0(t),S_1(t))`$ we have
$$\frac{d}{dt}F=\frac{F}{S_0},\dot{S_0}+\frac{F}{S_1},\dot{S_1}.$$
(3.6)
The two copies of the algebra are intertwined by the Poisson brackets. As a Poisson pair on $`M`$ we choose the following brackets
$$\begin{array}{cc}\hfill \{F,G\}_0& =A,[\frac{F}{S_0},\frac{G}{S_1}]+[\frac{F}{S_1},\frac{G}{S_0}]+S_1,[\frac{F}{S_0},\frac{G}{S_0}]\hfill \\ \hfill \{F,G\}_1& =A,[\frac{F}{S_1},\frac{G}{S_1}]S_0,[\frac{F}{S_0},\frac{G}{S_0}]\hfill \end{array}$$
(3.7)
The motivations can be found for instance in (see also ). Later on we shall see how to extend this definition to the case of an arbitrary number of copies.
###### Exercise 3.2
Check that equations (3.7) indeed define a Poisson pair. $`\mathrm{}`$
Let us now study the Hamiltonian vector fields. Those defined by the brackets $`\{,\}_0`$ have the form
$$\begin{array}{cc}\hfill \dot{S_0}& =[A,\frac{F}{S_1}]+[S_1,\frac{F}{S_0}]\hfill \\ \hfill \dot{S_1}& =[A,\frac{F}{S_0}].\hfill \end{array}$$
(3.8)
Those defined by the second bracket $`\{,\}_1`$ are
$$\begin{array}{cc}\hfill \dot{S_0}& =[S_0,\frac{F}{S_0}]\hfill \\ \hfill \dot{S_1}& =[A,\frac{F}{S_1}].\hfill \end{array}$$
(3.9)
It turns out that the Hamiltonian vector fields associated with the Poisson pencil are given by
$$\begin{array}{cc}\hfill \dot{S_0}& =[S_0+\lambda S_1,\frac{F}{S_0}][\lambda A,\frac{F}{S_1}]\hfill \\ \hfill \dot{S_1}& =[\lambda A,\frac{F}{S_0}]+[A,\frac{F}{S_1}].\hfill \end{array}$$
(3.10)
This computation allows to display an interesting property of these vector fields. If we multiply the second equation by $`\lambda `$ and add the result to the first equation we find
$$(\lambda ^2A+\lambda S_1+S_0)^{}=[\frac{F}{S_0},\lambda ^2A+\lambda S_1+S_0].$$
(3.11)
This is a Lax representation with Lax matrix $`L(\lambda )=\lambda ^2A+\lambda S_1+S_0.`$ It depends polynomially on the parameter of the pencil. We have thus ascertained that all the Hamiltonian vector fields relative to the Poisson pencil (3.10) admit a Lax representation. The converse, however, is not necessarily true. Indeed, it must be noticed that the single Lax equation (3.11) is not sufficient to completely reconstruct the Poisson pencil (3.10). Additional constraints on the matrix $`L(\lambda )`$ are required to make the problem well-posed. The kind of constraints to be set are suggested by the geometric theory of reduction which we shall now outline.
### 3.3 Geometric reduction
We herewith outline a specific variant of the reduction technique of Marsden and Ratiu for Poisson manifolds. This variant is particularly suitable for bi-Hamiltonian manifolds.
Among the geometric objects defined by a Poisson pair $`(P_0,P_1)`$ on a manifold $`M`$ we consider:
a symplectic leaf $`S`$ of one of the two Poisson bivectors, say $`P_0`$.
the annihilator $`(TS)^0`$ of the tangent bundle of $`S`$, spanned by the 1โforms vanishing on the tangent spaces to $`S`$.
the image $`D=P_1(TS)^0`$ of this annihilator according to the second Poisson bivector $`P_1`$. It is spanned by the Hamiltonian vector fields associated with the Casimir functions of $`P_0`$ by $`P_1`$.
the intersection $`E=DTS`$ of the distribution $`D`$ with the tangent bundle of the selected symplectic leaf $`S`$.
It can be show that $`E`$ is an integrable distribution as a consequence of the compatibility of the Poisson brackets. Therefor we can consider the space of leaves of this distribution, $`N=S/E`$. We assume $`N`$ to be a smooth manifold. By the MarsdenโRatiu theorem, $`N`$ is a reduced bi-Hamiltonian manifold.
The reduced brackets on $`N`$ can be computed by using the process of โprolongation of functionsโ from $`N`$ to $`M`$. Given any function $`f:N`$, we consider it as a function on $`S`$, invariant along the leaves of $`E`$. Then we choose any function $`F:M`$ which annihilates $`D`$ and coincides with $`f`$ on $`S`$. This function is said to be a prolongation of $`f`$. It is not unique, but this fact is not disturbing. It can be show that, if $`F`$ and $`G`$ are prolongations of $`f`$ and $`g`$, their bracket $`\{F,G\}_\lambda `$ is an invariant function along $`E`$. Therefore it defines a function on $`N`$ which is by definition the reduced bracket $`\{f,g\}_\lambda `$. The final result, of course, is independent of the particular choices of the prolongations $`F`$ and $`G`$.
### 3.4 An explicit example
According to the spirit of these lectures, rather than discussing the proof of the reduction theorem stated in Section 3.3, we prefer to illustrate it on a concrete example. Let us thus perform the reduction of the Poisson pencil defined on two copies of $`๐ค=๐ฐ๐ฉ(2)`$. The matrices $`S_0`$ and $`S_1`$ are traceless matrices whose entries we denote as follows:
$$S_0=\left(\begin{array}{cc}p_0& r_0\\ q_0& p_0\end{array}\right),S_1=\left(\begin{array}{cc}p_1& r_1\\ q_1& p_1\end{array}\right).$$
(3.12)
The space $`M`$ has dimension six, and the entries of $`S_0`$ and $`S_1`$ are global coordinates on it. In these coordinates the differential of a function $`F:M`$ is represented by the pair of matrices
$$\frac{F}{S_0}=\left(\begin{array}{cc}\frac{1}{2}\frac{F}{p_0}& \frac{F}{q_0}\\ \frac{F}{r_0}& \frac{1}{2}\frac{F}{p_0}\end{array}\right),\frac{F}{S_1}=\left(\begin{array}{cc}\frac{1}{2}\frac{F}{p_1}& \frac{F}{q_1}\\ \frac{F}{r_1}& \frac{1}{2}\frac{F}{p_1}\end{array}\right).$$
(3.13)
###### Exercise 3.3
Define the pairing $`{\displaystyle \frac{F}{S}},\dot{S}`$ on $`๐ค`$ as the trace of the product of the matrices $`{\displaystyle \frac{F}{S}}`$ and $`\dot{S}`$. Show that the matrices (3.13) verify the defining equation (3.6). $`\mathrm{}`$
The Hamiltonian vector fields (3.8) and (3.9) are consequently given by
$$\begin{array}{cc}\hfill \dot{p_0}& =r_1\frac{F}{r_0}q_1\frac{F}{q_0}\frac{F}{q_1}\hfill \\ \hfill \dot{q_0}& =q_1\frac{F}{p_0}2p_1\frac{F}{r_0}+\frac{F}{p_1}\hfill \\ \hfill \dot{r_0}& =2p_1\frac{F}{q_0}r_1\frac{F}{p_0}\hfill \\ \hfill \dot{p_1}& =\frac{F}{q_0}\hfill \\ \hfill \dot{q_1}& =\frac{F}{p_0}\hfill \\ \hfill \dot{r_1}& =0\hfill \end{array}$$
(3.14)
and by
$$\begin{array}{cc}\hfill \dot{p_0}& =r_0\frac{F}{r_0}+q_0\frac{F}{q_0}\hfill \\ \hfill \dot{q_0}& =2p_0\frac{F}{r_0}q_0\frac{F}{p_0}\hfill \\ \hfill \dot{r_0}& =2p_0\frac{F}{q_0}+r_0\frac{F}{p_0}\hfill \\ \hfill \dot{p_1}& =\frac{F}{q_1}\hfill \\ \hfill \dot{q_1}& =\frac{F}{p_1}\hfill \\ \hfill \dot{r_1}& =0\hfill \end{array}$$
(3.15)
respectively.
Step 1: The reduced space $`N`$.
First we notice that the Hamiltonian vector fields (3.14) verify the constraints
$$\dot{r_1}=0,(r_0+p_1^2+r_1q_1)^{}=0.$$
(3.16)
It follows that the submanifold $`SM`$ defined by the equations
$$r_1=1,r_0+p_1^2+r_1q_1=0$$
(3.17)
is a symplectic leaf of the first Poisson bivector $`P_0`$. Furthermore, it follows that the annihilator $`(TS)^0`$ is spanned by the exact 1โforms $`dr_1`$ and $`d(r_0+p_1^2+r_1q_1)`$. By computing the images of these forms according to the second Poisson bivector (3.15), we find the distribution $`D`$. It is spanned by the single vector field
$$\begin{array}{cc}\hfill \dot{p_0}& =r_0\hfill \\ \hfill \dot{q_0}& =2p_0\hfill \\ \hfill \dot{r_0}& =0\hfill \\ \hfill \dot{p_1}& =1\hfill \\ \hfill \dot{q_1}& =2p_1\hfill \\ \hfill \dot{r_1}& =0\hfill \end{array}$$
(3.18)
which verifies the five constraints
$$\begin{array}{c}(p_0r_0p_1)^{}=0,(q_0+2p_0p_1r_0p_1^2)^{}=0,\hfill \\ (q_1+p_1^2)^{}=0,\dot{r_0}=0,\dot{r_1}=0.\hfill \end{array}$$
(3.19)
They show that $`DTS`$, and therefore $`E=D`$. Moreover they yield that the leaves of $`E`$ on $`S`$ are the level curves of the functions
$$\begin{array}{cc}\hfill u_1& =q_1+p_1^2\hfill \\ \hfill u_2& =p_0+p_1q_1+p_1^3\hfill \\ \hfill u_3& =q_0+2p_0p_1+q_1p_1^2+p_1^4\hfill \end{array}$$
(3.20)
We conclude that:
* $`N=^3`$;
* $`(u_1,u_2,u_3)`$ are global coordinates on $`N`$;
* the canonical projection $`\pi :SS/E`$ is defined by equations (3.20).
Step 2: The reduced brackets.
Consider any function $`f:N`$. The function
$$F:=f(q_1+p_1^2,p_0+p_1q_1+p_1^3,q_0+2p_0p_1+q_1p_1^2+p_1^4)$$
(3.21)
is clearly a prolongation of $`f`$ to $`M`$, since it coincides with $`f`$ on $`S`$, and is invariant along $`D`$. We can thus use $`F`$ to compute the first component of the reduced Hamiltonian vector field on $`N`$ according to the following algorithm:
$$\begin{array}{cc}\hfill \dot{u_1}& \stackrel{\left(\text{3.20}\right)}{=}\dot{q_1}+2p_1\dot{p_1}\stackrel{\left(\text{3.14}\right)}{=}\frac{F}{p_0}p_1\frac{F}{q_0}\hfill \\ & \stackrel{\left(\text{3.21}\right)}{=}\left(\frac{f}{u_2}+2p_1\frac{f}{u_3}\right)2p_1\frac{f}{u_3}=\frac{f}{u_2}.\hfill \end{array}$$
(3.22)
The other components are evaluated in the same way. The final result is that the Hamiltonian vector fields associated with the reduced Poisson pencil on $`N`$ are defined by
$$\begin{array}{cc}& \dot{u_1}=(u_1+\lambda )\frac{f}{u_2}+2u_2\frac{f}{u_3}\hfill \\ & \dot{u_2}=(u_1+\lambda )\frac{f}{u_1}+(u_32\lambda u_1)\frac{f}{u_3}\hfill \\ & \dot{u_3}=2u_2\frac{f}{u_1}+(2\lambda u_1u_3)\frac{f}{u_2}\hfill \end{array}$$
(3.23)
With this reduction process we passed from a sixโdimensional manifold $`M`$ to a three dimensional manifold $`N`$. Later on, we shall see that this manifold coincides with the invariant submanifold $`M_3`$ of KdV, defined by the constraint
$$u_{xxx}6uu_x=0.$$
(3.24)
Step 3: the GZ hierarchy.
To compute the Casimir function of the pencil (3.23) we notice that these vector fields obey the constraint
$$(2\lambda u_1u_3)\dot{u_1}+2u_2\dot{u_2}(u_1+\lambda )\dot{u_3}=0.$$
(3.25)
Therefore, integrating this equation, we obtain that
$$C(\lambda )=\lambda (u_1^2u_3)+(u_2^2u_1u_3)=\lambda C_0+C_1$$
(3.26)
is the Casimir sought for. It fulfills the scheme of the GZ theorem, and it defines a โshortโ Lenard chain
$$P_0dC_0=0P_1dC_0=P_0dC_1=X_1P_1dC_1=0.$$
(3.27)
Therefore the GZ โhierarchyโ consists of the single vector field
$$X_1:\begin{array}{c}\dot{u_1}=2u_2\hfill \\ \dot{u_2}=u_3+2u_1^2\hfill \\ \dot{u_3}=4u_1u_2\hfill \end{array}$$
(3.28)
As a last remark, we notice that this vector field coincides with the restriction of the first equation $`{\displaystyle \frac{u}{t_1}}=u_x`$ of the KdV hierarchy on the invariant submanifold (3.24). Indeed, by the procedure explained in Section 1, the reduced equation written in the โCauchy data coordinatesโ $`(u,u_x,u_{xx})`$ is given by
$$\begin{array}{cc}& \frac{u}{t_1}=u_x\hfill \\ & \frac{u_x}{t_1}=u_{xx}\hfill \\ & \frac{u_{xx}}{t_1}=6uu_x\hfill \end{array}$$
(3.29)
We can now pass from (3.29) to (3.28) by the change of variables
$$u_1=\frac{1}{2}u,u_2=\frac{1}{4}u_x,u_3=\frac{1}{4}u_{xx}\frac{1}{2}u^2.$$
(3.30)
This remark shows that the simplest reduced KdV flow is bi-Hamiltonian. In the fifth lecture we shall see that this property is general, and we shall explain the origin of the seemingly โad hocโ change of variables (3.30).
### 3.5 A more general example
To deal with higherโorder reduced KdV flows, we have to extend the class of bi-Hamiltonian manifolds to be considered. We outline the case of three copies of the algebra $`๐ค`$. The formulas are similar to the ones of equation (3.7), albeit a little more involved. The brackets $`\{F,G\}_0`$ and $`\{F,G\}_1`$ are now given by
$$\begin{array}{cc}\hfill \{F,G\}_0& =A,[\frac{F}{S_0},\frac{G}{S_2}]+[\frac{F}{S_1},\frac{G}{S_1}]+[\frac{F}{S_2},\frac{G}{S_0}]\hfill \\ & +S_2,[\frac{F}{S_0},\frac{G}{S_1}]+[\frac{F}{S_1},\frac{G}{S_0}]\hfill \\ & +S_1,[\frac{F}{S_0},\frac{G}{S_0}]\hfill \end{array}$$
(3.31)
and
$$\begin{array}{cc}\hfill \{F,G\}_1& =A,[\frac{F}{S_1},\frac{G}{S_2}]+[\frac{F}{S_2},\frac{G}{S_1}]\hfill \\ & +S_2,[\frac{F}{S_1},\frac{G}{S_1}]\hfill \\ & S_0,[\frac{F}{S_0},\frac{G}{S_0}].\hfill \end{array}$$
(3.32)
The comparison of the two examples allows to infer by induction the general rule for the Poisson pair, holding in the case of an arbitrary (finite) number of copies of $`๐ค`$. The pencil (3.31)โ(3.32) can be reduced according to the procedure shown before. If $`๐ค=๐ฐ๐ฉ(2)`$ and $`A`$ is still given by (3.3), the final result of the process is the following: We start from a nineโdimensional manifold $`M`$ and, after reduction, we arrive at a fiveโdimensional manifold $`N`$. It fulfills the assumption of the GZ theorem. The GZ hierarchy consists of two vector fields, which are the reduced KdV flows given by (1.31) and (1.32).
###### Exercise 3.4
Perform the reduction of the pencil (3.31)โ(3.32)) for $`๐ค=๐ฐ๐ฉ(2)`$.
$`\mathrm{}`$
## 4 The KdV theory revisited
In this lecture we consider again the KdV theory, but from a new point of view. Our purpose is twofold. The first aim is to show that the KdV hierarchy is another example of GZ hierarchy. The second aim is to explain in which sense the KdV hierarchy can be linearized. The algebraic linearization procedure dealt with in this lecture was suggested for the first time by Sato (see also the developments contained in ), who exploited the soโcalled Lax representation of the KdV hierarchy in the algebra of pseudoโdifferential operators. Here we shall give a different description, strictly related to the Hamiltonian representation of the KdV hierarchy as a kind of infinite-dimensional GZ hierarchy. However, the presentation does not go beyond the limits of a simple sketch of the theory. We refer to for full details.
### 4.1 Poisson pairs on a loop algebra
In this section we consider the infinite-dimensional Lie algebra $`M`$ of $`C^{\mathrm{}}`$โmaps from the circle $`S^1`$ into $`๐ค=๐ฐ๐ฉ(2)`$. A generic point of this manifold is presently a $`2\times 2`$ traceless matrix
$$S=\left(\begin{array}{cc}p(x)& r(x)\\ q(x)& p(x)\end{array}\right),$$
(4.1)
whose entries are periodic functions of the coordinate $`x`$ running over the circle. The three functions $`(p,q,r)`$ play the role of โcoordinatesโ on our manifold. The scalar-valued functions $`F:M`$ to be considered are local functionals
$$F=_{S^1}f(p,q,r;p_x,q_x,r_x;\mathrm{})๐x.$$
(4.2)
As before, their differentials are given by the matrices
$$\frac{\delta F}{\delta S}=\left(\begin{array}{cc}\frac{1}{2}\frac{\delta f}{\delta p}& \frac{\delta f}{\delta q}\\ \frac{\delta f}{\delta r}& \frac{1}{2}\frac{\delta f}{\delta p}\end{array}\right),$$
(4.3)
whose entries are the variational derivatives of the Lagrangian density $`f`$ with respect to the functions $`(p,q,r)`$. The Poisson pencil is similar to the first one considered in the previous lecture (see equation (3.2)). It is defined by
$$\{F,G\}_\lambda =S+\lambda A,[\frac{\delta F}{\delta S},\frac{\delta G}{\delta S}]+\omega (\frac{\delta F}{\delta S},\frac{\delta G}{\delta S}).$$
(4.4)
It differs from the previous example by the addition of the nontrivial cocycle
$$\omega (a,b)=_{S^1}\frac{da}{dx}b๐x.$$
(4.5)
This term is essential to generate partial differential equations. It is responsible for the appearance of the partial derivative in the expansion of the Hamiltonian vector fields
$$\dot{S}=\left(\frac{\delta F}{\delta S}\right)_x+[S+\lambda A,\frac{\delta F}{\delta S}].$$
(4.6)
###### Exercise 4.1
Recall that a twoโcocycle on $`๐ค`$ is a bilinear skewsymmetric map $`\omega :๐ค\times ๐ค`$ which verifies the cyclic condition
$$\omega (a,[b,c])+\omega (b,[c,a])+\omega (c,[a,b])=0.$$
Using this identity and the periodic boundary conditions check that equation (4.6) defines a Poisson bivector. $`\mathrm{}`$
### 4.2 Poisson reduction
We apply the same reduction technique used in the previous lecture, avoiding to give all the details of the computations. They can be either worked out by exercise or found in
The first Poisson bivector $`P_0`$ is defined by
$$\dot{S}=[A,\frac{\delta F}{\delta S}],$$
(4.7)
where $`A`$ is still defined by equation (3.3). These Hamiltonian vector fields obey the only constraint $`\dot{r}=0`$. Therefore the submanifold $`๐ฎ`$ formed by the matrices
$$S=\left(\begin{array}{cc}p& 1\\ q& p\end{array}\right)$$
(4.8)
is a symplectic leaf of $`P_0`$. The annihilator $`(T๐ฎ)^0`$ is spanned by the differentials of the functionals $`F:M`$ depending only on the coordinate function $`r`$. Consequently, the distribution $`D`$ is spanned by the vector fields
$$\begin{array}{cc}& \dot{p}=\frac{\delta f}{\delta r}\hfill \\ & \dot{q}=\left(\frac{\delta f}{\delta r}\right)_x2p\frac{\delta f}{\delta r}\hfill \\ & \dot{r}=0\hfill \end{array}$$
(4.9)
The distribution $`D`$ is thus tangent to $`๐ฎ`$ and $`E`$ coincides with $`D`$. The vector field (4.9) verifies the constraint
$$\dot{q}+2p\dot{p}+\dot{p_x}=(q+p^2+p_x)^{}=0.$$
(4.10)
It follows that the leaves of the distribution $`E`$ are the level sets of the function
$$u=q+p^2+p_x.$$
(4.11)
Therefore the quotient space $`N`$ is the space of scalar functions $`u:S^1`$, and (4.11) is the canonical projection $`\pi :๐ฎ๐ฎ/E`$. We see that the manifold $`N`$ is (isomorphic to) the phase space of the KdV equation.
We use the projection (4.11) to compute the reduced Poisson bivectors. The scheme of the computation is always the same. First we prolong any functional $`=_{S^1}f(u,u_x,\mathrm{})๐x`$ on $`N`$ into the functional
$$F(p,q,r)=_{S^1}f(q+p^2+p_x,q_x+2pp_x+p_{xx};\mathrm{})๐x$$
(4.12)
on $`๐ฎ`$. Then we compute its differential at the points of $`๐ฎ`$,
$$\frac{\delta F}{\delta S}=\left(\begin{array}{cc}\frac{1}{2}\left(\frac{\delta F}{\delta u}\right)_x+p\frac{\delta F}{\delta u}& \frac{\delta F}{\delta u}\\ 0& \frac{1}{2}\left(\frac{\delta F}{\delta u}\right)_xp\frac{\delta F}{\delta u}\end{array}\right).$$
(4.13)
Finally, we evaluate the reduced Hamiltonian vector fields on $`N`$ according to the usual scheme:
$$\begin{array}{cc}\hfill \dot{u}& \stackrel{\left(4.11\right)}{=}\dot{q}+\dot{p_x}+2p\dot{p}\hfill \\ & \stackrel{\left(4.6\right)}{=}\left[\left(\frac{\delta f}{\delta u}\right)_x+(q+\lambda )\frac{\delta f}{\delta p}2p\frac{\delta f}{\delta r}\right]+\left[\frac{1}{2}\left(\frac{\delta f}{\delta p}\right)_x+\frac{\delta f}{\delta r}+(q+\lambda )\frac{\delta f}{\delta q}\right]_x\hfill \\ & +2p\left[\frac{1}{2}\left(\frac{\delta f}{\delta p}\right)_x+\frac{\delta f}{\delta r}+(q+\lambda )\frac{\delta f}{\delta q}\right]\hfill \\ & \stackrel{\left(4.13\right)}{=}(q+\lambda )\left[\left(\frac{\delta f}{\delta u}\right)_x+2p\frac{\delta f}{\delta u}\right]+\left[\frac{1}{2}\left(\frac{\delta f}{\delta u}\right)_{xx}+\left(p\frac{\delta f}{\delta u}\right)_x+(q+\lambda )\frac{\delta f}{\delta u}\right]_x\hfill \\ & 2p\left[\frac{1}{2}\left(\frac{\delta f}{\delta u}\right)_{xx}+\left(p\frac{\delta f}{\delta u}\right)_x+(q+\lambda )\frac{\delta f}{\delta u}\right]\hfill \\ & =\frac{1}{2}\left(\frac{\delta f}{\delta u}\right)_{xxx}+2(q+p_x+p^2+\lambda )\left(\frac{\delta f}{\delta u}\right)_x+(q_x+p_{xx}+2pp_x)\frac{\delta f}{\delta u}\hfill \\ & \stackrel{\left(4.11\right)}{=}\frac{1}{2}\left(\frac{\delta f}{\delta u}\right)_{xxx}+2(u+\lambda )\left(\frac{\delta f}{\delta u}\right)_x+u_x\frac{\delta f}{\delta u}.\hfill \end{array}$$
(4.14)
We obtain the Poisson pencil of the KdV equation. This pencil is therefore the reduction of the โcanonicalโ pencil (4.4) over a loop algebra.
### 4.3 The GZ hierarchy
The simplest way for computing the Casimir function of the above pencil is to use the Miura map. Since this map relates the pencil to the simple bivector of the mKdV equation, it is sufficient to compute the Casimir of the latter bivector, and to transform it back to the phase space of the KdV equation.
We notice that the Casimir function of the mKdV hierarchy (1.39) is given by
$$๐ง(h)=2z_{S^1}h๐x,$$
(4.15)
where the constant $`z`$ has been inserted for future convenience.
To obtain the Casimir function of the KdV equation, we must โinvertโ the Miura map by expressing $`h`$ as a function of $`u`$. To do that we exploit the dependence of the Miura map on the parameter $`\lambda =z^2`$ of the pencil. We know that in the finite-dimensional case the Casimir function can be found as a polynomial in $`\lambda `$. In the infinite-dimensional case, we expect the Casimir function to be represented by a series. It is then natural to look at $`h`$ in the form of a Laurent series in $`z`$,
$$h(z)=z+\underset{l1}{}h_lz^l,$$
(4.16)
whose coefficients $`h_l`$ are scalar-valued periodic functions of $`x`$. In this way we change our point of view on the Miura map. Henceforth it must be looked at as a relation between a scalar function $`u`$ and a Laurent series $`h(z)`$. This change of perspective deeply influences all the mKdV theory. It is a possible starting point for the Sato picture of the KdV theory, as we shall show later.
By inserting the expansion (4.16) into the Miura map $`h_x+h^2=u+z^2`$ and equating the coefficients of different powers of $`z`$, we easily compute recursively the coefficients $`h_l`$ as differential polynomial of the function $`u`$. The first ones are
$$\begin{array}{c}h_1=\frac{1}{2}u\hfill \\ h_2=\frac{1}{4}u_x\hfill \\ h_3=\frac{1}{8}(u_{xx}u^2)\hfill \\ h_4=\frac{1}{16}(u_{xxx}4uu_x)\hfill \\ h_5=\frac{1}{32}(u_{xxxx}6uu_{xx}5u_x^2+2u^3).\hfill \end{array}$$
(4.17)
One can notice (see ) that all the even coefficients $`h_{2l}`$ are total $`x`$โderivatives. This remark explains the โstrangeโ enumeration with odd times used for the KdV hierarchy in the first lecture.
To compute concretely the GZ vector fields, besides the Casimir function
$$๐ง(u,z)=2z\underset{l1}{}_{S^1}h_lz^l๐x,$$
(4.18)
we need its differential. To simplify the notation we set
$$\alpha :=\frac{\delta ๐ง}{\delta u}=1+\underset{l1}{}\alpha _lz^l.$$
(4.19)
Once again, the simplest way for evaluating this series is to use the Miura map. We notice that $`\beta =2z`$ is the differential of the Casimir of the mKdV equation. From the transformation law of 1-forms,
$$\mathrm{\Phi }_{h}^{}{}_{}{}^{}(\alpha )=\beta ,$$
(4.20)
we then conclude that that $`\alpha `$ solves the equation
$$\alpha _x+2\alpha h=2z.$$
(4.21)
As before, the coefficients $`\alpha _l`$ can be computed recursively. One finds a Laurent series in $`\lambda =z^2`$,
$$\alpha =1\frac{1}{2}u\lambda ^1+\frac{1}{8}(3u^2u_{xx})\lambda ^2+\mathrm{},$$
(4.22)
whose first coefficients have already appeared in (1.20). From $`\alpha `$ we can easily evaluate the Lenard partial sums $`\alpha ^{(j)}=\left(\lambda ^j\alpha \right)_+`$ and write the odd GZ equations in the form
$$\frac{u}{t_{2j+1}}=\left(\frac{1}{2}_{xxx}+2(u+\lambda )_x+u_x\right)\left(\alpha ^{(j)}\right).$$
(4.23)
The even ones are
$$\frac{u}{t_{2j}}=0.$$
(4.24)
The above equations completely and tersely define the KdV hierarchy from the standpoint of the method of Poisson pairs.
### 4.4 The Central System
We shall now pursue a little further the farโreaching consequences of the change of point of view introduced in the previous subsection. According to this new point of view, the mKdV hierarchy is defined in the space $``$ of the Laurent series in $`z`$ truncated form above. This affects the whole picture.
Let us consider again the basic formulas of the mKdV theory. They are the Miura map,
$$h_x+h^2=u+z^2,$$
(4.25)
the formula for the currents (1.40),
$$H^{(2j+1)}=\frac{1}{2}\alpha _x^{(j)}+\alpha ^{(j)}h,H^{(2j)}=0,$$
(4.26)
and the definition of the mKdV hierarchy
$$\frac{h}{t_j}=_xH^{(j)}.$$
(4.27)
They were obtained in the first lecture. Presently they are complemented by the information that $`h(z)`$ is a Laurent series of the form (4.16). We shall now investigate the meaning of the above formulas in this new setting.
We start form the series $`h(z)`$, and we associate with it a new family of Laurent series $`h^{(j)}(z)`$ defined recursively by
$$h^{(j+1)}=h_x^{(j)}+hh^{(j)},$$
(4.28)
starting from $`h^{(0)}=1`$. They form a moving frame associated with the point $`h`$ in the space of (truncated) Laurent series. The first three elements of this frame are explicitly given by
$$h^{(0)}=1,h^{(1)}=h,h^{(2)}=h_x+h^2.$$
(4.29)
We see the basic block $`h_x+h^2`$ of the Miura transformation appearing. We call $`_+`$ the linear span of the series $`\{h^{(j)}\}_{j0}`$. It is a linear subspace of $``$, attached to the point $`h`$. We can now interpret the three basic formulas of the mKdV theory as properties of this linear space:
* The Miura map (4.25) tells us that the linear space $`_+`$ is invariant with respect to the multiplication by $`\lambda `$,
$$\lambda (_+)_+.$$
(4.30)
* The formula (4.26) for the currents then entails that the currents $`H^{(j)}`$, for $`j`$, belong to $`_+`$:
$$H^{(j)}_+.$$
(4.31)
* Furthermore, in conjunction with equation (4.21), it entails that the asymptotic expansion of the currents $`H^{(j)}`$ has the form
$$H^{(j)}=z^j+\underset{l1}{}H_l^iz^l=z^j+O(z^1).$$
(4.32)
* Finally, the mKdV equations (4.27) can be seen as the commutativity conditions of the operators $`(_x+h)`$ and $`\left({\displaystyle \frac{}{t_j}}+H^{(j)}\right)`$:
$$[_x+h,\frac{}{t_j}+H^{(j)}]=0.$$
(4.33)
Used together, conditions (4.31) and (4.33) imply that the operators $`\left({\displaystyle \frac{}{t_j}}+H^{(j)}\right)`$ leave the linear space $`_+`$ invariant:
$$\left(\frac{}{t_j}+H^{(j)}\right)(_+)_+.$$
(4.34)
This is the abstract but simple form of the laws governing the time evolution of the currents $`H^{(j)}`$. These equations are the โtopโ of the KdV theory, and form the basis of the Sato theory. It is not difficult to give them a concrete form. By using the form of the expansion (4.32) it is easy to show that equations (4.34) are equivalent to the infinite system of Riccatiโtype equations on the currents $`H^{(j)}`$:
$$\frac{H^{(j)}}{t_k}+H^{(j)}H^{(k)}=H^{(j+k)}+\underset{l=1}{\overset{j}{}}H_l^kH^{(jl)}+\underset{l=1}{\overset{k}{}}H_l^jH^{(kl)}.$$
(4.35)
It will be called the Central System (CS).
###### Exercise 4.2
Prove formulas (4.31) and (4.32).
### 4.5 The linearization process
The first reward of the previous work is the discovery of a linearization process. The equations (4.35) of the Central System are not directly linearizable, but they can be easily transformed into a new system of linearizable Riccati equations by a transformation in the space of currents. This idea is realized once again by a โMiura mapโ. The novelty, however, is that this map is now operating on the space of currents rather than on the phase space of the KdV equation.
We simply give the final result. Let us consider a new family of currents $`\{W^{(k)}\}_{k0}`$ of the form
$$W^{(k)}=z^k+\underset{l1}{}W_l^kz^l,$$
(4.36)
and let us denote by $`๐ฒ_+`$ their linear span in $``$. We define (see also ) a new system of equations on the currents $`W^{(k)}`$ by imposing the โconstraintsโ
$$\left(\frac{}{t_k}+z^k\right)(๐ฒ_+)๐ฒ_+$$
(4.37)
on their linear span $`๐ฒ_+`$. It is easily seen that they take the explicit form
$$\frac{W^{(k)}}{t_j}+z^jW^{(k)}=W^{(j+k)}+\underset{l=1}{\overset{j}{}}W_l^kW^{(jl)}.$$
(4.38)
They will be called the Sato equations (on the โbig cell of the Sato Grassmannianโ). They are a system of linearizable Riccati equations. This can be seen either from the geometry of a suitable group action on the Grassmannian or by means of the following more elementary considerations. We write equations (4.38) in the matrix form
$$\frac{๐ถ}{t_j}+๐ถ{}_{}{}^{T}\mathrm{\Lambda }_{}^{k}\mathrm{\Lambda }^k๐ถ=๐ถ\mathrm{\Gamma }_k๐ถ,$$
(4.39)
where $`๐ถ=\left(W_l^k\right)`$ is the matrix of the components of the currents $`W^{(k)}`$, $`\mathrm{\Lambda }`$ is the infinite shift matrix
$$\mathrm{\Lambda }=\left[\begin{array}{ccccc}0& 1& 0& \mathrm{}& \\ 0& 0& 1& 0& \mathrm{}\\ \mathrm{}& & \mathrm{}& \mathrm{}& \\ \mathrm{}& & & \mathrm{}& \mathrm{}\\ \mathrm{}& & & & \mathrm{}\end{array}\right],$$
(4.40)
and $`\mathrm{\Gamma }^k`$ is the convolution matrix of level $`k`$,
$$\mathrm{\Gamma }_k=\left[\begin{array}{cccccc}0& \mathrm{}& & 1& 0& \mathrm{}\\ \mathrm{}& & 1& 0& \mathrm{}& \mathrm{}\\ & & & & & \\ 1& 0& & & & \\ \mathrm{}& & & & & \end{array}\right].$$
(4.41)
One can thus check that the matrix Riccati equation (4.39) is solved by the matrix
$$๐ถ=๐ต๐ด^1,$$
(4.42)
where $`๐ด`$ and $`๐ต`$ satisfy the constant coefficients linear system
$$\frac{}{t_k}๐ด={}_{}{}^{T}\mathrm{\Lambda }_{}^{k}๐ด\mathrm{\Gamma }_k๐ต,\frac{}{t_k}๐ต=\mathrm{\Lambda }^k๐ต.$$
(4.43)
The closing remark is that the Sato equations are mapped into the Central System (4.35) by the following algebraic Miura map:
$$H^{(j)}=\frac{_{l=0}^jW_{jl}^0W^{(l)}}{W^{(0)}}.$$
(4.44)
The outcome of this long chain of extensions and transformations is the following algorithm for solving the KdV equation:
First we solve the linear system (4.43), with a suitably chosen initial condition, which we do not discuss here;
Then we use the projective transformation (4.42) and the Miura map (4.44) to recover the currents $`H^{(j)}`$;
Finally, we extract the first current $`H^{(1)}=h`$, and we evaluate the first component $`h_1`$ of its Laurent expansion in powers of $`z^1`$.
The function
$$u(x,t_3,\mathrm{})=2h_1|_{t_1=x}$$
(4.45)
is then a solution of the KdV equation.
### 4.6 The relation with the Sato approach
The equations (4.27) make sense for an arbitrary Laurent series $`h`$ of the form (4.16), even if it is not a solution of the Riccati equation $`h_x+h^2=u+z^2`$. Hence they define, for every $`j`$, a system of PDEs for the coefficients $`h_l`$. We will show<sup>6</sup><sup>6</sup>6See also the papers . that these systems are equivalent to the celebrated KP hierarchy of the Kyoto school (see the lectures by Satsuma in these volume). The usual definition of the KP equations can be summarized as follows. Let $`\mathrm{\Psi }๐`$ be the ring of pseudodifferential operators on the circle. It contains as a subring the space $`๐`$ of purely differential operators. Let us denote with $`()_+`$ the natural projection from $`\mathrm{\Psi }๐`$ onto $`๐`$. Let $`Q`$ be a monic operators of degree $`1`$,
$$Q=\underset{j1}{}q_j^j.$$
(4.46)
The KP hierarchy is the set of Lax equations for $`Q`$
$$\frac{}{t_j}Q=[\left(Q^j\right)_+,Q].$$
(4.47)
The aim of this subsection is to show that such a Lax representation just arises as a kind of a Euler form of the equations (4.27). Before stating the next result, we must observe that the relations (4.28) can be solved backwards, in such a way to define the Faร di Bruno elements $`h^{(j)}`$ for all $`j`$.
###### Proposition 4.3
Suppose the series $`h`$ of the form (4.16) to evolve according to a conservation law,
$$\frac{h}{t}=_xH,$$
(4.48)
for an arbitrary $`H`$. Then the Faร di Bruno elements $`h^{(j)}`$, for $`j`$, evolve according to
$$\left(\frac{}{t}+H\right)h^{(j)}=\underset{k=0}{\overset{\mathrm{}}{}}\left(\genfrac{}{}{0pt}{}{j}{k}\right)(_x^kH)h^{(jk)},$$
(4.49)
where
$$\left(\genfrac{}{}{0pt}{}{j}{k}\right)=\frac{j(j1)\mathrm{}(jk+1)}{k!},\left(\genfrac{}{}{0pt}{}{j}{0}\right)=1.$$
Now we consider the map $`\varphi :\mathrm{\Psi }๐`$, from the space of Laurent series to the ring of pseudodifferential operators on the circle, which acts on the Faร di Bruno basis according to
$$\varphi (h^{(j)})=^j.$$
(4.50)
This map is then extended by linearity (with respect to multiplication by a function of $`x`$) to the whole space $``$.
###### Definition 4.4
We call Lax operator of the KP theory the image
$$Q=\varphi (z)$$
(4.51)
of the first element of the standard basis in $``$.
If the $`q_j`$ are the components of the expansion of $`z`$ on the Faร di Bruno basis,
$$z=h^{(1)}\underset{j1}{}q_jh^{(j)},$$
(4.52)
then we can write
$$Q=\underset{j1}{}q_j^j$$
(4.53)
according to the definition of the map $`\varphi `$. We note that equation (4.52) uniquely defines the coefficients $`q_j`$ as differential polynomials of the components $`h_j`$ of $`h(z)`$:
$$\begin{array}{c}q_1=h_1,q_2=h_2,q_3=h_3+h_1^2\hfill \\ q_4=h_4+3h_1h_2h_1h_{1x}\hfill \\ \mathrm{}\mathrm{}\hfill \end{array}$$
(4.54)
This is an invertible relation between the $`h_j`$ and the $`q_j`$, so that equation (4.52) may be seen as a change of coordinates in the space $``$.
###### Proposition 4.5
The map $`\varphi `$ has the following three properties:
Multiplying a vector of the Faร di Bruno basis by a power $`z^k`$ of $`z`$ yields
$$\varphi (z^kh^{(j)})=^jQ^k.$$
(4.55)
The evolution along a conservation law of the form
$$\frac{h}{t}=_x\left(\underset{k}{}H_kz^k\right)$$
translates into
$$\frac{}{t}\left(\varphi (h^{(j)})\right)=\underset{k}{}[^j,H_k]Q^k.$$
(4.56)
If $`\pi _+`$ and $`\mathrm{\Pi }_+`$ are respectively the projection onto the positive part $`_+`$ and $`๐\mathrm{\Psi }๐`$, then
$$\varphi \pi _+=\mathrm{\Pi }_+\varphi .$$
(4.57)
To obtain the Sato form of the equations (4.27), we derive the equation
$$z=h^{(1)}\underset{l1}{}q_lh^{(l)}$$
(4.58)
with respect to the time $`t_j`$, getting
$$\underset{l1}{}\frac{q_l}{t_j}h^{(l)}=\frac{h^{(1)}}{t_j}\underset{l1}{}q_l\frac{h^{(l)}}{t_j}.$$
(4.59)
Applying the map $`\varphi `$ to both sides of this equation we obtain
$$\underset{l1}{}\frac{q_l}{t_j}^l=\underset{k1}{}[,H_k^j]Q^k\underset{k1}{}q_l[^l,H_k^j]Q^k,$$
(4.60)
or
$$\frac{Q}{t_j}+\underset{k1}{}[Q,H_k^j]Q^k=0.$$
(4.61)
Finally, we introduce the operator
$$B^{(j)}=\varphi (H^{(j)})=\varphi \left(z^j+\underset{k1}{}H_k^jz^k\right)=Q^j+\underset{k1}{}H_k^jQ^k$$
(4.62)
associated with the current density $`H^{(j)}`$, and we note that
$$B^{(j)}=\varphi (\pi _+(z^j)))=(\varphi (z^j))_+=(Q^j)_+.$$
(4.63)
Thus we can write (4.61) in the final form
$$\frac{Q}{t_j}+[Q,\left(Q^j\right)_+]=0,$$
(4.64)
which coincides with equation (4.47).
## 5 Lax representation of the reduced KdV flows
In this lecture we want to investigate more accurately the properties of the stationary KdV flows, that is, of the equations induced by the KdV hierarchy on the finiteโdimensional invariant submanifolds of the singular points of any equation of the hierarchy. Examples of these reductions have already been discussed in the first lecture. In the third lecture we realized, in a couple of examples, that the reduced flows were still bi-Hamiltonian. Although not at all surprising, this property is somewhat mysterious, since it is not yet well understood how the Poisson pairs of the reductions are related to the original Poisson pairs of the KdV equation. Moreover, even if the subject is quite old and classical (see, e.g., ), it was still lacking in the literature a systematic and coordinate free proof that such reduced flows are bi-Hamiltonian (see, however, ). In this lecture we will not provide such a proof, which is contained in , but we will give a sufficiently systematic algorithm to compute the reduced Poisson pair. This algorithm is based on the study of the Lax representation of the reduced equations.
### 5.1 Lax representation
In this section we associate a Lax matrix (polynomially depending on $`\lambda `$) with each element $`H^{(j)}`$. This matrix naturally arises from a change of basis in the linear space $`_+`$ attached to the point $`h`$. So far we have introduced two bases:
The moving frame $`\{h^{(j)}\}`$;
The canonical basis $`\{H^{(j)}\}`$.
Presently we introduce a third basis by exploiting the constraint
$$\lambda (_+)_+,$$
(5.1)
characteristic of the KdV theory. The new basis is formed by the multiples $`\{\lambda ^jH^{(0)},\lambda ^jH^{(1)}\}`$ of the first two currents. Formally we define
$$\text{ iii) the Lax basis: }(\lambda ^j,\lambda ^jh).$$
The use of this basis leads to a new representation of the currents $`H^{(j)}`$, where each current is written as a linear combination of the first two, $`H^{(0)}=1`$ and $`H^{(1)}=h`$, with coefficients that are polynomials in $`\lambda `$. Let us consider a few examples:
$$\begin{array}{cc}\hfill H^{(0)}& =1+0h\hfill \\ \hfill H^{(1)}& =01+1h\hfill \\ \hfill H^{(2)}& =\lambda 1+0h\hfill \\ \hfill H^{(3)}& =h_21+(\lambda h_1)h\hfill \\ \hfill H^{(4)}& =\lambda ^21\hfill \\ \hfill H^{(5)}& =(\lambda h_2+h_1h_2h_4)1+(\lambda ^2\lambda h_1+h_1^2h_3)h.\hfill \end{array}$$
(5.2)
This new representation also affects our way of writing the action of the operators $`\left({\displaystyle \frac{}{t_j}}+H^{(j)}\right)`$. Let these operators act on $`H^{(0)}`$ and $`H^{(1)}`$. For the basic invariance condition (4.34), we get an element in $`_+`$ which can be represented on the Lax basis. As a result we can write
$$\left(\frac{}{t_j}+H^{(j)}\right)\left[\begin{array}{c}1\\ h\end{array}\right]=L^{(j)}(\lambda )\left[\begin{array}{c}1\\ h\end{array}\right],$$
(5.3)
where $`L^{(j)}(\lambda )`$ is the Lax matrix associated with the current $`H^{(j)}`$. We shall see below the explicit form of some of these matrices.
It becomes now very easy to rewrite the Central System in the form of equations on the Lax matrices $`L^{(j)}(\lambda )`$. We simply have to notice that the equations (4.35) entail the โexactness conditionโ
$$\frac{H^{(j)}}{t_k}=\frac{H^{(k)}}{t_j},$$
(5.4)
from which it follows that the operators $`\left({\displaystyle \frac{}{t_j}}+H^{(j)}\right)`$ and $`\left({\displaystyle \frac{}{t_k}}+H^{(k)}\right)`$ commute:
$$[\frac{}{t_j}+H^{(j)},\frac{}{t_k}+H^{(k)}]=0.$$
(5.5)
It is now sufficient to evaluate this condition on $`(H^{(0)},H^{(1)})`$ and to expand on the Lax basis to find the โzero curvature representationโ of the KdV hierarchy:
$$\frac{L^{(j)}}{t_k}\frac{L^{(k)}}{t_j}+[L^{(j)},L^{(k)}]=0.$$
(5.6)
Suppose now that we are on the invariant submanifold formed by the singular points of the $`j`$โth member of the KdV hierarchy. On this submanifold
$$\frac{L^{(k)}}{t_j}=0k,$$
(5.7)
and the zero curvature representation becomes the Lax representation
$$\frac{L^{(j)}}{t_k}=[L^{(k)},L^{(j)}].$$
(5.8)
We have thus shown that all the stationary reductions of the KdV hierarchy admit a Lax representation. As a matter of fact, this Lax representation coincides with the Lax representation of the GZ systems on LieโPoisson manifolds studied in Section 3. The latter are bi-Hamiltonian systems. Therefore, we end up stating that the stationary reductions of the KdV theory are bi-Hamiltonian, and we can construct the associated Poisson pairs. We shall now see a couple of examples.
### 5.2 First example
We study anew the simplest invariant submanifold of the KdV hierarchy, defined by the equation
$$u_{xxx}6uu_x=0.$$
(5.9)
In this example we consider the constraint from the point of view of the Central System. Since the constraint is the stationarity of the time $`t_3`$, we have to consider only the first three Lax matrices. As for the matrix $`L^{(1)}`$, the following computation,
$$\begin{array}{cc}\hfill \left(\frac{}{t_1}+H^{(1)}\right)& 1=01+1h\hfill \\ \hfill \left(\frac{}{t_1}+H^{(1)}\right)& H^{(1)}\stackrel{\left(4.35\right)}{=}H^{(2)}+2h_1\stackrel{\left(5.2\right)}{=}(\lambda +2h_1)1+0h,\hfill \end{array}$$
(5.10)
shows that
$$L^{(1)}=\left(\begin{array}{cc}0& 1\\ \lambda +2h_1& 0\end{array}\right).$$
(5.11)
Similarly, the computation
$$\begin{array}{cc}\hfill \left(\frac{}{t_3}+H^{(3)}\right)1& \stackrel{\left(5.2\right)}{=}h_21+(\lambda h_1)h\hfill \\ \hfill \left(\frac{}{t_3}+H^{(3)}\right)H^{(1)}& \stackrel{\left(4.35\right)}{=}H^{(4)}+h_1H^{(2)}+h_2H^{(1)}+h_3+H_1^3\hfill \\ & \stackrel{\left(5.2\right)}{=}(\lambda ^2+\lambda h_1+2h_3h_1^2)1+h_2h\hfill \end{array}$$
(5.12)
yields
$$L^{(3)}=\left(\begin{array}{cc}h_2& \lambda h_1\\ \lambda ^2+h_1\lambda +2h_3h_{1}^{}{}_{}{}^{2}& h_2\end{array}\right).$$
(5.13)
On the submanifold $`M_3`$ defined by equation (5.9) this matrix verifies the Lax equation
$$\frac{L^{(3)}}{t_1}=[L^{(1)},L^{(3)}].$$
(5.14)
This equation completely defines the time evolution of the first three components $`(h_1,h_2,h_3)`$ of the current $`H^{(1)}=h`$. These components play the role of coordinates on $`M_3`$. We get
$$\begin{array}{cc}\hfill \frac{h_1}{t_1}& =2h_2\hfill \\ \hfill \frac{h_2}{t_1}& =2h_3h_{1}^{}{}_{}{}^{2}\hfill \\ \hfill \frac{h_3}{t_1}& =4h_1h_2\hfill \end{array}$$
(5.15)
By the change of coordinates
$$h_1=\frac{1}{2}u,h_2=\frac{1}{4}u_x,h_3=\frac{1}{8}(u_{xx}u^2),$$
coming from the inversion (4.15) of the Miura map, these equations take the form
$$\frac{u}{t_1}=u_x,\frac{u_x}{t_1}=u_{xx},\frac{u_{xx}}{t_1}=6uu_x,$$
(5.16)
already encountered in Lecture 1. This shows explicitly the connection between the two points of view.
To find the connection between these equations and the GZ equations dealt with in the first example of Lecture 3, we compare the Lax matrix
$$L^{(3)}(\lambda )=\lambda ^2\left(\begin{array}{cc}0& 0\\ 1& 0\end{array}\right)+\lambda \left(\begin{array}{cc}0& 1\\ h_1& 0\end{array}\right)+\left(\begin{array}{cc}h_2& h_1\\ 2h_3h_1^2& h_2\end{array}\right)$$
with the Lax matrix
$$S(\lambda )=\lambda ^2\left(\begin{array}{cc}0& 0\\ 1& 0\end{array}\right)+\lambda \left(\begin{array}{cc}p_1& 1\\ q_1& p_1\end{array}\right)+\left(\begin{array}{cc}p_0& (q_1+p_1^2)\\ q_0& p_0\end{array}\right)$$
associated with the points of the symplectic leaf defined by (3.17). We easily identify $`L^{(3)}`$ with the restriction of $`S(\lambda )`$ to $`p_1=0`$ upon setting
$$p_0=h_2,q_1=h_1,q_0=2h_3h_1^2.$$
(5.17)
By comparing these equations with the projection (3.20), which allows to pass from the symplectic leaf $`S`$ to the quotient space $`N=S/E`$, we obtain the change of coordinates
$$u_1=h_1,u_2=h_2,u_3=2h_3h_1^2,$$
(5.18)
connecting the reduction (5.15) of the Central System to the GZ system (3.28) dealt with in the third Lecture. The latter was, by construction, a bi-Hamiltonian system. We argue that also the reduction of the Central System herewith considered is a bi-Hamiltonian vector field, and that its Poisson pair is obtained by geometric reduction. Basic for this identification is the property of the Lax matrix $`L^{(3)}`$ of being a section of the fiber bundle $`\pi :SS/E`$ appearing in the geometric reduction. It is this property which allows to set an invertible relation among the coordinates $`(u_1,u_2,u_3)`$, coming from the geometric reduction, and the coordinates $`(h_1,h_2,h_3)`$ coming from the reduction of the Central System.
### 5.3 The generic stationary submanifold
It is now not hard to give the general form of the matrices $`L^{(j)}`$ for an arbitrary odd integer $`2j+1`$. First we observe that
$$\left(\frac{}{t_{2j+1}}+H^{(2j+1)}\right)1=H^{(2j+1)}\stackrel{\left(\text{4.26}\right)}{=}\frac{1}{2}\alpha _x^{(j)}+\alpha ^{(j)}h.$$
(5.19)
Then we notice that
$$\begin{array}{cc}& \left(\frac{}{t_j}+H^{(j)}\right)h=H^{(j+1)}+\underset{l=1}{\overset{j}{}}h_lH^{(jl)}+H_1^j\hfill \\ & =\frac{1}{2}\left(\alpha _x^{(j+1)}+\underset{l=1}{\overset{j}{}}h_l\alpha _x^{(j1)}\right)+H_1^j+\left(\alpha ^{(j+1)}+\underset{l=1}{\overset{j}{}}h_l\alpha ^{(j1)}\right)h.\hfill \end{array}$$
(5.20)
Therefore
$$L^{(j)}=\left(\begin{array}{cc}\frac{1}{2}\alpha _x^{(j)}& \alpha ^{(j)}\\ \frac{1}{2}(\alpha _x^{(j+1)}+_{l=1}^jh_l\alpha _x^{(j1)})+H_1^j& \alpha ^{(j+1)}+_{l=1}^jh_l\alpha ^{(j1)}\end{array}\right).$$
(5.21)
By using the definition (4.21) of the Lenard series $`\alpha (z)`$ of which the polynomials $`\alpha ^{(j)}`$ are the partial sums, it is easy to prove that $`L^{(j)}`$ is a traceless matrix.
We leave to the reader to specialize the matrix $`L^{(5)}`$, and to write explicitly the Lax equations
$$\frac{L^{(5)}}{t_1}=[L^{(1)},L^{(5)}],\frac{L^{(5)}}{t_3}=[L^{(3)},L^{(5)}].$$
(5.22)
They should be compared with the reduced KdV equations (1.31) and (1.32) on the invariant submanifold defined by the constraint
$$u_{xxxxx}10uu_{xxx}20u_xu_{xx}+30u^2u_x=0.$$
(5.23)
They should also be compared with the GZ equations obtained via the geometric reduction process applied to the LieโPoisson pairs defined on three copies of $`๐ฐ๐ฉ(2)`$ by equations (3.31) and (3.32). We have not displayed explicitly these equations yet. We will give their form in the next lecture.
### 5.4 What more?
There is nothing โsacredโ with the KdV theory. As we know, it is related with the constraint
$$z^2(_+)(_+),$$
(5.24)
which defines an invariant submanifold of the Central System. Many other constraints can be considered. For instance, the constraint
$$z^3(_+)(_+)$$
(5.25)
leads to the soโcalled Boussinesq theory, and is studied in . What is remarkable is that the change of constraint does not affect the algorithm for the study of the reduced equations. All the previous reasonings are valid without almost no change. The only difference resides in the fact that the computations become more involved. This remark allows to better appreciate the meaning of the process leading from the KdV equation to the Central System. We have not only given a new formulation to known equations. We have actually found a much bigger hierarchy, possessing remarkable properties, which coincides with the KdV hierarchy on a (small) proper invariant subset. The integrability properties belong to the bigger hierarchy, and hold outside the KdV submanifold. Many other interesting equations can be found by other processes of reduction. There is some evidence that a very large class of evolution equations possessing some integrability properties can be eventually recovered as a suitable reduction of the Central System, or of strictly related systems. However, we shall not pursue this point of view further, since it would lead us too far away from our next topic, the separability of the reduced KdV flows.
## 6 DarbouxโNijenhuis coordinates and Separability
In this lecture we shall consider the reduced KdV flows from a different point of view. Our aim is to probe the study of the geometry of the Poisson pair which, as realized in the third and fifth lectures, is associated with these flows. The final goal is to show the existence of a suitable set of coordinates defined by and adapted to the Poisson pair. They are called DarbouxโNijenhuis coordinates. We shall prove that they are separation coordinates for the HamiltonโJacobi equations associated with the reduced flows.
To keep the presentation within a reasonable size, we shall mainly deal with a particular example, and we shall not discuss thoroughly the theoretical background, referring to for more details. We shall use the example to display the characteristic features of the geometry of the reduced manifolds. The reader is asked to believe that all that will be shown is general inside the class of the reduced stationary KdV manifolds, whose Poisson pencils are of maximal rank. A certain care must be used in trying to extend these conclusions to other examples like the Boussinesq stationary reductions, whose Poisson pencils are not of maximal rank. They will not be covered in these lecture notes. The example worked out is the reduction of the first and the third KdV equations on the invariant submanifold defined by the equation
$$u_{xxxxx}10uu_{xxx}20u_xu_{xx}+30u^2u_x=0,$$
(6.1)
a problem addressed at the end of Section 5.3.
### 6.1 The Poisson pair
As we mentioned several times, the invariant submanifold $`M_5`$ defined by equation (6.1) has dimension five. From the standpoint of the Central System, it is characterized by the two equations
$$z^2(_+)(_+),H^{(5)}h=\lambda ^3+\underset{l=1}{\overset{5}{}}h_lH^{(5l)}+H_1^5.$$
(6.2)
We recall that the first constraint means that, inside the big cell of the Sato Grassmannian, we are working on the special submanifold corresponding to the KdV theory. The second constraint means that, inside the phase space of the KdV theory, we are working on the set (6.1) of singular points of the fifth flow. The two constraints play the following roles. The first constraint sets up a relation among the currents $`H^{(j)}`$: All the currents are expressed as linear combinations (with polynomial coefficients) of the first two currents $`H^{(0)}=1`$ and $`H^{(1)}=h`$. So this constraint drastically reduces the number of the unknowns $`H_l^j`$ to the coefficients $`h_l`$ of $`h`$. The second constraint then further cuts the degrees of freedom to a finite number, by setting relations among the coefficients $`h_l`$. It can be shown that only the first five coefficients $`(h_1,h_2,h_3,h_4,h_5)`$ survive as free parameters. All the other coefficients can be expressed as polynomial functions of the previous ones. By a process of elimination of the exceeding coordinate, one proves that the restriction of the first and third flows of the KdV hierarchy are represented by the following differential equations:
$$\begin{array}{cc}& \frac{h_1}{t_1}=2h_2\hfill \\ & \frac{h_2}{t_1}=2h_3h_{1}^{}{}_{}{}^{2}\hfill \\ & \frac{h_3}{t_1}=2h_1h_22h_4\hfill \\ & \frac{h_4}{t_1}=2h_5h_{2}^{}{}_{}{}^{2}2h_1h_3\hfill \\ & \frac{h_5}{t_1}=4h_3h_2+2h_{1}^{}{}_{}{}^{2}h_24h_1h_4\hfill \end{array}$$
(6.3)
and
$$\begin{array}{cc}& \frac{h_1}{t_3}=2h_4+2h_1h_2\hfill \\ & \frac{h_2}{t_3}=2h_5+h_{2}^{}{}_{}{}^{2}+h_{1}^{}{}_{}{}^{3}\hfill \\ & \frac{h_3}{t_3}=2h_1h_4+4h_{1}^{}{}_{}{}^{2}h_22h_3h_2\hfill \\ & \frac{h_4}{t_3}=2h_{3}^{}{}_{}{}^{2}2h_2h_4+2h_1h_{2}^{}{}_{}{}^{2}+h_{1}^{}{}_{}{}^{4}+h_{1}^{}{}_{}{}^{2}h_3\hfill \\ & \frac{h_5}{t_3}=2h_{1}^{}{}_{}{}^{2}h_44h_3h_4+2h_{1}^{}{}_{}{}^{3}h_2\hfill \end{array}$$
(6.4)
They can also be seen as the Lax equations (5.22). However, for our purposes, it is more important to recognize that the above equations are the GZ equations of the Poisson pencil defined on $`M_5`$. This pencil can be computed according to the reduction procedure explained in the third lecture. The final outcome is that the reduced Poisson bivector is given by
$$\begin{array}{cc}\hfill \dot{h_1}& =2\frac{H}{h_2}+2(h_1\lambda )\frac{H}{h_4}+2h_2\frac{H}{h_5}\hfill \\ \hfill \dot{h_2}& =2\frac{H}{h_1}+2(\lambda 2h_1)\frac{H}{h_3}2h_2\frac{H}{h_4}+(4\lambda h_12h_3h_1^2)\frac{H}{h_5}\hfill \\ \hfill \dot{h_3}& =2(2h_1\lambda )\frac{H}{h_2}+(2h_3+2h_1^24\lambda h_1)\frac{H}{h_4}+2(h_4+h_1h_2)\frac{H}{h_5}\hfill \\ \hfill \dot{h_4}& =2(\lambda h_1)\frac{H}{h_1}+2h_2\frac{H}{h_2}(2h_3+2h_1^24\lambda h_1)\frac{H}{h_3}\hfill \\ & +(2h_56h_1h_3+h_2^2+2h_1^3+4\lambda h_3+2\lambda h_1^2)\frac{H}{h_5}\hfill \\ \hfill \dot{h_5}& =2h_2\frac{H}{h_1}+(2h_3+h_1^24\lambda h_1)\frac{H}{h_2}2(h_4+h_1h_2)\frac{H}{h_3}\hfill \\ & (2h_56h_1h_3+h_2^2+2h_1^3+4\lambda h_3+2\lambda h_1^2)\frac{H}{h_4}.\hfill \end{array}$$
(6.5)
The Casimir function of this pencil is a quadratic polynomial,
$$C(\lambda )=C_0\lambda ^2+C_1\lambda +C_2,$$
(6.6)
and the coefficients are
$$\begin{array}{c}C_0=h_{1}^{}{}_{}{}^{3}2h_1h_3+h_5\hfill \\ C_1=h_2h_4h_1h_5+\frac{3}{2}h_{1}^{}{}_{}{}^{2}h_3\frac{1}{2}h_1h_{2}^{}{}_{}{}^{2}\frac{1}{2}h_{3}^{}{}_{}{}^{2}\frac{1}{2}h_{1}^{}{}_{}{}^{4}\hfill \\ C_2=\frac{1}{2}h_3h_{2}^{}{}_{}{}^{2}h_3h_5+\frac{1}{2}h_{1}^{}{}_{}{}^{5}+h_1h_{3}^{}{}_{}{}^{2}h_1h_2h_4\frac{3}{2}h_{1}^{}{}_{}{}^{3}h_3+h_{1}^{}{}_{}{}^{2}h_5+\frac{1}{2}h_{4}^{}{}_{}{}^{2}\hfill \end{array}$$
(6.7)
The Lenard chain is
$$\begin{array}{ccc}P_0dC_0\hfill & =& 0\hfill \\ P_0dC_1\hfill & =& P_1dC_0=\frac{๐ก}{t_1}\hfill \\ P_0dC_2\hfill & =& P_1dC_1=\frac{๐ก}{t_3}\hfill \\ & & P_1dC_2=\mathrm{\hspace{0.25em}\hspace{0.25em}0},\hfill \end{array}$$
(6.8)
where $`๐ก`$ is the vector $`(h_1,h_2,h_3,h_4,h_5)`$. It shows that the reduced flows are bi-Hamiltonian. Finally, if one uses the coordinate change (4.17) from the coordinates $`(h_1,h_2,h_3,h_4,h_5)`$ to the coordinates $`(u,u_x,u_{xx},u_{xxx},u_{xxxx})`$, one can put the equations (6.3) and (6.4) in the form (1.31) and (1.32) considered in the first lecture.
### 6.2 Passing to a symplectic leaf
We aim to solve equations (6.3) and (6.4) by the HamiltonโJacobi method. This requires to set the study of such equations on a symplectic manifold. This can be easily accomplished by noticing that these vector fields are already tangent to the submanifold $`S_4`$ defined by the equation
$$C_0=E,$$
(6.9)
for a constant $`E`$. We know that this submanifold is symplectic since $`C_0`$ is the Casimir of $`P_0`$. The dimension of $`S_4`$ is four, and the variables $`(h_1,h_2,h_3,h_4)`$ play the role of coordinates on it.
For our purposes it is crucial to remark an additional property of $`S_4`$: It is a bi-Hamiltonian manifold. This means that also the second bivector $`P_1`$ induces, by a process of reduction, a Poisson structure on $`S_4`$ compatible with the natural restriction of $`P_0`$. This is not a general situation. It holds as a consequence of a peculiarity of the Poisson pencil (6.5). The property we are mentioning concerns the vector field
$$Z=\frac{}{h_5}.$$
(6.10)
One can easily check that:
$`Z`$ is transversal to the symplectic leaf $`S_4`$.
The functions which are invariant along $`Z`$ form a Poisson subalgebra with respect to the pencil.
In simpler terms, the Poisson bracket of functions which are independent of $`h_5`$ is independent on $`h_5`$ as well. Since they coincide with the functions on $`S_4`$ (by the transversality condition), this property allows us to define a pair of Poisson brackets also on $`S_4`$. The first bracket is associated with the symplectic 2โform $`\omega _0`$ on $`S_4`$. It can be easily checked that
$$\omega _0=h_1dh_1dh_2+\frac{1}{2}(dh_2dh_4+dh_5dh_1).$$
(6.11)
The second Poisson bracket can be represented in the form
$$\{f,g\}_1=\omega _0(NX_f,X_g),$$
(6.12)
where $`X_f`$ and $`X_g`$ are the Hamiltonian vector fields associated with the functions $`f`$ and $`g`$ by the symplectic 2โform $`\omega _0`$, and $`N`$ is a $`(1,1)`$โtensor field on $`S_4`$, called the Nijenhuis tensor associated with the pencil (see, e.g., ). In our example one obtains
$$\begin{array}{cc}\hfill N=& \left(h_1\frac{}{h_1}h_2\frac{}{h_2}+(h_33h_1^2)\frac{}{h_3}2h_1h_2\frac{}{h_4}\right)dh_1\hfill \\ & +(h_3h_1^2)\frac{}{h_4}dh_2+\left(\frac{}{h_1}+2h_1\frac{}{h_3}+h_2\frac{}{h_4}\right)dh_3\hfill \\ & +\left(\frac{}{h_2}+h_1\frac{}{h_4}\right)dh_4.\hfill \end{array}$$
(6.13)
Thus we arrive at the following picture of the GZ hierarchy considered in this lecture. It is formed by a pair of vector fields, $`X_1`$ and $`X_3`$, defined by (6.3) and (6.4). They are tangent to the symplectic leaf $`(S_4,\omega _0)`$ defined by equations (6.9) and (6.11). This symplectic manifold is still bi-Hamiltonian, and therefore there exists a Nijenhuis tensor field $`N`$, defined by equation (6.12). The vector fields $`X_1`$ and $`X_3`$ span a Lagrangian subspace which is invariant with respect to $`N`$. One finds that they obey the following โmodified Lenard recursion relationsโ
$$\begin{array}{ccccc}NX_1& =& X_3& +& (\text{Tr }N)X_1\\ NX_3& =& & +& (\text{det }N)X_1.\end{array}$$
(6.14)
From them we can extract the matrix
$$๐ฅ=\left(\begin{array}{cc}\text{Tr }N& 1\\ \text{det }N& 0\end{array}\right)$$
(6.15)
which represents the action of $`N`$ on the abovementioned Lagrangian subspace. It will play a fundamental role in the upcoming discussion of the separability of the vector fields.
###### Exercise 6.1
Compute the expression of the reduced pencil on $`S_4`$ and check the form of the Nijenhuis tensor, as well as the modified Lenard recursion relations (6.14).
### 6.3 DarbouxโNijenhuis coordinates
We are now in a position to introduce the basic tool of the theory of separability in the bi-Hamiltonian framework: The concept of DarbouxโNijenhuis coordinates on a symplectic bi-Hamiltonian manifold, like $`S_4`$.
Given a symplectic 2โform $`\omega _0`$ and a Nijenhuis tensor $`N`$ coming from a Poisson pencil defined on a $`2n`$โdimensional manifold $``$, under the assumption that the eigenvalues of $`N`$ are real and functionally independent, one proves the existence of a system of coordinates $`(\lambda _1,\mathrm{},\lambda _n;\mu _1,\mathrm{},\mu _n)`$ which are canonical for $`\omega _0`$,
$$\omega _0=\underset{i=1}{\overset{n}{}}d\mu _1d\lambda _i,$$
(6.16)
and which allows to put $`N^{}`$ (the adjoint of $`N`$) in diagonal form:
$$N^{}d\lambda _i=\lambda _id\lambda _i,N^{}d\mu _i=\lambda _id\mu _i.$$
(6.17)
The coordinates $`\lambda _i`$ are the eigenvalues of $`N^{}`$, and therefore can be computed as the zeroes of the minimal polynomial of $`N`$:
$$\lambda ^n+c_1\lambda ^{n1}+\mathrm{}+c_n=0.$$
(6.18)
The coordinates $`\mu _j`$ can be computed as the values that a conjugate polynomial
$$\mu =f_1\lambda ^{n1}+\mathrm{}+f_n$$
(6.19)
assumes on the eigenvalues $`\lambda _j`$, that is,
$$\mu _j=f_1\lambda _j^{n1}+\mathrm{}+f_n,j=1,\mathrm{},n.$$
(6.20)
The determination of this polynomial, which is not uniquely defined by the geometric structures present in the theory, requires a certain care. Although there is presently a sufficiently developed theory on the DarbouxโNijenhuis coordinates and on their computation, for the sake of brevity we shall not tackle this problem, but rather limit ourselves to display these polynomials in the example at hand. They are
$$\begin{array}{cc}& \lambda ^2h_1\lambda +(h_1^2h_3)=0\hfill \\ & \mu h_2\lambda +(h_1h_2h_4)=0\hfill \end{array}$$
(6.21)
The important idea emerging from the previous discussion is that the GZ equations are often coupled with a special system of coordinates related with the Poisson pair.
###### Exercise 6.2
Check that the polynomials (6.21) define a system of DarbouxโNijenhuis coordinates for the pair $`(\omega _0,N)`$ considered above.
### 6.4 Separation of Variables
We start from the classical Stรคckel theorem on the separability, in orthogonal coordinates, of the HamiltonโJacobi equation associated with the natural Hamiltonian
$$H(q,p)=\frac{1}{2}g^{ii}(q)p_i^2+V(q_1,\mathrm{},q_n)$$
(6.22)
on the cotangent bundle of the configuration space. According to Stรคckel, this Hamiltonian is separable if and only if there exists as invertible matrix $`S(q_1,\mathrm{},q_n)`$ and a vector $`U(q_1,\mathrm{},q_n)`$ such that $`H`$ is among the solutions $`(H_1,\mathrm{},H_n)`$ of the linear system
$$\frac{1}{2}p_i^2=U_i(q)+\underset{j=1}{\overset{n}{}}S_{ij}(q)H_j,$$
(6.23)
and $`S`$ and $`U`$ verify the Stรคckel condition:
The rows of $`S`$ and $`U`$ depend only on the corresponding coordinate.
This means for instance that the elements $`S_{1j}`$ and $`U_1`$ depend only on the first coordinate $`q_1`$, and so on. Such a matrix $`S`$ is called a Stรคckel matrix (and $`U`$ a Stรคckel vector).
The strategy we shall follow to prove the separability of the HamiltonโJacobi equations associated with the GZ vector fields $`X_1`$ and $`X_3`$ on the manifold $`S_4`$ considered above, is to show that the DarbouxโNijenhuis coordinates allow to define a Stรคckel matrix for the corresponding Hamiltonians.
The construction of the Stรคckel matrix starts from the matrix $`๐ฅ`$ which relates the vector field $`X_1`$ and $`X_3`$ to the Nijenhuis tensor $`N`$ (see equation (6.15)). One can prove that this matrix satisfies the remarkable identity
$$N^{}d๐ฅ=๐ฅd๐ฅ.$$
(6.24)
This is a matrix equation which must be interpreted as follows: $`d๐ฅ`$ is a matrix of 1โforms, and $`N^{}`$ acts separately on each entry of this matrix; $`๐ฅd๐ฅ`$ denotes the matrix multiplication of the matrices $`๐ฅ`$ and $`d๐ฅ`$, which amounts to linearly combine the 1โforms appearing in $`d๐ฅ`$. In our example, equation (6.24) becomes
$$\begin{array}{ccccc}N^{}d(\text{Tr}N)& =& d(\text{det}N)& +& (\text{Tr}N)d(\text{Tr}N)\\ N^{}d(\text{det}N)& =& & +& (\text{det}N)d(\text{Tr}N)\end{array}$$
(6.25)
###### Exercise 6.3
Check that this equations are verified by the Nijenhuis tensor (6.13).
We leave for a moment the particular case we are dealing with, and we suppose that, on a symplectic bi-Hamiltonian manifold fulfilling the conditions of Subsection 6.3, a family of $`n`$ vector fields $`(X_1,X_3,\mathrm{},X_{2n1})`$ is given. We assume that they are Hamiltonian with respect to $`P_0`$, say, $`X_{2i1}=P_0dC_i`$, and that there exists a matrix $`๐ฅ`$ such that
$$NX_{2i1}=\underset{j=1}{\overset{n}{}}๐ฅ_i^jX_{2j1}\text{for all }i.$$
(6.26)
Finally, we suppose that $`๐ฅ`$ satisfies condition (6.24). Then, from the matrix $`๐ฅ`$ we build up the matrix $`๐ณ`$ whose rows are the leftโeigenvectors of $`๐ฅ`$. In other words, we construct a matrix $`๐ณ`$ such that
$$๐ฅ=๐ณ^1\mathrm{\Lambda }๐ณ,$$
(6.27)
where $`\mathrm{\Lambda }=\text{diag}(\lambda _1,\mathrm{},\lambda _n)`$ is the diagonal matrix of the eigenvalues of $`๐ฅ`$, coinciding with the eigenvalues of $`N`$. The matrix $`๐ณ`$ is normalized by imposing that in each row there is a constant component. A suitable normalization criterion, for instance, is to set the entries in the last column equal to $`1`$.
###### Theorem 6.4
If the matrix $`๐ฅ`$ verifies condition (6.24) (as it is always true in our class of examples), then the matrix $`๐ณ`$ is a (generalized) Stรคckel matrix in the DarbouxโNijenhuis coordinates.
This theorem means that the rows of the matrix $`๐ณ`$ verify the following generalized Stรคckel condition: The entries of the first row of $`๐ณ`$ depend only on the canonical pair $`(\lambda _1,\mu _1)`$, those of the second row on $`(\lambda _2,\mu _2)`$, and so on. With respect to the classical case recalled at the beginning of this lecture, we notice that by generalizing the class of Hamiltonians considered, we have been obliged to extend a little bit the notion of Stรคckel matrix. However, this extension does not affect the theorem of separability. Indeed, as a consequence of the fact that the matrix $`๐ฅ`$ is defined by the vector fields $`(X_1,X_3,\mathrm{},X_{2n1})`$ themselves through equation (6.26), one can prove that $`๐ณ`$ is a Stรคckel matrix for the corresponding Hamiltonians $`(C_1,\mathrm{},C_n)`$.
###### Theorem 6.5
The column vector
$$๐=๐ณ๐,$$
(6.28)
where $`๐`$ is the column vector of the Hamiltonians $`(C_1,\mathrm{},C_n)`$, verifies the (generalized) Stรคckel condition in the DarbouxโNijenhuis coordinates. This means that the first component of $`๐`$ depends only on the pair $`(\lambda _1,\mu _1)`$, the second on $`(\lambda _2,\mu _2)`$, and so on.
We shall not prove these two theorems here, preferring to see them โat workโ in the example at hand. First we consider the matrix $`T`$. Due to the form (6.15) of the matrix $`๐ฅ`$, it is easily proved that
$$๐ณ=\left(\begin{array}{cc}\lambda _1& 1\\ \lambda _2& 1\end{array}\right).$$
(6.29)
Indeed, the equation $`\mathrm{๐ณ๐ฅ}=\mathrm{\Lambda }๐ณ`$ follows directly from the characteristic equation for the tensor $`N`$. It should be noted that the matrix $`๐ณ`$ has been computed without computing explicitly the eigenvalues $`\lambda _1`$ and $`\lambda _2`$. It is enough to use the first of equations (6.21), defining the DarbouxโNijenhuis coordinates. The matrix $`๐ณ`$ clearly possess the Stรคckel property (even in the classical, restricted sense).
The vector $`๐`$ can be computed as well without computing explicitly the coordinates $`(\lambda _j,\mu _j)`$. It is sufficient, once again, to use the equations (6.21). We now pass to prove that equation (6.28), in our example, has the particular form
$$\begin{array}{cc}\hfill \frac{1}{2}\mu _1^2\frac{1}{2}\lambda _1^2E\lambda _1^2& =\lambda _1C_1+C_2\hfill \\ \hfill \frac{1}{2}\mu _2^2\frac{1}{2}\lambda _2^2E\lambda _2^2& =\lambda _2C_1+C_2.\hfill \end{array}$$
(6.30)
We notice that proving this statement is tantamount to proving that the following equality between polynomials,
$$\mu (\lambda )^2\lambda ^5=2C(\lambda ),$$
(6.31)
is verified in correspondence of the eigenvalues of $`N`$. This can be done as follows. Let us write the polynomials defining the DarbouxโNijenhuis coordinates in the symbolic form
$$\begin{array}{c}\lambda ^2=e_1\lambda +e_2\hfill \\ \mu =f_1\lambda +f_2.\hfill \end{array}$$
(6.32)
The coefficients $`(e_j,f_j)`$ of these polynomials must be regarded as known functions of the coordinates on the manifold. By squaring the second polynomial and by eliminating $`\lambda ^2`$ by means of the first equation, we get
$$\mu ^2=f_1^2(e_1\lambda +e_2)+2f_1f_2\lambda +f_1f_2=(f_1^2e_1+2f_1f_2)\lambda +(f_1^2e_2+f_2^2).$$
(6.33)
In the same way we obtain
$$\begin{array}{cc}\hfill \lambda ^5=\lambda \lambda ^4& =\lambda [(e_1^3+2e_1e_2)\lambda +(e_1^2e_2+e_2^2)]\hfill \\ & =(e_1^4+3e_1^2e_2+e_2^2)\lambda +(e_1^3e_2+2e_1e_2^2).\hfill \end{array}$$
(6.34)
Finally,
$$C(\lambda )=C_0\lambda ^2+C_1\lambda +C_2=(C_0e_1+C_1)\lambda +(C_0e_2+C_2).$$
(6.35)
By inserting these expressions into equation (6.31), we see that the resulting equation splits into two parts, according to the โsurvivingโ powers of $`\lambda `$:
$$\begin{array}{cc}\hfill \lambda & :(f_1^2e_1+2f_1f_2)(e_1^4+3e_1^2e_2+e_2^2)=2(C_0e_1+C_1)\hfill \\ \hfill 1& :(e_1^2+e_2+e_2^2)(e_1^3e_2+2e_1e_2^2)=2(C_0e_2+C_2).\hfill \end{array}$$
(6.36)
This method allows to reduce the proof of the separability of the HamiltonโJacobi equation(s) to the procedure of checking that explicitly known functions identically coincide on the manifold.
We end our discussion of the separability at this point. Our aim was simply to introduce the method of Poisson pairs, and to show by means of concrete examples how it can be profitably used to define and solve special classes of integrable Hamiltonian equations. We hope that the examples discussed in these lectures might be successful in giving at least a feeling of the nature and the potentialities of this method. |
warning/0002/hep-ph0002238.html | ar5iv | text | # 1 Introduction
## 1 Introduction
Considerable attention has been devoted to double Higgs boson production at future $`e^+e^{}`$ and hadron colliders, both in the Standard Model (SM) and the MSSM (see Refs. for an incomplete list of references for SM $`e^+e^{}`$ and hadron colliders, respectively, and similarly for the MSSM). For the SM, detailed signal-to-background studies already exist for a LC environment , for both โreducibleโ and โirreducibleโ backgrounds , which have assessed the feasibility of experimental analyses. At the LHC, since here the typical SM signal cross sections are of the order of 10 fb , high integrated luminosities would be needed to generate a statistically large enough sample of double Higgs events. These would be further obscured by an overwhelming background, making their selection and analysis in a hadronic environment extremely difficult. Thus, in this contribution we will concentrate only on the case of the MSSM.
In the Supersymmetric (SUSY) scenario, the phenomenological potential of these reactions is two-fold. Firstly, in some specific cases, they can furnish new discovery channels for Higgs bosons. Secondly, they are all dependent upon several triple Higgs self-couplings of the theory, which can then be tested by comparing theoretical predictions with experimental measurements. This is the first step in the reconstruction of the Higgs potential itself<sup>1</sup><sup>1</sup>1The determination of the quartic self-interactions is also required, but appears out of reach for some time: see Refs. for some cross sections of triple Higgs production..
The Higgs Working Group (WG) has focused much of its attention in assessing the viability of these reactions at future TeV colliders. However, the number of such processes is very large both at the LHC and a LC , so only a few โreferenceโ reactions could be studied in the context of this Workshop. Work is in progress for the longer term, which aims to cover most of the double Higgs production and decay phenomenology at both accelerators .
These reference reactions were chosen to be $`gghh`$ for the LHC (see top of Fig. 1) and $`e^+e^{}hhZ`$ for the LC (see bottom of Fig. 1), where $`h`$ is the lightest of the MSSM scalar Higgs bosons. The reason for this preference is simple. Firstly, a stable upper limit exists on the value of $`m_h`$, of the order of 130 GeV, now at two-loop level , so that its detection is potentially well within the reach of both the LHC and a LC. In contrast, the mass of all other Higgs bosons of the MSSM may vary from the electroweak (EW) scale, $`๐ช(m_Z)`$, up to the TeV region. In addition, as noted in Ref. , the multi-$`b`$ final state in $`gghhb\overline{b}b\overline{b}`$, with two resonances and large transverse momenta, may be exploited in the search for the $`h`$ scalar in the large $`\mathrm{tan}\beta `$ and moderate $`m_A`$ region. This is a corner of the MSSM parameter space that has so far eluded the scope of the standard Higgs production and decay modes . (The symbol $`A`$ here denotes the pseudoscalar Higgs boson of the MSSM, and we reserve the notation $`H`$ for the heaviest scalar Higgs state of the model.) However, this paper will not investigate the LHC discovery potential in this mode, given the very sophisticated treatment of the background (well beyond the scope of this note) required by the assumption that no $`h`$ scalar state has been previously discovered (see below). This will be done in Ref. . Furthermore, the $`gghh`$ and $`e^+e^{}hhZ`$ modes largely dominate double Higgs production , at least for centre-of-mass (CM) energies of 14 TeV at the LHC and 500 GeV in the case of a LC, the default values of our analysis. (Notice that we assume no polarization of the incoming beams in $`e^+e^{}`$ scatterings.) Finally, when $`m_H\stackrel{>}{}2m_h`$, the two reactions are resonant, as they can both proceed via intermediate states involving $`H`$ scalars, through $`ggH`$ and $`e^+e^{}HZ`$, which in turn decay via $`Hhh`$ . Thus, the production cross sections are largely enhanced (up to two orders of magnitude above typical SM rates at the LHC ) and become clearly visible. This allows the possibility of probing the trilinear $`Hhh`$ vertex at one or both these colliders.
The dominant decay rate of the MSSM $`h`$ scalar is into $`b\overline{b}`$ pairs, regardless of the value of $`\mathrm{tan}\beta `$ . Therefore, the final signatures of our reference reactions always involve four $`b`$-quarks in the final state. (In the case of a LC environment, a further trigger on the accompanying $`Z`$ boson can be exploited.)
If one assumes very efficient tagging and high-purity sampling of $`b`$-quarks, the background to $`hh`$ events at the LHC is dominated by the irreducible QCD modes . Among these, we focus here on the cases $`q\overline{q},ggb\overline{b}b\overline{b}`$, as representative of ideal $`b`$-tagging performances. These modes consist of a purely QCD contribution of $`๐ช(\alpha _s^4)`$, an entirely EW process of $`๐ช(\alpha _{em}^4)`$ (with no double Higgs intermediate states) and an $`๐ช(\alpha _s^2\alpha _{em}^2)`$ component consisting of EW and QCD interactions. (Note that in the EW case only $`q\overline{q}`$ initiated subprocesses are allowed at tree-level.) For a LC, the final state of the signal is $`b\overline{b}b\overline{b}Z`$, with the $`Z`$ reconstructed from its decay products in some channel. Here, the EW background is of $`๐ช(\alpha _{em}^5)`$ away from resonances (and, again, contains no more than one intermediate Higgs boson), whereas the EW/QCD background is proportional to $`(\alpha _s^2\alpha _{em}^3)`$.
In general, EW backgrounds can be problematic due to the presence of $`Z`$ vectors and single Higgs scalars yielding $`b\overline{b}`$ pairs, with the partons being typically at large transverse momenta and well separated. In contrast, the QCD backgrounds involve no heavy objects decaying to $`b\overline{b}`$ pairs and are dominated by the typical infrared (i.e., soft and collinear) QCD behaviour of the partons in the final state. However, they can yield large production rates because of the strong couplings.
In this study, we investigate the interplay between the signal and background at both colliders, adopting detector as well as dedicated selection cuts. We carry out our analysis at both parton and hadron level. The plan of this note is as follows. The next Section details the procedure adopted in the numerical computation. Sect. 3 displays our results and contains our discussion. Finally, in the last section, we summarise our findings and consider possible future studies.
## 2 Calculation
For the parton level simulation, the double Higgs production process at the LHC, via $`gg`$ fusion, has been simulated using the program of Ref. to generate the interaction $`gghh`$, with the matrix elements (MEs) taken at leading-order (LO) for consistency with our treatment of the background. We then perform the two $`hb\overline{b}`$ decays to obtain the actual $`4b`$-final state. For double Higgs production at a LC, we use a source code for the signal derived from that already used in Ref. . At both colliders, amplitudes for background events were generated by means of MadGraph and the HELAS package . Note that interferences between signal and backgrounds, and between the various background contributions themselves, have been neglected. This is a good approximation for the interferences involving the signal because of the very narrow width of the MSSM lightest Higgs boson. Similarly, the various background subprocesses have very different topologies, and one would expect their interferences to be small in general.
The Higgs boson masses and couplings of the MSSM can be expressed at tree-level in terms of the mass of the pseudoscalar Higgs state, $`m_A`$, and the ratio of the vacuum expectation values of the two neutral fields in the two iso-doublets, $`\mathrm{tan}\beta `$. At higher order however, top and stop loop-effects can become significant. Radiative corrections in the one-loop leading $`m_t^4`$ approximation are parameterised by
$`ฯต{\displaystyle \frac{3G_Fm_t^4}{\sqrt{2}\pi ^2\mathrm{sin}^2\beta }}\mathrm{log}{\displaystyle \frac{m_S^2}{m_t^2}}`$ (1)
where the SUSY breaking scale is given by the common squark mass, $`m_S`$, set equal to $`1`$ TeV in the numerical analysis. If stop mixing effects are modest at the SUSY scale, they can be accounted for by shifting $`m_S^2`$ in $`ฯต`$ by the amount $`\mathrm{\Delta }m_S^2=\widehat{A}^2[1\widehat{A}^2/(12m_S^2)]`$ ($`\widehat{A}`$ is the trilinear common coupling). The charged and neutral CP-even Higgs boson masses, and the Higgs mixing angle $`\alpha `$ are given in this approximation by the relations:
$`m_{H^\pm }^2`$ $`=`$ $`m_A^2+m_Z^2\mathrm{cos}^2\theta _W,`$
$`m_{h,H}^2`$ $`=`$ $`\frac{1}{2}[m_A^2+m_Z^2+ฯต`$
$``$ $`\sqrt{(m_A^2+m_Z^2+ฯต)^24m_A^2m_Z^2\mathrm{cos}^22\beta 4ฯต(m_A^2\mathrm{sin}^2\beta +m_Z^2\mathrm{cos}^2\beta )}],`$
$`\mathrm{tan}2\alpha `$ $`=`$ $`\mathrm{tan}2\beta {\displaystyle \frac{m_A^2+m_Z^2}{m_A^2m_Z^2+ฯต/\mathrm{cos}2\beta }}\text{with}{\displaystyle \frac{\pi }{2}}\alpha 0,`$ (2)
as a function of $`m_A`$ and $`\mathrm{tan}\beta `$. The triple Higgs self-couplings of the MSSM can be parameterised in units of $`M_Z^2/v`$, $`v=246`$ GeV, as,
$`\lambda _{hhh}`$ $`=`$ $`3\mathrm{cos}2\alpha \mathrm{sin}(\beta +\alpha )+3{\displaystyle \frac{ฯต}{m_Z^2}}{\displaystyle \frac{\mathrm{cos}\alpha }{\mathrm{sin}\beta }}\mathrm{cos}^2\alpha ,`$
$`\lambda _{Hhh}`$ $`=`$ $`2\mathrm{sin}2\alpha \mathrm{sin}(\beta +\alpha )\mathrm{cos}2\alpha \mathrm{cos}(\beta +\alpha )+3{\displaystyle \frac{ฯต}{m_Z^2}}{\displaystyle \frac{\mathrm{sin}\alpha }{\mathrm{sin}\beta }}\mathrm{cos}^2\alpha ,`$
$`\lambda _{HHh}`$ $`=`$ $`2\mathrm{sin}2\alpha \mathrm{cos}(\beta +\alpha )\mathrm{cos}2\alpha \mathrm{sin}(\beta +\alpha )+3{\displaystyle \frac{ฯต}{m_Z^2}}{\displaystyle \frac{\mathrm{cos}\alpha }{\mathrm{sin}\beta }}\mathrm{sin}^2\alpha ,`$
$`\lambda _{HHH}`$ $`=`$ $`3\mathrm{cos}2\alpha \mathrm{cos}(\beta +\alpha )+3{\displaystyle \frac{ฯต}{m_Z^2}}{\displaystyle \frac{\mathrm{sin}\alpha }{\mathrm{sin}\beta }}\mathrm{sin}^2\alpha ,`$
$`\lambda _{hAA}`$ $`=`$ $`\mathrm{cos}2\beta \mathrm{sin}(\beta +\alpha )+{\displaystyle \frac{ฯต}{m_Z^2}}{\displaystyle \frac{\mathrm{cos}\alpha }{\mathrm{sin}\beta }}\mathrm{cos}^2\beta ,`$
$`\lambda _{HAA}`$ $`=`$ $`\mathrm{cos}2\beta \mathrm{cos}(\beta +\alpha )+{\displaystyle \frac{ฯต}{m_Z^2}}{\displaystyle \frac{\mathrm{sin}\alpha }{\mathrm{sin}\beta }}\mathrm{cos}^2\beta .`$ (3)
Next-to-leading order (NLO) effects are certainly dominant, though the next-to-next-to-leading order (NNLO) ones cannot entirely be neglected (especially in the Higgs mass relations). Thus, in the numerical analysis, the complete one-loop and the leading two-loop corrections to the MSSM Higgs masses and the triple Higgs self-couplings are included. The Higgs masses, widths and self-couplings have been computed using the HDECAY program described in Ref. , which uses a running $`b`$-mass in evaluating the $`hb\overline{b}`$ decay fraction. Thus, for consistency, we have evolved the value of $`m_b`$ entering the $`hbb`$ Yukawa couplings of the $`hb\overline{b}`$ decay currents of our processes in the same way.
For our analysis, we have considered $`\mathrm{tan}\beta =3`$ and $`50`$. For the LHC, high values of $`\mathrm{tan}\beta `$ produce a signal cross section much larger than the $`\mathrm{tan}\beta =3`$ scenario, over almost the entire range of $`m_A`$. However, this enhancement is due to the increase of the down-type quark-Higgs coupling, which is proportional to $`\mathrm{tan}\beta `$ itself, and serves only to magnify the dominance of the quark box diagrams of Fig. 1. Unfortunately, these graphs have no dependence on either of the two triple Higgs self-couplings entering the gluon-gluon process considered here, i.e., $`\lambda _{hhh}`$ and $`\lambda _{Hhh}`$. Thus, although the cross section is comfortably observable, all sensitivity to such vertices is lost. Therefore, the measurement of the triple Higgs self-coupling, $`\lambda _{Hhh}`$, is only feasible at the LHC for low $`\mathrm{tan}\beta `$ due to the resonant production of the heavy Higgs boson (see Fig. 5a of Ref. ).
In contrast, the cross section for double Higgs production at the LC is small for large $`\mathrm{tan}\beta `$ because there is no heavy Higgs resonance (see Fig. 8 of Ref. ). As soon as it becomes kinematically possible to decay the heavy Higgs into a light Higgs pair, the $`ZZH`$ coupling is already too small to generate a sizable cross section. Furthermore, the continuum MSSM cross section is suppressed with respect to the SM cross section since the MSSM couplings $`ZZH`$ and $`ZZh`$ vary with $`\mathrm{cos}(\beta \alpha )`$ and $`\mathrm{sin}(\beta \alpha )`$, respectively, with respect to the corresponding SM coupling. Notice that in this regime, at a LC, the $`\lambda _{hhh}`$ vertex could in principle be accessible instead, since $`\lambda _{hhh}\lambda _{Hhh}`$ (see Fig. 2 of Ref. ) and because of the kinematic enhancement induced by $`m_hm_H`$. Unfortunately, we will see that the size of the $`e^+e^{}hhZ`$ cross section itself is prohibitively small.
We assume that $`b`$-jets are distinguishable from light-quark and gluon jets and no efficiency to tag the four $`b`$-quarks is included in our parton level results. We further neglect considering the possibility of the $`b`$-jet charge determination. Also, to simplify the calculations, the $`Z`$ boson appearing in the final state of the LC process is treated as on-shell and no branching ratio (BR) is applied to quantify its possible decays. In practice, one may assume that it decays leptonically (i.e., $`Z\mathrm{}^+\mathrm{}^{}`$, with $`\mathrm{}=e,\mu ,\tau `$) or hadronically into light-quark jets (i.e., $`Zq\overline{q}`$, with $`qb`$), in order to avoid problems with $`6b`$-quark combinatorics. Furthermore, in the LC analysis, we have not simulated the effects of Initial State Radiation (ISR), beamstrahlung or Linac energy spread. Indeed, we expect them to affect signal and backgrounds rather similarly, so we can neglect them for the time being. Indeed, since a detailed phenomenological study, including both hadronisation and detector effects, already exists for the case of double Higgs-strahlung in $`e^+e^{}`$ , whose conclusions basically support those attained in the theoretical study of Ref. , we limit ourselves here to update the latter to the case of the MSSM.
So far only resonant production $`gg`$ $``$ $`H`$ $``$ $`hh`$ $``$ $`b\overline{b}b\overline{b}`$ has been investigated , with full hadronic and detector simulation and considering also the (large) QCD backgrounds, and a similar study does not exist for continuum double Higgs production at the LHC. (See Ref. for a detailed account of the $`ggH`$ $``$ $`hh`$ $``$ $`\gamma \gamma b\overline{b}`$ decay channel.) The event simulation has been performed by using a special version of PYTHIA , in which the relevant LO MEs for double Higgs production of Ref. have been implemented by M. El Kacimi and R. Lafaye. These MEs take into account both continuum and resonant double Higgs boson production and their interferences. (The insertion of those for $`e^+e^{}`$ processes is in progress.) The PYTHIA interface to HDECAY has been exploited in order to generate the MSSM Higgs mass spectrum and the relevant Higgs BRs, thus maintaining consistency with the parton level approach. As for the LHC detector simulation, the fast simulation package was used, with high luminosity (i.e., $`๐t=100`$ fb<sup>-1</sup>) parameters.
The motivation for our study is twofold. On the one hand, to complement the studies of Ref. by also considering the continuum production $`gghhb\overline{b}b\overline{b}`$ at large $`\mathrm{tan}\beta `$. On the other hand, to explore the possibility of further kinematic suppression of the various irreducible backgrounds to the resonant channel at small $`\mathrm{tan}\beta `$.
## 3 Results
### 3.1 The LHC analysis
In our LHC analysis, following the discussion in Sect. 2, we focus most of our attention on the case $`\mathrm{tan}\beta =3`$, with $`m_A=210`$ GeV, although other combinations of these two MSSM parameters will also be considered. We further set $`A=\mu =1`$ TeV and take all sparticle masses (and other SUSY scales) to be as large as $`1`$ TeV.
#### 3.1.1 $`gghhb\overline{b}b\overline{b}`$ at parton level
In our parton level analysis, we identify jets with the partons from which they originate (without smearing the momenta) and apply all cuts directly to the partons. We mimic the finite coverage of the LHC detectors by imposing a transverse momentum threshold on each of the four $`b`$-jets,
$$p_T(b)>30\mathrm{GeV}$$
(4)
and requiring their pseudorapidity to be
$$|\eta (b)|<2.5.$$
(5)
Also, to allow for their detection as separate objects, we impose an isolation criterium among $`b`$-jets,
$$\mathrm{\Delta }R(bb)>0.4,$$
(6)
by means of the usual cone variable $`\mathrm{\Delta }R(ij)=\sqrt{\mathrm{\Delta }\eta (ij)^2+\mathrm{\Delta }\varphi (ij)^2}`$, defined in terms of relative differences in pseudo-rapidity $`\eta _{ij}`$ and azimuth $`\varphi _{ij}`$ of the $`i`$-th and $`j`$-th $`b`$-jets.
As preliminary and very basic selection cuts (also to help the stability of the numerical integration), we have required that the invariant mass of the entire $`4b`$-system is at least twice the mass of the lightest MSSM Higgs boson (apart from mass resolution and gluon emission effects), e.g.,
$$m(bbbb)2m_h40\mathrm{GeV},$$
(7)
and that exactly two $`h`$-resonances are reconstructed, such that
$$|m(bb)m_h|<20\mathrm{GeV}.$$
(8)
In doing so, we implicitly assume that the $`h`$ scalar boson has already been discovered and its mass measured through some other channel, as we have already intimated in the Introduction.
After the above cuts have been implemented, we have found that the two $`4b`$-backgrounds proceeding through EW interactions are negligible compared to the pure QCD process. In fact, the constraints described in eqs. (7)โ(8) produce the strongest suppression, almost completely washing out the relatively enhancing effects that the cuts in (4)โ(6) have on the EW components of the backgrounds with respect to the pure QCD one, owning to the intermediate production of massive $`Z`$ bosons in the former. In the end, the production rates of the three subprocesses scale approximately as their coupling strengths: i.e., $`๐ช(\alpha _s^4)`$ : $`๐ช(\alpha _s^2\alpha _{em}^2)`$ : $`๐ช(\alpha _{em}^4)`$. Therefore, in the reminder of our analysis, we will neglect EW effects, as they represent not more than a 10% correction to the QCD rates, which are in turn affected by much larger QCD $`K`$-factors. As for the pure QCD background itself, it hugely overwhelms the double Higgs signal at this stage. The cross section of the former is about 7.85 pb, whereas that of the latter is approximately 0.16 pb.
To appreciate the dominance of the $`m_h`$ cuts, one may refer to Fig. 2, where the distributions in transverse momentum of the four $`p_T`$-ordered $`b`$-quarks (such that $`p_T(b_1)>\mathrm{}>p_T(b_4)`$) of both signal and QCD background are shown. Having asked the four $`b`$-jets of the background to closely emulate the $`gghhb\overline{b}b\overline{b}`$ kinematics, it is not surprising to see a โdegeneracyโ in the shape of all spectra. Clearly, no further background suppression can be gained by increasing the $`p_T(b)`$ cuts. The same can be said for $`\eta (b)`$ and $`\mathrm{\Delta }R(bb)`$. Others quantities ought to be exploited.
In Fig. 3, we present the signal and QCD background distributions in the minimum angle formed between the two $`b`$-quarks coming from the โsame Higgsโ (i.e., those fulfilling the cuts in (8)) in the $`4b`$-system rest frame (the plot is rater similar for the maximum angle, thus also on average). There, one can see a strong tendency of the two $`2b`$-pairs produced in the Higgs decays to lie back-to-back, reflecting the $`22`$ intermediate dynamics of Higgs pair production via $`gghh`$. Missing such kinematically constrained virtual state, the QCD background shows a much larger angular spread towards small $`\theta _{\mathrm{min}}(bb)`$ values, eventually tamed by the isolation cut (6).
The somewhat peculiar shape of the signal distribution is due to destructive interference. Recall that the signal contains not only diagrams proceeding via a heavy Higgs resonance (the upper-left hand graph of Fig. 1), which results in the large peak in Fig. 3, but also contains a continuum contribution mediated by box graphs (the upper-right hand graph of Fig. 1). These two contributions destructively interfere leading to the depletion of events between the large back-to-back peak and the small remaining โbumpโ of the continuum contribution as seen in Fig. 3.
In the end, a good criterium to enhance the signal-to-background ratio ($`S/B`$) is to require, e.g., $`\theta (bb)>2.4`$ radians, i.e., a separation between the $`2b`$-jets reconstructing the lightest Higgs boson mass of about 140 degrees in angle. (Incidentally, we also have investigated the angle that each of these $`2b`$-pairs form with the beam axis, but found no significant difference between signal and QCD background).
An additional consequence that one should expect from the presence of two intermediate massive objects in $`gghhb\overline{b}b\overline{b}`$ events is the spherical appearance of the jets in the final state, in contrast to the usual planar behaviour of the infrared QCD interactions. These phenomena can be appreciated in Fig. 4. Notice there the strong tendency of the background to yield high thrust configurations, again controlled by the separation cuts when $`T`$ approaches unity. On the contrary, the average value of the thrust in the signal is much lower, being the effect of accidental pairings of โwrongโ $`2b`$-pairs (the shoulder at high thrust values) marginal. An effective selection cut seems to be, e.g., $`T<0.85`$.
Furthermore, if the heavy Higgs mass is sufficiently well measured at the LHC then one can exploit the large fraction of $`4b`$-events which peak at $`m_H`$ in the signal, as dictated by the $`Hhh`$ decay, improving the signal-to-background ratio. This peak at $`m_H`$ can be clearly seen in the left hand plot of Fig. 5, where it dominates the QCD background, even for bins 13 GeV wide. In fact, not only could the QCD background be considerably suppressed but also those contributions to $`gghh`$ not proceeding through an intermediate $`H`$ state should be removed, this greatly enhancing the sensitivity of the signal process to the $`\lambda _{Hhh}`$ coupling. This can be seen in the right hand plot of Fig. 5 where the signal is shown on a logarithmic scale. The continuum contribution due to the box graphs (and its destructive interference with the heavy Higgs decay contribution) is now evident although one should note that it is considerably suppressed compared to the peak at $`m_H`$.
Now, if a less than 10% mass resolution can be achieved on the light and heavy Higgs masses, then one can tighten cut (8) to $`|m(bb)m_h|<10`$ GeV and introduce the additional cut $`|m(bbbb)m_H|<20`$ GeV. These cuts taken together with those in $`\theta (bb)`$ and $`T`$ already suggested, reduce the QCD background to the same level as the signal. In fact, we have found that the cross section of the background drops to approximately 174 fb whereas that of the signal remains as large as 126 fb, this yielding a very high statistical significance at high luminosity. Even for less optimistic mass resolutions the signal-to-background ratio is still significantly large. For example, selecting events with $`|m(bb)m_h|<20`$ GeV and $`|m(bbbb)m_H|<40`$ GeV, the corresponding numbers are approximately 102 fb for the signal and 453 fb for the background. Notice that the signal actually decreases as these Higgs mass windows are made larger. This is due to our insistence that exactly two $`b\overline{b}`$ pairs should reconstruct the light Higgs mass. As the light Higgs mass window is enlarged from $`m_h\pm 10`$ GeV to $`m_h\pm 20`$ GeV, it becomes more likely that accidental pairings reconstruct the light Higgs boson. Since one is then unable to unambiguously assign the $`b`$ quarks to the light Higgs bosons, the event is rejected and the signal drops.
Although we have discussed here an ideal situation which is difficult to match with more sophisticated hadronic and detector simulations, it still demonstrates that the measurement of the $`\lambda _{Hhh}`$ coupling could be well within the potential of the LHC, at least for our particular choice of MSSM parameters. Comforted by such a conclusion, we now move on to more realistic studies.
#### 3.1.2 $`gghhb\overline{b}b\overline{b}`$ at the LHC experiments
Although the LHC experiments will be the first where one can attempt to measure the Higgs self-couplings, the analysis is very challenging because of the smallness of the production cross sections. Even in the most favourable cases, the production rate is never larger than a few picobarns, already including one-loop QCD corrections, as computed in Ref. . The cross sections at this accuracy are given in Tab. 1, for the resonant process (case 1 with $`m_H=220`$ GeV) as well as three non resonant scenarios: one at the same $`\mathrm{tan}\beta `$ but with the $`Hhh`$ decay channel closed (case 2), a second at very large $`\mathrm{tan}\beta `$ and no visible resonance (case 3) and, finally, the SM option (case 4, where $`m_h`$ identifies with the mass of the standard Higgs state).
#### 3.1.3 LHC trigger acceptance
For $`4b`$-final states, possible LHC triggers are high $`p_T`$ electron/muons and jets. As an example, the foreseen ATLAS level 1 trigger thresholds on $`p_T`$ and acceptance for a $`4b`$-selection (with the four $`b`$-jets reconstructed in the detector) are given in Tab. 2, assuming the LHC to be running at high luminosity.
The overall trigger acceptance is at best 8โ9%, for cases 2,3,4. The very low efficiency for case 1 is clearly a consequence of the small value of the difference $`m_H2m_h`$, translating into a softer $`p_T(b)`$ spectrum with respect to the other cases (compare the left-hand with the right-hand side of Fig. 6). One can further see in the left-hand plot of Fig. 6 that the bulk of the signal lies below the lowest $`p_T(b)`$ threshold of Tab. 2 (i.e., $`90`$ GeV), so that adopting smaller trigger thresholds could result in a dramatic enhancement of our efficiency. Of course, this would also substantially increase the low transverse momentum QCD background, as we can see in the parton level analysis of Fig. 2.
For example, by lowering the thresholds to 180, 80 and 50 GeV for 1, 3 and 4 jets, respectively (compare to Tab. 2), the overall trigger acceptance on the signal goes up to 1.8%, i.e., by almost a factor of 4. Meanwhile though, the ATLAS level-1 jet trigger rates increase by a factor of 10 . Anyhow, even for our poor default value of $`ฯต(bbbb)`$ in Tab. 2, we will see that case 1 still yields a reasonable number of events in the end. Optimisations of the $`b`$-jet transverse momentum thresholds are in progress .
#### 3.1.4 LHC events selection for $`gghhb\overline{b}b\overline{b}`$
Jets are reconstructed merging tracks inside $`\mathrm{\Delta }R(bb)=0.4`$. Only jets with transverse energy/momentum greater than 30 GeV and with $`|\eta (b)|<2.5`$ are kept. (Thus, the same cuts as in the parton level analysis, now applied instead to jets.) The effect from pile up is included in the resolution. A jet energy correction is then applied.
The invariant masses of each jet pair can then be computed. Assuming that the lightest Higgs boson mass is known, events with $`m(bb)`$ sufficiently close to $`m_h`$ can efficiently be selected, see Fig. 7. Another cut on the $`\mathrm{\Delta }R(bb,bb)`$ between pairs of $`b`$-jets can also be applied to reduce the intrinsic combinatorial background, since the latter concentrates at large $`\mathrm{\Delta }R(bb,bb)`$ values, see Fig. 8.
For case 1, as already discussed, we can further impose that the invariant mass of the four $`b`$-jets should be the heavy Higgs mass, $`m_H`$, in order to select the $`Hhh`$ resonance, as confirmed by Fig. 9. In the other three cases, where the $`Hhh`$ splitting is no longer dominant (MSSM) or non-existent (SM), one can still insist that the $`4b`$-jet invariant mass should be higher than two times the lightest Higgs mass, see Fig. 10 and recall eq. (7). Finally, following Fig. 11, by constraining the $`b`$-jets pairs four-momenta around the known light Higgs mass value, $`m_h`$, one can further reject the intrinsic background by means of the $`m(bbbb)`$ spectrum.
#### 3.1.5 LHC $`b`$-tagging in $`gghhb\overline{b}b\overline{b}`$
The $`b`$-tagging efficiency at high luminosity is set to 50%, with $`p_T`$ dependent correction factors for jets rejection. An average rejection of 10 for $`c`$-jets and 100 for light-jets can be expected. We then studied the effect on the selection efficiency of requiring from one to four $`b`$-tags, although it is clear that, according to the parton level studies, the huge background rate demands four $`b`$-tags, leading to a 6% tagging efficiency overall.
#### 3.1.6 Event rates at the LHC
Taking into account all the efficiencies described above, and using the NLO normalisation of Tab. 1, one can extract the number of expected events per year at the LHC at high luminosity given in Tab. 3. The selection cuts enforced here are the following. For a start, we have kept configurations where $`|m(bb)m_h|<30`$ GeV (cases 1,3,4) or $`|m(bb)m_h|<20`$ GeV (case 2) and $`\mathrm{\Delta }R(bb,bb)<2.5`$ (all four cases). (If more than two $`m_h`$โs are reconstructed, the best two $`2b`$-pairs are selected according to the minimum value of $`\delta M^2=[m_hm(bb)]^2+[m_hm^{}(bb)^2]`$.) Then, a cut on $`m(bbbb)`$ is applied: in presence of the $`Hhh`$ resonance (case 1) we have kept events within an $`m_H`$ mass window of $`\pm 2\sigma `$ (about 82% of the total number survive); otherwise (cases 2,3,4) we have adjusted the $`m(bbbb)\stackrel{>}{}2m_h`$ cut so to keep 90% of the sample. In the end, one finds the numbers in Tab. 3, that are encouraging indeed.
In conclusion then, looking at the results in Tab. 3 and bearing in mind the potential seen in reducing the pure QCD background via $`gg๐ช(\alpha _s^4)b\overline{b}b\overline{b}`$ (see Figs. 35), one should be confident in the LHC having the potential to measure the $`\lambda _{Hhh}`$ coupling in resonant $`Hhh`$ events (case 1). To give more substance to such a claim, we have now initiated background studies at hadron and detector level, following the guidelines obtained by the parton level analysis . As for other configurations of the MSSM (such as case 2) or in the SM (case 4), the expectations are more pessimistic. Case 3 deserves further attention. In fact, notice the large number of events surviving and recall what mentioned in the Introduction concerning the potential of the non-resonant $`gghhb\overline{b}b\overline{b}`$ process as a discovery channel of the light Higgs boson of the MSSM in the large $`\mathrm{tan}\beta `$ region at moderate $`m_A`$ values, a corner of the parameter space where the $`h`$ coverage is given only by SM-like production/decay modes, thus not allowing one to access information on the MSSM parameters. Results on this topic too will be presented in Ref. .
### 3.2 The LC analysis
Here, we closely follow the selection procedure advocated in Ref. . In order to resolve the four $`b`$-jets as four separate systems inside the LC detector region, we impose the following cuts. First, that the energy of each $`b`$-jet is above a minimum threshold,
$$E(b)>10\mathrm{GeV}.$$
(9)
Second, that any $`b`$-jet is isolated from all others, by requiring a minimum angular separation,
$$\mathrm{cos}\theta (b,b)<0.95.$$
(10)
Similarly to the hadronic analysis, one can optimise $`S/B`$ by imposing the constraints ,
$$m(bbbb)2m_h10\mathrm{GeV},$$
(11)
$$|m(bb)m_h|<5\mathrm{GeV},$$
(12)
on exactly two combinations of $`2b`$-jets. Here, note that the mass resolution adopted for the quark systems is significantly better than in the LHC case, due to the cleanliness of the $`e^+e^{}`$ environment and the expected performance of the LC detectors in jet momentum and angle reconstruction . Thus, given such high mass resolution power from the LC detection apparatus, one may further discriminate between $`h`$ and $`Z`$ mass peaks by requiring that none of the $`2b`$-jet pairs falls around $`m_Z`$,
$$|m(bb)m_Z|>5\mathrm{GeV}.$$
(13)
Moreover, in the double Higgs-strahlung process $`e^+e^{}hhZ`$, the four $`b`$-quarks are produced centrally, whereas this is generally not the case for the background (see the discussion in Ref. ). This can be exploited by enforcing
$$|\mathrm{cos}\theta (bb,bbb,bbbb)|<0.75,$$
(14)
where $`\theta (bb,bbb,bbbb)`$ are the polar angles of all two-, three- and four-jet systems.
Fig. 12 shows the production and decay rates of the signal process, $`e^+e^{}hhZb\overline{b}b\overline{b}Z`$, as obtained at the partonic level, after the cuts (9)โ(10) have been implemented. The MSSM setup here includes some mixing, having adopted $`A=2.4`$ TeV and $`\mu =1`$ TeV, at both $`\mathrm{tan}\beta =3`$ and 50. Notice the onset of the $`Hhhb\overline{b}b\overline{b}`$ decay sequence in the Higgs-strahlung process $`e^+e^{}HZ`$ at low $`\mathrm{tan}\beta `$. The same does not occur for large values, as previously explained. The impact of the above jet selection cuts on the signal is marginal, as the $`b`$-quarks are here naturally isolated and energetic, being the decay products of heavy objects. In fact, the rates in Fig. 12 would only be 10โ20% higher if all the $`4b`$-quark phase space was allowed (the suppression being larger for smaller Higgs masses). At the height of the resonant peak around $`m_h104`$ GeV at $`\mathrm{tan}\beta =3`$, the signal rate is not large but observable, yielding more than one event every 1 fb<sup>-1</sup> of data. For a high luminosity 500 GeV TESLA design , this would correspond to more than 300 events per year. Given the very high efficiency expected in tagging $`b`$-quark jets, estimated at 90% for each pairs of heavy quarks , one should expect a strong sensitivity to the triple Higgs self-coupling. The situation at large $`\mathrm{tan}\beta `$ is much more difficult instead, being the production rates smaller by about a factor of 10.
In the left-hand side of Fig. 13 we present the EW background, after the constraints in (9)โ(10) have been enforced, in the form of the four dominant EW sub-processes. These four channels are the following.
1. $`e^+e^{}ZZZb\overline{b}b\overline{b}Z`$, first from the left in the second row of topologies in Fig. 3 of Ref. . That is, triple $`Z`$ production with no Higgs boson involved.
2. $`e^+e^{}h/HZZb\overline{b}b\overline{b}Z`$, first(first) from the left(right) in the fifth(fourth) row of topologies in Fig. 2 of Ref. (also including the diagrams in which the on-shell $`Z`$ is connected to the electron-positron line). That is, single Higgs-strahlung production in association with an additional $`Z`$, with the Higgs decaying to $`b\overline{b}`$. The cross sections of these two channels are obviously identical.
3. $`e^+e^{}h/HZZ^{}Z^{}Zb\overline{b}b\overline{b}Z`$, first from the right in the third row of topologies in Fig. 2 of Ref. . That is, single Higgs-strahlung production with the Higgs decaying to $`b\overline{b}b\overline{b}`$ via two off-shell $`Z^{}`$ bosons.
4. $`e^+e^{}Zh/Hb\overline{b}Z^{}Zb\overline{b}b\overline{b}Z`$, first(first) from the right(left) in the first(second) row of topologies in Fig. 2 of Ref. . That is, two single Higgs-strahlung production channels with the Higgs decaying to $`b\overline{b}Z`$ via one off-shell $`Z^{}`$ boson. Also the cross sections of these two channels are identical to each other, as in 2.
The $`๐ช(\alpha _s^2\alpha _{em}^3)`$ EW/QCD background is dominated by $`e^+e^{}ZZ`$ production with one of the two $`Z`$ bosons decaying hadronically into four $`b`$-jets. This subprocess corresponds to the topology in the middle of the first row of diagrams in Fig. 4 of Ref. . Notice that Higgs graphs are involved in this process as well (bottom-right topology in the mentioned figure of ). These correspond to single Higgs-strahlung production with the Higgs scalar subsequently decaying into $`b\overline{b}b\overline{b}`$ via an off-shell gluon. Their contribution is not entirely negligible, owing to the large $`ZH`$ production rates, as can be seen in the right-hand side of Fig. 13. The interferences among non-Higgs and Higgs terms are always negligible.
In performing the signal-to-background analysis, we have chosen two representative points only, identified by the two following combinations: (i) $`\mathrm{tan}\beta =3`$ and $`m_A=210`$ GeV (yielding $`m_h104`$ GeV and $`m_H220`$ GeV); (ii) $`\mathrm{tan}\beta =50`$ and $`m_A=130`$ GeV (yielding $`m_h120`$ GeV and $`m_H130`$ GeV). These correspond to the two asterisks in Fig. 12, that is, the maxima of the signal cross sections at both $`\mathrm{tan}\beta `$ values. The first corresponds to resonant $`Hhh`$ production, whereas the latter to the continuum case. If we enforce the constraints of eq. (11)โ(14), the suppression of both EW and EW/QCD is enormous, so that the corresponding cross sections are of $`๐ช(10^3)`$ fb, while the signal rates only decrease by a factor of four at most. This is the same situation that was seen for the SM case in Ref. . Indeed, in the end it is just a matter of how many signal events survive, the sum of the backgrounds representing no more than a 10% correction (see Fig. 11 of Ref. ). For example, after 500 fb<sup>-1</sup> of data collected, one is left with 156 and 15 events for case (i) and (ii), respectively. However, these numbers do not yet include $`b`$-tagging efficiency and $`Z`$ decay rates.
## 4 Summary
To summarise, the โdouble Higgs productionโ subgroup has contributed to the activity of the Higgs WG by assessing the feasibility of measurements of triple Higgs self-couplings at future TeV colliders. The machines considered were the LHC at CERN (14 TeV) and a future LC running at 500 GeV. In both cases, a high luminosity setup was assumed, given the smallness of the double Higgs production cross sections. In particular, the $`Hhh`$ resonant enhancement was the main focus of our studies, involving the lightest, $`h`$, and the heaviest, $`H`$, of the neutral Higgs bosons of the MSSM, in the kinematic regime $`m_H\stackrel{>}{}2m_h`$. This dynamics can for example occur in the following reactions: $`gghh`$ in the hadronic case and $`e^+e^{}hhZ`$ in the leptonic one, but only at low $`\mathrm{tan}\beta `$. These two processes proceed via intermediate stages of the form $`ggH`$ and $`e^+e^{}HZ`$, respectively, followed by the decay $`Hhh`$. Thus, they in principle allow one to determine the strength of the $`Hhh`$ vertex involved, $`\lambda _{Hhh}`$, in turn constraining the form of the MSSM Higgs potential itself. The signature considered was $`hhb\overline{b}b\overline{b}`$, as the $`hb\overline{b}`$ decay rate is always dominant.
We have found that several kinematic cuts can be exploited in order to enhance the signal-to-background rate to level of high significance, particularly at the $`e^+e^{}`$ machine. At the $`pp`$ accelerator, in fact, the selection of the signal is made much harder by the presence of an enormous background in $`4b`$ final states due to pure QCD. In parton level studies, based on the exact calculation of LO scattering amplitudes of both signals and backgrounds (without any showering and hadronisation effects but with detector acceptances), we have found very encouraging results. At a LC, the double Higgs signal can be studied in an essentially background free environment. At the LHC, the signal and the QCD background are in the end at the same level with detectable but not very large cross sections.
Earlier full simulations performed for the $`e^+e^{}`$ case had already indicated that a more sophisticated treatment of both signal and backgrounds, including fragmentation/hadronisation and full detector effects, should not spoil the results seen at the parton level. For the LHC, our preliminary studies of $`ggHhhb\overline{b}b\overline{b}`$ in presence of the $`gghhb\overline{b}b\overline{b}`$ continuum (and relative interferences) also point to the feasibility of the signal selection, after realistic detector simulation and event reconstruction. As for double $`h`$ production in the continuum, although not very useful for Higgs self-coupling measurements, this seems a promising channel, if not to discover the lightest MSSM Higgs boson certainly to study its properties and those of the Higgs sector in general (because of the large production and decay rates at high $`\mathrm{tan}\beta `$ and its sensitivity to such a parameter), as shown from novel simulations also presented in this study. (The discovery potential of this mode will eventually be addressed in Ref. .) Despite lacking a full background analysis in the LHC case, we have no reason to believe that a comparable degree of suppression of background events seen at parton level cannot be achieved also at hadron level. Progress in this respect is currently being made .
#### Acknowledgements
SM acknowledges financial support from the UK-PPARC. The authors thank Patrick Aurenche and the organisers of the Workshop for the stimulating environment that they have been able to create. DJM and MM thank Michael Spira for useful discussions. Finally, we all thank Elzbieta Richter-Was for many useful comments and suggestions. |
warning/0002/hep-th0002127.html | ar5iv | text | # 1 Introduction
## 1 Introduction
The four dimensional $`๐ฉ=2`$ supersymmetric $`SU(2)`$ Yang-Mills theory with $`N_f=4`$ is a theory with full $`SL(2,๐น๐น)`$ duality symmetry and global $`so(8)`$ symmetry . The duality symmetry has nontrivial implications for the dyonic spectrum of the theory. In fact the dyons fall into the $`\mathrm{๐}_๐ฏ,\mathrm{๐}_๐ฌ`$ or $`\mathrm{๐}_๐`$ representations according to their dyonic charges $`(p,q)`$. In addition, duality transformations permute such representations. This action can be characterized nicely as a homomorphism from the duality group $`SL(2,๐น๐น)`$ to the permutation group on three objects, the group of graph automorphisms of the $`so(8)`$ Dynkin diagram.
This quantum field theory can be viewed as living on a D3 brane probe of a 7-brane background . The 7-brane background includes the six 7-branes that compose the $`D_4`$ singularity of the Kodaira classification. Both the duality symmetry of the theory, and its interaction with the $`D_4`$ symmetry carried by the branes can be understood simply in terms of the 7-brane background . It is the purpose of the present paper to give a precise definition of the duality group associated to an arbitrary 7-brane configuration, and a description of its interaction with the Lie algebra carried by the configuration. Partial results were given in where the duality group of a brane configuration with monodromy $`KSL(2,๐น๐น)`$ was thought to be simply the stabilizer subgroup of $`K`$ in $`SL(2,๐น๐น)`$. While correct for many cases, this is not always true and further conditions must be satisfied.
The duality symmetry of the four-dimensional effective theory is the remnant of the $`SL(2,๐น๐น)`$ symmetry of IIB string theory. Consider the simplest case of a D3 in the vicinity of D7-branes. The three-brane is $`SL(2,๐น๐น)`$-symmetric while the D7-branes are left invariant by the subgroup generated by $`TSL(2,๐น๐น)`$; therefore the effective theory on the D3 still carries this subgroup as a duality symmetry (a trivial one because the theory does not have magnetic states). When 7-branes of different charges are involved only $`11`$ and $`11SL(2,๐น๐น)`$ leave each brane invariant. Those nontrivial transformations, however, which map the charges of the 7-branes to each other, do act as a duality symmetry of the effective theory because the transformed configuration is indistinguishable from the original one (although the position in the moduli space changes in general). The subtle task of identifying those $`SL(2,๐น๐น)`$ matrices which permute the charges of a given 7-brane configuration will be one subject of our present work. The difficulty arises because there is a large redundancy in characterizing a background by a list of 7-brane charges. To this end we will define the equivalence classes of configurations along the lines of : equivalent 7-brane setups are related to each other by the process of moving branch cuts of 7-branes through each other. Then we look for $`SL(2,๐น๐น)`$ transformations which map a given 7-brane configuration to another one in its equivalence class.
Having found the โunbroken partโ of $`SL(2,๐น๐น)`$ one should identify how this duality symmetry acts on the spectrum of the theory. As is well known from the Seiberg-Witten example, the duality group acts through the automorphism group of the root lattice of the algebra carried by the 7-branes. Such automorphisms are of two kinds: those related to the automorphisms of the Dynkin diagram, and those arising from Weyl reflections. In elucidating the general theory we have found it necessary to identify how Weyl symmetries are represented on the 7-brane configurations.
The states of D3 brane probe theories are strings or string junctions stretched between the D3 and (some of the) background 7-branes. When the 7-branes are on top of each other, the 7+1 dimensional low energy model is a Yang-Mills theory with some gauge algebra $`๐ข`$ while the spectrum of the 4-dimensional D3-theory furnishes a representation of $`๐ข`$ which is a global symmetry. When the 7-branes do not coincide (and many times they cannot) the states are no longer degenerate but $`๐ข`$ is still a useful spectrum generating algebra. In particular states constituting the orbits of the Weyl group of $`๐ข`$ are of the same multiplicity. Indeed, Weyl reflections derive from an ambiguity in the choice of a base of simple roots, and states related by Weyl reflections are really physically equivalent.
How do Weyl transformations act on junctions? Consider a configuration of D7-branes. Fig. 1(a) shows two strings, $`๐ฌ_1`$ and $`๐ฌ_2`$ which are in the fundamental representation of $`su(n)`$ and are mapped to each other by the Weyl reflection with the root corresponding to $`\alpha `$. This is the case because using familiar intersection rules and composition of open strings we see that $`๐ฌ_1+(๐ฌ_1๐ถ)๐ถ=๐ฌ_1๐ถ=๐ฌ_2`$. Physically, this Weyl reflection is implemented by the exchange process of the two D7-branes (see Fig. 1(b)), giving us an identical theory in which the two states $`๐ฌ_1`$ and $`๐ฌ_2`$ are replaced with each other. The cases involving mutually nonlocal 7-branes like the one in Fig. 1(c) are treated similarly except that one should be careful when exchanging 7-branes because the path along which they are moved is relevant. Anticipating a more extended discussion later in the paper, we note that the rules for transpositions of 7-branes imply that the transformation shown in Fig. 1(d) leaves the brane configuration invariant. In the process of transposition the two junctions $`๐_1`$ and $`๐_2=๐_1+๐ถ(๐_1๐ถ)`$ can be shown to transform into each other as a consequence of prong creation taking place when 7-branes move through strings. Thus $`๐_1`$ and $`๐_2`$, while not equivalent junctions, represent states that are physically indistinguishable. We show generally that Weyl reflections act on junctions as automorphisms of the junction lattice generated by exchanging 7-branes through a path along which they are mutually local.
Weyl and duality transformations are different in nature: the former is an automorphism of the junction lattice mapping states into states with the same asymptotic charges, while the latter, acting through an element of $`SL(2,๐น๐น)`$ can change the asymptotic charges of junctions. In addition to this action, duality transformations are characterized by the kind of transformations they generate on the Lie-algebraic data characterizing the junctions. We show that these amount to automorphisms of the root lattice $`Q`$ of the Lie algebra carried by the 7-branes. Sometimes duality transformations give Weyl transformations of the root lattice, in which case they relate states with possibly different asymptotic charges, but having essentially equivalent Lie-algebraic data. On the other hand, some duality transformations map to automorphisms that correspond to symmetries of the Dynkin diagram. Such transformations are never Weyl reflections, and relate states in different Weyl orbits or different representations.
We define the duality group $`๐`$ of a brane configuration as the subgroup of $`SL(2,๐น๐น)`$ whose effect on the brane configuration can be undone by crossing transformations. A 7-brane background carrying a Lie algebra $`๐ข`$ has a duality group $`๐`$ that interacts nontrivially with $`๐ข`$ when the symmetry group $`\mathrm{\Gamma }`$ of the Dynkin diagram of $`๐ข`$ is nontrivial. We will see that in general there is no canonical map $`๐\mathrm{\Gamma }`$ because of the surprising fact that invariant transpositions, transpositions that leave the brane configuration invariant, do not always act as Weyl transformations on the Lie algebraic data of the junction. Sometimes they may act as some outer automorphism of the root lattice. Since $`\mathrm{\Gamma }`$ precisely represents outer automorphisms, the lack of a canonical map simply results because after an $`SL(2,๐น๐น)`$ duality, the restoring crossing transformation is ambiguous up to invariant transpositions.
If the action of the set invariant transpositions on the Lie algebraic data induces a Weyl transformation then the above map is well defined. This is the case for all the configurations realizing finite algebras. Particularly interesting are the cases of configurations carrying affine exceptional symmetries . In these cases the Dynkin diagrams typically have nontrivial automorphism groups and therefore dualities interact rather nontrivially with the Lie algebra data. Also for these configurations not all invariant transpositions induce Weyl transformations on the Lie algebraic data of a junction. Therefore for these configurations the above map is only well defined with respect to a particular set of transpositions used to undo the $`SL(2,๐น๐น)`$ transformations. In this paper we compute the duality groups $`๐`$ and the maps to the symmetry groups $`\mathrm{\Gamma }`$ of Dynkin diagrams for all brane configurations that can be localized in a compactification of type IIB string theory (see ).
Let us mention two questions that we have not discussed in this paper. While the action of dualities on root junctions determines fully the action of duality on weight junctions for finite algebras, in the case of affine Kac-Moody algebras more work is necessary to understand how dualities act on junctions that represent general weight vectors. Thus, for the affine exceptional configurations discussed here, our results are restricted to junctions in the root lattice. Second, we have not characterized the duality action for the configuration giving the Lie algebra $`\widehat{E}_9`$. Since this is not an affine Kac-Moody algebra the automorphism group of its root lattice appears to be unfamiliar.
In this paper, as in our previous ones we have focused on the symmetry aspects of 7-brane configurations. Another line of works dealing mostly with tests of F-theory/heterotic duality has been reviewed in and some recent works of interest are .
The paper is arranged as follows. In section 2 we define equivalent 7-brane backgrounds and the duality group of a 7-brane configuration. Section 3 introduces the Weyl transformations in 7-brane language. In section 4 we discuss the homomorphism between the duality (effective) group and the automorphism group of the Dynkin diagram of the underlying algebra. In sections 5 and 6 we systematically compute duality groups and homomorphisms for 7-brane configurations corresponding to finite and affine algebras, respectively.
## 2 Brane configurations, duality and crossing transformations
The purpose of this section is to give a precise definition and begin the characterization of the duality group $`๐(w)`$ of a configuration $`w`$ of 7-branes. The duality group is a subgroup of $`SL(2,๐น๐น)`$ leaving the brane configuration invariant in a sense we describe in detail. Roughly, an $`SL(2,๐น๐น)`$ element belongs to the duality group when its action on the branes can be undone by a crossing transformation, that is successive operations of brane transpositions. We show that both the monodromy $`K(w)`$ of a configuration and the $`(11)`$ matrix belong to the duality group $`๐(w)`$. We also examine how crossing transformations act on junctions, and that leads to the understanding of how dualities act on junctions. In particular we show that the duality transformation $`K(w)`$ can be defined to leave invariant junctions having zero asymptotic charge. The general characterization of duality groups by their action on junctions is left for the next section.
Our presentation here will be formal for the sake of precision and brevity. Certainly the idea of the duality group of a brane configuration is not new. The duality groups of Seiberg-Witten $`N=2`$ SYM theories, in particular those of the $`N_f4`$ $`SU(2)`$ theories have been studied in detail and correspond to the duality groups of the 7-brane configurations where these four-dimensional theories arise on a D3 brane probe. A definition of the duality group for $`D_{N>4}`$ and $`E_6,E_7,E_8`$ was given in . The present one is a refined version of that definition, applicable to other brane configurations as well.
We define what we mean by a 7-brane configuration keeping the canonical picture in mind: the branes are arranged along the real axis with their branch cuts going downwards and the monodromies are listed in the order the branes appear from left to right .
### 2.1 Definitions and properties
* Definition: 7-brane configuration or simply a configuration is defined to be a word of $`SL(2,๐น๐น)`$ matrices conjugate to $`T^1`$. That is, a 7-brane is characterized by the monodromy of the axion-dilaton field $`\tau `$ around it. Via introduction of branch cuts $`\tau `$ is made single valued and $`[A]`$ stands for a 7-brane with a cut where $`\tau `$ jumps by $`ASL(2,๐น๐น)`$. To define a 7-brane background we list the branes with cuts going downwards from left to right. The set of all configurations consisting of $`n`$ 7-branes we denote as $`๐_n`$:
$`๐_nw`$ $`=`$ $`[A_1][A_2]\mathrm{}[A_n]A_i=g_iT^1g_i^1,g_iSL(2,๐น๐น),`$
$`K(w)`$ $``$ $`A_n\mathrm{}A_2A_1SL(2,๐น๐น).`$ (2.1)
$`K(w)`$ is the overall monodromy associated with the configuration $`w`$. The purpose of the square bracket is to distinguish words of matrices from products of them, i.e. $`[A_1A_2]๐_1`$ denotes a one-letter word of the matrix $`A=A_1A_2`$ while $`[A_1][A_2]๐_2`$ is a two letter word made of these two.
* Definition: $`SL(2,๐น๐น)`$ action on $`๐_n`$. $`SL(2,๐น๐น)`$ symmetry of IIB transforms $`\tau `$ and as a consequence the monodromy of a 7-brane is conjugated by this transformation $`gSL(2,๐น๐น)`$. The image of a configuration is simply the word of the transformed matrices:
$`SL(2,๐น๐น)g:๐_n๐_n[A_1][A_2]\mathrm{}[A_n][gA_1g^1][gA_2g^1]\mathrm{}[gA_ng^1].`$ (2.2)
* Definition: Transposition of 7-branes. The position of the branch cuts of the 7-branes are unphysical and they can be relocated by performing an immaterial $`SL(2,๐น๐น)`$ transformation on all physical parameters in a selected region. In particular the relative order of the cuts can be changed but the 7-branes through which a branch cut is moved are subject to the same $`SL(2,๐น๐น)`$ transformation. The elementary change of a configuration is when a cut of a brane is moved through its left hand ($`P`$) or right hand ($`P^1`$) neighbor. The transposition of the $`m`$-th and $`m+1`$-th letter of a word is thus performed by the following rule :
$`\begin{array}{ccc}P_m:\hfill & \mathrm{}[A_m][A_{m+1}]\mathrm{}\mathrm{}[A_{m+1}][A_{m+1}A_mA_{m+1}^1]\hfill & \\ & & \\ P_m^1:\hfill & \mathrm{}[A_m][A_{m+1}]\mathrm{}\mathrm{}[A_m^1A_{m+1}A_m][A_m]\mathrm{}.\hfill & \end{array}`$ (2.6)
Notice that $`A_m^1A_{m+1}A_m`$ is conjugate to $`T^1`$ if $`A_{m+1}`$ is, therefore the result is indeed in $`๐_n`$. Also note that in general $`P_m^2`$ is non trivial.
* Property 1: The transpositions satisfy the following Braid group relations:
$`P_mP_m^{}=P_m^{}P_m,\text{if}|mm^{}|>1,P_mP_{m+1}P_m`$ $`=`$ $`P_{m+1}P_mP_{m+1}.`$ (2.7)
* Definition: The Group of crossing transformations $`_N`$ (corresponding to branch cut moves) is defined by its action on $`๐_n`$. This group is generated by all transpositions $`\{P_m\}_{m=1}^{n1}`$ subject to the constraint (2.7).
* Property 2: The $`SL(2,๐น๐น)`$ transformations as defined in (2.2) commute with the transpositions. If $`g`$ denotes an $`SL(2,๐น๐น)`$ transformation and $`b`$ a series of transpositions, we have $`g(b(w))=b(g(w))`$. To prove this it suffices to examine the case when $`b`$ is a single transposition, say $`P_1`$:
$`\begin{array}{ccccc}[A_1][A_2]& & \stackrel{P_1}{}& & [A_2][A_2A_1A_2^1]\\ & & & & \\ g& & & & g\\ & & & & \\ [gA_1g^1][gA_2g^1]& & \stackrel{P_1}{}& & [gA_2g^1][gA_2A_1A_2^1g^1]\end{array}`$ (2.13)
* Definition: Equivalence group. As explained before, if two 7-brane configurations differ by either an overall $`SL(2,๐น๐น)`$ or by crossing transformations, they are physically identical. The need for distinguishing between configurations up to this equivalence motivates the following definition. The equivalence group of $`๐_n`$ is the direct product of the two groups: $`_nSL(2,๐น๐น)\times _n`$. The action of elements of $`_n`$ on $`๐_n`$ is well-defined due to the commutativity of the actions of the two factors, $`(g,b)w=g(b(w))=b(g(w))`$. The product on this set is defined by $`(g_1,b_1)(g_2,b_2)=(g_1g_2,b_1b_2)`$ where $`g_iSL(2,๐น๐น),b_i_n`$.
Acting with an element of $`๐_n`$ changes the 7-branes in general. There is however a subgroup of $`๐_n`$ which leaves the configuration invariant and thus acts as a symmetry of the D3 probe theory. Let us therefore define
* Definition: The symmetry group $`๐ฎ(w)_n`$ of the configuration $`w`$ is given by
$`๐ฎ(w)\{(g,b)_n|(g,b)w=w\}.`$ (2.14)
$`๐ฎ`$ is manifestly a group, indeed a very large one because for a given $`g`$ there typically are many choices of $`b`$ satisfying the condition $`(g,b)w=w`$. For a given $`g`$ the transformation $`b`$ is not unique because there are crossing transformations $`b`$ that leave $`w`$ invariant. These form a normal subgroup $`(w)`$ of $`๐ฎ(w)`$:
$`(w)\{(11,b)๐ฎ(w)\}.`$ (2.15)
The subgroup of the IIB duality group $`SL(2,๐น๐น)`$ which preserves $`w`$ consists of all those elements $`gSL(2,๐น๐น)`$ appearing in $`๐ฎ(w)`$. This means forgetting about the compensating transformation $`b`$ and thus leads one to define:
* Definition: The duality group $`๐(w)`$ of a configuration $`w`$ is defined as
$`๐(w)๐ฎ(w)/(w)\{gSL(2,๐น๐น)|b_n\text{such that}(g,b)w=w\}.`$ (2.16)
This is clearly the subgroup of $`SL(2,๐น๐น)`$ whose elements leave $`w`$ invariant after the action of a suitable crossing transformation.
* Proposition: $`K(w)๐(w)`$: the duality group necessarily contains the overall monodromy. In addition, $`11๐(w)`$.
* Proof: Consider the configuration $`w=[A_1][A_2]\mathrm{}[A_n],`$ with $`K_n=K(w)=A_n\mathrm{}A_1`$. Also define $`K_i=A_i\mathrm{}A_1`$, for $`1in`$. We will explicitly construct an element $`b`$ of $`_n`$ in terms of transpositions which satisfies $`(K_n,b)w=w`$. Let us perform the $`SL(2,๐น๐น)`$ transformation with $`K_n^1`$:
$`[A_1][A_2]\mathrm{}[A_n]`$ $`\stackrel{K_n^1}{}`$ $`[K_n^1A_1K_n][K_n^1A_2K_n]\mathrm{}[K_n^1A_nK_n].`$ (2.17)
As the first step, apply the product $`(P_1\mathrm{}P_{n1})_N`$ to the rhs of (2.17): this corresponds to moving the cut of the rightmost brane through the rest of them; then repeat this process $`n1`$ times. We claim that not only the order of the branes is restored but their charges are transformed back to the original ones:
$`\stackrel{P_1\mathrm{}P_{n1}}{}`$ $`[K_n^1A_nK_n][K_n^1A_nA_1A_n^1K_n]\mathrm{}[K_n^1A_nA_{n1}A_n^1K_n]`$ (2.18)
$`=[K_{n1}^1A_nK_{n1}][K_{n1}^1A_1K_{n1}]\mathrm{}[K_{n1}^1A_{n1}K_{n1}]`$
$`\stackrel{P_1\mathrm{}P_{n1}}{}`$ $`[K_{n1}^1A_{n1}K_{n1}][K_{n1}^1A_{n1}A_nA_{n1}^1K_{n1}]\mathrm{}[K_{n1}^1A_{n1}A_{n2}A_{n1}^1K_{n1}]`$
$`=[K_{n2}^1A_{n1}K_{n2}][K_{n2}^1A_nK_{n2}]\mathrm{}[K_{n2}^1A_{n2}K_{n2}]`$
$`\mathrm{}`$
$`\stackrel{P_1\mathrm{}P_{n1}}{}`$ $`[A_1][A_2]\mathrm{}[A_n].`$
Thus we proved that the overall monodromy is indeed in the duality group, that is
$`(K(w),(P_1P_2\mathrm{}P_{n1})^n)w=(K(w)^1,(P_1P_2\mathrm{}P_{n1})^n)w=w.`$ (2.19)
* For the transformation $`g=11`$ we simply note that this transformation does not change the word describing the configuration since for each brane $`[A_i][gA_ig^1]=[A_i]`$. This transformation is clearly in the duality group and since it plainly leaves the 7-brane configuration invariant it does not need to be accompanied by brane transpositions ($`b`$ can be taken to be the identity in (2.16)).
### 2.2 The action on invariant charges
Having seen how the $`SL(2,๐น๐น)`$ equivalence actions transform the 7-brane configurations we would like to know the fate of junctions ending on these 7-branes. A junction is characterized by its invariant charges; the effective number of prongs on each 7-brane of the configuration . The action of an overall $`SL(2,๐น๐น)`$ transformation is trivial: the charges of each string composing the junction transform as a doublet but the invariant charges do not change.
In general when one performs a crossing transformation, not only the 7-brane labels change but the invariant charges on those branes change as well. This is most easily seen from the active viewpoint: instead of relocating the cuts moving them through the 7-branes, we can move the 7-branes. Whenever a 7-brane in motion crosses a string segment, additional prongs on that brane might be created.
It suffices to determine how the charges change in the transposition of two consecutive 7-branes, more complicated cases are considered by successive transpositions. Consider therefore a junction on a pair of 7-branes $`[K_{๐ณ_1}][K_{๐ณ_2}][๐ณ_1][๐ณ_2]`$ with invariant charges $`[Q_1][Q_2]`$. Here we label the branes, as in with their charge vector $`๐ณ=(p,q)`$, in terms of which the corresponding $`SL(2,๐น๐น)`$ monodromy matrix is written as $`K_๐ณ11+\mathrm{๐ณ๐ณ}^TS`$. According to (2.6) the charges of the branes transform as ($`s_{12}๐ณ_1\times ๐ณ_2=p_1q_2q_1p_2`$):
$`\begin{array}{ccccccc}P\hfill & :\hfill & [๐ณ_1][๐ณ_2]\hfill & \hfill & [๐ณ_2][K_{๐ณ_2}๐ณ_1]\hfill & =\hfill & [๐ณ_2][๐ณ_1+s_{12}๐ณ_2]\hfill \\ P^1\hfill & :\hfill & [๐ณ_1][๐ณ_2]\hfill & \hfill & [K_{๐ณ_1}^1๐ณ_2][๐ณ_1]\hfill & =\hfill & [๐ณ_2+s_{12}๐ณ_1][๐ณ_1].\hfill \end{array}`$ (2.22)
The action on the invariant charges can be determined by looking at how many prongs are created/annihilated on each 7-brane in the canonical presentation , but charge conservation alone gives the answer too:
$`\begin{array}{ccccccc}P\hfill & :\hfill & [Q_1][Q_2]\hfill & \hfill & [Q_2s_{12}Q_1][Q_1]\hfill & & \\ P^1\hfill & :\hfill & [Q_1][Q_2]\hfill & \hfill & [Q_2][Q_1s_{12}Q_2].\hfill & & \end{array}`$ (2.25)
Appending the invariant charges as superscripts to the branes, the complete action for both transpositions is:
$`[๐ณ_1]^{Q_1}[๐ณ_2]^{Q_2}`$ $`\stackrel{P}{}`$ $`[๐ณ_2]^{Q_2s_{12}Q_1}[๐ณ_1+s_{12}๐ณ_2]^{Q_1}`$ (2.26)
$`[๐ณ_1]^{Q_1}[๐ณ_2]^{Q_2}`$ $`\stackrel{P^1}{}`$ $`[๐ณ_2+s_{12}๐ณ_1]^{Q_2}[๐ณ_1]^{Q_1s_{12}Q_2}.`$ (2.27)
The crossing transformation that restores the original brane configuration after application of the monodromy $`K`$ (see (2.19)) has an important property: it does not change the invariant charges of junctions with zero asymptotic charge. This is best seen by visualizing the motion of the 7-branes. In Fig. 2 we show the effect of the cyclic transformation $`P_1\mathrm{}P_{n1}`$: it corresponds to moving the rightmost 7-brane above the rest to the left. If a junction has no asymptotic charge, this 7-brane does not cross any string segment along the transformation and therefore the invariant charge of any of 7-branes does not change. Performing this transformation $`n`$ times rearranges the 7-branes in the original order and although their monodromies change, the invariant charges of any given junction without asymptotic charge remain the same.
We note in passing that while the $`g=11๐(w)`$ transformation does not change the list of branes, the invariant charges of any junction will change sign.
## 3 Implementing Lie algebra Weyl transformations
In section 2 we introduced the group $`(w)`$ of invariant crossing transformations. The elements of this group are crossing transformations $`b`$ that leave the brane configuration $`w`$ invariant, namely, $`bw=w`$. It is the purpose of this section to understand some of the structure of $`(w)`$.
We know that invariant crossing transformations act on junctions by shuffling their invariant charges, and therefore act on weight vectors of the Lie algebra carried by the branes. We will show that Weyl transformations of the Lie algebra can always be implemented at the level of junctions by invariant crossing transformations and this will take most of the work in the present section. Nevertheless, there are sometimes invariant transformations that do not correspond to Weyl transformations, but rather correspond to outer automorphisms of the root lattice of the algebra. This we found to be a surprise. We will give here a nontrivial example of this phenomenon.
It is useful to define the subgroup $`_W(w)`$ called the group of invariant crossing transformations of Weyl type. A transformation is said to belong to $`_W(w)`$ if its action on weights corresponding to junctions is a Weyl transformation. There are some special invariant transformations that do not change junctions at all. Such transpositions belong to $`_W(w)`$ since they imply a trivial identity Weyl transformation. We give an example of such transposition.
As we will see later in this paper, whenever $`_W(w)`$ coincides with $`(w)`$, namely, if all invariant transpositions are of Weyl type, the characterization of duality groups is very much simplified. This will be the case for finite algebras, but not the case for affine ones.
### 3.1 Weyl transformations as invariant crossing transformations
In this section our main objective is to prove that for any Weyl transformation of the Lie algebra $`๐ข(w)`$ carried by a 7-brane configuration $`w`$ there is a crossing transformation which implements this Weyl transformation at the level of junctions. This crossing transformation leaves the brane configuration invariant and simply acts on the invariant charges of junctions supported on the configuration. This action is such that the associated weight vectors undergo the desired Weyl transformation. We restrict our attention to the finite and affine Kac-Moody algebras that can be obtained on localizable brane configurations.
We begin by noting that Weyl transformations are generated by elementary reflections using the simple roots of the algebras in question. We also recall that each simple root junction of the $`๐^N,๐_{0N3},๐_{N0},๐_{6N8}`$ and $`\widehat{๐}_{0N8}`$ configurations can be represented by an open string $`๐ถ`$ ($`๐ถ^2=2`$) stretched between two possibly mutually non-local 7-branes.
We now claim that there is a rearrangement of the configuration, equivalent to a crossing transformation $`b`$, such that the root $`๐ถ`$ in question becomes an open string $`b(๐ถ)`$ stretched between two mutually local 7-branes. This rearrangement corresponds to moving one of the two 7-branes on which $`๐ถ`$ ends along the path of $`๐ถ`$ itself until the brane is just to the side of other brane. This motion, for a particular example, is shown in Fig. 3. One can imagine the branch cut of the moving brane staying vertical, and one can see that indeed this motion simply corresponds to a sequence of transpositions of branes. After this motion, however, many of the branes of the configuration may have changed identity, and therefore the configuration has not been left invariant.
We now claim that the elementary Weyl reflection using the root $`๐ถ`$ is obtained by first doing the transpositions in $`b`$, then doing the transposition $`P_i`$ that interchanges the two mutually local 7-branes supporting the now short open string $`b(๐ถ)`$, and then using the brane that sits where the first one ended to retrace backwards the path, this is simply done by applying $`b^1`$. Two facts should be noted. First, the original open string $`๐ถ`$ changes direction, thus becoming $`๐ถ`$ as one would expect for a Weyl transformation generated by $`๐ถ`$. Second, by retracing the path, all changes of brane labels that the first tracing of the path caused are compensated and the brane configuration is now left invariant. It is now left to show that this sequence of operations $`(b^1P_ib)`$ performs the expected Weyl reflection on arbitrary junctions supported on the configuration. That is,
$$\text{W}_\alpha =b^1P_ib:๐๐+(๐๐ถ)๐ถ.$$
(3.1)
We first show explicitly that this formula holds in the special case when $`๐ถ`$ is an open string stretching between two adjacent mutually local branes $`[๐ณ_๐ข]`$ and $`[๐ณ_{i+1}]`$. In this case $`๐ถ=๐ณ_๐ข๐ณ_{๐ข+\mathrm{๐}}`$, and $`W_\alpha =P_i`$, just the exchange of the two branes $`[๐ณ_๐ข]`$ and $`[๐ณ_{i+1}]`$. Consider now a general junction supported on the configuration
$$๐=Q_k๐ณ_k=๐ณ_๐ขQ_i+๐ณ_{๐ข+\mathrm{๐}}Q_{i+1}+\mathrm{}(๐\alpha )=Q_{i+1}Q_i.$$
(3.2)
The exchange $`P_i`$ of the mutually local branes $`[๐ณ_๐ข]`$ and $`[๐ณ_{๐ข+\mathrm{๐}}]`$ maps $`๐`$ to the new junction $`๐^{}`$ defined as
$$P_i:๐๐^{}=๐ณ_๐ขQ_{i+1}+๐ณ_{๐ข+\mathrm{๐}}Q_i+\mathrm{}=๐+(๐๐ถ)๐ถ,$$
(3.3)
where use was made of the explicit expression for $`๐ถ`$ and of equation (3.2). This confirms our claim for this special case.
Let us now return to the general problem and compute the action of $`b^1P_ib`$ on a general junction as follows
$$(b^1P_ib)(๐)=b^1\left[P_i(b(๐))\right]=b^1\left[b(๐)+(b(๐)b(๐ถ))b(๐ถ)\right],$$
(3.4)
where in the last step we used (3.3) where the role of $`๐ถ`$ is played here by the junction $`b(๐ถ)`$ extending between the two mutually local branes that are transposed. Since $`b^1b`$ equals the identity on any junction, and the intersection of two junctions is invariant under crossing transformations we find
$$(b^1P_ib)(๐)=๐+(๐๐ถ)๐ถ.$$
(3.5)
This completes our proof that Weyl transformations can be realized as crossing transformations that leave invariant the brane configuration. Our realization has been very specific, and while all such crossing transformations belong to $`_W(w)`$, they do not necessarily exhaust it, as we illustrate in section 3.2.
One can wonder if there is a useful notion of Weyl reflections of the junction lattice of a configuration that makes no reference to the Lie algebra carried by the configuration. For any junction $`๐ถ`$ such that $`(๐ถ๐ถ)=2`$ we could define
$`๐ฒ_๐ถ(๐)=๐+(๐๐ถ)๐ถ.`$ (3.6)
This transformation preserves intersection numbers and it is therefore an automorphism of the junction lattice. The transformations generated by junctions $`๐ถ`$ of zero asymptotic charge correspond to the Lie algebraic Weyl transformations since such junctions are roots. On the other hand a transformation $`๐ฒ_๐ถ`$ generated by a junction $`๐ถ`$ with asymptotic charge will change the asymptotic charge of the junction on which it acts. The significance of such transformations is unclear since they generically map BPS junctions to non-BPS junctions.<sup>1</sup><sup>1</sup>1Consider the $`\widehat{๐}_9`$ junction $`๐ถ=_{i=1}^8\mu _i\omega ^i+2\omega ^p\delta ^{(0,1)}`$, with the $`E_8`$ weight satisfying $`\mu ^2=2`$. This junction has self-intersection minus two and non-zero asymptotic charge. Consider now $`๐ฒ_๐ถ`$ acting on the BPS junction $`\delta ^{(0,1)}`$. One readily finds $`๐ฒ_๐ถ(\delta ^{(0,1)})=2_{i=1}^8\mu _i\omega ^i+4\omega ^p\delta ^{(0,1)}.`$ We know that a necessary condition for a junction $`๐`$ of asymptotic charge $`(p,q)`$ to be BPS is that $`๐๐2+\text{gcd}(p,q)`$ . Since $`(๐ฒ_๐ถ(\delta ^{(0,1)}))(๐ฒ_๐ถ(\delta ^{(0,1)}))=0<2+\text{gcd}(4,0),`$ it follows that $`๐ฒ_๐ถ(\delta ^{(0,1)})`$ cannot be BPS. Thus it may be that the only useful reflections of the junction lattice are those generated by roots of the Lie algebra carried by the configuration.
### 3.2 Further examples of invariant crossing transformations
Above we presented a group of invariant crossing transformations that act on junctions via Weyl transformations of their Lie-algebraic data. This group, however, does not contain all invariant transpositions $`(w)`$ (see (2.15))of the generic 7-brane configuration. Among the additional invariant crossing transformations there are ones that are of Weyl type as well as others which are not. In the following we give an example for each.
$`๐_๐`$: For simplicity, consider first the case of $`๐_\mathrm{๐}`$. This configuration has no root junction (a junction of self-intersection minus two and zero asymptotic charge) therefore there is no Weyl group of the usual kind, and there are no transpositions of the type discussed in the previous subsection. Nevertheless there is a nontrivial crossing transformation leaving the configuration invariant and thus belonging to $`(w)`$. This transformation leaves all charges unchanged, as will be shown now. Indeed,
$`๐^{Q_A}๐^{Q_B}๐^{Q_C}\stackrel{P_1}{}๐^{Q_B+Q_A}๐_{[\mathrm{๐},\mathrm{๐}]}^{}{}_{}{}^{Q_A}๐^{Q_C}\stackrel{P_2}{}๐^{Q_B+Q_A}๐^{Q_C+Q_A}๐^{Q_A}`$
$`\stackrel{P_2}{}๐^{Q_B+Q_A}๐^{Q_C}๐_{[\mathrm{๐},\mathrm{๐}]}^{}{}_{}{}^{Q_C+Q_A}\stackrel{P_1}{}๐^{Q_A+Q_CQ_B}๐_{[\mathrm{๐},\mathrm{๐}]}^{}{}_{}{}^{Q_B+Q_A}๐_{[\mathrm{๐},\mathrm{๐}]}^{}{}_{}{}^{Q_C+Q_A}`$ (3.7)
$`\stackrel{P_1P_2P_2P_1}{}๐^{Q_A}๐_{[\mathrm{๐},\mathrm{๐}]}^{}{}_{}{}^{Q_C}๐_{[\mathrm{๐},\mathrm{๐}]}^{}{}_{}{}^{2Q_CQ_B}\stackrel{P_2}{}๐^{Q_A}๐^{Q_B}๐^{Q_C}.`$
This transposition is obtained by first moving the $`๐`$-brane around the $`\mathrm{๐๐}`$ branes twice and then moving the transformed $`๐`$-brane through the cut of the $`๐`$-brane. This brane configuration has this particular invariance because $`T^2๐(w)`$ (section 5.3), and therefore the $`SL(2,๐น๐น)`$ transformation $`T^{2k}`$ induced on the $`\mathrm{๐๐}`$ system by the $`๐`$-brane can be undone by a transposition.
Now, we can turn to the case of $`๐_๐`$, where clearly the same transformation exists leaving the configuration and all invariant charges unchanged. This transformation is trivially of Weyl type in that it acts on the junctionsโs Lie algebraic data as the identity. On the other hand, it is also clear that this transformation cannot be obtained by composition of transformations that interchange branes connected by an open string. This is the case because all such open strings in $`๐_๐`$ join $`๐`$ branes and thus the $`๐`$ and $`๐`$ branes are never interchanged.
$`\widehat{\stackrel{~}{๐}}_\mathrm{๐}`$: This 7-brane configuration also does not admit any ordinary root junction. The only BPS junctions are multiples of the delta junction $`\delta `$ that encircles the configuration. Having no real roots we have no transformations generated by open strings. Surprisingly, there is an invariant transformation that actually changes the Lie algebraic data of junctions. In doing so it shows that elements of $`(w)`$ may generate in general nontrivial automorphisms that are not of Weyl type. It would be of interest to exhibit those explicitly for affine configurations having a nontrivial Weyl group, but we focus here our attention to the case of $`\widehat{\stackrel{~}{๐}}_\mathrm{๐}`$. The transformation in question is actually analogous to the one just considered above. We first perform the following transpositions:
$`๐^{Q_A}๐_{[2,1]}^{Q_1}๐_{[1,2]}^{Q_2}๐_{[1,1]}^{Q_3}\stackrel{P_1P_2P_3P_3P_2P_1}{}๐^{Q_AQ_1+2Q_2+Q_3}๐_{[3,1]}^{Q_1+Q_A}๐_{[3,2]}^{Q_2+Q_A}๐_{[0,1]}^{Q_3Q_A}.`$ (3.8)
This transposition consists of moving the $`๐`$-brane around the other 7-branes once in the counter clockwise direction. The linear transformation $`g`$ induced on the charges is given by
$`\left(\begin{array}{cc}Q_A& \\ Q_1& \\ Q_2& \\ Q_3& \end{array}\right)\left(\begin{array}{cccc}1& 1& 2& 1\\ 1& 1& 0& 0\\ 1& 0& 1& 0\\ 1& 0& 0& 1\end{array}\right)\left(\begin{array}{c}Q_A\\ Q_1\\ Q_2\\ Q_3\end{array}\right).`$ (3.21)
The effect on the charges of moving the $`๐`$-brane around the other 7-branes three times is given by the linear transformation $`g^3`$. We can restore the original configuration by transposition $`P_2P_3`$,
$`๐^{Q_A}๐_{[2,1]}^{Q_1}๐_{[1,2]}^{Q_2}๐_{[1,1]}^{Q_3}\stackrel{(P_1P_2P_3P_3P_2P_1)^3}{}๐^{Q_A^{\prime \prime \prime }}๐_{[5,1]}^{Q_1^{\prime \prime \prime }}๐_{[7,2]}^{Q_2^{\prime \prime \prime }}๐_{[2,1]}^{Q_3^{\prime \prime \prime }}\stackrel{P_2P_3}{}๐^{\widehat{Q}_A}๐_{[2,1]}^{\widehat{Q}_1}๐_{[1,2]}^{\widehat{Q}_2}๐_{[1,1]}^{\widehat{Q}_3}.`$
The transformed charges $`\widehat{Q}`$ are
$`\widehat{Q}_A=Q_A^{\prime \prime \prime },\widehat{Q}_1=Q_3^{\prime \prime \prime }+3Q_2^{\prime \prime \prime }3Q_1^{\prime \prime \prime },\widehat{Q}_2=Q_1^{\prime \prime \prime },\widehat{Q}_3=Q_2^{\prime \prime \prime },`$ (3.22)
where $`\stackrel{}{Q}^{\prime \prime \prime }`$ are obtained from $`\stackrel{}{Q}`$ by the linear transformation $`g^3`$. Using (3.21) we get
$`\widehat{Q}_A`$ $`=`$ $`Q_A3Q_1+6Q_2+3Q_3,`$
$`\widehat{Q}_1`$ $`=`$ $`3Q_A+3Q_2+2Q_3,`$
$`\widehat{Q}_2`$ $`=`$ $`3Q_A2Q_1+6Q_2+3Q_3,`$ (3.23)
$`\widehat{Q}_3`$ $`=`$ $`3Q_A+3Q_17Q_23Q_3.`$
Thus we see that $`(11,P_2P_3(P_1P_2P_3P_3P_2P_1)^3)(w)`$ has a nontrivial action on the charges. Since there are no real root junctions with support on this configuration this element of $`(w)`$ is not of Weyl type. The motion of the $`๐`$-brane around the other three branes has the effect equivalent to the action of a global $`SL(2,๐น๐น)`$ transformation by $`T^3`$. We will discuss this action in more detail in section 6.
## 4 Duality groups and Dynkin graph automorphisms
The results of the previous section allow us to find and characterize the duality group $`๐(w)`$ of a brane configuration $`w`$ with monodromy $`K(w)`$. Since crossing transformations cannot change the monodromy $`K(w)`$ of the configuration, any element of the duality group must leave the monodromy invariant and therefore is contained in Stab$`(K)`$, the stabilizer of $`K`$ in $`SL(2,๐น๐น)`$ ($`g\text{Stab}(K)gKg^1=K`$)
$$๐(w)\text{Stab}(K(w)).$$
(4.1)
The duality group $`๐(w)`$ of the configuration will be the subgroup of Stab$`(K)`$ that leaves the configuration invariant in the sense discussed in the previous section (see (2.16)). In any concrete case it is relatively straightforward to determine the group Stab$`(K)`$. Then one must select $`๐`$ by finding the subgroup of Stab$`(K)`$ for which there exist crossing transformations that restore the configuration.
We have seen that $`K๐`$ (Proposition, sect. 2.1) and therefore $`\{K\}`$, the group generated by $`K`$, is a subgroup of the duality group $`๐`$. Since $`\{K\}`$ is a normal subgroup of Stab$`(K)`$, it is also a normal subgroup of $`๐`$. We are thus led to define the quotient group
$$\overline{๐}๐/\{K\},$$
(4.2)
referred to as the reduced duality group, that will play an important role in the computations. Another general fact discussed before is that the transformation $`11SL(2,๐น๐น)`$, clearly contained in Stab$`(K)`$ is also an element in $`๐`$.
Consider now the action of an element $`(g,b)๐ฎ(w)`$ of the symmetry group of the brane configuration. The set of transpositions $`b`$ that restore the original configuration via $`(g,b)w=w`$ will shuffle the invariant charges of junctions and therefore this symmetry maps junctions to (typically) different junctions. One nevertheless gets an automorphism of the junction lattice $`\mathrm{\Lambda }_๐`$; namely for any two junctions $`๐_\mathrm{๐}`$ and $`๐_\mathrm{๐}`$ mapping to $`๐_1^{}`$ and $`๐_2^{}`$ one has $`๐_1๐_2=๐_1^{}๐_2^{}`$. We therefore have a map
$$๐ฎ\text{Aut}(\mathrm{\Lambda }_๐),$$
(4.3)
from the symmetry group of the brane configuration to the group of automorphisms of the junction lattice.
Let us focus for the moment on junctions representing roots of the Lie algebra associated to the brane configuration. In fact, more generally, consider junctions associated to elements in the root lattice $`Q`$ of the Lie algebra. Such junctions, as discussed in length in earlier papers have zero asymptotic charges. It follows that symmetry transformations in $`๐ฎ`$ will map these junctions among themselves. In addition, for such junctions $`๐_1๐_2=\lambda _1\lambda _2`$, where $`\lambda _1,\lambda _2Q`$ are the associated elements in the root lattice. Since duality elements map to automorphisms of the junction lattice, by restricting to junctions associated to $`Q`$ duality elements map to automorphisms of the root lattice $`Q`$. We therefore have:
$$๐ฎ\text{Aut}(Q).$$
(4.4)
On the other hand there is no canonical map $`๐`$ Aut$`(\mathrm{\Lambda }_๐)`$, nor there is a canonical map $`๐`$ Aut$`(Q)`$ from the duality group of the configuration. This is so because for duality elements, the compensating crossing transformation in $`(w)`$ used to restore the brane configuration is not uniquely defined. As discussed in section 3, a brane configuration typically admits crossing transformations $`\widehat{b}`$ that leave it invariant. In fact, any Weyl reflection of the root lattice is generated on junctions by a crossing transformation $`\widehat{b}`$ that leaves the configuration invariant. In the language of (2.16) the ambiguous action of a duality transformation $`g`$ on junctions arises because the computation of such action requires the choice of some $`b`$ such that $`(g,b)w=w`$. But $`b`$ is ambiguous, if $`b`$ satisfies this equation, $`b\widehat{b}`$ does as well.
If a configuration $`w`$ has the property that $`(w)=_W(w)`$, namely, every invariant transposition is of Weyl type, then the map
$`๐\text{Aut}(Q)/W.`$ (4.5)
is well defined. This is the case because all invariant transpositions map to the Weyl group. Since the duality elements $`\{K\}`$ always map to Weyl group we also have the well defined map
$`\overline{๐}\text{Aut}(Q)/W.`$ (4.6)
In the above homomorphisms the quotient group to the right is well defined since $`W`$ is a normal subgroup of $`\text{Aut}(Q)`$ (for this and other facts quoted below see, ). This quotient, for algebras of finite type, is simply the (graph) automorphism group $`\mathrm{\Gamma }`$ of the Dynkin diagram (and the $`\pm `$ above is not necessary). For infinite Kac-Moody algebras the above quotient includes, in addition to the graph automorphism $`\mathrm{\Gamma }`$, the generator $`(1)`$ which changes the sign of every vector in the root lattice, and this is never a Weyl transformation. For finite algebras the transformation $`(1)`$ of $`Q`$ is many times a Weyl transformation. From the list of finite algebras we consider, $`(1)`$ is not a Weyl element for the $`A_N`$ series and for $`E_6`$. In such cases, $`(1)`$ is equivalent, up to Weyl transformations, to a graph automorphism. Since the element $`11๐SL(2,๐น๐น)`$ precisely acts as $`(1)`$ on $`Q`$ this shows that this transformation is nontrivial (i.e. not Weyl) for the $`A_N`$ series, for $`E_6`$ and for affine algebras.
In the next section we will show that $`(w)=_W(w)`$ for 7-brane configurations realizing finite algebras. This condition, however, may not satisfied for configurations realizing affine algebras, as we illustrated in section 3.2. Therefore in these cases (4.5) and (4.6) are not well defined since there is no unique choice of invariant transposition and different choices can induce non-Weyl action on the roots.
Our strategy in this case would be to define the maps $`๐\pm \mathrm{\Gamma }`$ and $`\overline{๐}\pm \mathrm{\Gamma }`$ with respect to a fixed set of invariant transpositions that can undo the effect of an $`SL(2,๐น๐น)`$ transformation. The map is thus dependent on the choice of transpositions. An invariant characterization will require better understanding of the structure of $`(w)`$.
When defined, our interest is in the homomorphisms $`๐\pm \mathrm{\Gamma }`$ and $`\overline{๐}\pm \mathrm{\Gamma }`$, the latter capturing the interplay of duality transformations with Lie algebraic data. Dualities in $`๐`$ that map to nontrivial elements of $`\mathrm{\Gamma }`$ relate junctions appearing in different representations or junctions appearing as vectors in $`Q`$ that are not related by Weyl transformations. For the finite algebras, we shall find cases when the map $`\varphi :\overline{๐}\pm \mathrm{\Gamma }`$ is an isomorphism (for $`E_6`$, for example), and cases when it is onto but not one to one ($`D_4`$, for example).
We can readily find the implications for weight vectors in the case of configurations leading to finite algebras. In this situation $`๐_1๐_2=\lambda _1\lambda _2+f(p_1,q_1;p_2,q_2)`$, where $`f`$ is a quadratic form determined solely by the monodromy $`K`$ . Since duality transformations preserve $`K`$ the automorphisms of the junction lattice arising from dualities give automorphisms of the weight lattice $`\mathrm{\Lambda }`$. On the other hand for finite algebras Aut$`(Q)=`$ Aut $`(\mathrm{\Lambda })`$ since $`\mathrm{\Lambda }=Q^{}`$. It thus follows that the homomorphism $`๐\pm \mathrm{\Gamma }`$ carries information on how representations in different conjugacy classes are mapped into each other by duality transformations. We leave the question of duality action on junctions corresponding general affine weight vectors open.
## 5 Duality groups for finite-type configurations
In this section we calculate duality groups $`๐`$ and give the homomorphisms to the corresponding Dynkin-graph automorphism groups $`\mathrm{\Gamma }`$. These homomorphisms are characterized by $`\varphi :\overline{๐}\mathrm{\Gamma }`$, as discussed before. This section focuses on brane configurations of elliptic and parabolic monodromies realizing finite Lie algebras. The case of finite algebras is relatively simple to analyze. The 7-brane configurations realizing finite algebras have the property that $`(w)=_W(w)`$ i.e, all invariant transpositions are of Weyl type. To prove this consider the action of a invariant transposition on the weight vector of a junction with support on the 7-brane configuration. If the transposition is not of Weyl type it will induce an outer automorphism on the weight vector. It was shown in that for 7-brane configurations realizing finite algebras the conjugacy class of a weight vector corresponding to a junction is determined by the asymptotic charge of the junction. Since a transposition cannot change the asymptotic charge of a junction, it cannot change the conjugacy class of the corresponding weight vector. Therefore the action of the transposition on the weight vector cannot be an outer automorphism and hence must be a Weyl transformation. This simplification implies that the homomorphism $`\overline{๐}\mathrm{\Gamma }`$ is well defined for these cases. In table 1 we list the relation between the conjugacy classes and the asymptotic charge for various 7-brane configurations realizing finite algebras. In the case of $`E_8`$ since there are no outer automorphisms therefore every transpositions is trivially of Weyl type. This is consistent with the fact that there is a single conjugacy class for $`E_8`$.
Throughout this and the next section we denote by $`\{\mathrm{}\}`$ a group generated by the elements indicated by dots. In addition $`\{\mathrm{}|\mathrm{}\}`$ will denote the group generated by the elements to the left of the vertical bar, modulo the relations to the right of the bar.
### 5.1 $`๐_๐`$ configuration: $`๐^{๐+\mathrm{๐}}`$
This configuration is built from $`(N+1)`$ $`[1,0]`$ branes. The monodromy is $`K=T^{N1}`$ and $`\text{Stab}(K)=\{11,T\}`$. Since $`T`$ preserves the charges of the 7-branes, it belongs to the duality group just like $`11๐`$. Therefore $`๐(๐_๐)=\text{Stab}(K)=\{11,T\}=๐น๐น_2\times ๐น๐น`$, and $`\overline{๐}(๐_๐)=\{11,T|T^{N+1}=11\}=๐น๐น_2\times ๐น๐น_{N+1}`$.
On the other hand $`\mathrm{\Gamma }(A_N)=๐น๐น_2`$ for $`N2`$ and is generated by the transformation $`๐ช:(a_1,a_2,\mathrm{},a_n)(a_n,a_{n1},\mathrm{},a_1)`$ of the Dynkin labels. Since $`T`$ does not affect the invariant charges it leaves all Dynkin labels unchanged and therefore $`\varphi (T)=11\mathrm{\Gamma }(A_N)`$. The transformation $`11`$, however, changes the sign of all invariant charges and therefore of all the Dynkin labels
$$11:(a_1,a_2,\mathrm{},a_n)(a_1,a_2,\mathrm{},a_n).$$
(5.2)
The Weyl transformation $`W`$ which corresponds to rotating the 7-brane configuration by half a full turn can be seen to map
$$W:(a_1,a_2,\mathrm{},a_n)(a_n,a_{n1},\mathrm{},a_1).$$
(5.3)
We now recognize that the action of $`11SL(2,๐น๐น)`$ on the Dynkin labels is given by the composition of $`๐ช`$ and $`W`$. Therefore the homomorphism $`\varphi `$ from $`๐`$ to $`\mathrm{\Gamma }(A_N)`$ is defined as
$`\varphi (T)`$ $`=`$ $`+1\mathrm{\Gamma }(A_N)`$
$`\varphi (11)`$ $`=`$ $`๐ช\mathrm{\Gamma }(A_N).`$ (5.4)
In case of $`๐_\mathrm{๐}`$ the computation of the duality group is identical and therefore $`\overline{๐}(๐_\mathrm{๐})=๐น๐น_2\times ๐น๐น_2`$. On the other hand here $`\mathrm{\Gamma }(A_1)=1`$ and the $`11`$ transformation is simply a Weyl transformation, thus the homomorphism $`\varphi `$ is trivial.
### 5.2 $`๐_๐`$ configurations: $`๐^{๐+\mathrm{๐}}๐`$
Since $`๐_\mathrm{๐}=๐_\mathrm{๐}`$, and $`๐_{๐\mathrm{๐}}`$ have hyperbolic monodromies, we need only focus on the configurations $`๐_\mathrm{๐},๐_\mathrm{๐}`$ and $`๐_\mathrm{๐}`$.
* $`๐_\mathrm{๐}:`$ The monodromy $`K(๐_\mathrm{๐})(ST)^1`$ and $`\text{Stab}(K)=\{(ST)^1\}`$ ($`11\text{Stab}(K)`$ since $`(ST)^3=11`$). Since $`K๐`$, $`๐(๐_\mathrm{๐})=\text{Stab}(K)=๐น๐น_6`$, and $`\overline{๐}(๐_\mathrm{๐})=1`$. This configuration supports no junctions without asymptotic charges so there is no $`๐ข`$ associated to it.
* $`๐_\mathrm{๐}:`$ Here $`K(๐_\mathrm{๐})S^1`$ and $`\text{Stab}(K)=\{S\}`$. Since $`K๐`$, $`๐(๐_\mathrm{๐})=\text{Stab}(K)=๐น๐น_4`$ and $`\overline{๐}(๐_\mathrm{๐})=1`$. Since $`\mathrm{\Gamma }(A_1)=1`$, $`\varphi `$ is the trivial isomorphism.
* $`๐_\mathrm{๐}:`$ In this case $`K(๐_\mathrm{๐})ST`$ and $`\text{Stab}(K)=\{ST\}`$. Since $`(ST)^3=11`$, $`๐(๐_\mathrm{๐})=\text{Stab}(K)=๐น๐น_6`$ and $`\overline{๐}(๐_\mathrm{๐})=๐น๐น_2`$. Here $`\mathrm{\Gamma }(A_2)=๐น๐น_2`$, and it follows from (5.2) and (5.3) that the homomorphism $`\varphi `$ from $`\overline{๐}(๐_\mathrm{๐})`$ to $`\mathrm{\Gamma }(A_2)`$ is given by
$`\varphi (11)`$ $`=`$ $`+1\mathrm{\Gamma }(A_2),`$
$`\varphi (11)`$ $`=`$ $`๐ช\mathrm{\Gamma }(A_2),`$ (5.6)
where $`๐ช(a_1,a_2)=(a_2,a_1)`$.
### 5.3 $`๐_๐`$ configurations: $`๐^๐\mathrm{๐๐}`$
We begin with some general remarks applicable whenever $`N0`$. We will show that $`T๐`$. Indeed a $`T`$ transformation can be undone by taking the rightmost $`๐`$ brane on a round trip encircling $`\mathrm{๐๐}`$ branes:
$$๐^{๐\mathrm{๐}}๐\mathrm{๐๐}\stackrel{T}{}๐^{๐\mathrm{๐}}๐๐_{[\mathrm{๐},\mathrm{๐}]}๐_{[\mathrm{๐},\mathrm{๐}]}\stackrel{P_NP_{N+1}P_{N+1}P_N}{}๐^{๐\mathrm{๐}}๐\mathrm{๐๐}.$$
(5.8)
In doing this operation the invariant charges on the branes, denoted as $`Q_1,\mathrm{}Q_{N1}`$, for the inert $`๐`$ branes, $`Q_N`$ for the rightmost $`๐`$ brane, and $`Q_B,Q_C`$, transform as
$`Q_i`$ $`\stackrel{T}{}`$ $`Q_i,i=1\mathrm{}N1`$
$`Q_N`$ $`\stackrel{T}{}`$ $`Q_NQ_B+Q_C,`$
$`Q_B`$ $`\stackrel{T}{}`$ $`Q_B+Q_N,`$ (5.9)
$`Q_C`$ $`\stackrel{T}{}`$ $`Q_C+Q_N.`$
On the other hand, we know form eq.(6.27) of how Dynkin labels are given in terms of invariant charges:
$`a_i`$ $`=`$ $`Q_iQ_{i+1},i=1\mathrm{}N2`$
$`a_{N1}`$ $`=`$ $`Q_{N1}Q_N,`$ (5.10)
$`a_N`$ $`=`$ $`Q_{N1}+Q_N+Q_BQ_C.`$
One immediately deduces from the last two equations the action of the duality $`T`$ on the Dynkin labels
$$T:(a_1,\mathrm{},a_{N1},a_N)(a_1,\mathrm{},a_N,a_{N1}).$$
(5.11)
This exchange of the last two Dynkin labels is the familiar $`๐น๐น_2`$ automorphism of the $`D_N`$ Dynkin graph.
* $`๐_\mathrm{๐}:`$ The monodromy $`K(\mathrm{๐๐})=T^4`$ and $`\text{Stab}(K)=\{11,T\}`$. A junction of asymptotic charge $`(p,q)`$ on this configuration satisfies the condition $`p+q0(\text{mod}2)`$. After a transformation by $`T^k`$ if the branes can be brought back to the original ones by branch cut moves then the transformed asymptotic charge $`(p^{},q^{})=(p+kq,q)`$ must also satisfy the same condition. This implies that $`k0(\text{mod}2)`$. Indeed, a probe D3-brane in this background realizes $`๐ฉ=2`$ pure SW-theory whose BPS spectrum is not invariant under $`T`$ transformation . One can verify, however, that this configuration is invariant under transformation by $`T^2`$:
$$\mathrm{๐๐}\stackrel{T^2}{}\mathrm{๐๐}_{[\mathrm{๐},\mathrm{๐}]}\stackrel{P_1^1}{}\mathrm{๐๐}.$$
(5.12)
It then follows that $`๐(๐_\mathrm{๐})=\{11,T^2\}`$ and $`\overline{๐}(๐_\mathrm{๐})=\{T^2|T^4=11\}=๐น๐น_4`$. Since this configuration supports no junctions without asymptotic charges there is no $`๐ข`$ associated to it.
* $`๐_\mathrm{๐}:`$ Here $`K(๐_\mathrm{๐})=T^3`$ and $`\text{Stab}(K)=\{11,T\}`$. Since we have an $`๐`$ brane $`T๐`$. Thus $`๐(๐_\mathrm{๐})=\{11,T\}`$ and $`\overline{๐}(๐_\mathrm{๐})=\{T|T^3=11\}=๐น๐น_6`$. This configuration does not support any root, therefore there is no Dynkin diagram. Nevertheless, as is well known, it carries a $`u(1)`$ algebra, whose associated junction is the non-BPS junction $`\overline{}๐=2๐๐๐`$ . For an arbitrary junction $`๐`$ the corresponding $`u(1)`$ charge $`Q^{}`$ is proportional to $`๐\overline{๐}2Q_1+Q_BQ_C`$. Both $`11`$ and $`T`$ are checked to take $`Q^{}Q^{}`$.
* $`๐_\mathrm{๐}:`$ This configuration has $`K(๐_\mathrm{๐})=T^2`$, $`\text{Stab}(K)=\{11,T\}=๐(๐_\mathrm{๐})`$. Therefore $`\overline{๐}(๐_\mathrm{๐})=\{T|T^2=11\}=๐น๐น_4`$. The configuration supports two roots representing the $`A_1A_1`$ algebra, it corresponds to two disconnected Dynkin nodes, with Dynkin labels $`a_1`$ and $`a_2`$ correctly given by (5.3). The action of $`T`$ as given in (5.11) simply exchanges the two Dynkin labels. This is the non-trivial element of $`\mathrm{\Gamma }(A_1A_1)=๐น๐น_2`$. Therefore $`\varphi :\overline{๐}(๐_\mathrm{๐})=๐น๐น_4๐น๐น_2`$ via $`\varphi (T)=1.`$
* $`๐_\mathrm{๐}:`$ Here $`K(๐_\mathrm{๐})=T`$, $`\text{Stab}(K)=\{11,T\}=๐(๐_\mathrm{๐})`$, and $`\overline{๐}(๐_\mathrm{๐})=\{T|T=11\}=๐น๐น_2`$. Also $`T:(a_1,a_2,a_3)=(a_1,a_3,a_2),`$ is the non-trivial element of $`\mathrm{\Gamma }(A_3)=๐น๐น_2`$ (the labeling of nodes follows the $`D_N`$ conventions; node number one is in the middle). Thus the homomorphism $`\varphi :\overline{๐}(๐_\mathrm{๐})\mathrm{\Gamma }(A_3)=๐น๐น_2`$ is the isomorphism $`\varphi (T)=1.`$
* $`๐_\mathrm{๐}:`$ In this case the monodromy is $`11`$ and therefore $`\text{Stab}(K)=SL(2,๐น๐น)`$. Since invariance under $`T`$ has been already established, we show that $`๐(๐_\mathrm{๐})=\text{Stab}(K)=SL(2,๐น๐น)`$ by demonstrating the invariance of the configuration under $`S`$. Indeed,
$`๐^\mathrm{๐}\mathrm{๐๐}`$ $`\stackrel{S}{}`$ $`(๐_{[\mathrm{๐},\mathrm{๐}]})^\mathrm{๐}\mathrm{๐๐}\stackrel{(P_5P_4P_3P_2P_1P_1P_2P_3P_4)}{}๐^\mathrm{๐}\mathrm{๐๐}.`$ (5.13)
It is possible to anticipate the action of $`S`$ on the $`D_4`$ Dynkin labels. Recall from that the various conjugacy classes of $`so(8)`$ are correlated with asymptotic $`(p,q)`$ charges mod 2. In particular in eq. (6.26) of we see that $`\mathrm{๐}_๐ฏ`$ and $`\mathrm{๐}_๐ฌ`$ representations arise from $`(1,0)`$ and $`(0,1)`$ charges respectively (mod 2), while $`\mathrm{๐}_๐`$ arises from $`(1,1)`$. We see that mod 2, the action of $`S`$ on those asymptotic charges exchanges the ones corresponding to $`\mathrm{๐}_๐ฏ`$ and $`\mathrm{๐}_๐ฌ`$ while it leaves invariant that corresponding to $`\mathrm{๐}_๐`$. In our conventions, $`\mathrm{๐}_๐ฏ`$ and $`\mathrm{๐}_๐ฌ`$ are associated to the first and third nodes of the Dynkin diagram, and therefore we expect $`S`$ to act as the graph automorphism $`a_1a_3`$.
The transformations in (5.13) imply that under $`S`$, the invariant charges transform as:
$`Q_i`$ $`\stackrel{S}{}`$ $`Q_BQ_i+{\displaystyle Q_k}i=1\mathrm{}4`$
$`Q_B`$ $`\stackrel{S}{}`$ $`Q_C2Q_B2{\displaystyle Q_k}`$ (5.14)
$`Q_C`$ $`\stackrel{S}{}`$ $`Q_B{\displaystyle Q_k}.`$
The resulting action on the Dynkin labels (5.3) is given by
$`S:(a_1,a_2,a_3,a_4)`$ $``$ $`(a_1,a_2,a_3,a_1+2a_2+a_3+a_4).`$ (5.15)
A little calculation shows that $`S`$ is a composition of Weyl reflections and the expected graph automorphism:
$$S=๐ชW_{\alpha _1+\alpha _2+\alpha _3}W_{\alpha _2},๐ช(a_1,a_2,a_3,a_4)(a_3,a_2,a_1,a_4).$$
(5.16)
* $`๐_{๐\mathrm{๐}}:`$ Here $`K(๐_๐)=T^{4N}`$ and $`\text{Stab}(K)=\{11,T\}=๐(๐_๐)`$. Thus $`\overline{๐}(๐_๐)=\{T|T^{N4}=11\}=๐น๐น_{2(N4)}`$. Under $`T`$, using (5.11), we have $`\varphi :\overline{๐}(๐_๐)\mathrm{\Gamma }(D_N)=๐น๐น_2`$ is fixed by $`\varphi (T)=1`$. Since we have a homomorphism, $`\varphi (11)=\varphi (T^{N4})=[\varphi (T)]^{N4}=(1)^{N4}=(1)^N`$. Thus for $`๐_N`$ with $`N`$ even, the transformation $`11`$ maps to a Weyl transformation, while for $`N`$ odd, the transformation $`11`$ is equivalent to the nontrivial graph automorphism up to a Weyl transformation. This is as expected; a change of sign of all Dynkin labels in the $`D_N`$ algebras is a Weyl transformation only for $`N`$ even (Ref. , sect. 13).
### 5.4 $`๐_๐`$ configuration: $`๐^{๐\mathrm{๐}}\mathrm{๐๐๐}`$
* $`๐_\mathrm{๐}`$: Here $`K(๐_\mathrm{๐})(ST)^1`$ and $`\text{Stab}(K)=\{11,K\}=๐น๐น_6=๐(๐_\mathrm{๐})`$. Therefore $`\overline{๐}(๐_\mathrm{๐})=\{11\}=๐น๐น_2`$. In addition, $`\mathrm{\Gamma }(E_6)=๐น๐น_2`$, and its non-trivial element, up to a Weyl transformation, changes the sign of all the Dynkin labels (maps representations to their conjugates). It follows that the homomorphism $`\varphi :\overline{๐}(๐_\mathrm{๐})๐น๐น_2`$ is fixed by $`\varphi (11)=1`$. This is an isomorphism.
* $`๐_\mathrm{๐}`$: Here $`K(๐_\mathrm{๐})S`$ and $`\text{Stab}(K)=\{S\}=๐น๐น_4=๐(๐_\mathrm{๐})`$. It follows that $`\overline{๐}(E_7)=\{11\}`$. Since $`\mathrm{\Gamma }(E_7)`$ is also trivial the homomorphism $`\varphi `$ is trivial. Dualities will preserve Weyl orbits, and therefore representations. While the duality $`11=S^2`$ changes the sign of all Dynkin labels, this is simply a Weyl transformation of $`E_7`$.
* $`๐_\mathrm{๐}`$: Here $`K(๐_\mathrm{๐})(ST)^1`$ and $`\text{Stab}(K)=\{(ST)^1\}=๐น๐น_6=๐(๐_\mathrm{๐})`$. Just as in the case of $`E_7`$ we have $`\overline{๐}(๐_\mathrm{๐})=\{11\}`$, $`\mathrm{\Gamma }(E_8)=1`$ and a trivial homomorphism $`\varphi `$.
## 6 Duality groups for affine configurations
In this section we will try to extend the result of previous section to the case of affine exceptional configurations $`\widehat{\stackrel{~}{๐}}_0`$, $`\widehat{\stackrel{~}{๐}}_1`$, and the series $`\widehat{๐}_N`$ for $`1N8`$. These configurations are more interesting because of their relation with del Pezzo surfaces but at the same time more difficult to analyze since for these configurations $`(w)_W(w)`$ i.e, not all transpositions are of Weyl type. This means that there are elements in $`(w)`$ whose action on roots may be outer automorphisms of the root lattice. Our strategy in this case will be to find the map $`\overline{๐}(w)/\{11\}\pm \mathrm{\Gamma }`$ for a fixed set of transformations used to undo the $`SL(2,๐น๐น)`$ transformations in $`๐(w)`$.
We begin with some general remarks applicable to the the affine exceptional brane configurations $`\widehat{๐}_๐=๐^{๐\mathrm{๐}}\mathrm{๐๐๐๐}`$ with $`2N8`$. All such configurations have at least one $`๐`$ brane. They have monodromy $`K(\widehat{๐}_๐)=T^{9N}`$, and one readily finds that Stab$`(K)=\{11,T\}`$. We now show that $`T๐`$ for $`n>1`$ by an explicit calculation quite similar to that given in (5.8). We make the $`๐`$ brane do a counterclockwise round trip around the other branes:
$`\mathrm{๐๐๐๐๐}\stackrel{T}{}\mathrm{๐๐}_{[\mathrm{๐},\mathrm{๐}]}๐_{[\mathrm{๐},\mathrm{๐}]}๐_{[\mathrm{๐},\mathrm{๐}]}๐_{[\mathrm{๐},\mathrm{๐}]}\stackrel{P_1P_2P_3P_4P_4P_3P_2P_1}{}\mathrm{๐๐๐๐๐}.`$ (6.1)
Let $`\{Q_A,Q_B^1,Q_C^1,Q_B^2,Q_C^2\}`$ denote the invariant charges on the branes. The transformed invariant charges are found to be
$`Q_A`$ $`\stackrel{T}{}`$ $`Q_A+q,`$
$`Q_B^1`$ $`\stackrel{T}{}`$ $`Q_B^1+Q_A,`$
$`Q_C^1`$ $`\stackrel{T}{}`$ $`Q_C^1+Q_A,`$ (6.2)
$`Q_B^2`$ $`\stackrel{T}{}`$ $`Q_B^2Q_A,`$
$`Q_C^2`$ $`\stackrel{T}{}`$ $`Q_C^2Q_A.`$
Here, $`q=Q_C^1+Q_C^2Q_B^1Q_B^1`$ is the total $`q`$-charge of the junction $`๐`$. Since we are only interested in junctions that correspond to states in the root lattice, we set $`q=0`$. We see that the effect of the $`T`$ transformation on a junction of $`\widehat{E}_N`$ with zero $`q`$ charge is simply
$$๐๐+Q_A(๐)๐น.$$
(6.3)
Here, $`Q_A(๐)`$ is the invariant charge on the $`๐`$-brane used to undo the effect of the $`T`$ transformation, and $`๐น=๐_1+๐_2๐_1๐_2`$ . Indeed, when there is more than one $`๐`$ brane, any of them can be used to undo the effect of $`T`$. We have therefore shown that
$$๐(\widehat{๐}_๐)=\{11,T\},2N8.$$
(6.4)
$`\widehat{๐}_\mathrm{๐}`$ being the composition of two copies of $`D_4`$ has $`๐(\widehat{๐}_\mathrm{๐})=SL(2,๐น๐น)`$. In addition, since $`K(\widehat{๐}_๐)=T^{9N}`$ we also have
$$\overline{๐}(\widehat{๐}_๐)=\{11,T|T^{9N}=1\}=๐น๐น_2\times ๐น๐น_{9N},2N8.$$
(6.5)
The group Aut$`(Q)/W`$ of an affine algebra, written as $`\pm \mathrm{\Gamma }`$ in section 4, is more precisely written as
$$\text{Aut}(Q)/W=๐น๐น_2\times \mathrm{\Gamma }$$
(6.6)
where the element $`(1,e)`$, with $`e`$ the identity in $`\mathrm{\Gamma }`$, is the transformation $`QQ`$ reversing the sign of all the vectors in the root lattice, and thus reversing the sign of all Dynkin labels. An element of the form $`(0,h\mathrm{\Gamma })`$ simply acts by the graph automorphism $`h`$ of the Dynkin graph of the affine algebra. The duality $`11๐(\widehat{๐}_๐)`$ maps to $`(1,0)`$ in $`\overline{๐}`$ and then
$$\varphi :(1,0)\overline{๐}=๐น๐น_2\times ๐น๐น_{9N}(1,e)๐น๐น_2\times \mathrm{\Gamma }$$
(6.7)
Our computations will require finding how $`T๐`$ acts. For this we note that it maps to $`(0,1)\overline{๐}`$. We will find that
$$\varphi :(0,1)\overline{๐}=๐น๐น_2\times ๐น๐น_{9N}(0,h(T))๐น๐น_2\times \mathrm{\Gamma },$$
(6.8)
where $`h(T)`$ is a graph automorphism. This map respects the product structure of the groups involved.
To simplify the formulae we also introduce the following notation for Weyl transformations,
$$W_{i_1^{n_1}i_2^{n_2}\mathrm{}i_k^{n_k}}W_{n_1\alpha _{i_1}+n_2\alpha _{i_2}+\mathrm{}n_k\alpha _{i_k}}.$$
(6.9)
Let us now consider in detail the various configurations in the above series. We will show the brane configurations and indicate the simple root junctions. Then we select an $`๐`$ brane to undo the $`T`$ duality and use equation (6.3) to find the action on the simple roots. The final step is writing this action as the composition of a Weyl transformation and the action arising from a Dynkin graph automorphism. The answer is the graph automorphism $`h(T)\mathrm{\Gamma }`$ defined in equation (6.8).
* $`\widehat{๐}_\mathrm{๐}`$: The $`T`$ transformation acts trivially on the roots since the $`๐`$-brane supports no root. Therefore
$$h(T)=0.$$
(6.10)
* $`\widehat{๐}_\mathrm{๐}`$: Using the rightmost $`๐`$ brane of the figure, $`T`$ acts on the simple roots as follows,
$`T:(\alpha _0,\alpha _1,\alpha _2;\beta _0,\beta _1)`$ $``$ $`(\alpha _0\delta ,\alpha _1+\delta ,\alpha _2;\beta _0+\delta ,\beta _1\delta ).`$ (6.11)
Here, $`(\alpha _0,\alpha _1,\alpha _2)`$ are the roots of $`\widehat{A}_2`$ and $`(\beta _0,\beta _1)`$ are the roots of $`\widehat{A}_1`$. It is easy to verify that acting on simple roots
$$T=(W_{\alpha _1}W_{\alpha _2}๐ช,W_{\beta _0}๐ช^{}),$$
(6.12)
where the first and second terms indicate the action on the $`\widehat{A}_2`$ and $`\widehat{A}_1`$ roots respectively. The graph automorphisms $`๐ช`$ and $`๐ช^{}`$ are
$$๐ช(\alpha _0,\alpha _1,\alpha _2)=(\alpha _2,\alpha _0,\alpha _1),๐ช^{}(\beta _0,\beta _1)=(\beta _1,\beta _0).$$
(6.13)
These are elements of $`\mathrm{\Gamma }(\widehat{A}_2\widehat{A}_1)=๐_6\times ๐น๐น_2`$, where $`๐_6`$ denotes the symmetry group of the triangle, here formed by the $`\widehat{A}_2`$ simple roots, with $`๐ช`$ the elementary rotation. In addition, $`๐ช^{}`$ is the nontrivial element of $`๐น๐น_2`$, representing the exchange of the two simple roots of $`\widehat{A}_1`$. In summary;
$$h(T)=(๐ช,๐ช^{})๐_6\times ๐น๐น_2=\mathrm{\Gamma }(\widehat{A}_2\widehat{A}_1).$$
(6.14)
* $`\widehat{๐}_\mathrm{๐}`$: Using the rightmost $`๐`$ brane the nontrivial action of $`T`$ on the roots is given by
$$T:(\alpha _2,\alpha _4)(\alpha _2\delta ,\alpha _4+\delta ).$$
(6.15)
A calculation shows that
$$T=๐ชW_1W_2W_3W_0W_1W_2,๐ช(\alpha _0,\alpha _1,\alpha _2,\alpha _3,\alpha _4)=(\alpha _3,\alpha _4,\alpha _0,\alpha _1,\alpha _2).$$
(6.16)
Here $`๐ช`$ implements the transformation $`\omega ^3\mathrm{\Gamma }(\widehat{A}_4)=๐_{10}`$, where $`\omega `$ is a cyclic minimal rotation of the pentagon. Note that $`\omega ^3`$ is a generator for the $`๐น๐น_5`$ subgroup of rotations of $`๐_{10}`$. Thus, in summary
$$h(T)=๐ช=\omega ^3๐_{10},$$
(6.17)
Note that in $`\overline{๐}`$, $`T^5=11`$, and this is consistent with the map to $`๐_{10}`$.
* $`\widehat{๐}_\mathrm{๐}`$: Using the rightmost $`๐`$ brane, the action of $`T`$ is:
$$T:(\alpha _0,\alpha _5)(\alpha _0\delta ,\alpha _5+\delta ),$$
(6.18)
This time we find
$$T=๐ชW_2W_{1234}W_{35},๐ช(\alpha _0,\alpha _1,\alpha _2,\alpha _3,\alpha _4,\alpha _5)=(\alpha _4,\alpha _5,\alpha _3,\alpha _2,\alpha _1,\alpha _0)$$
(6.19)
Here $`๐ช\mathrm{\Gamma }(\widehat{D}_5)=๐_8`$, is a generator for the $`๐น๐น_4`$ subgroup of $`๐_8`$, consistent with $`T^4=11`$ in $`\overline{๐}`$. In summary:
$$h(T)=๐ช\mathrm{\Gamma }(\widehat{D}_5)=๐_8.$$
(6.20)
* $`\widehat{๐}_\mathrm{๐}`$: Using the leftmost brane, one finds that $`T`$ induces the transformations
$$T:(\alpha _0,\alpha _5)(\alpha _0\delta ,\alpha _5+\delta ).$$
(6.21)
A calculation shows that
$$T=๐ชW_{34}W_{1234}W_{1236}W_{23}W_{45},$$
(6.22)
where $`๐ช\mathrm{\Gamma }(\widehat{E}_6)=S_3`$ is the generator of the $`๐น๐น_3`$ subgroup of $`S_3`$ performing the rigid minimal rotation of the Dynkin diagram:
$$๐ช(\alpha _0,\alpha _1,\alpha _2,\alpha _3,\alpha _4,\alpha _5,\alpha _6)=(\alpha _1,\alpha _5,\alpha _4,\alpha _3,\alpha _6,\alpha _0,\alpha _2).$$
(6.23)
This is compatible with $`T^3=11`$ in $`\overline{๐}`$. In summary:
$$h(T)=๐ช\mathrm{\Gamma }(\widehat{E}_6)=S_3.$$
(6.24)
* $`\widehat{๐}_\mathrm{๐}`$: Using the leftmost $`๐`$ brane, $`T`$ acts as
$`T`$ $`:(\alpha _0,\alpha _1)(\alpha _0\delta ,\alpha _1+\delta )`$ (6.25)
$`T`$ $`=๐ชW_{45}W_{45^267}W_{123^24567}W_{12^23^24^257}W_{34},`$
where $`๐ช\mathrm{\Gamma }(\widehat{E}_7)=๐น๐น_2`$ is the nontrivial generator of the graph automorphism, and exchanges the two long branches of the Dynkin diagram:
$$๐ช(\alpha _0,\alpha _1,\alpha _2,\alpha _3,\alpha _4,\alpha _5,\alpha _6,\alpha _7)=(\alpha _1,\alpha _0,\alpha _6,\alpha _5,\alpha _4,\alpha _3,\alpha _2,\alpha _7).$$
(6.26)
This is compatible with $`T^2=11`$ in $`\overline{๐}`$. In summary:
$$h(T)=๐ช\mathrm{\Gamma }(\widehat{E}_7)=๐น๐น_2.$$
(6.27)
* $`\widehat{๐}_\mathrm{๐}`$: Using the leftmost $`๐`$ brane the $`T`$ action is:
$$T:(\alpha _0,\alpha _1)(\alpha _02\delta ,\alpha _1+\delta )$$
(6.28)
Since $`\mathrm{\Gamma }(\widehat{E}_8)=1`$, the above ought to be a pure Weyl transformation. Indeed,
$`T=W_0W_{123456}W_{123458}W_{34567}W_{2345^2678}W_6W_{345^2678}W_{234568}`$
$`W_{56}W_{1234567}W_{12}W_{12345678}W_{34568}W_{12345^268}.`$ (6.29)
Therefore $`h(T)=e`$ is the identity element in the trivial group $`\mathrm{\Gamma }(\widehat{E}_8)=1`$. Indeed $`T=11`$ in $`๐`$ as well.
We now consider in detail the special cases:
$``$ $`\widehat{\stackrel{~}{๐}}_\mathrm{๐}`$: This configuration consists of three 7-branes $`๐_{[\mathrm{๐},\mathrm{๐}]}๐_{[\mathrm{๐},\mathrm{๐}]}๐_{[\mathrm{๐},\mathrm{๐}]}`$ (eq(3.10) of with last two branes interchanged). Therefore $`K(\widehat{๐}_\mathrm{๐})=T^9`$ and $`\text{Stab}(K)=\{11,T\}`$. An important property of this configuration is that the asymptotic charge $`(p,q)`$ of every junction satisfies the condition $`pq0(\text{mod}\mathrm{\hspace{0.17em}3})`$. An argument similar to the one used for $`๐_\mathrm{๐}`$ proves that if $`T^k๐(w)`$ then $`k0(\text{mod}\mathrm{\hspace{0.17em}3})`$. By explicit computation one can show that $`T^3๐(w)`$:
$`๐_{[\mathrm{๐},\mathrm{๐}]}๐_{[\mathrm{๐},\mathrm{๐}]}๐_{[\mathrm{๐},\mathrm{๐}]}\stackrel{T^3}{}๐_{[\mathrm{๐},\mathrm{๐}]}๐_{[\mathrm{๐},\mathrm{๐}]}๐_{[\mathrm{๐},\mathrm{๐}]}\stackrel{P_2^1P_1^1}{}๐_{[\mathrm{๐},\mathrm{๐}]}๐_{[\mathrm{๐},\mathrm{๐}]}๐_{[\mathrm{๐},\mathrm{๐}]}.`$ (6.30)
Thus $`๐(\widehat{\stackrel{~}{๐}}_\mathrm{๐})=\{11,T^3\}`$ and $`\overline{๐}(\widehat{\stackrel{~}{๐}}_\mathrm{๐})=\{11,T^3|T^9=11\}=๐น๐น_2\times ๐น๐น_3`$. This configuration supports no algebra. The only junctions having no asymptotic charges (localized) are multiples of the delta junction
$$\delta =๐ฑ_{[\mathrm{๐},\mathrm{๐}]}+๐ฑ_{[\mathrm{๐},\mathrm{๐}]}๐ฑ_{[\mathrm{๐},\mathrm{๐}]},$$
(6.31)
which can be represented as a $`(1,0)`$ string circling the branes in the counterclockwise direction. The action of the $`T^3`$ duality transformation on a general junction is found as usual using the brane transpositions in (6.30). Denoting the invariant charges on the $`๐_{[\mathrm{๐},\mathrm{๐}]},๐_{[\mathrm{๐},\mathrm{๐}]},\text{and}๐_{[\mathrm{๐},\mathrm{๐}]}`$ branes by $`Q_1,Q_2`$ and $`Q_3`$ respectively, we find that
$$T^3:(Q_1,Q_2,Q_3)(Q_2,Q_3,Q_1+3(Q_2+Q_3)).$$
(6.32)
In this notation the delta junction is $`(1,1,1)`$, and one readily verifies that $`T^3`$ leaves it invariant. This was expected, since $`T^3`$ leaves invariant the $`(1,0)`$ string, and brane transpositions cannot affect a looping string.
$``$ $`\widehat{\stackrel{~}{๐}}_\mathrm{๐}`$: Here $`K(\widehat{๐}_\mathrm{๐})=T^8`$ and $`\text{Stab}(K)=\{11,T\}`$. Since this brane configuration has an $`๐`$brane a simple computation shows that $`T๐(w)`$:
$$\mathrm{๐๐}_{[\mathrm{๐},\mathrm{๐}]}๐_{[\mathrm{๐},\mathrm{๐}]}๐_{[\mathrm{๐},\mathrm{๐}]}\stackrel{T}{}\mathrm{๐๐}_{[\mathrm{๐},\mathrm{๐}]}๐_{[\mathrm{๐},\mathrm{๐}]}๐_{[\mathrm{๐},\mathrm{๐}]}\stackrel{P_1P_2P_3P_3P_2P_1}{}\mathrm{๐๐}_{[\mathrm{๐},\mathrm{๐}]}๐_{[\mathrm{๐},\mathrm{๐}]}๐_{[\mathrm{๐},\mathrm{๐}]}.$$
(6.33)
Thus $`๐(\widehat{\stackrel{~}{๐}}_\mathrm{๐})=\{11,T\}`$ and $`\overline{๐}(\widehat{\stackrel{~}{๐}}_\mathrm{๐})=\{11,T|T^8=11\}=๐น๐น_2\times ๐น๐น_8`$. A localized junction on this configuration is specified by the number of $`\delta `$ loops (the junction in (6.31)) and a $`u(1)`$ charge $`\overline{Q}`$. We thus write $`๐=m\delta +\overline{Q}\overline{๐}`$ where
$$\overline{๐}=3๐๐ฑ_{[\mathrm{๐},\mathrm{๐}]}๐ฑ_{[\mathrm{๐},\mathrm{๐}]}.$$
(6.34)
Under the $`T`$ transformation the junction $`๐(m,Q)`$ transforms as
$$T:(m,\overline{Q})(m+3\overline{Q},\overline{Q}),$$
(6.35)
which follows because the $`๐`$ brane that must be circled around has invariant charge $`3\overline{Q}`$.
Since the monodromy of the configuration is $`T^8`$, we see that
$$T^8:(m,\overline{Q})(m+24\overline{Q},\overline{Q}),$$
(6.36)
should be a transformation that can be generated simply by crossing transformations. Such a tranposition (actually its inverse) was discussed in section 3.2. It is not of Weyl type because there are no real root junctions for this brane configuration.
Indeed, the required transposition is the inverse of that discussed in section 3.2 for the case of $`\widehat{\stackrel{~}{๐}}_1`$. These transpositions are equivalent to first taking the $`๐`$ brane clockwise around the other three branes three times. This has the effect of changing the labels of the other three branes as if acted by $`T^3`$. Those branes are then restored to their original labels by performing the transpositions indicated in the second step of (6.30). The first step takes a junction $`๐=m\delta +\overline{Q}\overline{}๐`$ and adds to it $`(3Q_A)=9\overline{Q}`$ delta junctions. This step, while changing the labels of the three rightmost branes, it does not change the invariant charges they have; these are, in the notation used for $`\widehat{\stackrel{~}{๐}}_\mathrm{๐}`$, $`\overline{Q}(1,0,1)`$ (see (6.34)). Using (6.32) we see that under the restoring transposition: $`\overline{Q}(1,0,1)\overline{Q}(0,1,2)=\overline{Q}(1,0,1)+\overline{Q}\delta `$. Therefore, the complete series of transpositions $`\stackrel{~}{b}`$ adds $`(9+1)\overline{Q}=8\overline{Q}`$ delta junctions:
$$\stackrel{~}{b}:(m,\overline{Q})(m8\overline{Q},\overline{Q})$$
(6.37)
Comparing with (6.36) we see that indeed $`T^8`$ has the same effect as the transposition $`(\stackrel{~}{b})^3`$. Since a $`T`$ duality adds $`3\overline{Q}`$ delta junctions, and we can add or remove $`8\overline{Q}`$ delta junctions by transpositions, $`T^8`$ is the lowest power of $`T`$ that is equivalent to a transposition.
$``$ $`\widehat{๐}_\mathrm{๐}`$: Here $`K(\widehat{๐}_\mathrm{๐})=T^8`$. A junction of asymptotic charge $`(p,q)`$ with support on this configuration satisfies the condition $`p+q0(\text{mod}\mathrm{\hspace{0.17em}2})`$. Thus if $`T^k๐(w)`$ then $`k0(\text{mod}\mathrm{\hspace{0.17em}2})`$. By explicit computation one can show that $`T^2๐(w)`$:
$$\mathrm{๐๐๐๐}\stackrel{T^2}{}๐_{[\mathrm{๐},\mathrm{๐}]}๐_{[\mathrm{๐},\mathrm{๐}]}๐_{[\mathrm{๐},\mathrm{๐}]}๐_{\mathrm{๐},\mathrm{๐}]}\stackrel{P_3^1P_1^1}{}\mathrm{๐๐๐๐}.$$
(6.38)
Thus $`๐(\widehat{๐}_\mathrm{๐})=\{11,T^2\}`$ and $`\overline{๐}(\widehat{๐}_\mathrm{๐})=\{11,T^2|T^8=11\}=๐น๐น_2\times ๐น๐น_4`$. The $`T^2`$ transformation acts on the invariant charges in the following way,
$$T^2:(Q_{B_1},Q_{C_1},Q_{B_2},Q_{C_2})(Q_{C_1},Q_{B_1}+2Q_{C_1},Q_{C_2},Q_{B_2}+2Q_{C_2}).$$
(6.39)
The invariant charges of the root junctions and the delta junction $`\delta `$ are
$$๐ถ_\mathrm{๐}=(1,0,1,0),๐ถ_\mathrm{๐}=(0,1,0,1),\delta =๐ถ_0+๐ถ_1=(1,1,1,1).$$
(6.40)
Thus from (6.39) and (6.40) it follows that
$$T^2:(\alpha _0,\alpha _1)(\alpha _0\delta ,\alpha _1+\delta ).$$
(6.41)
$`\mathrm{\Gamma }(\widehat{A}_1)=๐น๐น_2`$ where the nontrivial element is $`๐ช(\alpha _0,\alpha _1)=(\alpha _1,\alpha _0)`$. One can verify that acting on the simple roots
$$T^2=๐ชW_0.$$
(6.42)
Therefore, $`\varphi :\overline{๐}(\widehat{๐}_\mathrm{๐})๐น๐น_2\times \mathrm{\Gamma }(\widehat{A}_1)`$, is determined by
$$h(T^2)=๐ช.$$
(6.43)
This completes our analysis of the duality groups of affine exceptional brane configurations.
## Acknowledgements
We wish to acknowledge useful conversations with Y. Imamura and V. Kac.
This research was supported in part by the US Department of Energy under contract #DE-FC02-94ER40818. |
warning/0002/math-ph0002001.html | ar5iv | text | # Finite Dimensional Representations of Quadratic Algebras with Three Generators and Applications 11footnote 1Talk given at the VI International Wigner Symposium, 16-22 August 1999, Instabul, Turkey
## 1 Introduction
In classical mechanics, integrable system is a system possessing more constants of motion in addition to the energy. A comprehencive review of the two-dimensional integrable classical systems is given by Hietarinta, where the space was assumed to be flat. The case of non flat space is under current investigation.
An interesting subset of the totality of integrable systems is the set of systems, which possess a maximum number of integrals, these systems are termed as superintegrable ones. The Coulomb and the harmonic oscillator potentials are the most familiar classical superintegrable systems, whose their quantum counterpart has nice symmetry properties, which are described by the $`so(N+1)`$ and $`su(N)`$ Lie algebras.
The Hamiltonian of the classical systems is a quadratic function of the momenta. In the case of the flat space all the known two dimensional superintegrable systems with quadratic integrals of motion are simultaneously separable in more than two orthogonal coordinate systems. The integrals of motion of a two dimensional superintegrable system in flat space close in a classical quadratic Poisson algebra. The study of the quadratic Poisson algebras is a matter under investigation related to several branches of physics as: the solution of the classical Yang - Baxter equation , the two dimensional superintegrable systems in flat space or on the sphere , the statistics or the case of โexactly solvableโ classical problems .
The quantization of classical integrable systems turns generally to quantum integrable systems, but sometime one has to add correction terms to the integrals of motion or to the Hamiltonian, these correction terms seem to be of order $`๐ช(\mathrm{}^2)`$ . The classical Poisson algebra is shifted to some quantum polynomial algebra, the same thing is true in the case of quadratic Poisson algebra corresponding to the Yang \- Baxter equation, which is turned to a quantum quadratic associative algebra. The same idea was discussed in ref., where the classical problems, which are expressed by a quadratic Poisson algebra are mapped to quantum ones described by the corresponding quantum operator quadratic algebra. The same shift is indeed true for the superintegrable systems, where the classical ones correspond to the quantum ones and the classical quadratic Poisson algebra is mapped to a quadratic associative algebra.
In this contribution we study the general form of the quadratic algebras, which are encountered in the case of the two dimensional quantum superintegrable systems, these algebras are called $`Qu(3)`$. In references was conjectured that, the energy eigenvalues correspond to finite dimensional representations of the latent quadratic algebras. Granovkii et al in studied the representations of the quadratic Askey - Wilson algebras $`QAW(3)`$. Using there the proposed ladder representation, the finite dimensional representations are calculated and this method was applied to several superintegrable systems . Another method for calculating the finite dimensional representations is the use of the deformed oscillator algebra and their finite dimensional version which are termed as generalized deformed parafermionic algebras.
## 2 The Qu(3) Algebra
Let consider the quadratic associative algebra generated by the generators $`\{A,B,C\}`$, which satisfy the commutation relations
$$\begin{array}{c}[A,B]=C\hfill \\ [A,C]=\alpha A^2+\beta B^2+\gamma \{A,B\}+\delta A+ฯตB+\zeta \hfill \\ [B,C]=aA^2+bB^2+c\{A,B\}+dA+eB+z\hfill \end{array}$$
(1)
After rotating the generators $`A`$ and $`B`$, we can always consider the case $`\beta =0`$.
The Jacobi equality for the commutator induces the relation
$$[A,[B,C]]=[B,[A,C]]$$
the following relations
$$b=\gamma ,c=\alpha \text{and}e=\delta $$
must be satisfied, and consequently the general form of the quadratic algebra (1) can be explicitly written as follows:
$$[A,B]=C$$
(2)
$$[A,C]=\alpha A^2+\gamma \{A,B\}+\delta A+ฯตB+\zeta $$
(3)
$$[B,C]=aA^2\gamma B^2\alpha \{A,B\}+dA\delta B+z$$
(4)
The Casimir of this algebra is given by:
$$\begin{array}{cc}\hfill K=& C^2\alpha \{A^2,B\}\gamma \{A,B^2\}+(\alpha \gamma \delta )\{A,B\}+\hfill \\ \hfill +& (\gamma ^2ฯต)B^2+(\gamma \delta 2\zeta )B+\hfill \\ \hfill +& \frac{2a}{3}A^3+(d+\frac{a\gamma }{3}+\alpha ^2)A^2+(\frac{aฯต}{3}+\alpha \delta +2z)A\hfill \end{array}$$
(5)
another useful form of the Casimir of the algebra is given by:
$$\begin{array}{cc}\hfill K=& C^2+\frac{2a}{3}A^3\frac{\alpha }{3}\{A,A,B\}\frac{\gamma }{3}\{A,B,B\}+\hfill \\ & +\left(\frac{2\alpha ^2}{3}+d+\frac{2a\gamma }{3}\right)A^2+\left(ฯต+\frac{2\gamma ^2}{3}\right)B^2+\hfill \\ & +\left(\delta +\frac{a\gamma }{3}\right)\{A,B\}+\left(\frac{2\alpha \delta }{3}+\frac{aฯต}{3}+\frac{d\gamma }{3}+2z\right)A+\hfill \\ & +\left(\frac{\alpha ฯต}{3}+\frac{2\delta \gamma }{3}2\zeta \right)B+\frac{\gamma z}{3}\frac{\alpha \zeta }{3}\hfill \end{array}$$
(6)
This quadratic algebra has many similarities to the Racah algebra $`QR(3)`$, which is a special case of the Askey - Wilson algebra $`QAW(3)`$. The algebra (24) does not coincide with the Racah algebra $`QR(3)`$, if $`a0`$ in the relation (4). We shall call this algebra $`Qu(3)`$ algebra. Unless this difference between $`Qu(3)`$ and $`QR(3)`$ algebra a representation theory can be constructed by following the same procedures as they were described by Granovskii, Lutzenko and Zhedanov in ref. . In this paper we shall give a realization of this algebra using the deformed oscillator techniques. The finite dimensional representations of the algebra $`Qu(3)`$ will be constructed by constructing a realization of the algebra $`Qu(3)`$ with the generalized parafermionic algebra introduced by Quesne.
## 3 Deformed Parafermionic Algebra
Let now consider a realization of the algebra $`Qu(3)`$, by using of the deformed oscillator technique, i.e. by using a deformed oscillator algebra $`\{b^{},b,๐ฉ\}`$, which satisfies the
$$[๐ฉ,b^{}]=b^{},[๐ฉ,b]=b,b^{}b=\mathrm{\Phi }\left(๐ฉ\right),bb^{}=\mathrm{\Phi }\left(๐ฉ+1\right)$$
(7)
where the function $`\mathrm{\Phi }(x)`$ is a โwell behavedโ real function which satisfies the the the boundary condition:
$$\mathrm{\Phi }(0)=0,\mathrm{\Phi }(x),\text{for}x>0$$
(8)
As it is well known this constraint imposes the existence a Fock type representation of the deformed oscillator algebra, which is bounded by bellow, i.e. there is a Fock basis $`|n>,n=0,1,\mathrm{}`$ such that
$$\begin{array}{c}๐ฉ|n>=n|n>\hfill \\ b^{}|n>=\sqrt{\mathrm{\Phi }\left(n+1\right)}|n+1>,n=0,1,\mathrm{}\hfill \\ b|0>=0\hfill \\ b|n>=\sqrt{\mathrm{\Phi }\left(n\right)}|n1>,n=1,2,\mathrm{}\hfill \end{array}$$
(9)
The Fock representation (9) is bounded by bellow.
In the case of nilpotent deformed oscillator algebras, there is a positive integer $`p`$, such that
$$b^{p+1}=0,\left(b^{}\right)^{p+1}=0$$
the above equations imply that
$$\mathrm{\Phi }(p+1)=0,$$
(10)
In that case the deformed oscillator (7) has a finite dimensional representation, with dimension equal to $`p+1`$, this kind of oscillators are called deformed parafermion oscillators of order $`p`$.
An interesting property of the deformed parafermionic algebra is that the existence of a faithfull finite dimensional representation of the algebra implies that:
$$๐ฉ\left(๐ฉ1\right)\left(๐ฉ2\right)\mathrm{}\left(๐ฉp\right)=0$$
(11)
The above restriction and the constraints (8) and (10) imply that the general form of the structure function $`\mathrm{\Phi }(๐ฉ)`$ has the general form:
$$\mathrm{\Phi }(๐ฉ)=๐ฉ(p+1๐ฉ)(a_0+a_1๐ฉ+a_2๐ฉ^2+\mathrm{}a_{p1}๐ฉ^{p1})$$
## 4 Oscillator realization of the algebra $`Qu(3)`$
We shall show, that there is a realization of the algebra $`Qu(3)`$, such that
$`A=A\left(๐ฉ\right)`$ (12)
$`B=b\left(๐ฉ\right)+b^{}\rho \left(๐ฉ\right)+\rho \left(๐ฉ\right)b`$ (13)
where the $`A[x],b[x]`$ and $`\rho (x)`$ are functions, which will be determined. In that case (2) implies:
$$C=[A,B]C=b^{}\mathrm{\Delta }A\left(๐ฉ\right)\rho \left(๐ฉ\right)\rho \left(๐ฉ\right)\mathrm{\Delta }A\left(๐ฉ\right)b$$
(14)
where
$$\mathrm{\Delta }A\left(๐ฉ\right)=A\left(๐ฉ+1\right)A\left(๐ฉ\right)$$
Using equations (12), (13) and (3) we find:
$$\begin{array}{cc}\hfill [A,C]=& [A\left(๐ฉ\right),b^{}\mathrm{\Delta }A\left(๐ฉ\right)\rho \left(๐ฉ\right)\rho \left(๐ฉ\right)\mathrm{\Delta }A\left(๐ฉ\right)b]=\hfill \\ \hfill =& b^{}\left(\mathrm{\Delta }A\left(๐ฉ\right)\right)^2\rho \left(๐ฉ\right)+\rho \left(๐ฉ\right)\left(\mathrm{\Delta }A\left(๐ฉ\right)\right)^2b=\hfill \\ \hfill =& \alpha A^2+\gamma \{A,B\}+\delta A+ฯตB+\zeta =\hfill \\ \hfill =& b^{}\left(\gamma \left(A\left(๐ฉ+1\right)+A\left(๐ฉ\right)\right)+ฯต\right)\rho \left(๐ฉ\right)+\hfill \\ & +\rho \left(๐ฉ\right)\left(\gamma \left(A\left(๐ฉ+1\right)+A\left(๐ฉ\right)\right)+ฯต\right)b+\hfill \\ & +\alpha A\left(๐ฉ\right)^2+2\gamma A\left(๐ฉ\right)b\left(๐ฉ\right)+\delta A\left(๐ฉ\right)+ฯตB\left(๐ฉ\right)+\zeta \hfill \end{array}$$
(15)
therefore we have the following relations:
$`\left(\mathrm{\Delta }A\left(๐ฉ\right)\right)^2=\gamma \left(A\left(๐ฉ+1\right)+A\left(๐ฉ\right)\right)+ฯต`$ (16)
$`\alpha A\left(๐ฉ\right)^2+2\gamma A\left(๐ฉ\right)b\left(๐ฉ\right)+\delta A\left(๐ฉ\right)+ฯตB\left(๐ฉ\right)+\zeta =0`$ (17)
while the function $`\rho \left(๐ฉ\right)`$ can be arbitrarily determined. In fact this function can be fixed, in order to have a polynomial structure function $`\mathrm{\Phi }(x)`$ for the deformed oscillator algebra (7).
The solutions of equation (16) depend on the value of the parameter $`\gamma `$, while the function $`b(๐ฉ)`$ is uniquely determined by equation (17) (provided that almost one among the parameters $`\gamma `$ or $`ฯต`$ is not zero). At this stage, the cases $`\gamma 0`$ or $`\gamma =0`$, should be treated separately. We can see that:
* $`\gamma 0`$
In that case the solutions of equations (16) and (17) are given by:
$$A\left(๐ฉ\right)=\frac{\gamma }{2}\left((๐ฉ+u)^21/4\frac{ฯต}{\gamma ^2}\right)$$
(18)
$$\begin{array}{cc}\hfill b\left(๐ฉ\right)=& \frac{\alpha \left((๐ฉ+u)^21/4\right)}{4}+\frac{\alpha ฯต\delta \gamma }{2\gamma ^2}\hfill \\ & \frac{\alpha ฯต^22\delta ฯต\gamma +4\gamma ^2\zeta }{4\gamma ^4}\frac{1}{\left((๐ฉ+u)^21/4\right)}\hfill \end{array}$$
(19)
* $`\gamma =0,ฯต0`$
The solutions of equations (16) and (17) are given by:
$$A(๐ฉ)=\sqrt{ฯต}\left(๐ฉ+u\right)$$
(20)
$$b(๐ฉ)=\alpha \left(๐ฉ+u\right)^2\frac{\delta }{\sqrt{ฯต}}\left(๐ฉ+u\right)\frac{\zeta }{ฯต}$$
(21)
The constant $`u`$ will be determined later.
Using the above definitions of equations $`A(๐ฉ)`$ and $`b(๐ฉ)`$, the left hand side and right hand side of equation (4) gives the following equation:
$$\begin{array}{c}2\mathrm{\Phi }(๐ฉ+1)\left(\mathrm{\Delta }A\left(๐ฉ\right)+\frac{\gamma }{2}\right)\rho (๐ฉ)2\mathrm{\Phi }(๐ฉ)\left(\mathrm{\Delta }A\left(๐ฉ1\right)\frac{\gamma }{2}\right)\rho (๐ฉ1)=\hfill \\ =aA^2\left(๐ฉ\right)\gamma b^2(๐ฉ)2\alpha A\left(๐ฉ\right)b(๐ฉ)+dA\left(๐ฉ\right)\delta b(๐ฉ)+z\hfill \end{array}$$
(22)
Equation (5) gives the following relation:
$$\begin{array}{cc}\hfill K=& \\ \hfill =& \mathrm{\Phi }(๐ฉ+1)\left(\gamma ^2ฯต2\gamma A\left(๐ฉ\right)\mathrm{\Delta }A^2\left(๐ฉ\right)\right)\rho (๐ฉ)+\hfill \\ & +\mathrm{\Phi }(๐ฉ)\left(\gamma ^2ฯต2\gamma A\left(๐ฉ\right)\mathrm{\Delta }A^2\left(๐ฉ1\right)\right)\rho (๐ฉ1)\hfill \\ & 2\alpha A^2\left(๐ฉ\right)b(๐ฉ)+\left(\gamma ^2ฯต2\gamma A\left(๐ฉ\right)\right)b^2(๐ฉ)+\hfill \\ & +2\left(\alpha \gamma \delta \right)A\left(๐ฉ\right)b(๐ฉ)+\left(\gamma \delta 2\zeta \right)b(๐ฉ)+\hfill \\ & +\frac{2}{3}aA^3\left(๐ฉ\right)+\left(d+\frac{1}{3}a\gamma +\alpha ^2\right)A^2\left(๐ฉ\right)+\hfill \\ & +\left(\frac{1}{3}aฯต+\alpha \delta +2z\right)A\left(๐ฉ\right)\hfill \end{array}$$
(23)
Equations (22) and (23) are linear functions of the expressions $`\mathrm{\Phi }\left(๐ฉ\right)`$ and $`\mathrm{\Phi }\left(๐ฉ+1\right)`$, then the function $`\mathrm{\Phi }\left(๐ฉ\right)`$ can be determined, if the function $`\rho (๐ฉ)`$ is given. The solution of this system, i.e. the function $`\mathrm{\Phi }\left(๐ฉ\right)`$ depends on two parameters $`u`$ and $`K`$ and it is given by the following formulae:
* $`\gamma 0`$
$$\rho (๐ฉ)=\frac{1}{32^{12}\gamma ^8(๐ฉ+u)(1+๐ฉ+u)(1+2(๐ฉ+u))^2}$$
and
$$\begin{array}{c}\mathrm{\Phi }(๐ฉ)=3072\gamma ^6K(1+2(๐ฉ+u))^2\hfill \\ 48\gamma ^6(\alpha ^2ฯต\alpha \delta \gamma +aฯต\gamma d\gamma ^2)\hfill \\ (3+2(๐ฉ+u))(1+2(๐ฉ+u))^4(1+2(๐ฉ+u))+\hfill \\ +\gamma ^8(3\alpha ^2+4a\gamma )(3+2(๐ฉ+u))^2(1+2(๐ฉ+u))^4(1+2(๐ฉ+u))^2+\hfill \\ +768(\alpha ฯต^22\delta ฯต\gamma +4\gamma ^2\zeta )^2+\hfill \\ +32\gamma ^4(1+2(๐ฉ+u))^2(112(๐ฉ+u)+12(๐ฉ+u)^2)\hfill \\ (3\alpha ^2ฯต^26\alpha \delta ฯต\gamma +2aฯต^2\gamma +2\delta ^2\gamma ^24dฯต\gamma ^2+8\gamma ^3z+4\alpha \gamma ^2\zeta )\hfill \\ 256\gamma ^2(1+2(๐ฉ+u))^2\hfill \\ (3\alpha ^2ฯต^39\alpha \delta ฯต^2\gamma +aฯต^3\gamma +6\delta ^2ฯต\gamma ^23dฯต^2\gamma ^2+2\delta ^2\gamma ^4+\hfill \\ +2dฯต\gamma ^4+12ฯต\gamma ^3z4\gamma ^5z+12\alpha ฯต\gamma ^2\zeta 12\delta \gamma ^3\zeta +4\alpha \gamma ^4\zeta )\hfill \end{array}$$
(24)
* $`\gamma =0,ฯต0`$
$$\rho (๐ฉ)=1$$
$$\begin{array}{c}\mathrm{\Phi }(๐ฉ)=\hfill \\ =\frac{1}{4}\left(\frac{K}{ฯต}\frac{z}{\sqrt{ฯต}}\frac{\delta }{\sqrt{ฯต}}\frac{\zeta }{ฯต}+\frac{\zeta ^2}{ฯต^2}\right)\hfill \\ \frac{1}{12}\left(3da\sqrt{ฯต}3\alpha \frac{\delta }{\sqrt{ฯต}}+3\left(\frac{\delta }{\sqrt{ฯต}}\right)^26\frac{z}{\sqrt{ฯต}}+6\alpha \frac{\zeta }{ฯต}6\frac{\delta }{\sqrt{ฯต}}\frac{\zeta }{ฯต}\right)(๐ฉ+u)\hfill \\ +\frac{1}{4}\left(\alpha ^2+da\sqrt{ฯต}3\alpha \frac{\delta }{\sqrt{ฯต}}+\left(\frac{\delta }{\sqrt{ฯต}}\right)^2+2\alpha \frac{\zeta }{ฯต}\right)(๐ฉ+u)^2\hfill \\ \frac{1}{6}\left(3\alpha ^2a\sqrt{ฯต}3\alpha \frac{\delta }{\sqrt{ฯต}}\right)(๐ฉ+u)^3+\frac{1}{4}\alpha ^2(๐ฉ+u)^4\hfill \end{array}$$
(25)
## 5 Finite dimensional representations of the algebra $`Qu(3)`$
Let consider a representation of the algebra $`Qu(3)`$, which is diagonal to the generator $`A`$ and the Casimir $`K`$. Using the parafermionic realization defined by equations (12) and (13), we see that this a representation diagonal to the parafermionic number operator $`๐ฉ`$ and the Casimir $`K`$. The basis of a such representation corresponds to the Fock basis of the parafermionic oscillator, i.e. the vectors $`|k,n>,n=0,1,\mathrm{}`$of the carrier Fock space satisfy the equations
$$๐ฉ|k,n>=n|k,n>,K|k,n>=k|k,n>$$
The structure function (24) (or respectively (24) ) depend on the eigenvalues of the of the parafermionic number operator $`๐ฉ`$ and the Casimir $`K`$. The vectors $`|k,n>`$ are also eigenvectors of the generator $`A`$, i.e.
$$A|k,n>=A(k,n)|k,n>$$
In the case $`\gamma 0`$ we find from equation (18)
$$A(k,n)=\frac{\gamma }{2}\left((n+u)^21/4\frac{ฯต}{\gamma ^2}\right)$$
In the case $`\gamma =0,ฯต0`$ we find from equation (20)
$$A(k,n)=\sqrt{ฯต}\left(n+u\right)$$
then the parameter $`u=u(k,p)`$ is a solution of the system of equations (26).
If the deformed oscillator corresponds to a deformed Parafermionic oscillator of order $`p`$ then the two parameters of the calculation $`k`$ and $`u`$ should satisfy the constrints (8) and (10) the system:
$$\begin{array}{c}\mathrm{\Phi }(0,u,k)=0\\ \mathrm{\Phi }(p+1,u,k)=0\end{array}$$
(26)
then the parameter $`u=u(k,p)`$ is a solution of the system of equations (26).
Generally there are many solutions of the above system, but a unitary representation of the deformed parafermionic oscillator is restrained by the additional restriction
$$\mathrm{\Phi }(x)>0,\text{for}0<x<p+1$$
We must point out that the system (26) corresponds to a representation with dimension equal to $`p+1`$.
## 6 Application of the case $`\gamma =0`$
In this section, we shall give an example of the calculation of eigenvalues of a superintegrable two-dimensional system, by using the methods of the previous section. The calculation by an empirical method was performed in and the solution of the same problem by using separation of variables was studied in . Here in order to show the effects of the quantization procedure we donโt use $`\mathrm{}=1`$ as it was considered in references and . That means that the following commutation relations are taken in consideration:
$$[x,p_x]=i\mathrm{},[x,p_x]=i\mathrm{}$$
The superintegrable Holt system corresponds to the Hamiltonian:
$$H=\frac{1}{2}\left(p_x^2+p_y^2\right)+\frac{\omega ^2}{2}\left(4x^2+y^2\right)+\frac{k^2\frac{1}{4}}{y^2}.$$
(27)
This superintegrable system with two integrals:
$$A=p_x^2+4\omega ^2x^2,\text{and}B=\{p_y,xp_yyp_x\}2\omega ^2xy^22(1/4k^2)\frac{x}{y^2}$$
From the above definitions we can verify that:
$$[H,A]=0,[H,B]=0,$$
and
$$[A,B]=C,[A,C]=16\mathrm{}^2\omega ^2B,$$
$$[B,C]=24h^2A^264\mathrm{}^2HA+8\mathrm{}^2\left(4H^2+\omega ^2(14k^2+3\mathrm{}^2)\right)$$
The above algebra is a quadratic algebra $`Qu(3)`$ of the form (24), corresponding to the following values of the coefficients:
$$\begin{array}{c}\alpha =0,\gamma =0,\delta =0,ฯต=16\mathrm{}^2\omega ^2,\zeta =0\\ a=24\mathrm{}^2,d=64\mathrm{}^2H,z=8\mathrm{}^2\left(4H^2+\omega ^2(14k^2+3\mathrm{}^2)\right)\end{array}$$
The value of the Casimir operator (5) is given by:
$$\begin{array}{cc}\hfill K=& C^2+\hfill \\ & +16\mathrm{}^2A^364\mathrm{}^2HA^216\mathrm{}^2\omega ^2B^2+\hfill \\ & +16\mathrm{}^2\left(4H^2+\omega ^2(14k^2+11\mathrm{}^2)\right)A=\hfill \\ \hfill =& 256\mathrm{}^4\omega ^2H\hfill \end{array}$$
The representation of the above algebra, which is diagonal to the Hamiltonian $`H`$ and to the integral of motion $`A`$, corresponds to the eigenvalues of the energy $`E`$ and the eigenvalue of the Casimir equal to $`256\mathrm{}^4\omega ^2E`$ In that case, equations (22) and (23), which determine the function $`\mathrm{\Phi }(x)`$, are respectively:
$$\begin{array}{c}32\mathrm{}^2E^28\mathrm{}\omega \mathrm{\Phi }(x)+8\mathrm{}\omega \mathrm{\Phi }(x+1)8\mathrm{}^2\omega ^2\hfill \\ 24\mathrm{}^4\omega ^2+32\mathrm{}^2k^2\omega ^2+256\mathrm{}^3E\omega (x+u)384\mathrm{}^4\omega ^2(x+u)^2=0\hfill \end{array}$$
$$\begin{array}{c}32\left(\mathrm{\Phi }(x)+\mathrm{\Phi }(x+1)\right)\mathrm{}^2\omega ^2+\hfill \\ +64\mathrm{}^3\omega \left(4E^2+\omega ^2+11\mathrm{}^2\omega ^24k^2\omega ^2\right)(x+u)\hfill \\ 1024\mathrm{}^4E\omega ^2(x+u)^2+1024\mathrm{}^5\omega ^3(x+u)^3=0\hfill \end{array}$$
The above equations can be solved and we find that:
$$\begin{array}{cc}\hfill \mathrm{\Phi }(x)=& \frac{\mathrm{}\left(4E^2+8\mathrm{}E\omega +\omega ^2+3\mathrm{}^2\omega ^24k^2\omega ^2\right)}{2\omega }\hfill \\ & +\frac{\mathrm{}\left(4E^2+16\mathrm{}E\omega +\omega ^2+11\mathrm{}^2\omega ^24k^2\omega ^2\right)(x+u)}{\omega }\hfill \\ & 8\mathrm{}^2\left(2E+3\mathrm{}\omega \right)(x+u)^2+16\mathrm{}^3\omega (x+u)^3\hfill \end{array}$$
The two parameters of this equation are the parameter $`u`$ and the eigenvalue $`E`$ of the energy $`H`$, therefore we can solve the system:
$$\mathrm{\Phi }(0)=0,\mathrm{\Phi }(p+1)=0$$
and we find two solutions:
$$u=\frac{1}{2},E=\frac{\omega }{2}\left(4\mathrm{}(1+p)\pm \sqrt{1+\mathrm{}^2+4k^2}\right)$$
## 7 Application of the case $`\gamma 0`$
Let consider the case of a potential on the two dimensional hyperboloid taken from ref (case of the potential $`V_1`$).
The two dimensional hyperpoloid is characterized by the cartesian coordinates $`\omega _0,\omega _1,\omega _2`$, which obey to the restriction $`\omega _0^2(\omega _1^2+\omega _1^2)=1`$. The Hamiltonian is given by
$$H=\frac{1}{2}\mathrm{}_{LB}+V$$
where $`\mathrm{}_{LB}`$ is the Laplace - Beltrami operator for the details see ref , where
$$V=\frac{\alpha ^2}{\omega _2^2}\frac{\gamma ^2}{(\omega _0\omega _1)^2}+\beta ^2\frac{\omega _0+\omega _1}{\omega _0\omega _1)^3}$$
The two integrals are
$$A=\left(\omega _0_{\omega _1}+\omega _1_{\omega _0}\right)^22\beta ^2\left(\frac{\omega _0+\omega _1}{\omega _0\omega _1}\right)^2+2\gamma ^2\frac{\omega _0+\omega _1}{\omega _0\omega _1}$$
and
$$B=\left(\omega _0_{\omega _2}+\omega _2_{\omega _0}\omega _1_{\omega _2}+\omega _2_{\omega _1}\right)^2\frac{2\beta ^2\omega _2^2}{(\omega _0\omega _1)^2}\frac{2\alpha ^2(\omega _0\omega _1)^2}{\omega _2^2}$$
The operators $`H,A`$ and $`B`$ satisfy the following commutation relations:
$$[H,A]=0,[H,A]=0$$
$$[A,C]=8\{A,B\}+16\gamma ^2A16B16\gamma ^2(12\alpha ^22H)$$
$$[B,C]=8B^232\beta ^2A16\gamma ^2B16\beta ^2(14\alpha ^2+4H)$$
The above algebra is a quadratic algebra $`Qu(3)`$ of the form (24), corresponding to the following values of the coefficients:
$$\begin{array}{c}\alpha =0,\gamma =8,\delta =16\gamma ^2,ฯต=16,\zeta =16\gamma ^2(12\alpha ^22H)\\ a=0,d=32\beta ^2,z=16\beta ^2(14\alpha ^2+4H)\end{array}$$
The value of the Casimir operator (5) is given by:
$$\begin{array}{cc}\hfill K=& C^2\hfill \\ & \frac{8}{3}\{A,B,B\}+\frac{176}{3}B^232\beta ^2A^216\gamma ^2\{A,B\}\hfill \\ & \left(64\alpha ^2\gamma ^2+64\gamma ^2H\frac{352}{3}\gamma ^2\right)B\left(\frac{352}{3}128\alpha ^2+128\beta ^2H\right)\beta ^2A\hfill \\ & \frac{128}{3}\beta ^2(14\alpha ^2+4H)=\hfill \\ \hfill =& 16(4\beta ^2+8\alpha ^2\beta ^2+8\alpha ^4\beta ^23\gamma ^4+\hfill \\ & +8\alpha ^2\gamma ^48\beta ^2H+16\alpha ^2\beta ^2H+8\beta ^2H^2)\hfill \end{array}$$
The representation of the above algebra, which is diagonal to the Hamiltonian $`H`$, to the integral of motion $`A`$, corresponds to the eigenvalues of the energy $`E`$ and the eigenvalue of the Casimir equal to above cited value.
In that case, equations (22) and (23), which determine the function $`\mathrm{\Phi }(x)`$, give the form (24):
$$\begin{array}{cc}\hfill \mathrm{\Phi }(x)=& 32^{34}(2\beta ^2\gamma ^48\beta ^2(u+x)+8\beta ^2(u+x)^2)\hfill \\ & (\alpha ^2+\alpha ^4+E+2\alpha ^2E+E^2+(1+4\alpha ^24E)(u+x)+\hfill \\ & +(54\alpha ^2+4E)(u+x)^28(u+x)^3+4(u+x)^4)\hfill \end{array}$$
The two parameters of this equation are the parameter $`u`$ and the eigenvalue $`E`$ of the energy $`H`$, therefore we can solve the system:
$$\mathrm{\Phi }(0)=0,\mathrm{\Phi }(p+1)=0$$
and we find the solution:
$$u=\frac{1}{2}\left(1\frac{\gamma ^2}{\sqrt{2}\beta }\right),E=\frac{1}{2}\left(2p+2+\sqrt{2\alpha ^2+1/4}\frac{\gamma ^2}{\sqrt{2}\beta }\right)^2+\frac{1}{8}$$
## 8 Discussion
From the above discussion, we have shown how to calculate finite dimensional representations of the $`Qu(3)`$ algebra and we have given an application of this method in the calculation of the energy eigenvalues of the superintegrable systems. The systematic algebraic study of all the known superintegrable systems is under investigation. |
warning/0002/hep-ph0002180.html | ar5iv | text | # A short course in effective Lagrangians. Lectures delivered at the VII Mexican Workshop on Particles and Fields, Merida, Yucatan, Mexico, 10-17 November, 1999.
## I Introduction.
When studying a physical system it is often the case that there is not enough information to provide a fundamental description of some of its properties. In such cases one must parameterize the corresponding effects by introducing new interactions with coefficients to be determined phenomenologically. Experimental limits or measurement of these parameters then (hopefully) provides the information needed to provide a more satisfactory description.
A standard procedure for doing this is to first determine the dynamical degrees of freedom involved and the symmetries obeyed, and then construct the most general Lagrangian, the effective Lagrangian for these degrees of freedom which respects the required symmetries. The method is straightforward, quite general and, most importantly, it works!
In following this approach one must be wary of several facts. Fist it is clear that the relevant degrees of freedom can change with scale (e.g. mesons are a good description of low-energy QCD, but at higher energies one should use quarks and gluons); in addition, physics at different scales may respect different symmetries (e.g. mass conservation is violated at sufficiently high energies). It follows that the effective Lagrangian formalism is in general applicable only for a limited range of scales. It is often the case (but no always!) that there is a scale $`\mathrm{\Lambda }`$ so that the results obtained using an effective Lagrangian are invalid for energies above $`\mathrm{\Lambda }`$.
The formalism has two potentially serious drawbacks. First, effective Lagrangian has an infinite number of terms suggesting a lack of predictability. Second, even though the model has an UV cutoff $`\mathrm{\Lambda }`$ and will not suffer from actual divergences, simple calculations show that is is a possible for this type of theories to generating radiative corrections that grow with $`\mathrm{\Lambda }`$, becoming increasingly important for higher and higher order graphs. Either of these problems can render this approach useless. It is also necessary verify that the model is unitary.
I will discuss below how these problems are solved, an provide several applications of the formalism. The aim is to give a flair of the versatility of the approach, not to provide an exhaustive review of all known applications.
## II Familiar examples
### A Euler-Heisenberg effective Lagrangian
This Lagrangian summarizes QED at low energies (below the electron mass) . At these energies only photons appear in real processes and the effective Lagrangian will be then constructed using the photon field $`A_\mu `$, and will satisfy a $`U(1)`$ gauge and Lorenz invariances. Thus it can be constructed in terms of the field strength $`F_{\mu \nu }`$ or the loop variables $`๐(\mathrm{\Gamma })=_\mathrm{\Gamma }A๐x`$. The latter are non-local, so that a local description would involve only $`F`$, namely <sup>*</sup><sup>*</sup>* There is no $`F\stackrel{~}{F}`$ terms since it is a total derivative.
$`_{\mathrm{eff}}`$ $`=`$ $`_{\mathrm{eff}}(F)`$ (1)
$`=`$ $`aF^2+bF^4+c(F\stackrel{~}{F})^2+dF^2(F\stackrel{~}{F})\mathrm{}`$ (2)
One can arbitrarily normalize the fields and so choose $`a=1/4`$. The constants $`b,c`$ and $`d`$ have units of mass<sup>-2</sup>.
Note that the term $`d`$ violates CP. Though we know QED respects C and P, it is possible for other interactions to violate these symmetries, there is nothing in the discussion above that disallows such terms and, in fact, weak effects will generate them. For this system we are in a privileged position for we know the underlying physics, and so we can calculate $`b,c,d,\mathrm{}`$. The leading effects come form QED which yields $`b,c1/(4\pi m_e)^2`$ at 1 loop . The parameters $`b`$ and $`c`$ summarize all the leading virtual electron effects. (see Fig. 1). Forgetting about this underlying structure we could have simply defined a scale $`M`$ and taken $`b,c1/M^2`$ (so that $`M=4\pi m_e`$), and while this is perfectly viable, $`M`$ is not relevant phenomenologically speaking as it does not corresponds of a physical scale. In order to extract information about the physics underlying the effective Lagrangian from a measurement of $`b`$ and $`c`$ we must be able to at least estimate the relation between these constants and the underlying scales.
In addition we also know that $`d\xi /(4\pi v)`$ with $`v246\text{GeV}`$ and $`\xi `$ is a very small constant proportional to the Jarlskog determinant . The effective Lagrangian can hold terms with radically different scales and limits on some constants cannot, in general, translate to others. In this case the terms are characterized by different CP transformation properties, and it is often the case that such global symmetries are useful in differentiating terms in the effective Lagrangian. The point being that a term violating a given global symmetry at scale $`\mathrm{\Lambda }`$ will generate all terms in the effective Lagrangian with the same symmetry properties through radiative corrections. The caveat in the argument being that the underlying theory might have some additional symmetries not apparent at low energies which might further segregate interactions and so provide different scales for operators with the same properties under all low energy symmetries.
When calculating with the effective Lagrangian the effects produced by the new terms proportional to $`b,c`$ are suppressed by a factor $`(E/4\pi m_e)^2`$, where $`E`$ is the typical energy on the process and $`Em_e`$. Thus the effects of these terms are tiny, yet they are noticeable because they generate a new effect: $`\gamma \gamma `$ scattering.
### B (Standard) Superconductivity
This is a brief summary of the very nice treatment provided by Polchinski . The system under consideration has the electron field $`\psi `$ as its only dynamical variable (the phonons are assumed to have been integrated out, generating a series of electron self-interactions), it respects $`U(1)`$ electromagnetic gauge invariance, as well as Galilean invariance and Fermion number conservation.
Assuming a local description, the first few terms in the effective Lagrangian expansion are (neglecting those containing photons for simplicity)
$$_{\mathrm{eff}}=_k\psi _๐ค^{}\left[i_te_๐ค+\mu \right]\psi _๐ค+\psi _๐ค^{}\psi _๐ฅ\psi _๐ช\psi _๐ฉ^{}\delta (๐ค๐ฅ๐ช+๐ฉ)V_{\mathrm{๐ค๐ฅ๐ช}}+\mathrm{}$$
(3)
In this equation the relation $`e_๐ค=\mu `$ determines the Fermi surface, while $`V\frac{(\text{electron-photon coupling})^2}{(\text{phonon mass})^2}`$ summarizes the virtual phonon effects. In order to determine the importance of the various terms we need the dimensions of the field $`\psi `$. A vector k lies on the Fermi Surface (FS) if $`e_๐ค=\mu `$, if p is near the FS one can write $`๐ฉ=๐ค+\mathrm{}\widehat{๐ง}`$ (with $`e_๐ค=\mu `$). Scaling towards the FS implies $`\mathrm{}s\mathrm{}`$ with $`s0`$. Then assuming $`\psi s^d\psi `$ the quadratic terms in the action will be scale invariant provided $`d=1/2`$. The quartic terms in the action then scales as $`s`$ and becomes negligible near the FS except when the pairing condition $`๐ช+๐ฅ=0`$ is obeyed. In this case the quartic term scales as $`s^0`$ and cannot be ignored. In fact this term determines the most interesting behavior of the system at low temperatures (see for full details).
### C Electroweak interactions
Again I will follow the general recipe. I will concentrate only on the (low energy) interactions involving lepton fields, which are then the degrees of freedom. Since I assume the energy to be well below the Fermi scale, the only relevant symmetries are $`U(1)`$ gauge and Lorenz invariances. In addition there is the question whether the heavy physics will respect the discrete symmetries $`C`$, $`P`$ or $`CP`$; using perfect hindsight I will retain terms that violate these symmetries
Assuming a local description I have
$$_{\mathrm{eff}}=\overline{\psi }_i(i\overline{)}Dm_i)\psi _i+f_{ijkl}\left(\overline{\psi }_i\mathrm{\Gamma }^a\psi _j\right)\left(\overline{\psi }_k\mathrm{\Gamma }_a\psi _l\right)+\mathrm{}$$
(4)
where the ellipsis indicate terms containing operators of higher dimension, or those involving the electromagnetic field. The matrices $`\mathrm{\Gamma }`$ are to be chosen among the 16 independent basis $`\mathrm{\Gamma }^a=\{1,\gamma _\mu ,\sigma _{\mu \nu },\gamma _\mu \gamma _5,\gamma _5\}`$
The coefficients for the first two terms are be fixed by normalization requirements. While a SM calculation gives $`fg^2/m_W^2=1/v^2`$ ($`v246`$GeV) and is generated by tree-level graphs (see Fig. 2) because of this the scale $`1/\sqrt{f}`$ is, in fact, the scale of the heavy physics and so the model is applicable at energies swell below $`v`$. The four fermion interactions summarize the leading virtual gauge boson effects. The contributions of the four-fermion operators to processes with typical energy $`E`$ are suppressed by a factor $`E^2/v^2`$. These can be observed (or bounded) despite the $`Ev`$ condition because they generate new effects: $`C`$ and $`P`$ (and some of them chirality) violation.
### D Strong interactions at low energies
In this case we are interested in the description of the interactions among the lightest hadrons, the meson multiplet. The most convenient parameterization of these degrees of freedom is in terms of a unitary field $`U`$ such that $`U=\mathrm{exp}(\lambda _a\pi ^a/F)`$ where $`\pi ^a`$ denote the eight meson fields, $`\lambda ^a`$ the Gell-Mann matrices and $`F`$ is a constant (related to the pion decay constant). The symmetries obeyed by the system are chiral $`SU(3)_L\times SU(3)_R`$, Lorenz invariance, $`C`$ and $`P`$.
With these constraints the effective Lagrangian takes the form
$$_{\mathrm{eff}}=a\text{tr}U^{}U+\left[b\text{tr}_\mu U^{}_\nu U^\mu U^{}^\nu U+\mathrm{}\right]+\mathrm{}$$
(5)
I can set $`aF^2`$ by properly normalizing the fields. In this case the leading term in the effective Lagrangian will determine all (leading) low-energy pion interactions in terms of the single constant $`F`$. The effects form the higher-order terms have been measured and the data requires $`b1/(4\pi )^2`$. This result is also predicted by the consistency of this approach which requires that radiative corrections to $`a`$, $`b`$, etc. should be at most of the same size as their tree-level values.
## III Basic ideas on the applicability of the formalism
Being a model with intrinsic an cutoff there are no actual ultraviolet divergences in most effective Lagrangian computations. Still there are interesting renormalizability issues that arise when doing effective Lagrangian loop computations.
Imagine doing a loop calculation including some vertices terms of (mass) dimension higher than the dimension of space-time. These must have coefficients with dimensions of mass to some negative power. The loop integrations will produce in general terms growing with $`\mathrm{\Lambda }`$ the UV cutoff which are polynomials in the external momenta Since a graph can be rendered convergent by taking sufficient number of derivatives with respect to the external momenta. and will preserve the symmetries of the model . Hence these terms which may grow with $`\mathrm{\Lambda }`$ correspond to vertices appearing in the most general effective Lagrangian and can be absorbed in a renormalization of the corresponding coefficients. They have no observable effects (though they can be used in naturality arguments .
Effective theories will also be unitary provided one stays within the limits of their applicability. Should one exceed them new channels will open (corresponding to the production of the heavy excitations) and unitarity violating effects will occur. This is not produced by real unitarity violating interactions, but due to our using the model beyond its range of applicability (e.g. it the typical energy of the process under consideration reaches of exceeds $`\mathrm{\Lambda }`$). One can, of course, extend the model, but this necessarily introduces ad-hoc elements and will dilute the generality gained using effective theories.
For example consider $`WWZ`$ interactions with an effective Lagrangian of the form
$$_{\mathrm{eff}}=\lambda (p,k)W_{\mu \nu }(k)W^{\nu \rho }(p)Z_\rho {}_{}{}^{\mu }(pk)+\mathrm{};$$
(6)
(where $`V_{\alpha \beta }=_\alpha V_\beta _\beta V_\alpha `$) One can then choose $`\lambda `$ to insure unitarity is preserved (at least in some processes), for example
$$\lambda (p,k)=\frac{\lambda _0}{(pk+\mathrm{\Lambda })^n}$$
(7)
which, for $`n`$ sufficiently large insures that the cross section for the reaction $`e^+e^{}ZWW`$ is unitary, since it behaves as $`s^{22n}`$ for a CM energy$`=s\mathrm{\Lambda }^2`$. But the very same effective vertex also modifies other reactions such as, for example $`u\overline{d}WZW`$ where the cross section now has a factor $`(s\mathrm{\Lambda }^2)^{2n}`$ and will exhibit resonant behavior if $`s\mathrm{\Lambda }^2`$. If one requires $`s\mathrm{\Lambda }^2`$ (as required by the consistency of the formalism) there are neither unitarity violations nor resonance effects. If, however, one uses the above Ansatz to extend the range of applicability to $`s\mathrm{\Lambda }^2`$ and beyond then very clear resonances should be observed in hadron colliders. Given that these have not been observed one must use for $`\mathrm{\Lambda }`$ a value significantly larger than the average CM energy for the hard $`W`$ pair production cross section.
## IV Using effective Lagrangians
Effective Lagrangians provide an efficient way of summarizing some (perhaps very complex) interactions. The idea is simply to include all the effective vertices produces by those excitations which are not directly observed.
For example given a real scalar field $`\varphi `$ and assume that all Fourier components above a scale $`\mathrm{\Lambda }`$ are not directly observable (ฤฑ.e. the available energies lie all below $`\mathrm{\Lambda }`$), then the effective Lagrangian is obtained by integrating over the variables observable at energies $`\mathrm{\Lambda }`$; writing $`\varphi =\varphi _0+\varphi _1`$, with
$$\varphi _0(๐ค):|๐ค|<\mathrm{\Lambda }\varphi _1(๐ค):\mathrm{\Lambda }|๐ค|<\mathrm{\Lambda }_1$$
(8)
then by definition
$`e^{iS_{\mathrm{eff}}}`$ $`=`$ $`{\displaystyle [d\varphi _1]e^{iS(\varphi _0,\varphi _1)}},S_{\mathrm{eff}}={\displaystyle d^nx_{\mathrm{eff}}}`$ (9)
where $`_{\mathrm{eff}}`$ is obtained by expanding $`S_{\mathrm{eff}}`$ in powers of $`\mathrm{\Lambda }`$ which gives an infinite tower of local operators.
Another common situation where effective Lagrangians appear occurs when some heavy excitations are integrated out. This can be illustrated by the following toy model Iโm cheating in order to get a closed form for the effective action, a more realistic model should include a term $`\varphi _1^4`$
$$S=d^nx\left[\overline{\psi }(i\overline{)}m)\psi +\frac{1}{2}(\varphi )^2\frac{1}{2}\mathrm{\Lambda }^2\varphi _1^2+f\varphi \overline{\psi }\psi \right]$$
(10)
where $`\varphi `$ is heavy. A simple calculation gives
$$S_{\mathrm{eff}}=d^nx\left[\overline{\psi }(i\overline{)}m)\psi +\frac{1}{2}f^2\overline{\psi }\psi \frac{1}{\text{ }\text{ }\text{ }\text{ }\text{ }+\mathrm{\Lambda }^2}\overline{\psi }\psi \right]$$
(11)
and
$$_{\mathrm{eff}}=\overline{\psi }(i\overline{)}m)\psi +\frac{f^2}{2\mathrm{\Lambda }^2}\underset{l=1}{\overset{\mathrm{}}{}}\overline{\psi }\psi \left(\frac{\text{ }\text{ }\text{ }\text{ }\text{ }}{\mathrm{\Lambda }^2}\right)^n\overline{\psi }\psi $$
(12)
Note that terms with large number of derivatives will be suppressed by a large power of the small factor $`(E/\mathrm{\Lambda })`$, if we are interested in energies $`E\mathrm{\Lambda }`$ the whole infinite set of vertices must be included in order to reproduce the $`\varphi `$ pole.
### A How to parameterize ignorance
If one knows the theory we can, in principle, calculate $`_{\mathrm{eff}}`$ (or do a full calculation). Yet there are many cases where the underlying theory is not known. In these cases an effective theory if obtained by writing all possible interactions among the light excitations. The model then has an infinite number of terms each with an unknown parameter, and these constants then parameterize all possible underlying theories. The terms which dominate are those usually called renormalizable (or, equivalently, marginal or relevant). The other terms are called non-renormalizable, or irrelevant, since their effects become smaller as the energy decreases
This recipe for writing effective theories must be supplemented with some symmetry restrictions. The most important being that the all the terms in the effective Lagrangian must respect the local gauge invariance of the low-energy physics (more technically, the one respected by the renormalizable terms in the effective action) . The reason is that the presence of a gauge variant term will generate all gauge variant interactions thorough renormalization group evolution.
##### a Gauge invariantizing
Using a simple argument it is possible to turn any theory into a gauge theory and so it appears that the requirement of gauge invariance is empty. That this is not the case is explained here. I first describe the trick which grafts gauge invariance onto a theory and then discuss the implications.
Consider an arbitrary theory with matter fields (spin 0 and 1/2) and vector fields $`V_\mu ^n`$, $`n=1,\mathrm{}N`$. Then
* Choose a (gauge) group $`G`$ with $`N`$ generators $`\{T^n\}`$. Define a covariant derivative $`D_\mu =_\mu +V_\mu ^nT^n`$ and assume that the $`V_\mu ^n`$ are gauge fields.
* Invent a unitary field $`U`$ transforming according to the fundamental representation of $`G`$ and construct the gauge invariant composite fields
$$๐ฑ_\mu ^n=\text{tr}T^nU^{}D_\mu U$$
(13)
Taking $`\text{tr}T^nT^m=\delta _{nm}`$, it is easy to see that in the unitary gauge $`U=1`$, $`๐ฑ_\mu ^n=V_\mu ^n`$.
Thus if simply replace $`V๐ฑ`$ in the original theory we get a gauge theory. Does this mean that gauge invariance irrelevant since it can be added at will? In my opinion this is not the case.
In the above process all matter fields are assumed gauge singlets (none are minimally coupled to the gauge fields).In the case of the standard model , for example, the universal coupling of fermions to the gauge bosons would be accidental in this approach. In order to recover the full predictive power commonly associated with gauge theories, the matter fields must transform non-trivially under $`G`$ which can be done only if there are strong correlations among some of the couplings. It is not trivial to say that the standard model group is $`SU(3)\times SU(2)\times U(1)`$ with left-handed quarks transforming as $`(3,2,1/6)`$, left-handed leptons as $`(1,2,1/2)`$, etc., as opposed to a $`U(1)^{12}`$ with all fermions transforming as singlets .
### B How to estimate ignorance
A problem which I have not addressed so far is the fact that effective theories have an infinite number of coefficients, with the (possible) problem or requiring an infinite number of data points in order to make any predictions. On the other hand, for example, if this is the case why is it that the Fermi theory of the weak interactions is so successful?
The answer to this question lies in the fact that not all coefficients are created equal, there is a hierarchy . As a result, given any desired level of accuracy, only a finite number of terms need to be included. Moreover, even though the effective Lagrangian coefficients cannot be calculated without knowing the underlying theory, they can still be bounded using but a minimal set of assumptions about the heavy interactions. It is then also possible to estimate the errors in neglecting all but the finite number of terms used.
As an example consider the standard model at low energies and calculate two processes: Bhaba cross section and the anomalous magnetic moment of the electron. For Bhaba scattering there is a contribution due the $`Z`$-boson exchange (see Fig. 2)
$$e^+e^{}Ze^+e^{}\text{generates}๐ช=\frac{1}{2m_Z^2}\left(\overline{e}\mathrm{\Gamma }\gamma ^\mu e\right)\left(\overline{e}\mathrm{\Gamma }\gamma _\mu e\right)$$
(14)
where $`\mathrm{\Gamma }=g_V+g_A\gamma _5`$. The coefficient of the effective operator $`๐ช`$ is then $`(\text{coupling}/\text{physical mass})^21/v^2`$
The electron anomalous magnetic moment receives contributions from virtual $`W`$, $`Z`$ and $`H`$ exchanges (see Fig. 3). The corresponding low-energy operator is
$$๐ช=\overline{e}\sigma _{\mu \nu }eF^{\mu \nu }$$
(15)
In this case the coefficient $`\{\text{coupling}/[4\pi (\text{physical mass})]\}^21/(4\pi v)^2`$ <sup>ยง</sup><sup>ยง</sup>ยงIn addition the coefficient is suppressed by a factor of $`m_e`$ since it violates chirality..
The point of this exercise is to illustrate the fact that, for weakly coupled theories, loop-generated operators have smaller coefficients than operators generated at tree level. Leading effects are produced by operators which are generated at tree level.
### C Coefficient estimates
In this section I will provide arguments which can be used to estimate (or, at least bound) the coefficients in the effective Lagrangian. These are order of magnitude calculations and might be off by a factor of a few; it is worth noting that no single calculation has provided a significant deviation from these results.
The estimate calculations should be done separately for weakly and strongly interacting theories. I will characterize the first as those where radiative corrections are smaller than the tree-level contributions. Strongly interacting theories will have radiative corrections of the same size at any order Should the radiative corrections increase with the order of the calculation, it is likely that the dynamic variables being used are not appropriate for the regime where the calculation is being done.
#### 1 Weakly interacting theories
In this case leading terms in the effective Lagrangian are those which can be generated at tree level by the heavy physics. Thus the dominating effects are produced by operators which have the lowest dimension (leading to the smallest suppression from inverse powers of $`\mathrm{\Lambda }`$) and which are tree-level generated (TLG) operators can be determined .
When the heavy physics is described by a gauge theory it is possible to obtained all TLG operators . The corresponding vertices fall into 3 categories, symbolically
* vertices with $`4`$ fermions.
* vertices with $`2`$ fermions and $`k`$ bosons; $`k=2,3`$
* vertices with $`n`$ bosons; $`n=4,6`$.
A particular theory may not generate one or more of these vertices, the only claim is that there is a gauge theory which does.
In the case of the standard model with lepton number conservation the leading operators have dimension 6 . Subleading operators are either dimension 8 and their contributions are suppressed by an additional factor $`(E/\mathrm{\Lambda })^2`$ in processes with typical energy $`E`$. Other subleading contributions are suppressed by a loop factor $`1/(4\pi )^2`$. Note that it is possible to have situations where the only two effects are produced by either dimension 8 TLG operators or loop generated dimension 6 operators. In this case the former dominates only when $`\mathrm{\Lambda }>4\pi E`$.
##### a Triple gauge bosons
The terms in the electroweak effective Lagrangian which describe the interaction of the $`W`$ and $`Z`$ bosons generated by some heavy physics underlying the standard model has received considerable attention recently . In terms of the $`SU(2)`$ and $`U(1)`$ gauge fields $`W`$ and $`B`$ and the scalar doublet $`\varphi `$ these interactions are
$`_{\mathrm{eff}}`$ $`=`$ $`{\displaystyle \frac{1}{\mathrm{\Lambda }^2}}\left(\alpha _W๐ช_W+\alpha _{BW}๐ช_{BW}\right)`$ (16)
$`๐ช_W`$ $`=`$ $`ฯต_{IJK}W_{\mu \nu }^IW^J{}_{}{}^{\nu }{}_{\lambda }{}^{}W_{}^{K}^{\lambda \mu }`$ (17)
$`๐ช_{WB}`$ $`=`$ $`\varphi ^{}\tau ^I\varphi W_{\mu \nu }^IB^{\mu \nu }`$ (18)
The above arguments inly that there is no TLG operator containing three gauge bosons. This means that all effective contributions to the $`WWZ`$ and $`WW\gamma `$ interactions are loop generated, so their coefficients necessarily take the form $`(\text{coupling constants})/(16\pi ^2)`$. Thus the parameters $`\kappa `$ and $`\lambda `$ commonly used to parameterize these interactions are of order $`5\times 10^3`$. Experiments providing limits significantly above this value provide no information about the heavy physics.
#### 2 Strongly interacting theories
I will imagine a theory containing scalars and fermions which interact strongly. Gauge couplings are assumed to be small and will be ignored. This calculation is useful for low energy chiral theories but not for low energy QCD .
A generic effective operator in this type of theories takes the form
$$๐ช_{abc}\lambda \mathrm{\Lambda }^4\left(\frac{\varphi }{\mathrm{\Lambda }_\varphi }\right)^a\left(\frac{\psi }{\mathrm{\Lambda }_\psi }^{3/2}\right)^b\left(\frac{}{\mathrm{\Lambda }}\right)^c$$
(19)
Then the condition that these dynamic variables appropriately describe the physics below $`\mathrm{\Lambda }`$ implies that radiative corrections to the couplings are at most as large as the tree-level values, namely $`\delta _{\mathrm{rad}}\lambda \lambda `$. A straightforward estimate (including a factor of $`1/(16\pi ^2)`$ for each loop) shows that this condition is satisfied only if
$$\mathrm{\Lambda }_\psi =\frac{1}{(4\pi )^{2/3}}\mathrm{\Lambda },\mathrm{\Lambda }_\varphi =\frac{1}{4\pi }\mathrm{\Lambda },\lambda =\frac{1}{16\pi ^2}$$
(20)
In terms of $`U\mathrm{exp}(\varphi /\mathrm{\Lambda }_\varphi )`$, the operators take the form
$$๐ช_{abc}=\frac{1}{(4\pi )^{2b}}\mathrm{\Lambda }^{4c3b/2}^cU^a^{}\psi ^b$$
(21)
In particular the coefficient of the two derivative operators $`\text{tr}U^{}U`$ is $`\mathrm{\Lambda }_\varphi ^2`$.
For the case where $`\varphi `$ represents the interpolating field for the lightest mesons PCAC implies $`\mathrm{\Lambda }_\varphi =f_\pi `$ . Then
$$\psi ^4\frac{1}{f_\pi ^2}^4U^4\frac{1}{16\pi ^2}\psi ^2^2U^2\frac{1}{4\pi f_\pi }$$
(22)
(note that these are upper bounds). The extensive data on low energy meson reactions can be used to gauge the validity of these predictions, they are indeed satisfied. In particular the $`(U)^4`$ terms have coefficients $`1/(16\pi ^2)`$.
For the case of the standard model the field $`U`$ can be used to provide masses for the $`W`$ and $`Z`$ bosons without a physical Higgs being present (the price is that the model breaks down at energies $`4\pi v=3`$TeV). In this case the gauge fields are introduced minimally and it is the term $`(DU)^2`$ gives a mass to the $`W`$ and $`Z`$ which fixes $`\mathrm{\Lambda }_\varphi =v=246\text{GeV}`$ whence $`\mathrm{\Lambda }=3\text{TeV}`$; as before, the model makes no sense beyond this scale Tough it is conceivable that a full non-perturbative calculation would show that the theory cures itself and can be extended beyond this scale, there is no indication that this miracle occurs. In addition, when the gauge fields are reintroduced, the terms with 4 derivatives will generate triple-vector boson couplings, again leading to the estimates $`\lambda ,\kappa 5\times 10^3`$ .
### D Radiative corrections
Despite the presence of higher-dimensional operators radiative corrections can be calculated in the usual way. As an example imagine calculating the corrections to the cross section for the reaction $`e^+e^{}e^+e^{}`$ using the standard model with the addition of a 4-fermion interaction
$$_{\mathrm{eff}}=_{\mathrm{eff}}^{\mathrm{SM}}+\frac{f}{\mathrm{\Lambda }^2}\left(\overline{\psi }\gamma ^\mu \psi \right)\left(\overline{\psi }\gamma _\mu \psi \right)+\mathrm{}$$
(23)
where $`\psi `$ denotes the electron field.
The calculation is illustrated in Fig. 4 where the loops involving the 4-fermion operator are cut-off at a scale $`\mathrm{\Lambda }`$. The SM and new physics (NP) contributions are, symbolically,
SM: $`{\displaystyle \frac{1}{v^2}}\left[1+{\displaystyle \frac{g^2}{16\pi ^2}}+\mathrm{}\right]`$ (24)
NP: $`{\displaystyle \frac{f}{\mathrm{\Lambda }^2}}\left[1+{\displaystyle \frac{f}{16\pi ^2}}+\mathrm{}\right]`$ (25)
Note that this consistent behavior (that the new physics effects disappear as $`\mathrm{\Lambda }\mathrm{}`$) results form having the physical scale of new physics $`\mathrm{\Lambda }`$ in the coefficient of the operator. Had we used $`f^{}/v^2`$ instead of $`f/\mathrm{\Lambda }^2`$ the new physics effects would appear to be enormous, and growing with each new loop. It is not that the use of $`f^{}/v^2`$ is wrong, it is only that it is misleading to believe $`f^{}`$ can be of order one; it must be suppressed by the small factor $`(v/\lambda )^2`$.
Using these results we see that this reaction is sensitive to $`\mathrm{\Lambda }`$ provided $`f(v/\mathrm{\Lambda })^2>`$sensitivity. If the sensitivity is, say 1% this corresponds to $`\mathrm{\Lambda }/\sqrt{f}>2.5`$TeV .
This perturbative calculation is manageable provided $`f<16\pi ^2`$, otherwise the underlying physics is strongly coupled. It is still possible in that case to provide estimates of the new physics contributions, though these are less reliable, these estimates imply that $`1+f/(4\pi )^2+\mathrm{}1`$ when $`f16\pi ^2`$.
## V Applications to electroweak physics
With the above results one can determine, for any given process, the leading contributions (as parameterized by the various effective operator coefficients). Using then the coefficient estimates one can provide the expected magnitude of the new physics effects with only $`\mathrm{\Lambda }`$ as an unknown parameter, and so estimate the sensitivity to the scale of new physics.
It is important to note that this is sometimes a rather involved calculation as all contributing operators must be included. For example, in order to determine the heavy physics effects on the oblique parameters one must calculate not only these affecting the vector boson polarization tensors, but also this which modify the Fermi constant, the fine structure constant, etc. as these quantities are used when extracting $`S`$, $`T`$ and $`U`$ from the data .
### A Effective lagrangian
In the following I will assume that the underlying physics is weakly coupled and derive the leadingoperators that can be expected form the existence of heavy excitations at scale $`\mathrm{\Lambda }`$.
The complete list of dimension 6 operators was cataloged a long time ago for the case where the low energy spectrum includes a single scalar doublet <sup>\**</sup><sup>\**</sup>\**More complicated scalar sectors have also been studied , though not exhaustively.. It is then straightforward to determine the subset of operators which can be TLG, they are
* Fermions: $`\left(\overline{\psi }_i\mathrm{\Gamma }^a\psi _j\right)\left(\overline{\psi }_k\mathrm{\Gamma }^a\psi _l\right)`$
* Scalars: $`|\varphi |^6,(|\varphi |^2)^2`$
* Scalars and fermions: $`|\varphi |^2\times \text{Yukawa term}`$
* Scalars and vectors: $`|\varphi |^2|D\varphi |^2,|\varphi ^{}D\varphi |^2`$
* Fermions, scalars and vectors: $`\left(\varphi ^{}T^nD^\mu \varphi \right)\left(\overline{\psi }_iT^n\gamma _\mu \psi _j\right)`$
where $`T`$ denotes a group generator and $`\mathrm{\Gamma }`$ a product of a group generator and a gamma matrix.
Observables affected by the operators in this list provide the highest sensitivity to new physics effects provided that the standard model effects are themselves small (or that the experimental sensitivity is large enough to observe small deviations). I will illustrate this with two (incomplete) examples
### B b-parity
This is a proposed method for probing new flavor physics . Its virtue lies in the fact that it is very simple and sensitive (though it does not provide the highest sensitivity for all observables). The basic idea is based on the observation that the standard model acquires an additional global $`U(1)_b`$ symmetry in the limit $`V_{ub}=V_{cb}=V_{td}=V_{ts}=0`$ (given the experimental values $`0.002<|V_{ub}|<0.005`$, $`0.036<|V_{cb}|<0.046`$, $`0.004<|V_{td}|<0.014`$, $`0.034<|V_{ts}|<0.046`$, deviations form exact $`U(1)_b`$ invariance will be small). Then for any standard model interaction a reaction to the type
$$n_ib\text{jet}+Xn_fb\text{jet}+Y$$
(26)
will obey
$$(1)^{n_i}=(1)^{n_f}$$
(27)
to very high accuracy. The number $`(1)^{\mathrm{\#}\mathrm{of}b\mathrm{jets}}`$ defines the b-parity of a state (it being understood that the top quarks have decayed).
The standard model is then b-parity even, and the idea is to consider a lepton collider <sup>โ โ </sup><sup>โ โ </sup>โ โ In hadron colliders there are sea-$`b`$ quarks which foul-up the argument and simply count the number of $`b`$ jets in the final state; new physics effects will show up as events with odd number of $`b`$ jets.
The standard model produces no measurable irreducible background, yet there are significant reducible backgrounds which reduced the sensitivity to $`\mathrm{\Lambda }`$. To estimate these effects I define
* $`ฯต_b=b\text{jet}`$ tagging efficiency
* $`t_c=c\text{jet}`$ mistagging efficiency (probability of mistaking a $`c\text{jet}`$ jet for a $`b\text{jet}`$
* $`t_j=`$light-jet mistagging efficiency (probability of mistaking a light-jet for a $`b\text{jet}`$
so that the measured cross section with $`k`$-b-jets is
$$\overline{\sigma }_k=\underset{u+v+w=k}{}\left[\left(\genfrac{}{}{0pt}{}{n}{u}\right)ฯต_b^u(1ฯต_b)^{nu}\right]\left[\left(\genfrac{}{}{0pt}{}{m}{v}\right)t_c^v(1t_c)^{mv}\right]\left[\left(\genfrac{}{}{0pt}{}{\mathrm{}}{w}\right)t_j^w(1tj)^\mathrm{}w\right]\sigma _{nm\mathrm{}}$$
(28)
where $`\sigma _{nm\mathrm{}}`$ denotes the cross section for the final state with $`n`$ b-jets, $`m`$ c-jets, and $`\mathrm{}`$ light jets. Note that $`\left[\left(\genfrac{}{}{0pt}{}{n}{u}\right)ฯต_b^u(1ฯต_b)^{nu}\right]`$ is the probability of tagging $`u`$ and missing $`nu`$ b-jets out of the $`n`$ available.
As an example consider
$$_{\mathrm{eff}}=_{\mathrm{sm}}+\frac{f_{ij}}{\mathrm{\Lambda }^2}\left(\overline{\mathrm{}}\gamma ^\mu \mathrm{}\right)\left(\overline{q}_i\gamma _\mu q_j\right)$$
(29)
where $`ij`$ denote family indices. Taking $`m_H=100\text{GeV}|f|=1t_c=t_j=0`$ the sensitivity to $`\mathrm{\Lambda }`$ is summarized by the following table
| Limits from $`e^+e^{}t\overline{c}+\overline{t}c+b\overline{s}+\overline{b}s1b\mathrm{jet}+X`$ | | | | |
| --- | --- | --- | --- | --- |
| $`\sqrt{s}`$ | $`L`$ | $`ฯต_b=50\%`$ | $`ฯต_b=60\%`$ | $`ฯต_b=70\%`$ |
| 200 GeV | 2.5 $`fb^1`$ | 1.4 TeV | 1.5 TeV | 1.6 TeV |
| 500 GeV | 75 $`fb^1`$ | 5.0 TeV | 5.2 TeV | 5.5 TeV |
| 1000 GeV | 200 $`fb^1`$ | 9.5 TeV | 10.0 TeV | 10.7 TeV |
These results are promising yet they will be degraded in a realistic calculation. First one must include the effects of having $`t_{c,j}0`$. In addition there are complications in using inclusive reactions such as $`e^+e^{}b+X`$ since the contributions form events with large number of jets can be very hard to evaluate (aside from the calculational difficulties there are additional complications when defining what a jet is). A more realistic approach is to restrict the calculation to a sample with a fixed number of jets ($`2`$ and $`4`$ are the simplest) and determine the sensitivity to $`\mathrm{\Lambda }`$ for various choices of $`ฯต_b`$ and $`t_j`$ using this population only.
### C CP violation
Just as for b-parity the CP violating effects are small within the standard model and so precise measurements of CP violating observable might be very sensitive to new physics effects.
In order to study CP violations it is useful to first define what the CP transformation is. In order to do this in general denote the Cartan group generators by $`H_i`$ and the root generators by $`E_\alpha `$, then it is possible to find a basis where all the group generators are real and, in addition, the $`H_i`$ are diagonal . Define then CP transformation by Transformations
$`\psi `$ $``$ $`C\psi ^{}\text{(fermions)}`$ (30)
$`\varphi `$ $``$ $`\varphi ^{}\text{(scalars)}`$ (31)
$`A_\mu ^{\left(i\right)}`$ $``$ $`A_\mu ^{\left(i\right)},(i:\text{ Cartan generator})`$ (32)
$`A_\mu ^{\left(\alpha \right)}`$ $``$ $`A_\mu ^{\left(\alpha \right)},(\alpha :\text{ root})`$ (33)
it is easy to see that the field strengths and currents transform as $`A_\mu `$, while $`D\varphi (D\varphi )^{}`$. It then follows that in this basis the whole gauge sector of any gauge theory is CP conserving; CP violation can arise only in the scalar potential and fermion-scalar interactions using this basis.
In order to apply this to electroweak physics I will need the list of TLG operators of dimension 6 which violate CP, they are given by <sup>โกโก</sup><sup>โกโก</sup>โกโกThe notation is the following: $`\mathrm{}`$ and $`q`$ denote the left-handed lepton and quark doublets; $`u`$, $`d`$ and $`e`$ denote the right handed quark and charged lepton fields. $`lambda`$ denote the Gell Mann matrices, $`\tau `$ the Pauli matrices, and $`ฯต=i\tau ^2`$. $`D`$ represents the covariant derivatives and $`\varphi `$ the scalar doublet.
$`\left(\overline{\mathrm{}}e\right)\left(\overline{d}q\right)\text{h.c.}\left(\overline{q}u\right)\epsilon \left(\overline{q}d\right)\text{h.c.}\left(\overline{q}\lambda ^Au\right)\epsilon \left(\overline{q}\lambda ^Ad\right)\text{h.c.}`$ (34)
$`\left(\overline{\mathrm{}}e\right)\epsilon \left(\overline{q}u\right)\text{h.c.}\left(\overline{\mathrm{}}u\right)\epsilon \left(\overline{q}e\right)\text{h.c.}|\varphi |^2\left(\overline{\mathrm{}}e\varphi \text{h.c.}\right)`$ (35)
$`|\varphi |^2\left(\overline{q}u\stackrel{~}{\varphi }\text{h.c.}\right)|\varphi |^2\left(\overline{q}d\varphi \text{h.c.}\right)|\varphi |^2_\mu \left(\overline{\mathrm{}}\gamma ^\mu \mathrm{}\right)`$ (36)
$`|\varphi |^2_\mu \left(\overline{e}\gamma ^\mu e\right)|\varphi |^2_\mu \left(\overline{q}\gamma ^\mu q\right)|\varphi |^2_\mu \left(\overline{u}\gamma ^\mu u\right)`$ (37)
$`|\varphi |^2_\mu \left(\overline{d}\gamma ^\mu d\right)`$ (38)
$`๐ช_1=\left(\varphi ^{}\tau ^I\varphi \right)D_\mu ^{IJ}\left(\overline{\mathrm{}}\gamma ^\mu \tau ^J\mathrm{}\right)`$ (39)
$`๐ช_2=\left(\varphi ^{}\tau ^I\varphi \right)D_\mu ^{IJ}\left(\overline{q}\gamma ^\mu \tau ^Jq\right)`$ (40)
$`๐ช_3=\left(\varphi ^{}\epsilon D_\mu \varphi \right)\left(\overline{u}\gamma ^\mu d\right)\text{h.c}`$ (41)
All operators except $`๐ช_{1,2,3}`$ violate chirality and their coefficients are strongly bounded by their contributions to the strong CP parameter $`\theta `$; in addition some chialiry violating operators contribute to meson decays (which again provide strong bounds for fermions in the first generation) and, finally, in natural theories some contribute radiatively to fermion masses and will be then suppressed by the smaller of the corresponding Yukawa couplings. For these reasons I will not consider them further. Moreover, since I will be interested in limits that can be obtained using current data, I will ignore operators whose only observable effects involve Higgs particles.
With these restrictions only $`๐ช_{1,2,3}`$ remain; their terms not involving scalars are
$`๐ช_1`$ $``$ $`{\displaystyle \frac{igv^2}{\sqrt{2}}}\left(\overline{\nu }_L\overline{)}W^+e_L\text{h.c.}\right)`$ (42)
$`๐ช_2`$ $``$ $`{\displaystyle \frac{igv^2}{\sqrt{2}}}\left(\overline{u}_L\overline{)}W^+d_L\text{h.c.}\right)`$ (43)
$`๐ช_3`$ $``$ $`{\displaystyle \frac{igv^2}{\sqrt{8}}}\left(\overline{u}_R\overline{)}W^+d_R\text{h.c.}\right)`$ (44)
The contributions from $`๐ช_{1,2}`$ can be absorbed in a renormalization of standard model coefficients whence only $`๐ช_3`$ produces observable effects, corresponding to a right-handed quark current. Existing data (from $`\tau `$ decays and $`m_W`$ measurements) implies $`\mathrm{\Lambda }>\mathrm{\hspace{0.17em}500}\text{GeV}`$
One can also determine the type of new interactions which might be probed using these operators . The heavy physics which can generate $`๐ช_3`$ at tree level is described in Fig. 5. If the underlying theory is natural we conclude that there will be no super-renormalizable couplings; in this case $`๐ช_3`$ will be generated by heavy fermion exchanges only <sup>\**</sup><sup>\**</sup>\**It is true that vertices involving light fermions, light scalars and heavy fermions produce mixings between the light and heavy scales, but this occurs at the one loop level. In contrast cubic terms of order $`\mathrm{\Lambda }`$ in the scalar potential would shift $`v`$ at tree level.
Note finally that these arguments are only valid for weakly coupled heavy physics. For strongly coupled theories other CP violating operators can be important, e.g.
$$\frac{f}{\mathrm{\Lambda }^2}B^{\mu \nu }\left(\overline{e}\gamma _\mu D_\nu e\text{h.c}\right)$$
(45)
since $`|f|1`$.
## VI Other applications
The effective Lagrangian approach can be applied in many other situations such as gravity and high temperature field theory. I will briefly consider the latter.
### A Large temperatures
It is a well-known fact that the thermodynamics of a system with Hamiltonian $`H`$ can be derived form the partition function $`\text{tr}e^{\beta H}`$. This resembles closely the (trace of the) quantum evolution operator $`e^{iHt}`$ hence we can obtain the thermodynamics of a system by the replacement $`it\beta `$: non-zero temperature field theory corresponds to Euclidean field theory on a cylinder of perimeter$`=\beta `$, I will denote the corresponding Euclidean time by $`\tau `$
Since the time direction is finite the fields are expanded in a Fourier series. For bosons one obtains
$$\varphi =\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}\frac{d^3k}{(2\pi )^3}\varphi _n(๐ค)e^{i(2n\pi T\tau +๐ค๐ซ)}$$
(46)
and the corresponding free propagator is given by
$$\frac{1}{(2n\pi T)^2+๐ฉ^2+m^2}$$
(47)
The field is periodic in $`\tau `$ due to the commutativity of the variables in the functional integral (there is a much more physical reason, called the Kubo-Martin-Schwinger condition) .
Note that the $`n0`$ modes become heavy as $`T\mathrm{}`$ so that in this limit only the $`n=0`$ modes remain and the theory reduces to a 3-D Euclidean field theory (there might be some subtleties involved, see below).
Fro fermions the expansion is in odd Fourier modes since the corresponding integration variables anticommute. explicitly
$$\psi =\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}\frac{d^3k}{(2\pi )^3}\psi _n(๐ค)e^{i(2n+1)\pi T\tau +๐ค๐ซ)}$$
(48)
with fee propagator
$$\frac{1}{\left[i(2n+1)\pi T+\mu \right]\gamma ^0๐ค\gamma m}$$
(49)
which shows that all modes become heavy as $`T\mathrm{}`$. There will be then no fermions in the spectrum at very large temperatures. Note that this occurs independently of the fermion mass .
Despite the absence of heavy fermions and scalars (effective mass $`T`$) at large temperatures, we can still ask what is their effect on the scalar modes that survive in this regime. To this end we can construct the corresponding effective theory. I will illustrate the procedure using a simple example.
Consider the following scalar theory
$$^{\left(4\right)}=\frac{1}{2}\left(\varphi ^2\right)\frac{1}{2}m^2\varphi ^2\frac{\lambda }{4!}\varphi ^4$$
(50)
Then the excitations which survive at large $`T`$ are
$$\phi (๐ฑ)=\sqrt{T}_0^\beta ๐\tau \varphi (๐ฑ,\tau )$$
(51)
where $`\phi `$ is the dynamical variable of a 3 dimensional Euclidean field theory (in 3 dimensions the scalar fields have units of $`\sqrt{\text{mass}}`$ which explains the $`\sqrt{T}`$ factor). The only symmetry (aside form Euclidean invariance) is the reflection symmetry $`\phi \phi `$. The scale of the new theory is set by $`\mathrm{\Lambda }=T`$, but in this case the model is supposed to describe physics above $`\mathrm{\Lambda }`$
With these considerations we can write the effective theory for $`\phi `$,
$$_{\mathrm{eff}}=\frac{1}{2}\left(\phi \right)^2+\frac{1}{2}a\phi ^2+\frac{1}{4!}b\phi ^4+\frac{c}{6!}\phi ^6+O(1/T)$$
(52)
note that $`b`$ is a super-renormalizable coupling and may lead to infrared problems.
The coefficients $`a`$, $`b`$, $`c`$, etc. can be calculated from the original theory. At one loop one obtains
$`a={\displaystyle \frac{\lambda T^2}{24}}b={\displaystyle \frac{m}{2\pi }}\left({\displaystyle \frac{\lambda T}{4m}}\right)^2c={\displaystyle \frac{1}{4\pi }}\left({\displaystyle \frac{\lambda T}{4m}}\right)^3`$ (53)
But this calculation has some potential problems. Consider the $`2k`$ point function at zero external momentum; the corresponding graphs are given in Fig. 6 A simple estimate (verified by explicit calculation) shows that
$$\text{Graph}\underset{\mathrm{prefactors}+\mathrm{dim}.\mathrm{analysis}}{\underset{}{\frac{\lambda ^k}{m^{2k4}}}}\times \underset{\mathrm{integral}+\mathrm{sums}}{\underset{}{\left(\frac{T}{m}\right)^{k+1}}}$$
(54)
which corresponds to the operator
$$๐ช^{\left(k\right)}\frac{\lambda ^k}{m^{2k4}}\left(\frac{T}{m}\right)^{k+1}\left(\sqrt{T}\phi \right)^{2k}\frac{1}{T}=m^3\left(\frac{\sqrt{\lambda }T\phi }{m^{3/2}}\right)^{2k}$$
(55)
whose coefficient has positive powers of $`\mathrm{\Lambda }(=T)`$ and are not suppressed at large temperatures. In fact, should this be correct the, effective theory expansion would be useless.
The solution to this infrared problem (diverging effective coefficients as $`m0`$ is well known for this type of theories : the propagator for the $`n=0`$ mode gets dressed and in so doing the $`m^2`$ gets shifted by an amount $`T^2`$. Explicitly, the graphs in Fig. 7 shift
$$m^2m^2+\frac{\lambda }{24}T^2$$
(56)
so that the previous expression for the effective operator coefficient becomes
$$๐ช^{\left(k\right)}\frac{\lambda ^{(3k)/2}}{T^k}\phi ^{2k}$$
(57)
which vanishes as $`T\mathrm{}`$. Note that there is still a remnant of the infrared properties of the theory in that the coefficients still diverge as $`\lambda 0`$.
##### a QCD at high temperatures
The previous arguments can be applied to the case of gauge theories. Just as for the scalar field, the gauge field is periodic in $`\beta `$ and can be expanded in Fourier modes. At high temperatures, all but the $`n=0`$ modes are heavy with masses $`T`$. The remaining light modes are
$$๐_{n=0}^A๐^AA^0{}_{n=0}{}^{A}\phi ^A$$
(58)
leaving a 3-D Euclidean $`SU(3)`$ model with gauge fields $`๐^A`$ and with a scalar octet (the $`\phi ^A`$). The 3-D gauge coupling constant is $`g\sqrt{T}`$ (where $`g`$ denotes the QCD gauge coupling)
The simplest infrared divergences are cured by the dressing the gluon propagator at one loop ; the $`\varphi ^A`$ propagator at large $`T`$ then becomes
$$\frac{1}{p^2}\frac{1}{p^2+cg^2T^2}$$
(59)
for some numerical constant $`c`$. But this effect is not extended to the $`๐^A`$ for the corresponding vacuum polarization obeys $`\mathrm{\Pi }_{ii}(p0)0`$ .
The fact that the a remain massless leads to various interesting problems. For example the higher order corrections to the free energy, provided by graphs in Fig. 8. Suppose that the gauge bosons have a (dynamically generated) mass $`m`$. In this case a graph with $`\mathrm{}`$ loops behaves as
$$g^6T^4(g^2T/m)^\mathrm{}3(\mathrm{}>3)$$
(60)
For the case where internal lines correspond to $`A^0`$ (or $`\phi ^A`$) $`mgT`$ and the graph is well behaved, $`g^{\mathrm{}+3}T^4`$. On the other hand when the internal lines represent $`A^i`$ (or, equivalently, $`๐^A`$) propagator a problem will arise unless $`mg^2T`$ is generated (we already know there is no $`O(g)`$ correction to $`m`$). This so-called magnetic mass has not been obtained perturbatively though it is widely believed to be generated.
Additional problems arise since the gauge coupling constant in the 3-d theory has dimensions of $`\sqrt{\text{mass}}`$ leading to super-renormalizable interactions with the related infrared divergences .
## VII Conclusions
In these lectures I have provided a review of some of the very many aspects and properties of effective theories, as well as some of their application. Despite this drawback I hope it does give a flair for the strength of the approach.
Effective theories will be used in deriving the implications of new data on the properties of the physics which underlies the standard model , but in addition it can be applied to a wide variety of phenomena ranging form QCD to superconductivity. It is this flexibility which makes the formalism so attractive. |
warning/0002/astro-ph0002310.html | ar5iv | text | # THE MASS-TO-LIGHT FUNCTION: ANTIBIAS AND ฮฉ_๐
## 1 INTRODUCTION
One of the oldest - and simplest - techniques for estimating the mass density of the universe is the mass-to-light method. In this method, the average ratio of the observed mass to light of the largest possible systems is used; assuming it is a fair sample, it can then be multiplied by the total luminosity density of the universe to yield the universal mass density. When the method is applied to rich clusters of galaxies โ the largest virilized systems for which a mass has been reliably determined โ the total mass density of the universe adds up to only $`\mathrm{\Omega }0.2`$(where $`\mathrm{\Omega }`$ is the mass density in units of the critical density) (Zwicky 1957, Abell 1965, Ostriker, Peebles & Yahil 1974, Bahcall 1977, Faber & Gallagher 1979, Trimble 1987, Peebles 1993, Bahcall, Lubin & Dorman 1995, Carlberg *et al.* 1996, 1997, and references therein). A fundamental assumption in this determination, however, is that the mass-to-light ratio ($`M/L`$) of clusters is a fair representation of the universal value. If the mass-to-light ratio of clusters is larger or smaller than the universal mean, then the resulting $`\mathrm{\Omega }`$ will be an over- or under- estimate, respectively. It is not clear whether this classic assumption of an unbiased representation by clusters is correct. More generally, if mass follows light (i.e., galaxies) on large scales โ thus $`M/L`$ constant โ the galaxy distribution is considered to be unbiased with respect to mass; if mass is distributed more broadly than light, as is generally believed, then the galaxy distribution is biased (i.e., more clustered) with respect to mass, and the above determination of $`\mathrm{\Omega }`$ is an underestimate. We investigate these questions of cluster representation and bias, and the impact they have on the measurement of $`\mathrm{\Omega }`$.
Observations of galaxies, groups and clusters of galaxies suggest that $`M/L`$ increases as a function of scale up to scales of hundreds of kiloparsecs (Schwarzschild 1954, Rubin & Ford 1970, Roberts & Rots 1973, Ostriker *et al.* 1974, Einasto *et al.* 1974, Davis *et al.* 1980, Trimble 1987, Gramann 1990, Zaritzky *et al.* 1993, Fischer *et al* $`1999`$), but then flattens out and remains approximately constant on larger scales (Bahcall, Lubin & Dorman 1995). In the modern context we normally interpret this fact as indicating that luminous galaxies are more concentrated in peak density regions than the dark matter because baryons are dissipational. The shape and amplitude of the mass-to-light function โ that is, the dependence of $`M/L`$ on scale, $`(M/L)(R)`$ โ can place powerful constraints on the amount and distribution of dark matter in the universe, as well as on the amount of bias and its dependence on scale. The $`M/L`$ function thus provides a direct, model-independent census of the total mass density of the universe.
What is the expected dependence of $`M/L`$ on scale? In this paper we investigate this question using large-scale, high resolution hydrodynamic cosmological simulations that contain dark matter and gas, and compare the results with observations. We find an excellent agreement between models and observations in the shape of the $`M/L`$ function; both data and models show an increase on small scales (hundreds of kpc) and a flattened $`(M/L)(R)`$ distribution on large scales. We use the comparison between data and simulations to determine the mass density of the universe. The amount of bias and its dependence on scale are also revealed. We find that clusters of galaxies are mildly *antibiased,* in the sense that mass is more concentrated than light on average. Previous determinations of $`\mathrm{\Omega }`$ using clusters of galaxies have thus *overestimated* $`\mathrm{\Omega }`$ due to this unaccounted antibias. The present investigation attempts to provide an unbiased determination of $`\mathrm{\Omega }`$ using, for the first time, the entire observed mass-to-light function. The above results do not disagree with previous estimates that the mass density of galaxies is unbiased or positively biased with respect to the total mass density in the high density regions; it is the light density that is shown here to be antibiased.
## 2 OBSERVATIONS
The observed mass-to-light ratio of galaxies, groups and clusters as a function of scale, $`(M/L_\mathrm{B})(R)`$, is taken from Bahcall, Lubin and Dorman (1995, hereafter BLD). In these data, masses are determined using different methods including velocity dispersion, gravitational lensing, and X-ray gas temperature. The luminosity L<sub>B</sub> throughout this paper refers to the *total* blue luminosity, corrected for both Galactic and internal extinction. The data for rich clusters (at $`R=1.5h^1`$Mpc) and for groups ($`R20`$kpc to $`1h^1`$Mpc), shown in the figures below, represent median $`M/L_\mathrm{B}`$ values of large samples, as does the $`M/L_\mathrm{B}`$ ratio for the luminous parts of typical $`L_{}`$ elliptical and spiral galaxies (see BLD for details). More recent observations of rich clusters from the CNOC cluster survey (Carlberg *et al. 1996, 1997*) yield consistent results. Based on the available data, BLD find that the $`M/L_\mathrm{B}`$ ratio of galaxy systems increases linearly with scale up to the scale of very large galactic halos ($`R0.2h^1`$Mpc), but then flattens on larger scales; they suggest that $`M/L_\mathrm{B}`$ does not increase significantly with scale beyond $`0.2h^1`$Mpc. Furthermore, BLD show that $`M/L_\mathrm{B}`$ of elliptical galaxies is approximately three times larger than that of spirals (at the same radius); both increase linearly with scale up to $`R0.2h^1`$Mpc. The total mass of groups and clusters can then be accounted for by the combined mass of their elliptical and spiral galaxy members, including their large halos, plus the intracluster gas. The large halos are likely to be stripped off in the dense environments of clusters, but their mass still remains in the clusters. Observations of weak gravitational lensing by foreground galaxies using the Sloan Digital Sky Survey (Fischer *et al*. 1999) find consistent results indicating large halos around galaxies.
Recently, the first determination of the mass and mass-to-light ratio of a large supercluster ($`6h^1`$Mpc), MS0302, was obtained using weak gravitational lensing (Kaiser *et al.*1999). The mass and $`M/L_\mathrm{B}`$ ratio of three individual clusters in the supercuster as well as the mass and $`M/L_\mathrm{B}`$ of the supercluster itself were all determined from the weak lensing observations. The results show, quite remarkably, the *same* $`M/L_\mathrm{B}`$ ratio for both the individual clusters and the large supercluster ($`260\pm 40`$ and $`280\pm 40hM_{}/L_{}`$, respectively) thus directly confirming a flat $`(M/L_\mathrm{B})(R)`$ function on large scales, as suggested by BLD. This new supercluster result is added in the figures below (converted to our standard L<sub>B</sub> system by adding the 30% contribution to the luminosity from spiral galaxies (Kaiser *et al.* 1999), correcting for passive luminosity evolution from $`z=0.42`$ to $`z0`$ following $`L_\mathrm{B}(1+z)`$, and correcting for internal extinction ($``$10%; BLD); the net correction factor is 1$`\pm `$0.2). We also show (for illustration only) the $`M/L_\mathrm{B}`$ ratio determined from the Least Action Method at $`30h^1`$Mpc by Tully, Shaya and Peebles (1994) and the observed range of Virgo Infall measurements (see BLD). While these provide less direct measures of mass than the supercluster weak lensing result (and are thus not included in our fits), they are all consistent with each other and with the observed flattening of $`M/L_\mathrm{B}`$ with scale. The data are presented in the figures below.
## 3 SIMULATIONS
We investigate the expected behavior of $`M/L_\mathrm{B}`$ as a function of scale using two sets of cosmological simulations which include both dark matter and gas: a large-scale, $`100h^1`$Mpc box simulation to study the large-scale behavior of $`M/L_\mathrm{B}`$, and a smaller, higher resolution simulation with a box size of $`11.1h^1`$Mpc, to investigate smaller scales. The large-scale hydrodynamic simulation, described by Cen and Ostriker (1999), uses the shock-capturing Total Variation Diminishing method on a Cartesian grid for gas dynamics (Ryu *et al,* 1993). A Particle-Mesh (PM) code is used for dark matter particles. An FFT is used to solve Poissonโs equation. In addition, the code accounts for cooling processes including metal cooling and heating and incorporates a heuristic galaxy formation scheme described by Cen and Ostriker (1999) (see also below). The cosmological model used is a flat Cold-Dark-Matter (CDM) model, with mass density $`\mathrm{\Omega }=0.37,`$ cosmologial constant density $`\mathrm{\Omega }_\mathrm{\Lambda }=0.63,`$ baryon density $`\mathrm{\Omega }_\mathrm{b}=0.039,`$ and a Hubble constant $`h=0.7`$ (where $`H_0=100h`$ kms s<sup>-1</sup>Mpc<sup>-1</sup>). A power-spectrum slope of $`n=0.95`$ and normalization $`\sigma _8=0.8`$ (the mass rms fluctuations on $`8h^1`$Mpc scale at $`z=0`$ ) were used, consistent with the cluster abundance normalization and the COBE microwave background fluctuations (White *et al.* 1993, Ostriker and Steinhardt 1995, Bahcall and Fan 1998). This model fits well current observational data (e.g., Ostriker and Steinhardt 1995, Krauss *et al.* 1995, Bahcall *et al.* 1999). A periodic box of $`100h^1`$Mpc on a side is used, with $`512^3`$ fluid cells and $`256^3`$ dark matter particles. The dark matter mass resolution is $`6\times 10^9h^1M_{}`$ and the grid cell size is $`0.2h^1`$Mpc. We consider only scales with radii $`R1h^1`$Mpc in this simulation, which is considerably larger than the cell size; on these scales, the relevant gravitational and hydrodynamical physics are accurately computed. On smaller scales we use the smaller, higher-resolution simulation described below.
Galaxies are โidentifiedโ in the simulation by the procedure described in Cen and Ostriker (1999): if a cellโs mass is higher than the Jeanโs mass, and if the cooling time of the gas in it is shorter than its dynamical time, and if the flow around the cell is converging, then it will have stars forming inside that cell. The code turns the baryonic fluid component into collisionless stellar particles (โgalaxy particlesโ) at a rate proportional to $`m_b/t_{\mathrm{dyn}}`$, where $`m_\mathrm{b}`$ is the mass of gas in the cell and $`t_{\mathrm{dyn}}`$ is the local dynamical time. These galaxy particles subsequently contribute to metal production, SN energy feedback and the background ionizing UV radiation. This algorithm is essentially the same as in Cen and Ostriker (1992) and also used by Katz, Hernquist and Weinberg (1996), Gnedin (1996) and Steinmetz (1996). The masses of the galaxy particles range from $`10^6`$ to $`10^9M_{}`$; thus many galaxy particles are contained in a single luminous galaxy in the real universe. Rather than group the particles into galaxies, we simply use the galaxy particles themselves, which makes the results less dependent on resolution.
Luminosities (in the relevant bands) are assigned to each cell following the Bruzual and Charlot (1993, 1998; hereafter BC) model; we use their instantaneous star-formation model, which best fits observations (Nagamine, Cen & Ostriker 1999). We also analyze our results using other BC models; the main conclusions are insensitive to the specific star-formation model used. The luminosities determined for each cell are summed over the galaxy particles in the cell and evolve with time as given by the BC model. The simulated luminosities are in excellent agreement with the observed luminosity density in the universe at different redshifts (Nagamine, Cen & Ostriker 1999; see also below).
With the above information we can now determine the mass-to-light ratio, $`M/L_\mathrm{B}`$ (where $`L_\mathrm{B}`$ is the light in the blue band) at different locations in the simulation volume and study it as a function of scale. In order to minimize possible uncertainties due to model luminosities, we normalize all luminosities - and thus $`M/L_\mathrm{B}`$ \- to the *observed* luminosity density of the universe, as discussed below; this ensures that our results are largely independent of the specific luminosity model used.
The behavior of $`M/L_\mathrm{B}`$ on small scales is determined in a similar manner using smaller ($`11.1h^1`$Mpc box), higher-resolution ($`5h^1`$kpc) Tree SPH simulations (see Dave *et al.* 1999). This simulation uses a similar cosmological model ($`\mathrm{\Omega }=0.4`$, $`\mathrm{\Omega }_{}=0.6`$, $`h=0.65`$, $`\sigma _8=0.8`$); the small difference between the models is adjusted in the final normalization of $`M/L_\mathrm{B}`$, but is insignificant. An $`\mathrm{\Omega }=1`$ CDM model, tilted with $`n=0.8`$, is also investigated using this simulation size. Galaxies are identified using SKID (see Katz, Weinberg & Hernquist 1996), and luminosities are assigned to each galaxy using the same BC model described above.
## 4 DEFINITION OF BIAS
The term โbiasโ has been used with different explicit and implicit definitions, so it is essential that we be clear. Kaiser (1984) introduced bias as the difference between the amplitude of the correlation function of high density regions (such as galaxies and clusters) relative to that of the mass in order to explain the exceptionally strong correlation function observed for rich clusters of galaxies (Bahcall & Soneria 1983, Klypin & Kopylov 1983). Similarly, Davis et al. (1985) introduced bias as the proportionality constant between the observed fluctuations in the number density of galaxies and the mass fluctuations found in simulations: $`(\mathrm{\Delta }N/N)_{\mathrm{gal}}b_{\mathrm{gal}}(\mathrm{\Delta }\rho /\rho )_m`$. Since some smoothing scale ($`R`$) must be utilized to calculate either side of the equation, bias must be explicitly a function of scale, $`b_{\mathrm{gal}}(R)`$. Implicit were observational criteria limiting the counted galaxies to be above a certain luminosity and surface brightness. If (and it is a substantial assumption) one identifies the number density of halos in simulations with the number density of galaxies then good dark matter simulations (e.g., Jenkins *et al.* 1998, Kravtsov & Klypin 1999, Colin *et al.* 1999), which can compute $`(\mathrm{\Delta }N/N)_{\mathrm{halo}}=b_{\mathrm{halo}}(R)(\mathrm{\Delta }\rho /\rho )_m`$, provide useful information and indicate low bias at large scales ($`b_{\mathrm{halo}}1`$), significant positive bias (b$`{}_{\mathrm{halo}}{}^{}>1`$) at intermediate scales, and antibias ($`b_{\mathrm{halo}}<1`$) on small ($`R<5h^1`$Mpc) scales due to merging of halos.
Hydrodynamic simulations which seek to identify the site of galaxy formation and estimate the formation rate can compute the mass overdensity in galaxies, $`(\mathrm{\Delta }\rho /\rho )_{\mathrm{gal}}`$, although poor resolution limits their ability to identify individual objects and to compute the galaxy number overdensity $`(\mathrm{\Delta }N/N)_{\mathrm{gal}}`$. Recent papers by Katz *et al* (1999) and Cen & Ostriker (2000) find significant positive bias on intermediate scales. Blanton *et al* (1999) discuss in detail the physical origin of this bias, and its dependence on scale. Einasto et al. (1999) discuss the physical origin of bias in terms of the fraction of mass that exists in the voids. The โsemi-analyticโ approach seeks to combine in simplified form elements of the physical approach utilized in the hydrodynamic modelling with the detailed resolution obtainable from pure N-body work and has produced suggestive and most useful comparisons with observations (Cole *et al.* 1994, Kauffmann *et al.* 1997). All of the above work find positive (but small) bias on large scales so it is important to understand the sense in which we will identify antibias in this work.
The best way of comparing simulations to observations is neither through $`(\mathrm{\Delta }N/N)_{\mathrm{halo}}`$ nor $`(\mathrm{\Delta }\rho /\rho )_{\mathrm{gal}}`$, but via $`(\mathrm{\Delta }j/j)_{\mathrm{gal}}`$, where $`j`$ is the light emitted by the galaxies in some band (here we use $`j_\mathrm{B}`$ in the blue band). We then compare $`(\mathrm{\Delta }j_\mathrm{B}/j_\mathrm{B})_{\mathrm{gal}}`$ with the same observed quantity (after correction for obscuration).
Figure 1 shows in three panels (for top-hat smoothing scales $`1.5,5,10h^1`$Mpc) the average and the dispersion of $`(\mathrm{\Delta }\rho /\rho )_{\mathrm{gal}}`$ and $`(\mathrm{\Delta }j_\mathrm{B}/j_\mathrm{B})_{\mathrm{gal}}`$ versus the total mass overdensity $`(\mathrm{\Delta }\rho /\rho )_m`$. Points above the diagonal line are positively biased and those below the line are antibiased. The $`(\mathrm{\Delta }\rho /\rho )_{\mathrm{gal}}`$ curves are similar to those shown in Cen & Ostriker (1992, 2000) as well as Blanton *et al* (1999) and indicate positive bias on all scales in our dense regions (approaching no bias in the highest density regions on these scales), and negative bias (antibias) in underdense regions (i.e., little or no galaxies in the โvoidsโ).
But we see that the light density $`j_\mathrm{B}`$ is antibiased, both relatively and absolutely, in the highest density regions at 1.5, 5 and $`10h^1`$Mpc scales: $`(\mathrm{\Delta }j_\mathrm{B}/j_\mathrm{B})/(\mathrm{\Delta }\rho /\rho )_m<1`$. The effect is small but real and easily understood. At low redshift the highest density regions typically represent rich clusters and superclusters (for large smoothing scales of $`1`$ to $`10h^1`$Mpc); the stars and galaxies in such regions tend to be old. This well-known observational fact is clearly seen in the simulations (Blanton *et al* 1999; Cen & Ostriker 2000); after clusters form (at $`z12`$) the member galaxies reside within a hot medium ($`T=10^710^8`$K) for which cooling is inefficient and further star formation is inhibited. In such old dense regions massive young blue stars are rare, and the light diminishes sharply with increasing time, especially in the blue. Our Bruzual-Charlot (BC) models, which incorporate standard stellar evolution, thus show relatively low blue light levels in the highest density regions at the present time; the age effect overcomes the slight bias to bring the typical values of $`(\mathrm{\Delta }j_\mathrm{B}/j_\mathrm{B})`$ below $`(\mathrm{\Delta }\rho /\rho )_m`$ and yield a small antibias in the highest density regions. The large amount of intracluster gas in these systems may also contribute to the antibias.
For observers who, in general, have no direct access to the ordinate in Figure 1, $`(\mathrm{\Delta }\rho /\rho )_m`$, it is interesting to consider the ratio of the bias in the high density regions (rich clusters and superclusters) to the bias in the more normal regions where most galaxies live at moderate overdensities. Since these latter have a significant positive bias, the relative antibias of high density regions as compared to low density regions is a factor of $``$2-4 (depending on scale).
## 5 THE MASS-TO-LIGHT FUNCTION
We now turn to the determination of the expected mass-to-light ratio of galaxy systems as a function of scale by investigating $`M/L_\mathrm{B}`$ for different size volumes in the simulations. In the $`100h^1`$Mpc box, we investigate volumes with radii ranging from $`R1h^1`$Mpc to $`62h^1`$Mpc (the volume-equivalent radius of the full box). For each volume of radius $`R`$ we determine the total mass $`M`$ and the total light $`L_\mathrm{B}`$ within the volume, and hence $`M/L_\mathrm{B}`$. The volumes are centered on randomly selected โgalaxiesโ in the box, for proper comparison with observations (i.e., we center on random cells with total galaxy particle mass exceeding 10<sup>11</sup> or 10<sup>12</sup> M; the results are insensitive to the specific threshold). For a given radius, a large number of volumes are selected; these random volumes represent a wide range of mass overdensities. Rich clusters and superclusters of galaxies populate the highest overdensity regions (at their respective scales), while loose groups and other galactic systems correspond to regions of lower overdensities.
The first questions we ask are: How does $`M/L_\mathrm{B}`$ depend on scale and on the local overdensity - does it flatten and become constant on large scales? And, does it vary with overdensity (at a given scale)?
The results are presented in Figure 2, together with the observational data discussed in ยง2. The immediately apparent result is that $`(M/L)(R)`$ increases with scale on small scales and flattens on large scales, as seen in the observations. Each of the $`(M/L_\mathrm{B})(R)`$ curves for $`R0.9h^1`$Mpc represents the simulation results for the mean of all volumes with overdensity above a given threshold (at any given scale, as indicated in Fig. 2). The highest overdensities are selected to correspond to observed rich clusters of galaxies ($`\mathrm{\Delta }\rho /\rho `$190 and $``$ 250 at $`R=1.5h^1`$Mpc, where $`\mathrm{\Delta }\rho /\rho `$ is the total mass overdensity; this corresponds approximately to richness class $``$ 0 and $``$ 1 clusters; Abell 1958); these are shown by the top solid and dashed curves. The lower overdensity regions are presented by the dot-dashed curves; these are typical for loose groups of galaxies at $`R1h^1`$Mpc. To illustrate the trend of $`(M/L_\mathrm{B})(R)`$ with overdensity, we scale the density thresholds with radius (from $`R=1.5h^1`$Mpc) assuming a density profile f $`\rho (r)r^{2.4}`$, as suggested by observations (e.g., Bahcall 1977, 1999, Peebles 1993, Carlberg *et al.* 1997); the results are similar for other reasonable extrapolations.
Voids, which contain little or no light (galaxies) but do contain some mass, exhibit very large $`M/L_\mathrm{B}`$ ratios (e.g., Figure 1); their contribution is of course included in the total $`(M/L_\mathrm{B})_{\mathrm{box}}`$ and $`\mathrm{\Omega }`$ values discussed below since these values refer to the entire amount of mass in the box.
The solid curve marked $`\mathrm{\Omega }`$ = 0.37 represents the mean $`M/L_\mathrm{B}`$ function for $`\mathrm{\Delta }\rho /\rho (R1.5h^1\mathrm{Mpc})190`$ for the $`\mathrm{\Omega }`$ = 0.37 simulation (converted to h = 1 for comparison with the data). The same solid line is then scaled up and down to $`\mathrm{\Omega }`$ = 1 and $`\mathrm{\Omega }`$ = 0.16 respectively (the latter, as shown below, is our best-fit value), using linear scaling with $`\mathrm{\Omega }`$, as expected (see below). The entire set of $`(M/L_\mathrm{B})(R)`$ curves for different overdensities is presented only once, for clarity, for $`\mathrm{\Omega }`$ = 0.16.
The shape of the $`(M/L_\mathrm{B})(R)`$ function is nearly independent of the specific model luminosities used; all models, including models with different but observationally acceptable initial mass function (eg., Salpeter 1955, Miller & Scalo 1979, Scalo 1986, for $`0.1M_{}`$), yield essentially the same function shape. In order to be independent of possible uncertainties also in the normalization of the model luminosities, we normalize $`L_\mathrm{B}`$ of the entire simulationโand thus $`M/L_\mathrm{B}`$ of the full boxโ to the *observed* luminosity density of the universe. The local luminosity density of the universe (in total B band luminosity, corrected for extinction) is observed to be $`j_\mathrm{B}=(2\pm 0.4)10^8hL_{(\mathrm{B})}\mathrm{Mpc}^3`$ (Efstathiou *et al.* 1988, Lin *et al.* 1996, Carlberg *et al.* 1997, Ellis 1997, Small *et al.* 1998 and references therein). Since the mass density of the universe is $`\rho =3\mathrm{\Omega }H_0^2/8\pi G=\mathrm{\Omega }\rho _{\mathrm{crit}}=2.78\times 10^{11}\mathrm{\Omega }h^2M_{}\mathrm{Mpc}^3`$, the universal mass-to-light can be expressed as $`M/L_\mathrm{B}\rho /j_\mathrm{B}=(1400\pm 280)\mathrm{\Omega }hM_{}/L_{(\mathrm{B})}`$, where $`L_\mathrm{B}`$ is the total, extinction corrected blue luminosity at $`z0`$. We normalize our simulation box to have the observed luminosity density of the universe, $`j_\mathrm{B}`$, as listed above; the $`M/L_\mathrm{B}`$ of the full box is thus fixed at $`(M/L_\mathrm{B})_{\mathrm{box}}=518h`$ (for $`\mathrm{\Omega }=0.37`$). Our results are therefore independent of the absolute value of the simulated luminosities. In fact, the direct simulation yields $`M/L_\mathrm{B}=520h`$ for the box, strongly supporting the appropriateness of the luminosity model used. Similarly, for $`\mathrm{\Omega }`$=1, $`M/L_\mathrm{B}`$ is normalized to be $`M/L_\mathrm{B}=1400h`$ ($`\mathrm{\Omega }`$=1), as required.
On scales smaller than $`0.9h^1`$Mpc, the smaller, higher-resolution simulation is used (ยง3) to determine $`(M/L_\mathrm{B})(R)`$ from $`R20`$ kpc to $`6h^1`$Mpc. Since the box is small, no high-density regions such as rich clusters are found (since these are rare objects). The $`(M/L_\mathrm{B})(R)`$ presented in Figure 2 represents the average of typical bright galaxies (corresponding approximately to overdensities above the threshold indicated by the dot-dash curve, as extrapolated to the smaller radii). The results are presented for $`\mathrm{\Omega }=0.16`$ (scaled down from $`\mathrm{\Omega }=0.4`$). The two sets of simulations agree well with each other in the overlap region of $``$1 to $`6h^1`$Mpc, thus strongly supporting these independent results.
The results of Figure 2 show that the simulated $`(M/L_\mathrm{B})(R)`$ function increases on small scales and then flattens on large scales as suggested by observations (Bahcall *et al.* 1995); the data and simulations exhibit the same overall shape of the $`(M/L_\mathrm{B})(R)`$ function. This result is independent of the specific luminosity model used; all models yield the same basic $`(M/L_\mathrm{B})(R)`$ shape. Even though $`M/L_\mathrm{B}`$ flattens to a constant value on large scales, a clear dependence of $`(M/L_\mathrm{B})(R)`$ on the local overdensity (within a given radius R) is apparent; high overdensity regions exhibit higher $`M/L_\mathrm{B}`$ ratios than lower density regions. The results indicate that high density regions (such as rich clusters and superclusters) are *antibiased* with respect to the mean, exhibiting higher $`M/L_\mathrm{B}`$ ratios than average; this implies that mass is more concentrated than light in the high density regions. This effect, as noted in the previous section, is likely caused by the age effect: high density clusters and superclusters are old systems, with low recent star-formation (and thus lower than average blue luminosity); the old galaxies that dominate these system have significantly reduced luminosities at this late time in their evolution. Since all measures of $`\mathrm{\Omega }`$ that utilize the $`M/L_\mathrm{B}`$ method use clusters and superclusters of galaxies โ which are shown here to overestimate the mean $`M/L_\mathrm{B}`$ of the universe โ these measures also overestimate $`\mathrm{\Omega }`$.
We can now determine an unbiased $`\mathrm{\Omega }`$ by properly matching the simulated $`(M/L_\mathrm{B})(R)`$ function to the data. As illustrated in Figure 2, both $`\mathrm{\Omega }=1`$ and $`\mathrm{\Omega }=0.37`$ greatly overestimate the observed $`M/L_\mathrm{B}`$ ratio of groups, clusters, and superclusters, on all scales, by a factor of $`6`$ (for $`\mathrm{\Omega }`$ = 1) and $``$2 (for $`\mathrm{\Omega }`$ = 0.37). This overestimate is seen for the *entire* observed range of the $`M/L_\mathrm{B}`$ function, not just for the classical case of clusters at $`1h^1`$Mpc. By fitting the entire observed and simulated mass-to-light function - properly matching to the relevant overdensities - we can determine an unbiased measure of $`\mathrm{\Omega };`$ we discuss this below (ยง6).
In Figure 3 we compare the observed $`(M/L_\mathrm{B})(R)`$ data with the simulated results for the relevant high- and low- overdensity regions. The high overdensity region (represented by the higher of the two bands at $`R1h^1`$Mpc) corresponds to typical rich clusters and superclusters of galaxies (at $`1.5h^1`$ and $`520h^1`$Mpc respectively; see specific overdensities listed in Figure 3). The low density region reflects environments typical of looser groups and other galaxy systems. The results are presented for both $`\mathrm{\Omega }`$ = 0.16 and $`\mathrm{\Omega }`$ = 1, as scaled from the $`\mathrm{\Omega }`$ = 0.37 simulation. On small scales, $`R`$ 20 kpc to $`6h^1`$Mpc, the results from the high-resolution simulation reflect the full $`M/L_\mathrm{B}`$ range obtained for individual galaxies and groups. These results are in full agreement with the large-scale simulations; the two independent results merge nicely with each other in the overlap region. The $`\mathrm{\Omega }=1`$ results on small scales ($`6h^1`$Mpc) are obtained directly from the $`\mathrm{\Omega }`$ = 1 high resolution simulations; these direct simulation results agree well with the scaled-up results from low $`\mathrm{\Omega }`$ thus supporting the linear scaling of $`M/L_\mathrm{B}`$ with $`\mathrm{\Omega }`$ on large scales.
## 6 DETERMINING $`\mathrm{\Omega }`$
The results presented in Fig. 3 provide a powerful illustration that an $`\mathrm{\Omega }`$ = 1 model significantly overestimates $`M/L_\mathrm{B}`$ on *all scales*. On large scales, the high overdensity band that represents typical rich clusters (at $`R1.5h^1`$Mpc, $`\mathrm{\Delta }\rho /\rho 250`$) overestimates the observed $`M/L_\mathrm{B}`$ value for clusters by a factor of $``$6โ a familiar result. A similar overestimate is seen for smaller groups of galaxies, for individual galaxies, and for superclusters. Even an $`\mathrm{\Omega }`$ = 0.37 model appears to overestimate $`(M/L_\mathrm{B})(R)`$, by a factor of $``$2 (with lower significance).
To determine the best fit value of $`\mathrm{\Omega }`$, we use two methods. In the first method, we use the observed $`M/L_\mathrm{B}`$ ratio of rich clusters of galaxies, and correct it to the proper global universal value (i.e., correct for the cluster antibias) by using the simulationโs ratio of $`M/L_\mathrm{B}`$ for the entire box to that of rich clusters. This ratio, $`b_{\mathrm{cl}}^{^{M/L_\mathrm{B}}}[(M/L_\mathrm{B})_{\mathrm{box}}/<M/L_\mathrm{B}>_{\mathrm{cl}}]_{\mathrm{sim}}`$, is the bias factor (in $`M/L_\mathrm{B}`$) of clusters. For rich clusters (richness class $``$ 1) at $`R1.5h^1`$Mpc, we find
$$b_{\mathrm{cl}}^{^{M/L_\mathrm{B}}}\left[\frac{<M/L_\mathrm{B}>_{\mathrm{box}}}{<M/L_\mathrm{B}>_{\mathrm{cl}}}\right]_{\mathrm{sim}}=0.75\pm 0.15.$$
(1)
The universal $`M/L_\mathrm{B}`$ value is thus given by $`<M/L_\mathrm{B}>_{\mathrm{cl}}\times b_{\mathrm{cl}}`$; rich clusters overestimate the mean value by a factor of $`1/b_{\mathrm{cl}}1.3`$. The error-bar in (1) reflects the rms scatter among the simulated cluster $`M/L_\mathrm{B}`$ values and the scatter among the different luminosity models investigated. Since only the relative ratio between the simulated $`(M/L_\mathrm{B})_{\mathrm{box}}`$ and $`<M/L_\mathrm{B}>_{\mathrm{cl}}`$ is used in this method, the luminosity normalization is unimportant. The mass density of the universe can be determined from the mean observed $`M/L_\mathrm{B}`$ of rich clusters (richness $``$ 1) at $`R1.5h^1`$Mpc, $`<M/L_\mathrm{B}>_{\mathrm{cl}}^{\mathrm{obs}}=300\pm 70hM_{}/L_{}`$ (BLD; Carlberg *et al.* 1997; with $`L_\mathrm{B}`$ in our standard system, at $`z=0`$), and the observed luminosity density of the universe, $`j_\mathrm{B}`$,
$$\rho _m=<M/L_\mathrm{B}>_o\times j_\mathrm{B}=<M/L_\mathrm{B}>_{\mathrm{cl}}^{\mathrm{obs}}\times b_{\mathrm{cl}}^{^{M/L_\mathrm{B}}}\times j_\mathrm{B}$$
(2)
where $`<M/L_\mathrm{B}>_o`$ is the universal value. Therefore
$$\mathrm{\Omega }\frac{\rho _m}{\rho _{\mathrm{crit}}}=\frac{<M/L_\mathrm{B}>_{\mathrm{cl}}^{\mathrm{obs}}\times b_{\mathrm{cl}}^{M/L_\mathrm{B}}}{(M/L_\mathrm{B})_{\mathrm{crit}}},$$
(3)
where $`(M/L_\mathrm{B})_{\mathrm{crit}}\rho _{\mathrm{crit}}/j_\mathrm{B}`$ is the value required for a critical density universe ($`\mathrm{\Omega }`$ = 1; see ยง5). Recent observations of the local galaxy luminosity function, corrected to the standard system of luminosity used here, yield $`j_\mathrm{B}=(2\pm 0.4)\times 10^8hL_{}\mathrm{Mpc}^3`$ and thus $`(M/L_\mathrm{B})_{\mathrm{crit}}=1400\pm 280hM_{}/L_{}`$ (Lin *et al.* 1996, Carlberg *et al.* 1997, Ellis 1997, Small *et al.* 1998). The conservative error-bar used above reflects the scatter among the different measurements as well as their uncertainties. We thus find
$$\mathrm{\Omega }=\frac{(300\pm 70)(0.75\pm 0.15)}{1400\pm 280}=0.16\pm 0.06.$$
(4)
The representative $`M/L_\mathrm{B}`$ value of the universe is $`<M/L_\mathrm{B}>_o`$ =225 $`\pm `$ 70, as given by the numerator of (4).
A second method of determining $`\mathrm{\Omega }`$ is fitting the entire observed $`M/L_\mathrm{B}`$ function of galaxies, groups, clusters, and superclusters (MS0302) to the simulated function, for the relevant overdensities. Here we use the high $`\mathrm{\Delta }\rho /\rho `$ band (Fig. 3) for rich clusters, the lower bound of this band for the MS0302 supercluster, and the low $`\mathrm{\Delta }\rho /\rho `$ band for groups (the upper sub-band is used since it best matches the group overdensities). The small-scale $`R<1h^1`$Mpc band is used for fitting the observed galaxies and small groups of galaxies (at $`0.5h^1`$Mpc). Fitting the observed to simulated $`M/L_\mathrm{B}`$ function has a single free parameter: $`\mathrm{\Omega }`$; the best $`\chi ^2`$ fit yields $`\mathrm{\Omega }`$ = 0.16 $`\pm `$ 0.02. Since the box normalization is fixed at the observed value of $`j_\mathrm{B}=(2\pm 0.4)10^8h`$, corresponding to $`(M/L_\mathrm{B})_{\mathrm{box}}=(1400\pm 280)\mathrm{\Omega }h`$ (ยง5), the result is essentially independent of the luminosity models. The result does depend however on the observed normalization $`j_\mathrm{B}`$; therefore $`\mathrm{\Omega }=0.16\pm 0.02(j_\mathrm{B}/(2\pm 0.4)10^8h`$), or equivalantly, $`\mathrm{\Omega }=0.16\pm 0.02[(1400\pm 280)h/(M/L_\mathrm{B})_{\mathrm{crit}}]`$. Allowing for the normalization uncertainty as well as for uncertainties in the overdensities and in model luminosities, we find
$$\mathrm{\Omega }=0.16\pm 0.05.$$
(5)
This value is consistent with the one obtained earlier using clusters of galaxies alone. Additional systematic uncertainties, while difficult to accurately determine, may contribute an additional $``$ 20%($`\pm `$ 0.03) to the above uncertainty (see below). The $`M/L_\mathrm{B}`$ function for this best-fit value, plotted in Fig. 3, reproduces well the entire observed $`M/L_\mathrm{B}`$ function, from galaxies to superclusters.
The error-bars given in (4,5) above may not include all possible systematic uncertainties. For example, if low surface brightness galaxies contribute significantly to the total luminosity density of the universe (over and above the extrapolated luminosity function), but not to the luminosity in clusters, this will increase $`j_\mathrm{B}`$ (thus decrease $`(M/L_\mathrm{B})_{\mathrm{crit}}`$) from the value used, therefore increasing $`\mathrm{\Omega }`$. However, if such galaxies exist also in groups, clusters, and superclusters โ this effect will cancel out. The effect, if exists, is expected to be small, and is at least partially covered by the large uncertainty adopted for $`j_\mathrm{B}`$ and $`(M/L_\mathrm{B})_{\mathrm{crit}}`$. Similarly, a diffuse intracluster light, which may account for $``$ 15% of the total cluster luminosity (Feldmeier et al. 2000), is not included in the observed cluster luminosity (it may in fact compensate for the contribution of low surface brightness galaxies in the field). If included, this will lower $`<M/L_\mathrm{B}>_{\mathrm{cl}}^{\mathrm{obs}}`$ and thus lower $`\mathrm{\Omega }`$ (by $``$ 15%). Systematic uncertainties in the simulations may also contribute โ but only if they are scale dependent (since the overall normalization is independent of the simulations). It is unlikely that significant shape changes exist on the scales considered here. While difficult to accurately determine such possible systematic uncertainties, we estimate that they may contribute an additional $``$ 20% uncertainty to $`\mathrm{\Omega }`$.
Figure 2 and 3 illustrate that the $`(M/L_\mathrm{B})(R)`$ function of high density regions increases with scale to an above-average peak at a clusters-superclusters scale of few Mpc, then decreases to the mean universal value. Conversely, low density regions reveal lower $`M/L_\mathrm{B}`$ values. This is consistent with observations of rich clusters (high density) versus groups (low density); groups exhibit lower $`M/L_\mathrm{B}`$ ratios than typical rich clusters by a factor of nearly two, as seen in both data and simulations. Based on the present results we also expect that observations of weak lensing in the โfieldโ, which are currently underway, will reveal lower $`M/L_\mathrm{B}`$ ratios than seen in clusters or superclusters of galaxies by a factor of up to $``$2, depending on the specific overdensities.
Our best-fit $`\mathrm{\Omega }`$ (eq. 5) is lower than previous estimates due to the antibias discussed above as well as the more robust use of the entire $`M/L`$ function โ not just clusters โin constraining $`\mathrm{\Omega }`$. A mass-density of $`\mathrm{\Omega }0.35,`$ frequently regarded as a current โmost popularโ value, appears to overestimate the entire observed $`M/L`$ function, *on all scales*, for galaxies, groups, clusters and superclusters.
The above analysis uses overdensities selected in the $`\mathrm{\Omega }=0.37`$ simulation (keeping the same overdensities for the different $`\mathrm{\Omega }`$โs). The actual overdensities (of groups, clusters, superclusters) in the lower $`\mathrm{\Omega }0.16`$ universe are of course twice as large, which can further reduce $`\mathrm{\Omega }`$ by $`20\%`$, to $`\mathrm{\Omega }0.13`$. However, this effect, which is caused mainly by the earlier cluster formation time in lower $`\mathrm{\Omega }`$ models is minimized by the fact that all objects form earlier in such models. If so, the overdensities need not be re-scaled. If they are re-scaled, the best-fit $`\mathrm{\Omega }`$ may be lower than given above (by $`20`$%).
## 7 ELLIPTICAL AND SPIRAL GALAXIES
On small scales, the data show that $`M/L_\mathrm{B}`$ of elliptical galaxies is larger than that of spirals by a factor of $``$3 (BLD; see also Tully and Shaya 1998); this is mostly due to lower blue luminosity in the older ellipticals, but could also be partially due to higher elliptical mass. To test this observation in the simulations, we identify old and young galaxies (thus mostly ellipticals and spirals respectively) by selecting galactic systems based on their redshift of formation. For example, in the large simulation box we define regions of โoldโ galaxies as those where the total galactic particle mass formed at high redshift (e.g., $`z>1.9`$) exceeds that which formed at low redshift (e.g., $`z<0.6`$) by a factor of five. Thus regions dominated by old galaxies satisfy: $`M_{\mathrm{gal}}(z<0.6)/M_{\mathrm{gal}}(z>1.9)<0.2`$. Similarly, regions dominated by โyoungโ galaxies satisfy $`M_{\mathrm{gal}}(z<0.6)/M_{\mathrm{gal}}(z>1.9)>0.2`$. Varying the specific redshift cuts and the fractional threshold (0.2, 0.4, 0.6) does not affect the final results discussed below.
In Figure 4 we present the $`M/L_\mathrm{B}`$ function for the old and young galaxies as discussed above (for $`R0.9h^1`$Mpc; solid and dashed curves). These curves are superimposed on the high and low overdensity bands from Fig. 3. The results show a strong correlation: the old galaxy $`(M/L_\mathrm{B})(R)`$ function traces remarkably well the high overdensity regions (such as clusters and superclusters), while the young galaxies trace well the low overdensity regions. No re-normalization has been applied, and the results are insensitive to reasonable changes in the redshift and threshold definitions of the young and old regions. This result is consistent with observations in the sense that high density regions are indeed best traced by old galaxies. The difference between $`M/L_\mathrm{B}`$ of the old and young galactic regions is approximately a factor of 2 to 3, consistent with observations. Extending the results to smaller scales of individual galaxies, we select old and young galaxies in the high-resolution simulations based on their colors: $`BV>0.65`$ (old) and $`BV<0.65`$ (young). We plot in Fig. 4 the mean 10% highest and lowest $`(M/L_\mathrm{B})(R)`$ for galaxies in these respective color cuts, for $`R20`$ kpc to 6 Mpc. The results depend only slightly on the exact cuts. The results are consistent with those obtained from the large simulation; they merge with each other in the overlap regions. The simulated results are consistent with the data for the entire $`(M/L_\mathrm{B})(R)`$ function if - and only if - $`\mathrm{\Omega }`$ $``$0.16, as shown in Figs. 2-4.
## 8 CONCLUSIONS
We use large-scale cosmological simulations to determine the expected mass-to-light ratio of galaxy systems and its dependence on scale. The $`(M/L_\mathrm{B})(R)`$ function is investigated from small scales of galaxies ($`R20`$ kpc) to large scales ($`R60h^1`$Mpc), and compared with observations of galaxies, groups, clusters, and superclusters. We use the results to evaluate the amount of bias on different scales (i.e., how mass traces light), and use the comparison with observations to determine the mass density of the universe, $`\mathrm{\Omega }`$.
We find the following results:
1. In high density regions the galaxy blue light is antibiased (i.e., lower) relative to the total mass density (while the galaxy mass density is not). This is due to the old age of the high density systems which leads to a relative decrease in their present-day luminosity, especially in the blue band that traces recent star formation.
2. The shape of the simulated $`(M/L_\mathrm{B})(R)`$ function is in excellent agreement with observations. The simulated $`M/L_\mathrm{B}(R)`$ function increases with scale on small scales and flattens on large scales, where $`M/L_\mathrm{B}`$ reaches a constant value, as observed. The mean flattening of $`(M/L_\mathrm{B})(R)`$ on large scales indicates that, on average, mass follows light on large scales (i.e., $`ML`$).
3. Even though $`M/L_\mathrm{B}`$ is approximately constant on large scales, we find that the actual value of $`M/L_\mathrm{B}`$ depends on the local mass overdensity, $`\mathrm{\Delta }\rho /\rho (<R)`$, at a given scale. High overdensity regions exhibit higher $`M/L_\mathrm{B}`$ ratios than lower density regions. The difference can typically be a factor of 2 to 3, consistent with observations of groups and clusters of galaxies (representing low and high density regions, respectively). The dependence of $`M/L_\mathrm{B}(R)`$ on overdensity indicates that high density regions such as rich clusters and superclusters are relatively *antibiased* \- they exhibit higher than average $`M/L_\mathrm{B}`$ values, implying that mass is more concentrated than light in these regions (see 1 above). In the blue luminosity band, the cluster $`M/L_\mathrm{B}`$ antibias is $`b_{\mathrm{cl}}^{^{M/L_\mathrm{B}}}=<M/L_\mathrm{B}>_o/<M/L_\mathrm{B}>_{\mathrm{cl}}=0.75\pm 0.15`$.
4. We find that the $`(M/L_\mathrm{B})(R)`$ function of high density regions is traced well by $`(M/L_\mathrm{B})(R)`$ of old (elliptical) galaxies; low density regions are traced well by young (spiral) galaxies. These results are consistent with observations.
5. We determine the mass density of the universe by fitting the simulated $`(M/L_\mathrm{B})(R)`$ function to observations. The best fit $`\mathrm{\Omega }`$ is lower than previous estimates based on cluster $`M/L`$ values because of the antibias discussed above as well as the more robust use of the entire $`M/L`$ function โ not just clusters โ in constraining $`\mathrm{\Omega }`$. We find a best-fit value of $`\mathrm{\Omega }=0.16\pm 0.05`$ (with an additional estimated uncertainty of $`\pm 0.03`$ for possible additional systematics); this value provides a remarkably good match to the data for galaxies, groups, clusters, and superclusters. The results are independent of the details of the models and provide a powerful measure of $`\mathrm{\Omega }`$. The only significant uncertainty we are aware of is due to the possibility that current observations may systematically underestimate the global mean luminosity density of the universe. This will produce a corresponding underestimate in our computation of $`\mathrm{\Omega }`$ unless there was also a corresponding underestimate in the luminosity of groups, clusters, and superclusters of galaxies.
###### Acknowledgements.
We thank J. Peebles, D. Spergel, P. Steinhardt, and M. Strauss for helpful discussions. This work was supported by NSF grants AST-9803137 and ASC-9740300. |
warning/0002/hep-ex0002052.html | ar5iv | text | # 1 Introduction
## 1 Introduction
At sea level, together with neutrinos, muons are the most abundant particles originated by the interactions of primary cosmic rays at the top of the atmosphere. Due to their relative stability and small cross sections, these particle are able to arrive deep underground and/or deep underwater. As a consequence, their study covers many aspects of cosmic ray physics.
Recently, atmospheric muon flux measurements received attention the context of the atmospheric neutrino anomaly . Because of the close relation between muon and neutrino production, it follows that the evaluation of the atmospheric muon flux can provide an important cross check on the atmospheric neutrino flux. Moreover, measurements of muon flux at all geomagnetic latitudes are crucial for the normalizazion of the calculated neutrino flux.
Measurements performed at different altitudes (sea level, at mountain level or in balloon born experiments) offer various advantages. First of all, a different set of measurements at different altitude values can provide informations about the longitudinal development of the muon component in cosmic ray showers. Moreover, the interpretation of data collected at the top of the atmosphere are not affected by the uncertainties inherent in particle production and propagation, since the muons (and the corresponding neutrinos) generated in the first stages of the cascade. Finally, the knowledge of high altitude muon data is crucial for sea level sub-GeV neutrinos: the corresponding sub-GeV muons originated in the same decay processes cannot reach the sea level, considering that the average muon energy loss in the atmosphere is of the order of 2 GeV; we are thus forced to take data at high altitudes.
On the other hand, measurements performed at ground level offer the advantage of a high stability, large collecting factor and a long exposure time due to the relatively favourable experimental conditions. They however suffer of an intrinsic difficulty in interpretation, since the muons that arrive at sea level are the last stage of a multi-step cascade process. This is true, in particular, for the measurements at high zenith angles, near the horizon, where the intermediate and high energy regions of the spectrum can be analysed ($`p_\mu `$= 10-100 GeV/c). Nevertheless, for this reason, sea level data offer the possibility to perform a robust check of the reliability of existing Monte Carlo codes.
Finally, underground measurements offer the possibility to extend the energy range of the muon spectrum beyond 1 TeV. Such measurements are of an indirect type, but their link with the direct low-energy observations gives the possibility to complete the picture of muon spectra measurements and to cross-check the validity of the global set of data.
Most of the experiments devoted to the measurement of the muon momentum spectra and intensity have been carried in the โ70s. The problem is that the results are often in disagreements with one another; the discrepancies are significantly larger than the experimental errors reported. Recently new instruments, mainly designed for balloon experiments, have been developed; they are able to give detailed information on the muon flux at different altitudes in the atmosphere . Also new measurements deep underground or by EAS arrays have added new information at very high energies.
In this paper we summarize the observations of the muon flux at sea level and deep underground and discuss some of the systematics connected with such measurements. For more complete discussion one can refer to the recent papers and to books .
## 2 Atmospheric muon production and propagation
Secondary muons are mainly produced in the decays of secondary mesons, mostly $`\pi ^\pm `$ and $`K^\pm `$. The most important decay channels, and their respective decay probabilities, are:
$$a)\pi ^\pm \mu ^\pm \nu _\mu 100\%$$
$$b)K^\pm \mu ^\pm \nu _\mu 63.5\%$$
in which the produced muons take on the average 79% and 52% of the energy of the $`\pi ^\pm `$ and $`K^\pm `$, respectively.
The contribution of $`K`$ decays to muon production is a function of the energy and ranges from $``$ 5% at low energies to an asymptotic value of $``$ 27% for E $`\stackrel{>}{_{}}1`$ TeV. At very high energies a small contribution arises from charmed particles. The analytical form of the muon production spectrum at a given height in the atmosphere can be derived by folding the two-body decay kinematics of the parent mesons with their production spectrum. The latter is generally expressed in terms of the so called โspectrum weightedโ moments
$$Z_{p\pi ^\pm }=_0^1x^\gamma \frac{dN_{p\pi ^\pm }}{dx}๐x$$
(1)
where $`dN_{p\pi ^\pm }/dx`$ is the pion production spectrum ($`x=E_\pi /E_p`$ and $`\gamma `$ is the differential primary spectral index). A similar expression can be obtained for kaons. In general, the development of the meson and muon components in the atmosphere depends on the energy range we are considering. The competition between interaction and decay of the particles plays a crucial role and the relative importance of the two processes depends on the energy. We can distinguish three different energy regions in the muon spectrum:
a) $`E_\mu ฯต_{\pi ,K}`$ , where $`ฯต_\pi `$ = 115 GeV and $`ฯต_K`$ = 850 GeV are the critical energy beyond which meson reinteractions cannot be neglected. This is the typical muon energy range studied by underground detectors or by ground based experiments looking at high inclined directions. In this case, the meson production spectrum have the same power law dependence of the primary cosmic rays, but the rate of their decay has an extra $`E^1`$ dependence with respect to the primary and meson spectrum (a consequence of the Lorentz time dilatation). The muon (and hence neutrino) flux takes the form: $`dN/dE_\mu =E_\mu ^{(\gamma +1)}`$, and a zenith dependence $`dN/dcos\theta (cos\theta )^1`$. It should be noted that, in this energy region, the enhancement of the $`K^\pm `$ contribution to the secondary lepton production is particularly important in the neutrino flux calculation, as a consequence of the two-body decay kinematics . This last remark does not hold for muons, for which the limited knowledge of meson production in this energy range is not so crucial as for neutrinos.
b) $`ฯต_\mu \stackrel{<}{_{}}E_\mu \stackrel{<}{_{}}ฯต_{\pi ,K}`$ , where $`ฯต_\mu `$ 1 GeV. In this energy range, almost all the mesons decay, and the muon flux has a power law dependence with the same spectral index of the parent mesons (and hence of the primary cosmic ray, in the assumption of complete Feynman scaling validity) and is almost independent on the zenith angle. A compact form which expresses the low and high energy regions is :
$$\frac{dN_\mu }{dE_\mu }(E_\mu ,\theta )0.14E_\mu ^{2.7}\left[\frac{1}{1+\frac{1.1E_\mu cos\theta }{ฯต_\pi }}+\frac{0.054}{1+\frac{1.1E_\mu cos\theta }{ฯต_K}}\right]$$
(2)
c) $`E_\mu \stackrel{<}{_{}}ฯต_\mu `$ . In this case, muon decays and the energy losses in the atmosphere cannot be neglected. Moreover, geomagnetic latitude and solar modulation now play an important role being the primary cosmic ray energy $`E_p<20`$ GeV.
We stress again the relevance of muon flux measurements for the knowledge of the neutrino flux. In principle, sea level neutrino flux computation can be derived directly from muon flux measurements high in the atmosphere ($`X<`$ 37 $`g/cm^2`$) . This approach gives good results, but only a complete Monte Carlo simulation can take into account second order effects. The main ingredients in Monte Carlo calculations (and the main sources of systematics) of atmospheric lepton production are the input primary cosmic ray spectrum and a detailed description of secondary multiparticle production in the atmosphere. The primary cosmic ray composition plays an important role only in the very high energy range; the composition is dominated by protons and $`\alpha `$ particles at energies below 100 GeV.
Among various Monte Carlo codes now available, we recall and which are the ones used to interpret the atmospheric neutrino anomaly, and the new code based on the FLUKA interaction model which takes into account 3-dimentional effects of secondary propagation in the atmosphere.
The comparison between the Monte Carlo evaluation of the muon flux at different altitudes and at sea level with the existing muon flux measurements constitutes one of the most powerful benchmark to assess the validity of the simulations.
## 3 Atmospheric muon flux measurements
Measurements of the absolute intensity, energy spectrum and positive-to-negative ratio of muons have been carried out many times in the past. Most of these observations were made at sea level and few at different mountain altitudes with counter telescopes separated by absorbers (Pb, Fe) and magnetic spectrometers. More recently, with the development of superconducting magnet, it has been possible to operate spectrometers also on board of balloons which led to accurate measurements at different levels in the atmosphere.
Here we will consider mainly ground-based measurements and those made with detectors on balloons near the ground level or very close to it. The relevant quantities that can be directly measured and will be discussed here, are:
\- absolute muon intensity
\- muon momentum spectrum
\- charge ratio
### 3.1 Absolute intensity measurements
The vertical muon intensity at sea level is a quantity which varies with the geomagnetic latitude, altitude, solar activity and atmospheric conditions.
The geomagnetic field tends to prevent low energy cosmic rays from penetrating through the magnetosphere down to the Earthโs atmosphere. At any point on the Earth one can define a threshold or cut-off rigidity, Pc, for cosmic rays arriving at a particular zenith and azimuth angle .
Primary nuclei having lower rigidity are excluded by the action of the geomagnetic field and do not contribute to production of secondaries in the atmosphere. The cut-off values range from less than 1 GV near the geomagnetic poles to about 16 GV for vertical particles near the equator . It results that geomagnetic effects are important for sea level muons up to about $``$ 5 GeV (Fig. 1). The effect is larger at higher altitudes; Conversi found that the vertical flux of muons with momentum around 0.33 GeV/c at latitude 60 deg was 1.8 times higher with respect to the flux at the equator.
Moreover, as cosmic ray primaries are predominantly positively charged particles, the flux and spectra in the East and West directions differ up to energies of about 100 GeV; the intensity from the West is stronger than that from the East. This effect increases with altitude.
In addition, the primary cosmic ray spectrum at the top of the atmosphere changes with the 11 year solar cycle as the configuration of the Interplanetary Magnetic Field (IMF) varies. It results that the cosmic ray flux is significantly โmodulatedโ up to energies of about 20 GeV (Fig. 2).
In order to estimate how these changes in the primary spectrum influence the counting rate of a muon detector it is necessary to know the โdifferential response curveโ . Its shape varies significantly with the depth of observations, see Fig. 3. Their detailed calculations depend on the properties of nuclear cascades in the atmosphere; more precise descriptions can be found in At the standard momentum of 1 GeV/c and at high latitudes the modulation is 7% and 4.5% for the differential and the integral fluxes, respectively .
So in making a comparison of muon observations at low energies (less than 20 GeV) it is very important to know the year and the location the measurements were made. Figure 4 shows the neutron monitor counting rate recorded by a middle latitude station since 1953. No continuous recording of the same kind exists for muon monitors. By the comparison of the peak to peak variations during the interval 1965-1972 one can estimate that the total muon flux changes are usually a facto 3 to 5 smaller than the observed neutron flux variations
Finally, changes in pressure and, particularly, temperature above the instrument up to the point of muon production by pions and kaons, produce variations of different amplitude in different energy range. The most conspicuous for muons at higher energies is the seasonal variation for which the results reported in the following have not been corrected for.
Classical definition of the hard component is related to penetration characteristics, it is to say the capability of crossing 167 $`g/cm^2`$, equivalent to roughly 15 cm of Pb. As a matter of fact this component is made of of muon with momenta $`p_\mu >0.32`$ GeV/c and less than 1% are protons, neutrons, electrons and pions.
Let us distinguish between vertical and horizontal integral muon fluxes. These latter are made in order to extend the range of the former beyond several tens of GeV/c, but usually they do not give absolute values of the vertical muon intensity. For this reason we will not discuss them here.
The first measurement of the integral vertical intensity was made by Greisen at latitude $`50^o`$ and altitude 259 m a.s.l. (corresponding to 1007 $`g/cm^2`$) who found the value: $`0.83\times 10^2\pm 1\%`$ $`cm^2s^1sr^1`$. Rossi noticed that this value needed to be corrected in order to account for showering and scattering of particles inside the apparatus. Successive measurements led to higher values. The observations made at different latitudes and during different years are presented in Fig. 5. Most measurements were made at high latitudes , and only the few at low latitudes ; were corrected for the geomagnetic effect. No corrections have been made for the solar modulation effects; the measurements are essentially grouped in the period 1967-1977, with one in 1998 . The agreement between the measurements is fairly good (all the data within 10%) and one has to take into account that the largest contribution to the deviations are the systematic errors due to incorrect knowledge of the acceptance, efficiency of the counters and correction for the multiple scattering.
For energies $`E_\mu >1`$ TeV direct measurments of the muon flux were made at highly inclined directions using large magnetic spectrometers , large emulsion chambers and EAS arrays . For the discussion of the results we refer to the recent paper .
### 3.2 Momentum spectra
These spectra have been measured many times for moments up to $``$ 100 TeV/c. Magnetic spectrometers are mainly used at low and intermediate energies, while observations at high energies are made close to the horizontal directions at ground level or deep underground. The latter are indirect measurements, since the ground level spectra have to be extracted from underground data. We will consider here only ground level and underground observations.
#### 3.2.1 Ground level measurements
Direct measurements of momentum spectra for $`p_\mu <1`$ TeV/c are important for the comparison of nuclear cascade models with available data. Furthermore by extending the model results to higher energies one can hope to be able to evaluate prompt muon production and/or charm production. In the momentum region 10 GeV/c - 1 TeV/c: a) the production spectrum of the charged pions cannot be represented by a power law but has a maximum at an energy that depends on both the altitude and the latitude b) the energy loss and the decay of muons must be properly considered. In order to join low to high momentum spectra it is important to have single experiments that cover the widest energy range.
To better see the differences between the sea-level spectra we plot in Fig. 6 the percentage deviations of the data from the best fit spectrum obtained by . Notice that even if individual errors are small (however increasing with momentum due to decreasing number of detectable particles and to the maximum detectable momentum), deviations up to $`\pm `$ 20% are observed probably because of systematic effects.
#### 3.2.2 Underground measurements
From underground muon intensity measurements, informations about sea-level muon spectra can be obtained using different procedures (see for example ). Here and in the following, we assume a standard procedure applied form large area underground experiment . The vertical muon intensity, for a given direction $`\theta ,\varphi `$ and a corresponding rock slant depth $`h`$ can be expressed as:
$$I_\mu ^V(h,\theta ,\varphi )=\left(\frac{1}{\mathrm{\Delta }T}\right)\frac{_iN_im_i}{_j\mathrm{\Delta }\mathrm{\Omega }_jA_jฯต_j/cos\theta _j}$$
(3)
where $`\mathrm{\Delta }T`$ is the total livetime of the experiment, $`N_i`$ is the number of detected events with multiplicity $`m_i`$ in the angular bin $`\mathrm{\Delta }\mathrm{\Omega }_j`$, $`A_j`$ and $`ฯต_j`$ are, respectively, the geometrical and intrinsic acceptance of the detector. The relation between the measured $`I_\mu ^V(h)`$ and the sea-level muon spectrum can be expressed as:
$$I_\mu ^V(h)=_0^{\mathrm{}}\frac{dN_\mu }{dE_\mu d\mathrm{\Omega }}P(E_\mu ,h)๐E_\mu ,$$
(4)
where $`P(E,h)`$ is the muon survival probability function determined via Monte Carlo. Assuming for the sea-level muon spectrum an expression of the form (2), leaving as free parameters the muon spectral index and a normalization constant, it is possible it is possible to unfold sea level muon spectrum from the measured absolute muon intensity.
In Fig. 7 are reported the results of the fit of MACRO data together with LVD , MSU and Baksan data. Data are presented multiplied by factor $`p^3`$ to better observe the variation of the spectrum in the whole energy region and to strengthen a possible flattening in the tail of the spectrum due to charm production. The statistics is still too poor to allow any definite assessment on the existence of this effect at energies $`>`$ 10 TeV/c.
In indirect measurements, accurate estimates of the systematic errors are needed. The main sources of systematics in (4) are the knowledge of the rock density overburden and the treatment of hard processes in the energy loss of muons in the rock. In the MACRO fit, for example, their overall contribution has been estimated to be $``$ 5% and 3% in the determination of the normalization constant and muon spectral index respectively.
### 3.3 Charge ratio
In the primary cosmic rays there is an excess of positively charged particles (protons) with respect to the total number of nucleons. This excess is transmitted via nuclear interactions to pions and further to muons. By assuming that the primary composition is constant in the energy range considered, this ratio will remain constant with the exception of high energies, where the contribution from kaons starts to become sizeable. The muon charge ratio is expected to increase also with zenith angle as the depth is increasing and likewise the energy of the primaries that produce muons of a given momentum at ground. This quantity is important to study nucleon-nucleon interactions, composition and kaon contribution. Magnetic spectrographs are used for determining this ratio. Because of systematic effects in the momentum measurement the values are usually much spread out. Moreover limited statistics at high energy makes it difficult to appreciate the energy dependence. We report in Fig. 8 only the recent data from Mass at lower momenta, and the two compilations made by . It is clear more measurements with longer exposures are still needed.
## 4 Conclusions
Since atmospheric muons and neutrinos are generated in the same processes, the accuracy of the neutrino flux calculation can be improved by forcing the poorly known input parameters of the cascade model to fit the data on the muon flux.
However, the data are still not sufficient for this purpose, since several sea level measurements of the vertical muon flux are in poor agreement with one another, even though each experiment has typically very good statistics.
Disagreement between the results of different experiments are present even if the quoted errors are relatively small in the majority of the experiments. It indicates the existence of significant systematic errors in some experiments by as much as 30-35% at momenta from 10 to 1000 GeV/c.
## 5 References |
warning/0002/hep-ph0002167.html | ar5iv | text | # CP VIOLATION AND QUARK MIXING
## 1 Introduction
Within the standard model (SM), CP violation is due to the presence of a nonzero complex phase in the Cabibbo-Kobayashi-Maskawa (CKM) quark mixing matrix $`V`$. A particularly useful parametrization of the CKM matrix, due to Wolfenstein, follows from the observation that the elements of this matrix exhibit a hierarchy in terms of $`\lambda `$, the Cabibbo angle. In this parametrization the CKM matrix can be written approximately as
$$V\left(\begin{array}{ccc}1\frac{1}{2}\lambda ^2& \lambda & A\lambda ^3\left(\rho i\eta \right)\\ \lambda (1+iA^2\lambda ^4\eta )& 1\frac{1}{2}\lambda ^2& A\lambda ^2\\ A\lambda ^3\left(1\rho i\eta \right)& A\lambda ^2& 1\end{array}\right).$$
(1)
The allowed region in $`\rho `$$`\eta `$ space can be elegantly displayed using the so-called unitarity triangle (UT). The unitarity of the CKM matrix leads to the following relation:
$$V_{ud}V_{ub}^{}+V_{cd}V_{cb}^{}+V_{td}V_{tb}^{}=0.$$
(2)
Using the form of the CKM matrix in Eq. (1), this can be recast as
$$\frac{V_{ub}^{}}{\lambda V_{cb}}+\frac{V_{td}}{\lambda V_{cb}}=1,$$
(3)
which is a triangle relation in the complex plane (i.e. $`\rho `$$`\eta `$ space). With the experimental precision expected in future $`B`$ (and $`K`$) decays, it may become necessary to go beyond leading order in $`\lambda `$ in the Wolfenstein parametrization given above. To this end, we follow here the prescription of Buras et al.: defining $`\overline{\rho }\rho (1\lambda ^2/2)`$ and $`\overline{\eta }\eta (1\lambda ^2/2)`$, we have
$$V_{us}=\lambda ,V_{cb}=A\lambda ^2,V_{ub}=A\lambda ^3(\rho i\eta ),V_{td}=A\lambda ^3(1\overline{\rho }i\overline{\eta })$$
(4)
The key point here is that the matrix elements $`V_{us},V_{cb}`$ and $`V_{ub}`$ remain unchanged, but $`V_{td}`$ is renormalized in going from leading order (LO) to next-to-leading order (NLO). The apex of the UT is now defined by $`(\overline{\rho },\overline{\eta })`$.
Constraints on $`\overline{\rho }`$ and $`\overline{\eta }`$ come from a variety of sources. For example, $`|V_{cb}|`$ and $`|V_{ub}|`$ can be extracted from semileptonic $`B`$ decays, and $`|V_{td}|`$ is at present probed in $`B_d^0`$$`\overline{B_d^0}`$ mixing. The interior CP-violating angles $`\alpha `$, $`\beta `$ and $`\gamma `$ can be measured through CP asymmetries in $`B`$ decays. Additional constraints come from CP violation in the kaon system ($`|ฯต|`$), as well as $`B_s^0`$$`\overline{B_s^0}`$ mixing.
A profile of the unitarity triangle was presented by us in early 1999. This analysis was done at NLO precision, taking into account the state-of-the-art calculations of the hadronic matrix elements from lattice QCD and data available at that time. Subsequently, an improved lower limit $`\mathrm{\Delta }M_s>14.3(\mathrm{ps})^1`$ was reported at the Lepton-Photon symposium in the summer of 1999. In this report, we update the results of our 1999 CKM-unitarity fits by incorporating this new limit on $`\mathrm{\Delta }M_s`$. As we shall see here, this measurement tightens the constraints on the CKM parameters. The other new ingredient in our fits is that we now use the improved Wolfenstein parametrization given in Eq. (4). We also compare our results with two other recent fits in which the new $`\mathrm{\Delta }M_s`$-limit has been incorporated, but which differ from us in details which we shall specify below.
If new physics (of any type) is present, the principal way in which it can enter in flavour physics is via new contributions, possibly with new phases, to $`K^0`$$`\overline{K^0}`$, $`B_d^0`$$`\overline{B_d^0}`$ and $`B_s^0`$$`\overline{B_s^0}`$ mixing. The tree decay amplitudes, being dominated by virtual $`W`$ exchange, remain essentially unaffected by new physics. Thus, even in the presence of new physics, the measured values of $`|V_{cb}|`$ and $`|V_{ub}|`$ correspond to their true SM values, so that two sides of the UT are unaffected. However, the third side, which depends on $`|V_{td}|`$, will in general be affected by new physics. Furthermore, the measurements of $`|ฯต|`$ and $`B_s^0`$$`\overline{B_s^0}`$ mixing, which provide additional constraints on the UT, will also be affected. If Nature is kind, the unitarity triangle, as constructed from direct measurements of $`\alpha `$, $`\beta `$ and $`\gamma `$, will be inconsistent with that obtained from independent measurements of the sides. If this were to happen, it would be clear evidence for the presence of physics beyond the SM, and would be most exciting. In such a case, the new physics is also expected to modify the decay rates and distributions of rare $`B`$-decays such as $`BX_s\gamma `$, $`BX_s\mathrm{}^+\mathrm{}^{}`$ and $`BX_s\nu \overline{\nu }`$, and of related exclusive decays. (Similarly, the corresponding decays dominated by the $`bd`$ transitions may also be affected.)
One type of new physics which has been extensively studied is supersymmetry (SUSY). A hint suggesting that SUSY might indeed be around the corner is the gauge-coupling unification: a supersymmetric grand unified theory does better than its non-supersymmetric counterpart. A great deal of effort has gone into a systematic study of the pattern of flavour violation in SUSY, in particular in the flavour-changing neutral-current processes in $`K`$ and $`B`$ decays. We shall concentrate here on the minimal supersymmetric standard model (MSSM), and update the anticipated profile of the UT and CP-phases which we presented earlier. Of particular interest here is the scenario of minimal supersymmetric flavour violation, which involves, in addition to the SM degrees of freedom, charged Higgs bosons, a light stop (assumed right-handed) and a light chargino, with all other degrees of freedom assumed heavy and hence effectively integrated out. This scenario can be embedded in supergravity (SUGRA) models with gauge-mediated supersymmetry breaking, in which the first two squark generations and the gluinos are assumed heavy. Regardless of which variant is used, the key assumption in our analysis is that there are no new phases in the couplings โ although there are many new contributions to meson mixing and rare decays, all are proportional to the same combination of CKM matrix elements as found in the SM. As explained above, in this class of models measurements of the CP phases will yield the true SM values for these quantities. However, measurements of meson mixing and rare decays will be affected by the presence of this new physics.
In Section 2, we discuss the profile of the unitarity triangle within the SM. We describe the input data used in the fits and present the allowed region in $`\rho `$$`\eta `$ space, as well as the presently-allowed ranges for the CP angles $`\alpha `$, $`\beta `$ and $`\gamma `$. We turn to supersymmetric models in Section 3. We review several variants of the MSSM, in which the new CP-violating phases are essentially zero. We also discuss the NLO corrections in such models and show that the SUSY contributions to $`K^0`$$`\overline{K^0}`$, $`B_d^0`$$`\overline{B_d^0}`$ and $`B_s^0`$$`\overline{B_s^0}`$ mixing are of the same form and can be characterized by a single parameter $`f`$. We compare the profile of the unitarity triangle in SUSY models, for various values of $`f`$, with that of the SM. We conclude in Section 4.
## 2 Unitarity Triangle: SM Profile
### 2.1 Input Data
We briefly describe below the experimental and theoretical data which constrain the CKM parameters. (For more details, we refer the reader to.) A summary can be found in Table 1.
* The CKM parameters $`\lambda `$, $`A`$, $`\rho `$ and $`\eta `$ are directly constrained through measurements of the CKM elements $`|V_{us}|=\lambda `$, $`|V_{cb}|`$ and $`|V_{ub}/V_{cb}|`$. In our fits we ignore the small error on $`\lambda `$. Also, the error on $`|V_{ub}/V_{cb}|`$ includes some theoretical model dependence.
* $`|ฯต|,\widehat{B}_K`$: In the standard model, $`|ฯต|`$ is essentially proportional to the imaginary part of the box diagram for $`K^0`$$`\overline{K^0}`$ mixing and is given by
$`|ฯต|`$ $`=`$ $`{\displaystyle \frac{G_F^2f_K^2M_KM_W^2}{6\sqrt{2}\pi ^2\mathrm{\Delta }M_K}}\widehat{B}_K\left(A^2\lambda ^6\overline{\eta }\right)(y_c\{\widehat{\eta }_{ct}f_3(y_c,y_t)\widehat{\eta }_{cc}\}`$ (5)
$`+\widehat{\eta }_{tt}y_tf_2(y_t)A^2\lambda ^4(1\overline{\rho })),`$
where $`y_im_i^2/M_W^2`$, and the functions $`f_2`$ and $`f_3`$ can be seen in. Here, the $`\widehat{\eta }_i`$ are QCD correction factors, calculated at next-to-leading order in Refs. ($`\widehat{\eta }_{cc}`$), ($`\widehat{\eta }_{tt}`$) and ($`\widehat{\eta }_{ct}`$). The theoretical hadronic uncertainty in the expression for $`|ฯต|`$ is in the renormalization-scale independent parameter $`\widehat{B}_K`$. In Table 1, the $`|ฯต|`$ entry is taken from Ref., while that for $`\widehat{B}_K`$ is based on lattice QCD methods, summarized in Ref..
* $`\mathrm{\Delta }M_d,f_{B_d}^2\widehat{B}_{B_d}`$: The mass difference $`\mathrm{\Delta }M_d`$ is calculated from the $`B_d^0`$$`\overline{B_d^0}`$ box diagram, which is dominated by $`t`$-quark exchange:
$$\mathrm{\Delta }M_d=\frac{G_F^2}{6\pi ^2}M_W^2M_B\left(f_{B_d}^2\widehat{B}_{B_d}\right)\widehat{\eta }_By_tf_2(y_t)|V_{td}^{}V_{tb}|^2,$$
(6)
where, using Eq. (1), $`|V_{td}^{}V_{tb}|^2=A^2\lambda ^6[\left(1\overline{\rho }\right)^2+\overline{\eta }^2]`$. Here, $`\widehat{\eta }_B=0.55`$ is the QCD correction, calculated in the $`\overline{MS}`$ scheme. Consistency requires that the top quark mass be rescaled from its pole (mass) value of $`m_t=175\pm 5`$ GeV to the value $`\overline{m_t}(m_t(pole))=165\pm 5`$ GeV in the $`\overline{MS}`$ scheme. The slight dependence of $`\widehat{\eta }_B`$ on $`\overline{m_t}(m_t(pole))`$ in the range given here is ignored. The entry for $`\mathrm{\Delta }M_d`$ in Table 1 is taken from Ref..
For the $`B`$ system, the hadronic uncertainty is given by $`f_{B_d}^2\widehat{B}_{B_d}`$, analogous to $`\widehat{B}_K`$ in the kaon system. Present estimates of this quantity are summarized in Ref., yielding $`f_{B_d}\sqrt{\widehat{B}_{B_d}}=(190\pm 23)`$ MeV in the quenched approximation. The effect of unquenching is not yet understood completely. Taking the MILC collaboration estimates of unquenching would increase the central value of $`f_{B_d}\sqrt{\widehat{B}_{B_d}}`$ by $`21`$ MeV. The range of $`f_{B_d}\sqrt{\widehat{B}_{B_d}}`$ given in Table 1 is chosen to take all these considerations into account.
* $`\mathrm{\Delta }M_s,f_{B_s}^2\widehat{B}_{B_s}`$: Mixing in the $`B_s^0`$$`\overline{B_s^0}`$ system is quite similar to that in the $`B_d^0`$$`\overline{B_d^0}`$ system. The $`B_s^0`$$`\overline{B_s^0}`$ box diagram is again dominated by $`t`$-quark exchange, and the mass difference between the mass eigenstates $`\mathrm{\Delta }M_s`$ is given by a formula analogous to that of Eq. (6):
$$\mathrm{\Delta }M_s=\frac{G_F^2}{6\pi ^2}M_W^2M_{B_s}\left(f_{B_s}^2\widehat{B}_{B_s}\right)\widehat{\eta }_{B_s}y_tf_2(y_t)|V_{ts}^{}V_{tb}|^2.$$
(7)
Using the fact that $`|V_{cb}|=|V_{ts}|`$ (Eq. (1)), it is clear that one of the sides of the unitarity triangle, $`|V_{td}/\lambda V_{cb}|`$, can be obtained from the ratio of $`\mathrm{\Delta }M_d`$ and $`\mathrm{\Delta }M_s`$,
$$\frac{\mathrm{\Delta }M_s}{\mathrm{\Delta }M_d}=\frac{\widehat{\eta }_{B_s}M_{B_s}\left(f_{B_s}^2\widehat{B}_{B_s}\right)}{\widehat{\eta }_{B_d}M_{B_d}\left(f_{B_d}^2\widehat{B}_{B_d}\right)}\left|\frac{V_{ts}}{V_{td}}\right|^2.$$
(8)
The only real uncertainty in this quantity is the ratio of hadronic matrix elements $`f_{B_s}^2\widehat{B}_{B_s}/f_{B_d}^2\widehat{B}_{B_d}`$. It is now widely accepted that the ratio $`\xi _s(f_{B_s}\sqrt{\widehat{B}_{B_s}})/(f_{B_d}\sqrt{\widehat{B}_{B_d}})`$ is probably the most reliable of the lattice-QCD estimates in $`B`$ physics. The value given Table 1 is based on Ref..
The present lower bound on $`\mathrm{\Delta }M_s`$ is: $`\mathrm{\Delta }M_s>14.3\text{(ps)}^1`$ (at $`95\%`$ C.L.). This bound has been established using the so-called โamplitude methodโ, and we follow this method in including the current information about $`B_s^0`$$`\overline{B_s^0}`$ mixing in the fits.
Referring to Table 1, we see that the quantities with the largest errors are $`\widehat{\eta }_{cc}`$ (28%), $`\widehat{B}_K`$ (16%), $`|V_{ub}/V_{cb}|`$ (15%) and $`f_{B_d}\sqrt{\widehat{B}_{B_d}}`$ (19%). Of these, the latter three are extremely important in defining the allowed $`\rho `$$`\eta `$ region (the large error on $`\widehat{\eta }_{cc}`$ does not affect the fit very much). The errors on two of these quantities โ $`\widehat{B}_K`$ and $`f_{B_d}\sqrt{\widehat{B}_{B_d}}`$ โ are purely theoretical in origin, and the error on $`|V_{ub}/V_{cb}|`$ has a significant theoretical component (model dependence). Thus, the present uncertainty in the shape of the unitarity triangle is due in large part to theoretical errors. Reducing these errors will be quite important in getting a precise profile of the unitarity triangle and the CP-violating phases.
There are two other measurements which should be mentioned here. First, the KTEV collaboration at Fermilab and the NA48 collaboration at CERN have reported in 1999 new measurements of direct CP violation in the $`K`$ sector through the ratio $`ฯต^{}/ฯต`$. Their results, together with those of the earlier experiments NA31 and E731 are as follows:
$`\mathrm{Re}(ฯต^{}/ฯต)`$ $`=`$ $`(28.0\pm 4.1)\times 10^4[\text{KTEV โ99}],`$
$`=`$ $`(18.5\pm 7.3)\times 10^4[\text{NA48 โ99}],`$
$`=`$ $`(23.0\pm 6.5)\times 10^4[\text{NA31 โ93}],`$
$`=`$ $`(7.4\pm 5.9)\times 10^4[\text{E731 โ93}],`$
yielding the present world average $`\mathrm{Re}(ฯต^{}/ฯต)=(21.2\pm 4.6)\times 10^4`$. This combined result excludes the superweak model.
A great deal of theoretical effort has gone into calculating this quantity at next-to-leading order accuracy in the SM. However, numerical estimates require a number of non-perturbative parameters, which are at present poorly known, yielding an theoretical uncertainty which is larger than an order of magnitude. Thus, whereas $`ฯต^{}/ฯต`$ represents a landmark measurement, removing the superweak model of Wolfenstein from further consideration, its impact on CKM phenomenology, particularly in constraining the CKM parameters, is marginal as $`ฯต^{}/ฯต`$ is dominated by non-perturbative uncertainties.
Second, the CDF collaboration has recently made a measurement of $`\mathrm{sin}2\beta `$. In the Wolfenstein parametrization, $`\beta `$ is the phase of the CKM matrix element $`V_{td}`$. From Eq. (1) one can readily find that
$$\mathrm{sin}(2\beta )=\frac{2\overline{\eta }(1\overline{\rho })}{(1\overline{\rho })^2+\overline{\eta }^2}.$$
(10)
Thus, a measurement of $`\mathrm{sin}2\beta `$ would put a strong contraint on the parameters $`\overline{\rho }`$ and $`\overline{\eta }`$. However, the CDF measurement gives
$$\mathrm{sin}2\beta =0.79_{0.44}^{+0.41},$$
(11)
or $`\mathrm{sin}2\beta >0`$ at 93% C.L. As we will see in the next section, this constraint is quite weak โ the indirect measurement (reported here) already constrains $`0.53\mathrm{sin}2\beta 0.93`$ at the 95% C.L. in the SM. In view of this, and given that it is not clear how to combine the above measurement (which allows for unphysical values of $`\mathrm{sin}2\beta `$) with the other data, we have not included this measurement in our fits.
### 2.2 SM Fits
In order to find the allowed region in $`\overline{\rho }`$$`\overline{\eta }`$ space, i.e. the allowed shapes of the unitarity triangle, the computer program MINUIT is used to fit the parameters to the constraints described above. In the fit, we allow ten parameters to vary: $`\overline{\rho }`$, $`\overline{\eta }`$, $`A`$, $`m_t`$, $`m_c`$, $`\eta _{cc}`$, $`\eta _{ct}`$, $`f_{B_d}\sqrt{\widehat{B}_{B_d}}`$, $`\widehat{B}_K`$, and $`\xi _s`$. The $`\mathrm{\Delta }M_s`$ constraint is included using the amplitude method. The allowed (95% C.L.) $`\overline{\rho }`$$`\overline{\eta }`$ region is shown in Fig. 1. The triangle drawn is to facilitate our discussions, and corresponds to the central values of the fits, $`(\alpha ,\beta ,\gamma )=(93^{},24^{},63^{})`$.
The CP angles $`\alpha `$, $`\beta `$ and $`\gamma `$ can be measured in CP-violating rate asymmetries in $`B`$ decays. These angles can be expressed in terms of $`\overline{\rho }`$ and $`\overline{\eta }`$. Thus, different shapes of the unitarity triangle are equivalent to different values of the CP angles. Referring to Fig. 1, the allowed ranges at 95% C.L. are given by
$$75^{}\alpha 121^{},16^{}\beta 34^{},38^{}\gamma 81^{},$$
(12)
or, equivalently,
$$0.88\mathrm{sin}2\alpha 0.50,0.53\mathrm{sin}2\beta 0.93,0.38\mathrm{sin}^2\gamma 0.98.$$
(13)
Of course, the values of $`\alpha `$, $`\beta `$ and $`\gamma `$ are correlated, i.e. they are not all allowed simultaneously. After all, the sum of these angles must equal $`180^{}`$. We illustrate these correlations in Figs. 2 and 3. In both of these figures, the SM plot is labelled by $`f=0`$. Fig. 2 shows the allowed region in $`\mathrm{sin}2\alpha `$$`\mathrm{sin}2\beta `$ space allowed by the data. And Fig. 3 shows the allowed (correlated) values of the CP angles $`\alpha `$ and $`\gamma `$. This correlation is roughly linear, due to the relatively small allowed range of $`\beta `$ \[Eq. (12)\].
The allowed ranges for the CKM-parameters obtained from our unitarity fits can be compared with those obtained by other groups. For example, concentrating on $`\mathrm{sin}2\alpha `$ and $`\mathrm{sin}2\beta `$, Plaszczynski and Schune get the following (95% C.L.) ranges:
$$0.95\mathrm{sin}2\alpha 0.50,0.50\mathrm{sin}2\beta 0.85,$$
(14)
which are very similar to the ranges obtained by us for these quantities \[Eq. (13)\]. While there are smallish differences in the input parameters from experimental measurements, the real difference in the two fits lies in the incorporation of the theoretical uncertainties. We have treated theoretical and experimental errors on the same footing. On the other hand, Plaszczynski and Schune have scanned over a โreasonable rangeโ of theoretical parameters, determined the allowed contours for fixed values of these parameters and taken the envelope of all the individual contours obtained in the allowed range. Of course, the size of the resulting envelope depends on the assumed theoretical range, so that a certain amount of subjectivity is already embedded. Given that the parametric input in the present analysis and in are similar, the closeness of the two fits implies that they do not depend sensitively on the prescription for handling theoretical errors.
In fact, one can turn the argument around: with improved limits on (or an actual measurement of) $`\mathrm{\Delta }M_s`$, the theoretical errors on $`f_{B_d}\sqrt{\widehat{B}_{B_d}}`$ and $`\widehat{B}_K`$ can be effectively reduced. To quantify these remarks, we examine the presently-allowed correlation in the parameters $`\widehat{B}_K`$ and $`f_{B_d}\sqrt{\widehat{B}_{B_d}}`$ which follows from our fits in the SM. Recall that the theoretically-allowed ranges for these quantities are $`f_{B_d}\sqrt{\widehat{B}_{B_d}}=215\pm 40`$ MeV and $`\widehat{B}_K=0.94\pm 0.15`$. Rather than present the 95% c.l. region in the $`\overline{\rho }`$$`\overline{\eta }`$ plane (Fig. 1), we use the fits to find the allowed 95% c.l. region in the $`\widehat{B}_K`$$`f_{B_d}\sqrt{\widehat{B}_{B_d}}`$ plane. The results are shown in Fig. 4, where we have allowed the hadronic parameters to vary in the range $`135\text{MeV}f_{B_d}\sqrt{\widehat{B}_{B_d}}295`$ MeV and $`0.64\widehat{B}_K1.24`$, which corresponds to allowing a $`\pm 2\sigma `$ uncertainty on each. (Note that there appears to be some structure near the solid line on the left-hand side of the figure. This is a numerical artifact due to the binning of the $`\mathrm{\Delta }M_s`$ data, and can be ignored. Only the solid line is important.) Only values of $`f_{B_d}\sqrt{\widehat{B}_{B_d}}`$ and $`\widehat{B}_K`$ which lie between the two solid lines in Fig. 4 are allowed at the 95% C.L. Note that present data do not allow a value $`f_{B_d}\sqrt{\widehat{B}_{B_d}}165`$ MeV. Likewise, values of $`f_{B_d}\sqrt{\widehat{B}_{B_d}}`$ in excess of $`230`$ MeV are highly correlated with the value of $`\widehat{B}_K`$. Thus, no values in excess of $`230`$ MeV are allowed for $`f_{B_d}\sqrt{B_{B_d}}`$ if $`\widehat{B}_K0.6`$ in the SM. This is very similar (though not identical) to the correlation shown in Ref..
Of course, one obtains more stringent constraints on the CKM parameters if significantly reduced errors are assumed for the input parameters. For example, a recent fit, assuming $`\delta V_{ub}/V_{cb}=\pm 8.8\%`$ (compared to $`\delta V_{ub}/V_{cb}=\pm 15\%`$ used here), and $`f_{B_d}\sqrt{\widehat{B}_{B_d}}=220\pm 28\text{MeV}`$ (as opposed to $`f_{B_d}\sqrt{\widehat{B}_{B_d}}=215\pm 40\text{MeV}`$ in Table 1), leads to a more precise determination of the apex of the unitarity triangle. In turn, this yields the following 95% C.L. ranges for the CP asymmetries:
$$0.73\mathrm{sin}2\alpha 0.26,0.63\mathrm{sin}2\beta 0.81,0.51\mathrm{sin}^2\gamma 0.93.$$
(15)
## 3 Unitarity Triangle: A SUSY Profile
In this section we update the profile of the unitarity triangle in supersymmetric (SUSY) theories. In general, minimal supersymmetric standard models (MSSM) have three physical phases, apart from the QCD vacuum parameter $`\overline{\theta }_{QCD}`$ which we shall take to be zero. The three phases are: (i) the CKM phase represented here by the Wolfenstein parameter $`\eta `$, (ii) the phase $`\theta _A=\mathrm{arg}(A)`$, and (iii) the phase $`\theta _\mu =\mathrm{arg}(\mu )`$. The last two phases, residing in the soft SUSY-breaking terms and in the scalar superpotential, are peculiar to SUSY models and their effects must be taken into account in a general supersymmetric framework. In particular, the CP-violating asymmetries which result from the interference between mixing and decay amplitudes can produce non-standard effects. Concentrating here on the $`\mathrm{\Delta }B=2`$ amplitudes, two new phases $`\theta _d`$ and $`\theta _s`$ arise, which can be parametrized as follows:
$$\theta _{d,s}=\frac{1}{2}\mathrm{arg}\left(\frac{B_{d,s}|_{eff}^{SUSY}|\overline{B}_{d,s}}{B_{d,s}|_{eff}^{SM}|\overline{B}_{d,s}}\right),$$
(16)
where $`^{SUSY}`$ is the effective Hamiltonian including both the SM degrees of freedom and the SUSY contributions. Thus, CP-violating asymmetries in $`B`$ decays would involve not only the phases $`\alpha `$, $`\beta `$ and $`\gamma `$, defined previously, but additionally $`\theta _d`$ or $`\theta _s`$. In other words, the SUSY contributions to the real parts of $`M_{12}(B_d)`$ and $`M_{12}(B_s)`$ are no longer proportional to the CKM matrix elements $`V_{td}V_{tb}^{}`$ and $`V_{ts}V_{tb}^{}`$, respectively. If $`\theta _d`$ or $`\theta _s`$ were unconstrained, one could not make firm predictions about the CP asymmetries in SUSY models. In such a case, an analysis of the profile of the unitarity triangle in such models would be futile.
However, the experimental upper limits on the electric dipole moments (EDMs) of the neutron and electron do provide a constraint on the phase $`\theta _\mu `$. In supergravity (SUGRA) models with a priori complex parameters $`A`$ and $`\mu `$, the phase $`\theta _\mu `$ is strongly bounded with $`\theta _\mu <0.01\pi `$.
As for the phase $`\theta _A`$, it can be of $`O(1)`$ in the small $`\theta _\mu `$ region, as far as the EDMs are concerned. However, in both the $`\mathrm{\Delta }S=2`$ and $`\mathrm{\Delta }B=2`$ transitions, and for low-to-moderate values of $`\mathrm{tan}\upsilon `$ <sup>1</sup><sup>1</sup>1In supersymmetric jargon, the quantity $`\mathrm{tan}\beta `$ is used to define the ratio of the two vacuum expectation values (vevs) $`\mathrm{tan}\beta v_u/v_d`$, where $`v_d(v_u)`$ is the vev of the Higgs field which couples exclusively to down-type (up-type) quarks and leptons. (See, for example, the review by Haber in Ref.). However, in discussing flavour physics, the symbol $`\beta `$ is traditionally reserved for one of the angles of the unitarity triangle. To avoid confusion, we will call the ratio of the vevs $`\mathrm{tan}\upsilon `$., it has been shown that $`\theta _A`$ does not change the phase of either the matrix element $`M_{12}(K)`$ or of $`M_{12}(B)`$. Hence, in SUGRA models, $`\mathrm{arg}M_{12}(B)|_{SUGRA}=\mathrm{arg}M_{12}(B)|_{SM}=\mathrm{arg}(\xi _t^2)`$, where $`\xi _t=V_{td}^{}V_{tb}`$. Likewise, the phase of the SUSY contribution in $`M_{12}(K)`$ is aligned with the phase of the $`t\overline{t}`$-contribution in $`M_{12}(K)`$, given by $`\mathrm{arg}(V_{td}V_{ts}^{})`$.
Thus, in SUGRA models, one can effectively set $`\theta _d0`$ and $`\theta _s0`$, so that the CP-violating asymmetries give information about the SM phases $`\alpha `$, $`\beta `$ and $`\gamma `$. Hence, an analysis of the UT and CP-violating phases $`\alpha `$, $`\beta `$ and $`\gamma `$ can be carried out in a very similar fashion as in the SM, taking into account the additional contributions to $`M_{12}(K)`$ and $`M_{12}(B)`$.
### 3.1 NLO Corrections to $`\mathrm{\Delta }M_d`$, $`\mathrm{\Delta }M_s`$ and $`ฯต`$ in Minimal SUSY Flavour Violation
A number of SUSY models share the features mentioned in the previous subsection, and the supersymmetric contributions to the mass differences $`M_{12}(B)`$ and $`M_{12}(K)`$ have been analyzed in a number of papers, following the pioneering work of Ref.. The SUSY contributions to $`\mathrm{\Delta }M_d`$, $`\mathrm{\Delta }M_s`$ and $`|ฯต|`$ in supersymmetric theories can be incorporated in a simple form:
$`\mathrm{\Delta }M_d`$ $`=`$ $`\mathrm{\Delta }M_d(SM)[1+f_d(m_{\chi _2^\pm },m_{\stackrel{~}{t}_R},m_{H^\pm },\mathrm{tan}\upsilon )],`$
$`\mathrm{\Delta }M_s`$ $`=`$ $`\mathrm{\Delta }M_s(SM)[1+f_s(m_{\chi _2^\pm },m_{\stackrel{~}{t}_R},m_{H^\pm },\mathrm{tan}\upsilon )],`$
$`|ฯต|`$ $`=`$ $`{\displaystyle \frac{G_F^2f_K^2M_KM_W^2}{6\sqrt{2}\pi ^2\mathrm{\Delta }M_K}}\widehat{B}_K\left(A^2\lambda ^6\overline{\eta }\right)(y_c\{\widehat{\eta }_{ct}f_3(y_c,y_t)\widehat{\eta }_{cc}\}`$ (17)
$`+\widehat{\eta }_{tt}y_tf_2(y_t)[1+f_ฯต(m_{\chi _2^\pm },m_{\stackrel{~}{t}_2},m_{H^\pm },\mathrm{tan}\upsilon )]A^2\lambda ^4(1\overline{\rho })).`$
The quantities $`f_d`$, $`f_s`$ and $`f_ฯต`$ can be expressed as
$$f_d=f_s=\frac{\widehat{\eta }_{2,S}(B)}{\widehat{\eta }_B}R_{\mathrm{\Delta }_d}(S),f_ฯต=\frac{\widehat{\eta }_{2,S}(K)}{\widehat{\eta }_{tt}}R_{\mathrm{\Delta }_d}(S),$$
(18)
where $`R_{\mathrm{\Delta }_d}(S)`$ is defined as
$$R_{\mathrm{\Delta }_d}(S)\frac{\mathrm{\Delta }M_d(SUSY)}{\mathrm{\Delta }M_d(SM)}(LO)=\frac{S}{y_tf_2(y_t)}.$$
(19)
The supersymmetric function $`S`$ is given in Ref., and the NLO functions $`\widehat{\eta }_{2,S}(B)`$ and $`\widehat{\eta }_{2,S}(K)`$ can be found in Ref.. The functions $`f_i`$, $`i=d,s,ฯต`$ are all positive definite, i.e. the supersymmetric contributions add constructively to the SM contributions in the entire allowed supersymmetric parameter space. The two QCD correction factors appearing in Eq. (18) are numerically very close to one another, with $`\widehat{\eta }_{2,S}(B)/\widehat{\eta }_B\widehat{\eta }_{2,S}(K)/\widehat{\eta }_{tt}=0.93`$. Thus, to an excellent approximation, one has $`f_d=f_s=f_ฯตf`$.
How big can $`f`$ be? This quantity is a function of the masses of the top squark, chargino and the charged Higgs, $`m_{\stackrel{~}{t}_R}`$, $`m_{\stackrel{~}{\chi }_2^\pm }`$ and $`m_{H^\pm }`$, respectively, as well as of $`\mathrm{tan}\upsilon `$. The maximum allowed value of $`f`$ depends on the model (minimal SUGRA, non-minimal SUGRA, MSSM with constraints from EDMs, etc.). From the published results we conclude that typically $`f`$ can be as large as $`0.45`$ in non-minimal SUGRA models for low $`\mathrm{tan}\upsilon `$ (typically $`\mathrm{tan}\upsilon =2`$), and approximately half of this value in minimal SUGRA models. Relaxing the SUGRA mass constraints, admitting complex values of $`A`$ and $`\mu `$ but incorporating the EDM constraints, and imposing the constraints mentioned above, $`f`$ could be larger. In all cases, the value of $`f`$ decreases with increasing $`\mathrm{tan}\upsilon `$ or increasing $`m_{\stackrel{~}{\chi }_2^\pm }`$ and $`m_{\stackrel{~}{t}_R}`$, as noted above.
### 3.2 SUSY Fits
For the SUSY fits, we use the same program as for the SM fits, except that the theoretical expressions for $`\mathrm{\Delta }M_d`$, $`\mathrm{\Delta }M_s`$ and $`|ฯต|`$ are modified as in Eq. (17). We compare the fits for four representative values of the SUSY function $`f`$ โ 0, 0.2, 0.4 and 0.75 โ which are typical of the SM, minimal SUGRA models, non-minimal SUGRA models, and non-SUGRA models with EDM constraints, respectively.
The allowed 95% C.L. regions for these four values of $`f`$ are all plotted in Fig. 5. As is clear from this figure, there is still a considerable overlap between the $`f=0`$ (SM) and $`f=0.75`$ regions. However, there are also regions allowed for one value of $`f`$ which are excluded for another value. Thus a sufficiently precise determination of the unitarity triangle might be able to exclude certain values of $`f`$ (including the SM, $`f=0`$).
From Fig. 5 it is clear that a measurement of the CP angle $`\beta `$ will not distinguish among the various values of $`f`$. Rather, it is the measurement of $`\gamma `$ or $`\alpha `$ which has the potential to rule out certain values of $`f`$. As $`f`$ increases, the allowed region moves slightly down and towards the right in the $`\overline{\rho }`$$`\overline{\eta }`$ plane, corresponding to smaller values of $`\gamma `$ (or equivalently, larger values of $`\alpha `$). We illustrate this in Table 2, where we present the allowed ranges of $`\alpha `$, $`\beta `$ and $`\gamma `$, as well as their central values (corresponding to the preferred values of $`\overline{\rho }`$ and $`\overline{\eta }`$), for each of the four values of $`f`$. From this Table, we see that the allowed range of $`\beta `$ is largely insensitive to the model. Conversely, the allowed values of $`\alpha `$ and $`\gamma `$ do depend somewhat strongly on the chosen value of $`f`$. Note, however, that one is not guaranteed to be able to distinguish among the various models: as mentioned above, there is still significant overlap among all four models. Thus, depending on what values of $`\alpha `$ and $`\gamma `$ are obtained, we may or may not be able to rule out certain values of $`f`$.
For completeness, in Table 3 we present the corresponding allowed ranges for the CP asymmetries $`\mathrm{sin}2\alpha `$, $`\mathrm{sin}2\beta `$ and $`\mathrm{sin}^2\gamma `$. Again, we see that the allowed range of $`\mathrm{sin}2\beta `$ is largely independent of the value of $`f`$. On the other hand, as $`f`$ increases, the allowed values of $`\mathrm{sin}2\alpha `$ become increasingly negative, while those of $`\mathrm{sin}^2\gamma `$ become smaller.
The allowed (correlated) values of the CP angles for various values of $`f`$ can be clearly seen in Figs. 2 and 3. As $`f`$ increases from 0 (SM) to 0.75, the change in the allowed $`\mathrm{sin}2\alpha `$$`\mathrm{sin}2\beta `$ (Fig. 2) and $`\alpha `$$`\gamma `$ (Fig. 3) regions is quite significant.
## 4 Conclusions
In the very near future, CP-violating asymmetries in $`B`$ decays will be measured at $`B`$-factories, HERA-B and hadron colliders. Such measurements will give us crucial information about the interior angles $`\alpha `$, $`\beta `$ and $`\gamma `$ of the unitarity triangle. If we are lucky, there will be an inconsistency in the independent measurements of the sides and angles of this triangle, thereby revealing the presence of new physics.
An interesting possibility, from the point of view of making predictions, are models which contribute to $`B^0`$$`\overline{B^0}`$ mixings and $`|ฯต|`$, but without new phases. One type of new physics which does just this is supersymmetry (SUSY). There are some SUSY models which do contain new phases, but they suffer from a lack of predictivity. However, there is also a large class of SUSY models with no new phases. In these models, there are new, supersymmetric contributions to $`K^0`$$`\overline{K^0}`$, $`B_d^0`$$`\overline{B_d^0}`$ and $`B_s^0`$$`\overline{B_s^0}`$ mixing. The key ingredient in our analysis is the fact that these contributions, which add constructively to the SM, depend on the SUSY parameters in essentially the same way. That is, so far as an analysis of the unitarity triangle is concerned, there is a single parameter, $`f`$, which characterizes the various SUSY models within this class of models ($`f=0`$ corresponds to the SM).
We have therefore updated the profile of the unitarity triangle in both the SM and some variants of the MSSM. We have used the latest experimental data on $`|V_{cb}|`$, $`|V_{ub}/V_{cb}|`$, $`\mathrm{\Delta }M_d`$ and $`\mathrm{\Delta }M_s`$, as well as the latest theoretical estimates (including errors) of $`\widehat{B}_K`$, $`f_{B_d}\sqrt{\widehat{B}_{B_d}}`$ and $`\xi _sf_{B_d}\sqrt{\widehat{B}_{B_d}}/f_{B_s}\sqrt{\widehat{B}_{B_s}}`$. In addition to $`f=0`$ (SM), we considered three SUSY values of $`f`$: 0.2, 0.4 and 0.75, representing the minimal SUGRA models, non-minimal SUGRA models, and non-SUGRA models with EDM constraints, respectively.
We first considered the profile of the unitarity triangle in the SM, shown in Fig. 1. We then compared the SM with the different SUSY models. The result can be seen in Fig. 5. As $`f`$ increases, the allowed region moves slightly down and to the right in the $`\overline{\rho }`$$`\overline{\eta }`$ plane. The main conclusion from this analysis is that the measurement of the CP angle $`\beta `$ will not distinguish among the SM and the various SUSY models โ the allowed region of $`\beta `$ is virtually the same in all these models. On the other hand, the allowed ranges of $`\alpha `$ and $`\gamma `$ do depend on the choice of $`f`$. For example, larger values of $`f`$ tend to favour smaller values of $`\gamma `$. The dependence of the CP angles on the value of $`f`$ is illustrated clearly in Tables 2 and 3. Thus, with measurements of $`\gamma `$ or $`\alpha `$, we may be able to rule out certain values of $`f`$ (including the SM, $`f=0`$). However, we also note that there is no guarantee of this happening โ at present there is still a significant region of overlap among all four models.
Acknowledgements:
One of us (A.A.) would like to thank the organizers of the DAPHNE โ99 workshop, in particular Giorgio Capon and Gino Isidori, for their kind hospitality. The work of D.L. was financially supported by NSERC of Canada. |
warning/0002/cond-mat0002191.html | ar5iv | text | # Transition Temperature and Magnetoresistance in Double-Exchange Compounds with Moderate Disorder
## I Introduction
Interest has revived recently in the perovskite manganese oxides $`\mathrm{A}_{1\mathrm{x}}\mathrm{B}_\mathrm{x}\mathrm{MnO}_3`$ (where A is a trivalent and B is a trivalent atom), which were first investigated in the 1950โs. As the doping $`x`$ and the temperature $`T`$ are varied, these manganese oxides show a rich variety of phases. Particularly interesting is the doping region $`0.1x0.3`$, where the compounds undergo a transition from either insulating or very high resistance metallic, paramagnetic phase at high temperatures to a ferromagnetic phase at low temperatures. Near the transition, the resistivity of the compounds changes by orders of magnitude. The application of a strong magnetic field substantially reduces this effect, thus giving rise to a very large negative magnetoresistance. The physical mechanism, responsible for this behaviour, has been recently the subject of much discussion and controversy. It was initally suggested, that the CMR in manganese oxides can be explained within the framework of the Double-Exchange Model (DE). In this model it is assumed that the on-site direct repulsion $`U`$ is the largest energy, followed in order by the Hundโs rule energy $`J`$ and the hybridization energy $`t`$ between Mn-orbitals at neighboring sites. The basic conduction step is then the interchange of valence between neighboring Mn : $`[\mathrm{Mn}^{+3}\mathrm{Mn}^{+4}\mathrm{Mn}^{+4}\mathrm{Mn}^{+3}]`$. The basic physical idea of the DE mechanism is that this electron conduction is largest when the initial and the final states are degenerate. The latter requirement corresponds to an alignment of the spins of the manganese ions. In the opposite case, the conduction rate is suppressed by a factor of $`t/J`$. As a result, a transition from a paramagnetic to a ferromagnetic state leads to a dramatic increase of the conductivity of the compound. Using a Dynamical Mean Field calculation (DMFT), the double-exchange mechanism was successfully used for a quantitative description of the experimental data in $`\mathrm{LaSrMnO}_3`$ compounds. A later study claimed that the agreement with the experiment found in Ref. was caused by an unphysical choise of the density of states (DOS), and was accidental. But a calculation by Furukawa with several different choises of the local DOS confirmed the results of his earlier work.
Subsequently a calculation carried out by the authors of Ref. , concluded that the double exchange model alone could not explain the experimental data for the manganese oxides. There were two objections: (i) that the double-exchange model gave a transition temperature an order of magnitude larger than experiments and (ii) that the often observed insulating-like resistivity (resistivity increasing with decreasing temperature) could not be explained by the double-exchange model. It was proposed in Ref. , that for the description of mangenese oxides, one should take into account a continuation to the metallic state of the Jahn-Teller distortion found for the insulating antiferromagnetic end-member ($`x0`$) in these compounds into the $`x`$-range of interest min some kind of dynamic fashion. As shown by a simple calculation , objection (i) turns out to be due to an inadequate appreciation of the energetics of the double-exchange process. The transition temperature is related to the difference in the electronic cohesive energy of the ferromagnetic and paramagnetic phases, and is not given by the transition temperature of a spin model as in Ref. . As regards the proposal of the effects of a possible Jahn-Teller distortion, substantial theoretical effort has failed to produce any results which can be compared to experiments for the resistivity.
Meanwhile, there has been further progress experimentally. It was only recently pointed out, that the manganese oxides at similar elecron densities show two qualitatively different types of behaviour: (i) a metal-insulator transition near $`T_c`$, which in this case has relatively low values $`280`$ K and (ii) a metallic behaviour both below (a good metal) and above (incoherent metal with the absolute value of the resistivity near the Mottโs limit) the critical temperature, which is comparably high ($`380`$ K). The difference appears to be the amount of disorder. This would tend to remove the possibility that the Jahn-Teller effects, were they to occur, have much to do with the resistivity behavior. Instead the question to ask is why disorder so dramatically modifies the temperature dependence of the conductivity in the paramagnetic phase, while simultaneously reducing the transition temperature. The relation of the resitivity to the magnetization $`M(T)`$ in the ferromagnetic phase also depends on disorder. For small disorder the temperature dependent part is proportional to the $`M(T)^2`$, while for large disorder a much stronger dependence is found.
It was suggested that disorder effects-due to spin-disorder, lattice polarons due to the 30% difference in the volume of the $`\mathrm{Mn}^{+3}`$ and the $`\mathrm{Mn}^{+4}`$ ions as well as extrinsic disorder acting in concert might be responsible for the resisitivity in the paramangetic phase. The possibility suggested that spin-disorder alone my be sufficient turns out, as shown by recent numerical studies, not to be correct. Additional randomness due to subsitution disorder has been used in calculations to explain the experimental data. Isoelectronic $`\mathrm{La}_{0.7\mathrm{x}}\mathrm{R}_\mathrm{x}\mathrm{Ca}_{0.3}\mathrm{MnO}_3`$ shows enormous descrease of the critical temperature when $`\mathrm{R}`$ is $`\mathrm{Y}`$ compared to when $`\mathrm{R}`$ is $`\mathrm{Pr}`$. Note that the ionic radius of $`\mathrm{La}^{3+}`$ is $`1.02\mathrm{A}`$, of $`\mathrm{Pr}^{3+}`$ is $`1.01\mathrm{A}`$, and of $`\mathrm{Y}^{3+}`$ is $`0.89\mathrm{A}`$. Note, that the substituion with $`\mathrm{Y}`$ besides changing the average bond angle introduces disorder.
Also very interesting is the fact that not only does spin-disorder disappear for $`TT_c`$, lattice disorder does as well. This is evidenced by the remarkable variation of the Debye-Waller factor with temperature below and above $`T_c`$. It is clear that spin and lattice disorder act in concert and quite unusual ways. Further that quenched lattice disorder generates extra lattice disorder which is annealed in the ferromagnetic phase.
If indeed the difference in the properties of the CMR materials is caused by the effect of the substitutional disorder, then it might be possible to account for the main features of the behavior of the โparamagnetic-metallicโ compounds using the โpureโ double-exchange model. To address this questions one of the main objectives of the present paper. We also consider the effect of the substitutional disorder, and show that it leads to a substantial decrease of the critical temperature of the para- to ferro-magnetic transition, in agreement with the observed difference in $`T_c`$ in different CMR materials. In a future paper we hope to address the more subtle issues connected with cation and other disorder in the mixed-valent compounds.
The paper is organized as follows. In the next section, we develop the variation mean field theory for the double-echange Hamiltonian. This is a systematization of the ansatz used in ref.. We calculate the spin distribution function, and the critical temperature of the ferromagnetic transtition. In the third section, we study the effect of the substitution disorder on this phase transition. In Section IV we develop a semiclassical transport theory for the CMR compounds, and calculate the magnetic field- and temperature- dependence of the resistivity. We close with a summary and discussion of future directions.
## II The Variational Mean Field Theory
In the semiclassical limit of large spin $`S`$ of the manganese ions, the effective electron Hamiltonian in the double-exchange model can be expressed as
$`H_{\mathrm{eff}}`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{ij}{}^{}}t_0\mathrm{cos}(\theta _{ij}/2)c_i^{}c_j+{\displaystyle \underset{i}{}}\left[v_ic_i^{}c_i\mu _BSB\mathrm{cos}\vartheta _i\right]`$ (1)
where the first sum includes hopping only between the nearest-neighbour manganese ions of different valencies, the angle $`\theta _{ij}`$ is defined as the angle between the ion spins $`๐_i`$ and $`๐_j`$, $`v_i`$ represents the effect of the substitutional disorder, $`B`$ is the magnetic field, and $`\mu _B`$ is the Bohr magneton. The angle $`\vartheta `$ is the angle between the spin $`๐_i`$ and the magnetic field. It is important to note that the assumption of large $`U`$ and $`J`$ compared to $`t`$ makes the charge carriers effectively spin-less.
Neglecting the correlations in the orientations of the neighbour spins, we represent the free energy $`F`$ of the system in terms of the single spin orientation distribution function $`P_\mathrm{\Omega }\left(๐\right)`$. In the mean-field approximation, the distribution function depends only on the angle $`\vartheta `$ between the local spin and the external magnetic field $`B`$:
$`P_\mathrm{\Omega }`$ $`=`$ $`{\displaystyle \frac{1}{2\pi }}P_\vartheta \left(\vartheta \right)`$ (2)
In the semiclassical limit, the corresponding spin entropy is then
$`S_{\mathrm{spins}}`$ $`=`$ $`{\displaystyle ๐\vartheta \mathrm{sin}\vartheta P_\vartheta \left(\vartheta \right)\mathrm{log}\left[P_\vartheta \left(\vartheta \right)\right]}+S_{\mathrm{spins}}^0\left(S\right)`$ (3)
where the function $`S_{\mathrm{spins}}^0\left(S\right)`$ does not depend on $`P_\vartheta `$, and is related to our choise of the normalization of the spin distribution function $`d(\mathrm{cos}\vartheta )P_\vartheta =1`$. This semiclassical approximation is discussed in detail in Appendix A.
The calculation of the energy for a given spin distribution is more complicated. When the transfer integral between the neighbouring sites $`i`$ and $`j`$ is equal to a constant value $`\stackrel{~}{t}`$, and the effects of the substitution disorder can be ignored, the electron energy is given by
$`E_t\left[\stackrel{~}{t}\right]`$ $`=`$ $`{\displaystyle _{\mathrm{}}^\mu }๐\epsilon \rho _0(\stackrel{~}{t};\epsilon )\epsilon `$ (4)
where $`\rho `$ is the electron density of states (DOS) corresponding to the Hamiltonian (1) with no diagonal disorder ($`v_i=0`$) and constant transfer intergral $`t_{ij}=t`$. To account for the effects of the substitution disorder $`v_i`$, we introduce an effective averaged DOS defined as
$`\rho (t,\epsilon )=\rho _0(t,\epsilon v)_v`$ (5)
where the average is performed over the distribution of the $`v_i`$โs.
To obtain the total energy, in the mean field approximation we average $`E_t[t]`$ over the distribution of the transfer integrals
$`E`$ $`=`$ $`{\displaystyle ๐tP_t\left(t\right)E_t\left[t\right]}`$ (6)
The transfer integral $`t_{ij}`$ can then be expressed in terms of the polar angles $`(\varphi _i,\vartheta _i)`$ and $`(\varphi _j,\vartheta _j)`$, which define the orientations of the corresponding spins, since they uniquely define the relative angle $`\theta `$. Therefore the integration over $`t`$ in Eq. (6) can be converted to the integration over the polar angles. Using the procedure discussed in detail in Appendix B, we derive the effective free energy functional, and by a direct minimization obtain the following integral equation for the spin distribution function:
$`P_\vartheta \left(\vartheta \right)`$ $`=`$ $`\mathrm{exp}\left[2{\displaystyle _0^{2\pi }}{\displaystyle \frac{d\varphi _1}{2\pi }}{\displaystyle _0^{2\pi }}{\displaystyle \frac{d\varphi _2}{2\pi }}{\displaystyle _0^\pi }๐\vartheta _1\mathrm{sin}\vartheta _1P_\vartheta \left(\vartheta _1\right){\displaystyle _{\mathrm{}}^\mu }๐\epsilon {\displaystyle \frac{\epsilon \mu }{T}}\rho (t_0\mathrm{cos}\left({\displaystyle \frac{\theta }{2}}\right);\epsilon ){\displaystyle \frac{\zeta }{T}}+{\displaystyle \frac{\mu _BSB}{T}}\mathrm{cos}\vartheta \right]`$ (7)
Here the parameter $`\zeta =\zeta (T,B)`$ accounts for the proper normalization of the distribution function $`P_\vartheta `$. The last term accounts for the energy $`\mu _BSB\mathrm{cos}\vartheta `$ of the spin, tilted at the angle $`\vartheta `$ with respect to the direction of the external magnetic field. Finally, the first term in the exponential of Eq. (7) represent the energy of the electron gas, which depends on the spin distribution via the effective โlocalโ bandwidth $`W\mathrm{cos}\theta /2`$, determined by the relative orientation of the near spins. Note, that this term depends nontrivially on $`\vartheta `$ via the relative angle $`\theta =\theta (\varphi _1,\vartheta _1;\varphi ,\vartheta )`$. This nonlinear integral equation allows a straightforward numerical solution by iterations. In Fig. 1 we plot the distribution $`P_\vartheta `$ for different values of the scaled temperature $`\tau T/t_0`$ and dimensionless magnetic field $`bB/t_0`$.
When the magnetization of the system is small, and the spin distribution is close to uniform (e.g. when the system is in the paramagnetic phase in a small external field), then the distribution function
$`P_\vartheta \left(\vartheta \right)`$ $`=`$ $`{\displaystyle \frac{1}{2}}+\delta p_\vartheta \left(\vartheta \right),\delta p_\vartheta 1`$ (8)
Expanding the exponential in the right hand side of Eq. (7) in $`\delta p_\vartheta `$, and keeping the terms up to the first order in $`\delta p_\vartheta `$, yields
$`\delta p_\vartheta `$ $`=`$ $`{\displaystyle \frac{3}{2}}M(T,B)\mathrm{cos}\vartheta +๐ช\left(M^2\right),`$ (9)
where $`M(T,B)=\chi \left(T\right)B+๐ช\left(B^2\right)`$ is the magnetization of the system. The susceptibility $`\chi `$ is then given by
$`\chi \left(T\right)={\displaystyle \frac{1}{3}}{\displaystyle \frac{\mu _B^2S^2}{TT_c}}`$ (10)
where the critical temperature $`T_c`$ is given by
$`T_c`$ $`=`$ $`{\displaystyle _0^\pi }๐\vartheta \mathrm{sin}\vartheta \mathrm{cos}\vartheta {\displaystyle _{\mathrm{}}^\mu }๐\epsilon \left(\mu \epsilon \right)\rho (t_0\mathrm{cos}{\displaystyle \frac{\vartheta }{2}};\epsilon )`$ (11)
It is worthwhile noting the relation of this theory to an earlier mean-field variational description of the manganese oxides of Ref. . There, in addition to the mean-field approximation, a specific functional form of the probability distribution of the angle between different spins was assumed, with the system magnetization being the variational parameter. This should be contrasted to the method of the present paper, when the functional form of the single-spin distribution function is derived variationally. The general dependence of the distribution derived here turns out to be quite similar to the one assumed earlier. But the results of the variational procedure developed in this section should be more accurate besides being on firmer ground. Another advantage of the present method is that it can be used for the description of the effects of the substitution disorder - something, which is hard to characterize within the framework of Ref. .
As follows from Eq. (11), the critical temperature explicitly depends on the density of states, and the resulting value is in fact sensitive to the actual shape of DOS. However, this has only a marginal effect on the dependence of $`T_c`$ on the concentration $`x`$. To illustrate this behaviour, in Fig. 2 we plot the critical temperature as a function of the charge carrier concentration $`x`$ for a rectangular (blue curve) and Gassian (red curve) densities of states. For comparison, we also plot the $`x(1x)`$ dependence (black line), obtained in an earlier work, and the experimental data of Ref. . The model densities of states are plotted in Fig. 3 and compared to the DOS $`\rho _t`$, corresonding to the hamiltonian (1) with constant transfer integral, and no diagonal disorder. The effective bandwidth of the model desities of states is chosen to accurately reproduce the second moment $`\epsilon ^2`$. Note, how accurately the Gaussian density of states fits the profile of $`\rho _t`$.
As follows from Fig. 2, a reasonable choise of the bandwidth $`W=1.8\mathrm{eV}`$ consistent with the calculations in the local density approximation, leads to a good agreement with the experimental data.
As explained earlier, the transition temperature is determined essentially by the difference in the cohesive energy of the ferromagnet and the paramagnet by the entropy of the paramagnet. The larger bandwidth of the ferromagnet by about 20% is the essential aspect of the energetics in the double-exchange problem.
Consider now the effect of substitutional disorder. Substitutional disorder increases the electron-bandwidth for the paramagnet. The removal of spin-disorder is then expected to decrease the change in the bandwidth on becoming a ferromagnet. This is explicitly borne out by the theory here.
Since the critical temperature is directly related to the effective DOS, it is sensitive to the substitution disorder in the system. Assuming the Gaussian distribution of the disorder strength $`v_i`$ with the standard deviation $`V_0`$ and Gaussian โbareโ DOS $`\rho _g\mathrm{exp}\left[\epsilon ^2/\left(3t^2\right)\right]`$, we obtain:
$`T_c`$ $`=`$ $`{\displaystyle _0^\pi }๐\vartheta {\displaystyle _{\mathrm{}}^\mu }๐\epsilon \left(\mu \epsilon \right)\rho [t_{\mathrm{eff}}\left(\vartheta \right),\epsilon ]\mathrm{sin}\vartheta \mathrm{cos}\vartheta `$ (12)
where the effective transfer integral $`t_{\mathrm{eff}}`$ is defined by the equation
$`{\displaystyle \frac{1}{t_{\mathrm{eff}}^2}}`$ $`=`$ $`{\displaystyle \frac{1}{t_0^2\mathrm{cos}^2\frac{\vartheta }{2}}}+{\displaystyle \frac{3}{V_0^2}}`$ (13)
The dependence (12) is shown in Fig. 4 for different electron concentrations. As extra disorder makes the ferromagnetic phase less favourable, the critical temperature goes down with an increase of $`V_0`$.
It might be tempting to attribute the difference in critical temperature between the โtype-Iโ and โtype-IIโ compounds to the effect of the substitutional disorder. In such model, an effective disorder strength of $`V_00.7W`$, would fully account for only $`30\%`$ difference in the critical temperatures of โtype-Iโ and โtype-IIโ compounds. We suspect that large enough lattice disorder, in concert with lattice disorder localises electronic states in the paramagnetic phase. New considerations then eneter in to the determination of the transition temperature. These will be discussed separately. Also missing from the discussion above is the effect of the formation of spin-polarons which must occur in the paramagnetic phase . They would tend to decrease $`T_c`$ but the number of spins in the polarons is rather small and only a modest numerical effect on the transition temperature is expected. They are however quite important for the dynamics near the transition.
## III Resistivity without lattice disorder: Semiclassical Treatment
In the mean-field approximation developed in the previous section, each spin independently fluctuates around the averaged value defined by the magnetization of the system. From the point of view of the semiclassical transport theory, that would correspond to effective independent โscatterersโ located at each point of the lattice. However, in the ferromagnetic phase, when the spin fluctuations are small compared to the averaged value, the corresponding electron mean free path may be substantially larger that the (Mn) lattice spacing. In this limit, in order to estimate the resistivity of the system, we can use the standard semiclassical transport theory.
We introduce the average transfer integral $`\overline{t}t_0\mathrm{cos}\left(\theta _{\alpha \beta }/2\right)`$, so that the the corresponding unperturbed Hamiltonian is defined as
$`H_0`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{<\alpha \beta >}{}}\overline{t}c_\alpha ^+c_\beta `$ (14)
and rest of $`H`$ is treated as the โperturbationโ
$`V`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{<ij>}{}}\delta t_{ij}c_i^+c_j`$ (15)
The standard plane-wave diagonalization of $`H_0`$ yields the dispersion law
$`ฯต_๐ค=\overline{t}\left[\mathrm{cos}\left(k_xa\right)+\mathrm{cos}\left(k_ya\right)+\mathrm{cos}\left(k_za\right)\right]`$ (16)
which describes โholesโ near $`๐ค=0`$ and โelectronsโ near e.g. $`๐ค=\pi /a(1,1,1)`$. When the Fermi energy is located near the bottom or near the top of the band, one can define the effective mass for the electrons and the holes respectively, $`m_{}=\frac{2\overline{t}a^2}{\mathrm{}^2}`$.
The kinetic equation for the electron distribution function $`f_๐ค`$ is , :
$`e๐{\displaystyle \frac{f^0}{๐ค}}={\displaystyle \frac{2\pi }{\mathrm{}}}{\displaystyle \underset{๐ค}{}}\left|๐ค\left|V\right|๐ค^{}\right|^2\delta \left(ฯต_๐คฯต_๐ค^{}\right)\left(f_๐คf_๐ค^{}\right)`$ (17)
where $`f^0(ฯต_๐ค)`$ is the equilibrium distribution function, and $`๐ค\left|V\right|๐ค^{}`$ is the matrix element of the โperturbationโ (15).
The kinetic equation (17) has the solution:
$`f_๐ค=f^0\left(ฯต_๐ค\right)eE\tau \left(๐ค\right){\displaystyle \frac{ฯต_๐ค}{k_z}}{\displaystyle \frac{f^0\left(ฯต_๐ค\right)}{ฯต_๐ค}}`$ (18)
where the relaxation time is defined by the following equation:
$`{\displaystyle \frac{1}{\tau \left(๐ค\right)}}`$ $`=`$ $`{\displaystyle \frac{3a^3}{2\pi ^2\mathrm{}}}\overline{\delta t^2}{\displaystyle ๐๐ค^{}\left(1+\frac{ฯต_{๐ค+๐ค^{}}}{6\overline{t}}\right)\delta \left(ฯต_๐คฯต_๐ค^{}\right)}`$ (19)
$``$ $`{\displaystyle \frac{a^3}{2\pi \mathrm{}}}\overline{\delta t^2}{\displaystyle ๐๐ค^{}\mathrm{sin}^2\left(k_z^{}a\right)\delta \left(ฯต_๐คฯต_๐ค^{}\right)}`$ (20)
Assuming a uniform dispersion $`ฯต_๐ค=ฯต\left(\left|๐ค\right|\right)`$, this expression reduces to the standard result for the transport relaxation time
$`{\displaystyle \frac{1}{\tau \left(k\right)}}`$ $`=`$ $`{\displaystyle ๐๐ค^{}W_{\mathrm{๐ค๐ค}^{}}\left(1\mathrm{cos}\theta _{\mathrm{๐ค๐ค}^{}}\right)\delta \left(ฯต_๐คฯต_๐ค^{}\right)}`$ (21)
where the scattering rate $`W_{\mathrm{๐ค๐ค}^{}}=3a^3\left(2\pi ^2\mathrm{}\right)^1\overline{\delta t^2}\left(1+ฯต_{๐ค+๐ค^{}}/\left(6\overline{t}\right)\right)`$, and $`\theta _{\mathrm{๐ค๐ค}^{}}`$ is the angle between the vectors $`๐ค`$ and $`๐ค^{}`$.
Near the top and the bottom of the band, the integrals in (20) allow a straightforward analytical evaluation. For example, for the holes we obtain:
$`\tau ^1\left(๐ค\right)`$ $`=`$ $`{\displaystyle \frac{6}{\pi \mathrm{}}}ka\left(1{\displaystyle \frac{2}{9}}\left(ka\right)^2\right){\displaystyle \frac{\overline{\delta t^2}}{\overline{t}}}`$ (22)
As we pointed out before, the semiclassical approach developed in the present section is appropriate only when the charge carrier mean free path $`\mathrm{}a`$. Using (22), for the ratio of the mean free path to the lattice spacing near the top of the band we obtain:
$`{\displaystyle \frac{\mathrm{}}{a}}`$ $`=`$ $`{\displaystyle \frac{\overline{t}^2}{\overline{\delta t^2}}}{\displaystyle \frac{3}{\pi \left(1\frac{2}{9}\left(ka\right)^2\right)}}`$ (23)
As follows from Eq. (20),(23), in the absence of substitution disorder, $`\mathrm{}/a`$ is always greater than $`(3/\pi )(\overline{t}^2/\overline{\delta t^2})`$. The ratio $`\overline{t}^2/\overline{\delta t^2}`$ is a monotocially decreasing function of temperature in the ferromagnetic phase, and constant above the $`T_c`$, where $`\overline{t}^2/\overline{\delta t^2}=8`$. Therefore, since the mean free path due to the spin disorder is substantially larger than the effective lattice spacing, we expect that in the relevant concentration range $`x0.10.3`$ such a โpure DEโ system would generally show the metallic behavior . Indeed, in a typical โtype-IIโ compound $`\mathrm{La}_{0.7}\mathrm{Sr}_{0.3}\mathrm{MnO}_3`$, the resistivity does show the metallic behaviour $`d\rho /dT>0`$ both above and below the transition.
The effect of nonzero magnetization (caused either by the transtion to the ferromagnetic phase, or by external magnetic field) on the conductivity is twofold: first, it suppresses the fluctuations in transfer integrals thus decreasing the corresponding scattering rate; second, the increase of the average transfer integral caused by the magnetization leads to a decrease of the effective mass $`m_{}1/\overline{t}`$. Both these factors lead to a decrease of the resistivity $`\rho `$. For a small magnetization,
$`\rho (M)=\rho _0(1\kappa \left(M/M_{\mathrm{max}}\right)^2)`$ (24)
where in the effective mass approximation the coefficient $`\kappa `$ is equal to $`9/5`$. For weakly disordered manganites, the resitivity indeed follows Eq. (23). As seen in the inset in Fig.5(b) the experimental value is about 2 and slowly varies with the electron density. Taking into account the band nonparabolicity leads to a weak dependence of $`\kappa `$ on the concentration $`x`$, but does not fully account for an increase of $`\kappa `$. The variation of the resistivity in the whole range of the sample magnetization $`0<M<M_{\mathrm{max}}`$ is shown in Fig. 5.
## IV Substitution disorder: the effect on resistivity
As we pointed out in the section II of the present paper, the $`30\%`$ difference in the critical temperatures of the โtype-Iโ and โtype-IIโ compounds implies, that in the โdisorderedโ compounds the effective scattering potential is of the order of the electron bandwidth. In such conditions, the localization effects can become important, and the semiclassical treatment of the previous section is no longer appropriate.
It has been proposed, that in the โstrongly disorderedโ (โtype-Iโ) compounds, the ferro- to paramagnetic phase transition drives the metal-insulator transition. In the paramagnetic phase, the โcombined effortโ of the substitution and spin disorder is sufficient to localize the charge carriers, while in the ferromagnetic phase, due to larger electron bandwidth and weaker spin disorder, the mobility edge is below the Fermi energy.
One might be tempted to think that this mechanism of the colossal magnetoresistance of the โdisorderedโ manganites reduces the problem to an Anderson-type transition as a function of disorder alone, where the spin disorder is a function of the magnetization. This is not correct since the magnetic entropy is essential to the transition which occurs at a finite temperature unlike the Anderson transition which occurs at $`T=0`$.
An important question however is wheather the resistivity near the transition can be expressed uniquely as a function of magnetization. If the phase transition (with or without โdiagonalโ disorder) is characterized by a divergent magnetic correlation length scale, this may be possible. It should be remembered however that resistivity depends on fluctuations at large momentum transfer. In a solid with lattice disorder, the ferromagnetic correlation length does not uniquely characterize the important disorder at short length scales even though it may be coupled to the magnetization as appears to be the case in the manganites. It is also possible that for sufficiently strong disorder, the ferromagnetic transition is replaced by a cross-over and there is no divergent correlation length. These are probably the reasons why no clear indication of scaling behavior expected in a continuous quantum phase transition were found in the recent esistivity measurements. These considerations, however, go beyond the mean-field type theory, developed in the present paper.
In any case, we believe that in order to fully understand the physics underlying the colossal magnetoresistance in doped mangnites, one also has to take into account the lattice disorder, which is coupled to the spin disorder, via their influence on the charge carriers. Indeed, the strong coupling between the spin and lattice disorder was recently demonstrated in two independent experiments. Clearly, the effective lattice disorder has itโs own nontrivial temperature dependence, and, being coupled to the charge carriers, therefore obviously leads to substantial deviations from the standard picture of the โstaticโ Anderson metal-insulator transition. However, at this point we defer the further description of this effect.
## V Conclusions
As shown in the first part of the present paper, the phase transition from paramagnetic to ferromagnetic phase in relatively pure manganese oxides can be successfully described by a variational approximation on the double-exchange Hamiltonian. The results obtained for the critical temperature and its evolution with doping and the chemical composition of the compound are consistent with the experimental data. The decrease of $`T_c`$ with modest disorder is also understood.
The resistivity of the โtype-IIโ manganese compounds can also be successfully described using the DE model. We showed, that e.g. the theoretical dependence of the resistivity on the sample magnetization is in a quantitative agreement with the expetimental data. Our calculations show the robustness of the results to the particular choice of the electron density of states, as should be obvious since the transition temperature depends on the difference of the cohesive energy of the paramagnetic and the ferromagnetic phases.
The principal problems of the manganites left unanswered in this paper concern the properties of the โtype-Iโ compounds and the remarkable effects of disorder in them in both the dynamic and static properties. These are also the more subtle problems. Especially interesting is the fact that extrinsic disorder appears to promote some additional disorder in the paramagnetic phase which is swept away togather with the spin-disorder in the ferromagnetic phase. We hope to provide an answer to these questions separately.
## A The Spin Entropy in the Semiclassical Mean-Field Approximation
In the mean field approximation, when the spin density matrix of the whole system $`\rho ^\mathrm{\Sigma }\left(S_z^{(i)}\right)`$ is represented as a product of diagonal density matrices $`\rho _i=\rho ^{(1)}\left(S_z^{(i)}\right)`$ of the individual spins
$`\rho ^\mathrm{\Sigma }`$ $`=`$ $`\mathrm{\Pi }_{i=1}^N\rho _i`$ (A1)
Then the total spin entropy
$`S_{\mathrm{spins}}^\mathrm{\Sigma }`$ $`=`$ $`\mathrm{Tr}\left\{\rho ^\mathrm{\Sigma }\mathrm{log}\left[\rho ^\mathrm{\Sigma }\right]\right\}`$ (A2)
is represented by
$`S_{\mathrm{spins}}=N{\displaystyle \underset{S_z=S}{\overset{S}{}}}\rho ^{(1)}\left(S_z\right)\mathrm{log}\left[\rho ^{(1)}\left(S_z\right)\right]`$ (A3)
In the semiclassical approximation $`S1`$, the summation over $`S_z`$ can be replaced by integration. Introducing new variable $`\vartheta \mathrm{arccos}(S_z/S)`$, we obtain:
$`S_{\mathrm{spins}}=NS{\displaystyle _1^1}d\mathrm{cos}\vartheta \rho ^{(1)}\left(S\mathrm{cos}\vartheta \right)\mathrm{log}\left[\rho ^{(1)}\left(S\mathrm{cos}\vartheta \right)\right]`$ (A4)
where $`\rho ^{(1)}\left(S\mathrm{cos}\vartheta \right)`$ is normalized as follows:
$`1`$ $`=`$ $`{\displaystyle \underset{S_z=S}{\overset{S}{}}}\rho ^{(1)}\left(S_z\right)=S{\displaystyle _1^1}d\mathrm{cos}\vartheta \rho ^{(1)}\left(S\mathrm{cos}\vartheta \right)`$ (A5)
We now define the spin orientation distribution function $`P_\vartheta \rho ^{(1)}\left(S\mathrm{cos}\vartheta \right)`$, normalized as
$`{\displaystyle _1^1}d\mathrm{cos}\vartheta \rho ^{(1)}\left(S\mathrm{cos}\vartheta \right)`$ $`=`$ $`1`$ (A6)
As follows from Eqns. (A5), (A6), the spin orientation distribution function
$`P_\vartheta `$ $`=`$ $`{\displaystyle \frac{1}{S}}\rho ^{(1)}\left(S\mathrm{cos}\vartheta \right)`$ (A7)
Therefore, the semiclassical spin entropy
$`S_{\mathrm{spins}}=N{\displaystyle _1^1}d\mathrm{cos}\vartheta P_\vartheta \mathrm{log}\left[P_\vartheta \right]+N\mathrm{log}\left[S\right]`$ (A8)
For example, in the paramgagnetic phase, when there are no external fields, and the spind orienataion distribution is uniform, $`P_\vartheta =1/2`$, the semiclassical spin entropy is equal to $`N\mathrm{log}(2S)`$, which is consistent with the exact result $`N\mathrm{log}(2S+1)`$ for $`S1`$. Note, that the main contribution to the semiclassical spin entropy comes actually from the distribution-independent term in Eq. (A8).
The semiclassical description, however, fails for large magnetization, when the spin system is almost completely polarized, and the distribution function starts to change substantially on the scale of $`\delta \vartheta 1/S`$. In this case, the original expression, Eq. (A3), should be used for the calculation of the spin entropy.
## B The Variational Free Energy Functional
In the present Apendix, we calculate the variational free energy functional for the double-exchange model. Using Eqns. (4),(6), for the electron energy we obtain:
$`E_e`$ $`=`$ $`{\displaystyle _0^{2\pi }}{\displaystyle \frac{d\varphi _1}{2\pi }}{\displaystyle _0^{2\pi }}{\displaystyle \frac{d\varphi _2}{2\pi }}{\displaystyle _0^\pi }๐\vartheta _1\mathrm{sin}\vartheta _1{\displaystyle _0^\pi }๐\vartheta _2\mathrm{sin}\vartheta _2P_\vartheta \left(\vartheta _1\right)P_\vartheta \left(\vartheta _2\right){\displaystyle _{\mathrm{}}^\mu }๐\epsilon \epsilon \rho (t_0\mathrm{cos}\left({\displaystyle \frac{\theta (\varphi _1,\vartheta _1;\varphi _2,\vartheta _2)}{2}}\right);\epsilon )`$ (B1)
while the extra spin energy
$`E_s=B{\displaystyle _0^{2\pi }}{\displaystyle \frac{d\varphi _1}{2\pi }}{\displaystyle _0^\pi }๐\vartheta _1\mathrm{sin}\vartheta _1\mathrm{cos}\vartheta _1P_\vartheta \left(\vartheta _1\right)`$ (B2)
The free energy can be obtained by the subsituting these expressions and the entropy (3) into the standard definition of the free energy
$`F`$ $`=`$ $`E_e+E_sTS`$ (B3)
In order to find the single-spin distribution $`P_\vartheta `$, one has to minimize the effective free energy, taking into account the constraints of normalization. Using the standard Lagrange multiplier method, for the effective free energy functional we obtain:
$`\stackrel{~}{F}[P_\vartheta ;\mu ,\lambda ,\zeta ]`$ $`=`$ $`{\displaystyle _0^{2\pi }}{\displaystyle \frac{d\varphi _1}{2\pi }}{\displaystyle _0^{2\pi }}{\displaystyle \frac{d\varphi _2}{2\pi }}{\displaystyle _0^\pi }d\vartheta _1\mathrm{sin}\vartheta _1{\displaystyle _0^\pi }d\vartheta _2\mathrm{sin}\vartheta _2P_\vartheta \left(\vartheta _1\right)P_\vartheta \left(\vartheta _2\right){\displaystyle _{\mathrm{}}^\mu }d\epsilon (\epsilon \lambda )\rho (t_0\mathrm{cos}\left({\displaystyle \frac{\theta }{2}}\right));\epsilon )`$ (B4)
$`+`$ $`T{\displaystyle _0^{2\pi }}{\displaystyle \frac{d\varphi _1}{2\pi }}{\displaystyle _0^\pi }๐\vartheta _1\mathrm{sin}\vartheta _1P_\vartheta \left(\vartheta _1\right)\mathrm{log}\left[P_\vartheta \left(\vartheta _1\right)\right]B{\displaystyle _0^{2\pi }}{\displaystyle \frac{d\varphi _1}{2\pi }}{\displaystyle _0^\pi }๐\vartheta _1\mathrm{sin}\vartheta _1\mathrm{cos}\vartheta _1P_\vartheta \left(\vartheta _1\right)`$ (B5)
$`+`$ $`\zeta {\displaystyle _0^\pi }๐\vartheta \mathrm{sin}\vartheta p_\vartheta \left(\vartheta \right)+\overline{F}[x,\lambda ,\zeta ]`$ (B6)
where the โconstantโ $`\overline{F}`$ represents the spin distriibution-independent part of the free energy. Here, the Lagrange multiplier $`\zeta `$ accounts for the normalization of the distribution function $`P_\vartheta `$, while the Lagrange multiplier $`\lambda `$ represents the constraint of having a fixed concentration of mobile electrons in the system.
It is straightforward to show by a direct calculation, that at the extremum of the functional (B6) $`\lambda =\mu `$. This has a clear physical meaning - the Lagrange multiplier $`\lambda `$ corresponds to the electron number conservation, and therefore shoul be equal to the electron electrochemical potential. Replacing $`\lambda `$ by $`\mu `$ in (B6), we finally obtain the effective mean field free energy functional:
$`\stackrel{~}{F}[P_\vartheta ;\mu ,\lambda ,\zeta ]`$ $`=`$ $`{\displaystyle _0^{2\pi }}{\displaystyle \frac{d\varphi _1}{2\pi }}{\displaystyle _0^{2\pi }}{\displaystyle \frac{d\varphi _2}{2\pi }}{\displaystyle _0^\pi }๐\vartheta _1\mathrm{sin}\vartheta _1{\displaystyle _0^\pi }๐\vartheta _2\mathrm{sin}\vartheta _2P_\vartheta \left(\vartheta _1\right)P_\vartheta \left(\vartheta _2\right){\displaystyle _{\mathrm{}}^\mu }๐\epsilon \left(\epsilon \mu \right)\rho (t_0\mathrm{cos}\left({\displaystyle \frac{\theta }{2}}\right);\epsilon )`$ (B7)
$`+`$ $`T{\displaystyle _0^{2\pi }}{\displaystyle \frac{d\varphi _1}{2\pi }}{\displaystyle _0^\pi }๐\vartheta _1\mathrm{sin}\vartheta _1P_\vartheta \left(\vartheta _1\right)\mathrm{log}\left[P_\vartheta \left(\vartheta _1\right)\right]B{\displaystyle _0^{2\pi }}{\displaystyle \frac{d\varphi _1}{2\pi }}{\displaystyle _0^\pi }๐\vartheta _1\mathrm{sin}\vartheta _1\mathrm{cos}\vartheta _1P_\vartheta \left(\vartheta _1\right)`$ (B8)
$`+`$ $`\zeta {\displaystyle _0^\pi }๐\vartheta \mathrm{sin}\vartheta p_\vartheta \left(\vartheta \right)+\overline{F}`$ (B9)
Taking the functional derivative of (B9) with respect to $`P_\vartheta `$, for the distribution we obtain:
$`P_\vartheta \left(\vartheta \right)`$ $`=`$ $`\mathrm{exp}\left[2{\displaystyle _0^{2\pi }}{\displaystyle \frac{d\varphi _1}{2\pi }}{\displaystyle _0^{2\pi }}{\displaystyle \frac{d\varphi _2}{2\pi }}{\displaystyle _0^\pi }๐\vartheta _1\mathrm{sin}\vartheta _1P_\vartheta \left(\vartheta _1\right){\displaystyle _{\mathrm{}}^\mu }๐\epsilon {\displaystyle \frac{\epsilon \mu }{T}}\rho (t_0\mathrm{cos}\left({\displaystyle \frac{\theta }{2}}\right);\epsilon ){\displaystyle \frac{\zeta }{T}}+{\displaystyle \frac{B}{T}}\mathrm{cos}\vartheta \right]`$ (B10)
Note, that the exponential on the right hand side nontrivially depends on $`\vartheta `$ via the angle $`\theta =\theta (\varphi _1,\vartheta _1;\varphi ,\vartheta )`$. |
warning/0002/quant-ph0002021.html | ar5iv | text | # On balance of information in bipartite quantum communication systems: entanglement-energy analogy
## I Introduction
It is astonishing that just lately after over sixteen years quantum formalizm reveals us new possibilities due to entanglement processing being a root of new quantum phenomena such as quantum cryptography with Bell theorem , quantum dense coding , quantum teleportation , quantum computation . It shows how important is to recognize not only the structure of the formalizm itself but also potential possibilites encoded in it.
In spite of many beautiful experimental and theoretical results on entanglement there are still difficulties in understanding its many faces. It seems to be a reflection of basic difficulties in the interpretation of quantum formalizm as well as quantum-clasical hybridism in our perception of Nature. To overcome the latter the existence of unitary information field being a necessary condition of any communication (or correlation) has been postulated as well as the information interpretation of quantum wavefunction has been considered . It rests on the generic information paradigm according to which the notion of information represents a basic category and it can be defined independently of probability itself . It implies that Nature is unbroken entity. However, according to double, hylemorphic nature of the unitary information field, there are two mutually coupled levels of physical reality in Nature: logical (informational) due to potential field of alternatives and energetic due the field of activities (events) . Then from the point of view of the generic information paradigm, quantum formalism is simply a set of extremely useful informational algorithms describing the above complementary aspects of the same, really existing unitary information field. It leads in a natural way to analogy between information (entanglement) and energy being nothing but a reflection of unity of Nature.
Following this route, one attempts to find some useful analogies in the quantum communication domain. Namely, physicists believe that there should exist the laws governing entanglement processing in quantum communication systems, that are analogous to those in thermodynamics.
Short history of this view has its origin in the papers by Bennett et al. who announced a possible irreversibility of the entanglement distillation process . Popescu and Rohrlich have pointed out analogy between distillation-formation of pure entangled states and Carnot cycle, and they have shown that entanglement is extensive quantity. The authors formulated principle of entanglement processing analogous to the second principle of thermodynamics: โEntanglement cannot increase under local quantum operations and classical communicationโ. Vedral and Plenio have considered the principle in detail and pointed out that there is some (although not complete) analogy between efficiency of distillation and efficiency of Carnot cycle. In Refs entanglement-energy analogy has been developed and conservation of information in closed quantum systems has been postulated in analogy with the first principle of thermodynamics: Entanglement of compound system does not change under unitary processes on one of the subsystems. Then an attempt to formulate the counterpart of the second principle in a way consistent with the above principle has been done (since in the original Popescu-Rohrlich formulation entanglement was not conserved).
The main purpose of the paper is to support entanglement-energy analogy by demonstration that in the closed bipartite quantum communication system the information is conserved. The paper is organised as follows. In section II we describe closed quantum communication bipartite system. The next section contains formal description of balance of quantum information involving notions of physical and logical work. In section IV we introduce the concept of useful logical work in quantum communication. In next section we present balance of information in teleportation. In section VI we discuss entanglement-analogy in the context of the Gibbs-Hemholtz-like equation connecting entanglement of formation, distillable entanglement and bound entanglement. In section VII we present the balance of information in the process of distillation. In the last section we discuss the objectivity of quantum information in context of information interpretation of quantum states and alghoritmic complexisity..
## II Closed quantum communication system: the model
Consider closed quantum communication (QC) system $`U`$ composite of system $`S`$, measuring system $`M`$ and environment $`R`$
$$U=S+M+R$$
(1)
where each system is split into Alice and Bob parts $`S_X,M_X,R_X`$; $`X=A,B`$.
It is assumed that Alice and Bob can control the system $`S_X`$ which does not interact with environment $`R_X`$. The $`M_X`$ system consists of $`m_X`$ qubits and cotinuously interacts with environment $`R_X`$. In result the system $`M_X`$, palying the role of โancillaโ, is measured in distinguished basis $`|x_1x_2\mathrm{}`$, $`x_i=0,1`$ . In this sense the measurement is understand here as the process of irreversible entanglement with some environment and the system $`R_X`$ is to ensure this ireversibility. Note that in the above approach the evolution of the system is unitary i. e. abandon the von Neumann projection postulate which leads to violation of energy-momentum conservation . Then acting on one part of entangled system, we have no way to annihilate entanglement. The latter can change only by means of interacting of the both entangled subsystems. It may be objected that we can destroy entanglement e.g. by randomizing the relative phases on the subsystems of interest. However, if the reduction of wave packet is not regarded to be a real physical process, then the above operation must be considered as entangling the subsystem with some other system by means of a unitary transformation. Then the entanglement will not vanish but it will spread over all the three subsystems.
The operations Alice and Bob can perform in our QC system are:
* Quantum communication: Alice and Bob can exchange particles from the system $`S_X`$.
* โClassical communicationโ: Alice and Bob can exchange particles from the system $`M_X`$
Note that the number of qubits of the systems $`S_A`$ and $`S_B`$ can change but the total number of qubits of the system $`M`$ is conserved (similarly for $`S`$). Besides Alice and Bob can perform unitary transformation over the system $`M_X+S_X`$, $`X=A,B`$.
We would like to stress one more that in our approach the measurement represents an irreversible entanglement rather than the โprojectionโ of the state. To see it consider the case when Alice and Bob share a singlet state and Alice performs a measurement on it. The the initial state of the system $`M_A+S_A+S_B`$ ($`M_A`$ represents the Aliceโs ancilla while $`S_A`$, $`S_B`$ correspond to the particles forming a singlet state) is
$`|\mathrm{\Psi }_{M_AS_AS_B}=|0_{M_A}|\mathrm{\Psi }_A^{singlet}=`$ (2)
$`|0_{M_A}{\displaystyle \frac{1}{\sqrt{2}}}(|0_{S_A}|1_{S_B}|1_{S_A}|0_{S_B})`$ (3)
Then Alice performs the unitary operation $`U`$ on subsystem $`M_A+S_A`$. This operation corresponds to the interaction between $`M_A`$ and $`S_A`$ and can be represented by C-NOT gate. As a result the whole system is in the state
$$|\mathrm{\Psi }^{}_{M_AS_AS_B}=\frac{1}{\sqrt{2}}(|0_{M_A}|0_{S_A}|1_{S_B}|1_{M_A}|1_{S_A}|0_{S_B})$$
(4)
Further $`M_A`$ can be irreversibly entangled with environment system $`R_A`$ (which models the irreversibility of the measurement). But $`R_A`$ is still on Alice side, hence we have entanglement between systems $`(R_A+M_A+S_A)`$ and $`S_B`$ unchanged and equal to $`E=1`$ e-bit.
Of course, there are some interpretational problems if one imagines that Alice โreads outโ the result of the measurement as then we encounter problems coming from possible extension of the model by the projection postulate. However that for practical reasons (i. e. as far as quantum information qualitative description is concerned) the informational processes like e. g. quantum teleportation do not require reading the data. Moreover, it must be noted that at the absence of the projection postulate the above model can be veiwed as consistent with โmany worldsโ interpretation .
## III Conservation of quantum information: formal description
To determine balance of information in the closed system $`U`$ we adopt two basic postulates
* Entanglement is a form of quantum information corresponding to internal energy.
* Sending qubits corresponds to work.
In accordance with the postulate 1, the information is physical quantity that, in particular, should be conserved in closed quantum systems, similarly as energy. The second postulate allows to deal with communication processes (in thermodynamics work is a functional of process). To obtain the balance we must define our โenergyโ and โworkโ quantitatively. To this end consider system $`X`$ described in the Hilbert space $``$, $`\mathrm{dim}=d`$ being in a state $`\varrho _X`$. We define informational content $`I_X`$ of the state $`\varrho _X`$ as follows (cf. ):
$$I_X=\mathrm{log}\mathrm{dim}S(\varrho _X)$$
(5)
where $`\mathrm{dim}=d`$, $`S(\varrho _X)S(X)`$ are the dimension of the Hilbert space and the von Neumann entropy of the system state. Note that $`I_X`$ satisfies the inequality $`0=I_X^{min}I_XI_XI_X^{max}=\mathrm{log}\mathrm{dim}`$ where $`I_X^{min}`$ and $`I_X^{max}`$ are the information content of the maximal mixed state and pure state respectively. Thus it is well defined quantity which measures informational content of the system $`\varrho _X`$.
The formula (5) needs some comment as usually one interprets the von Neumann entropy as a measure of information. In fact there is no contradiction. Imagine for a moment that we admitt projection postulate i. e. Alice knows the concrete result of the measurement. Then the von Neumann entropy measures the information gain after the measurement while the formula (5) corresponds to the information prior the measurement and this information, in particular, is maximal if the system is in pure state. This is the reason while we use the name informational content as it has actual rather than potential (i. e. related to the future measurement) character. Below we shall see that, after we abandon the projection postulate, the above formula allows to perform a balance of quantum information in a consistent way. Note that the Hilbert space dimension used in formula (5) is present also in definitions of other notions (see below), in particular in the case of useful logical work (sec. IV). It plays, to some extent, the role similar to the one in channels capacities theory or error correction codes where dimension of โerror freeโ subspace is a central notion.
Consider now the QC system $`U`$, being in the initial pure state $`\psi _{in}`$, described by general Alice-Bob Hilbert space scheme as follows
$$\begin{array}{c}_A\\ \\ _A^{}\end{array}_B\}\psi _{in},$$
(6)
where $`_A_A^{}`$, $`_B`$ are the Hilbert spaces of the $`S_A+M_A+R_A`$ and $`S_B+M_B+R_A`$ respectively. Then in accordance with (5) the information contents of the Alice and Bob subsystems are defined as follows
$$I_A=\mathrm{log}\mathrm{dim}(_A_A^{})S(A+A^{});$$
(7)
$$I_B=\mathrm{log}\mathrm{dim}_BS(B),$$
(8)
where $`\mathrm{dim}(_A_A^{})`$ and $`\mathrm{dim}_B`$ are the dimensions the corresponding Hilbert spaces while $`S(A+A^{})`$, $`S(B)`$ are the von Neumann entropies of the subsystems.
Now, after transmission of the system $`A^{}`$ to receiver (Bob) the Alice-Bob Hilbert space scheme is given by
$$_A\begin{array}{c}_B\\ \\ _A^{}\end{array}\}\psi _{out}$$
(9)
and the total sytem $`U`$ is in the final state $`\psi _{out}`$.
Now, in accordance with the above โsending qubits โ workโ postulate we consider physical work performed over the system $`U`$ being a physical transmission of particles. Consequently, we define $`W_p`$ as a number of sent qubits of the system $`A^{}`$
$$W_p=\mathrm{log}\mathrm{dim}_A^{}.$$
(10)
Note that after transmission of the system $`A^{}`$ to the Bob, there is increase of the information content of his subsystem. Then we say that the system $`U`$ performed the logical work $`W_l`$ that is defined as increase of the informational content of the Bob (in general - receiver) system.
$$W_l=I_{out}^BI_{in}^B$$
(11)
where $`I_{in}^B=I^B`$, $`I_{out}^B=I^{B+A^{}}`$. Then one can regard the physical work as sending โmatterโ while the logical work โ sending โformโ that is consistent with the assumed hylemorphic nature of the information field. Subsequently we can define initial and final entanglement of the system $`U`$ as
$$E_{in}=S(B)=S(A+A^{});E_{out}=S(A)=S(B+A^{}),$$
(12)
where obvious relations between the entropies of the subsystems hold. Now, in accordace with the first postalate, $`E_{in}`$ and $`E_{out}`$ are simply initial and final potential informations contained in the total system. Having so defined quantities it is not hard to obtain the following information balance equations
$$E_{in}+W_p=E_{out}+W_l$$
(13)
or equivalently
$$I_{in}^A+I_{in}^B+2E_{in}=I_{out}^A+I_{out}^B+2E_{out}=const.$$
(14)
Note that the latter equation is compatible with the principle of information conservation expressed in the following form (equivalent to the one in the Introduction): For a compound quantum system a sum of information contained in the subsystems and information contained in entanglement is conserved in unitary processes .
To see how the above formalism works, consider two simple examples with ideal quantum transmission. Suppose, Alice sends an unentangled qubit of the system $`S`$ to Bob. Then the physical work $`W_p`$ is equal to 1 qubit. In result the informational content of Bobโs system increases by 1, thus also the logical work $`W_l`$ amounts to one qubit. Of course, in this case both โinโ and โoutโ entanglement are 0.
Suppose now that Alice sends maximally entangled qubit to Bob. Here, again, physical work is 1 qubit, and there is no initial entanglement. However the final entanglement is one ebit and logical work is 0, because the state of the Bob system is now completely mixed.
Now we see that, according to the balance equation (13) the difference $`W_pW_l`$ between the physical and logical work is due to entanglement. Indeed, as in the above example, sending particle may result in increase of entanglement rather than performing nonzero logical work.
## IV Useful logical work: quantum communication
The basic question arises in the context of quantum communication. Does the balance (13) distinguish between quantum and โclassicalโ communication in our model? It follows from definition that the physical work does not distinguish between these types of communication. But what about logical work? Suppose that Alice sent to Bob a particle of the system $`M_A`$ in a pure state $`|0`$. But in our model such state does not undergo decoherence. Then the logical work $`W_l`$ is equal to one qubit. Needless to say it is not quantum communication. Hence the logical work is not โusefulโ in this case.
In quantum communication we are usually interested in sending faithfully any superpositions without decoherence. Therefore it is convenient to introduce the notion of useful logical work as follows.
Definition. Useful work is amount of qubits of the system $`S`$ transmitted without decoherence
$$W_u=\mathrm{log}\mathrm{dim},$$
(15)
where $``$ is the Hilbert space transmitted asymptotically faithfully. The latter means that any state of this space would be transmitted with asymptotically perfect fidelity. We see that the work performed in previous example was not useful, since in result of the process, only the states $`|0`$ or $`|1`$ can be transmitted faithfully.
## V Balance of information in teleportation
To see how the above formalism works, consider the balance of quantum information in teleportation . Now the system $`S_A`$ consists of a particle in unkown state and one particle from maximally entangled pair, whereas the second particle from the pair represents $`S_B`$ system. The system $`M_A`$ consists of two qubits that interact with environment $`R_A`$ (Fig. 1).
The latter is only to ensure effective irreversibility of the measurement and it is demonstrable that its action is irrelevant to the information balance in the case of teleportation. As one knows, the initial state can be written in the following form
$$\psi _{in}\psi _0=\psi _{S_A^{}}^{unknown}\psi _{S_A^{\prime \prime }S_B}^{singlet}|00_{M_A},$$
(16)
where $`\psi _{S_A^{}}^{unknown}`$ is the state to be teleported, $`\psi _{S_A^{\prime \prime }S_B}^{singlet}`$ is the singlet state of entangled pair and $`|00_{M_A}`$ is the initial state of the measuring system. It is easy to check that the initial entanglement $`E_{in}`$ of the initial state is equal to one $`e\text{bit}`$. Now Alice performs โmeasurementโ being local unitary transformation on her joint system $`S_A^{}+S_{A^{\prime \prime }}+M_A`$. In result $`\psi _{in}`$ transforms to
$$\psi _1=\frac{1}{2}\underset{i=0}{\overset{3}{}}\psi _{S_A^{}S_A^{\prime \prime }}^i\psi _B^{i(unknown)}|i_{M_A},$$
(17)
where $`\psi _{S_A,S_{A^{\prime \prime }}}^i`$ constitute Bell basis, $`\psi _B^{i(unknown)}`$ is rotated $`\psi _{S_A}^{unknown}`$, $`|i_{M_A}`$ is the state of the system $`M_A`$ indicating the result of the measurement ($`i`$-th Bell state obtained). Since the Aliceโs operation is unitary one, it does not change initial asymptotic entanglement. Subsequently, Alice sends the two particles of the system $`M_A`$ to Bob. In accordance with definition (6), it corresponds to two qubits $`W_p=2`$ of work performed over the system. At the same time the state $`\psi _1`$ transforms to $`\psi _2`$ of the form
$$\psi _2=\frac{1}{2}\underset{i=0}{\overset{3}{}}\psi _{S_A^{}S_A^{\prime \prime }}^i\psi _B^{i(unknown)}|i_{M_B}.$$
(18)
Finally Bob decouples the system $`S_B`$ from other ones by unitary transformation that of course does not change the asymptotic entanglement.
After classical communication from Alice entanglement of the total system increased to the value $`E_{out}=2`$ e-bits. Indeed Alice sends two particles of system $`M_A`$ to Bob which are entangled with particles $`S_A^{}`$, $`S_A^{\prime \prime }`$). On the other hand, the logical work performed by the system in the above process amounts to $`W_u=1`$. One can see the balance equation (13) is satisfied, and is of the following form
$$(E_{in}=1)+(W_p=2)=(E_{out}=2)+(W_u=1)$$
(19)
One easily recoginzes the result of the logical work in the transmission of the unknown state to Bob. Since it is faithfully transmitted independently of its particular form, we obtain that also useful logical work $`W_u`$ is equal to 1 qubit. Hence in the process of teleportation all the work performed by the system is useful, and represents quantum communication.
## VI Thermodynamical entanglement-energy analogy. Gibbs-Hemholtz-like equation
So far we have considered balance of information in closed QC system. For open system (being, in general, in mixed state) the situation is much more complicated being a reflection of fundamental irreversibility in the asymptotic mixed state entanglement processing . Namely it has been shown there is a discontinuity in the structure of noisy entanglement. It appeared that there are at least two quantitively different types of entanglement: free - useful for quantum communication, and bound - nondistillable, very weak and peculiar type of entanglement. In accordance with entanglement-energy analogy this new type of entanglement was defined by equality
$$E_F=E_{bound}+E_D,$$
(20)
where $`E_F`$ and $`E_D`$ are asymptotic entanglement of formation and distillable entanglement respectively. Note that for pure entangled states $`|\mathrm{\Psi }\mathrm{\Psi }|`$ we have always $`E_F=E_D`$, $`E_{bound}=0`$ . Then, in this case the whole entanglement can be converted into the useful quantum work (see Fig. 2a) with $`EE_F(|\mathrm{\Psi }\mathrm{\Psi }|)`$). For bound entangled mixed states we have $`E_D=0`$, $`E_F=E_{bound}`$. It is quite likely that $`E_F>0`$ (so far we know only that $`E_f>0`$ ). Here $`E_f`$ is entanglement of formation defined in Ref. . and then all prior nontrivial entanglement of formation would be completely lost. Thus in any process involving only separable or bound entangled states useful quantum work is just zero. In general, hovewer, it can happen that the state contains two different types of entanglement. Namely there are cases where $`E_{bound}`$ is strictly positive i. e. we have
$$E_{bound}=E_FE_D>0.$$
(21)
This reveals fundamental irreversibility in the domain of quantum asymptotic information processing . It can be viewed as an analogue to irreversible thermodynamical processes where only the free energy (which is not equal to the total energy) can be converted to useful work. This supports the view according to which the equation (20) can be regarded as quantum information counterpart of the thermodynamical Gibbs-Hemholtz equation $`U=F+TS`$ where quantities $`E_F`$, $`E_D`$, $`E_{bound}`$ correspond to internal energy $`U`$, free energy $`F`$ and bound energy $`TS`$ respectively ($`T`$ and $`S`$ are the temperature and the entropy of the system).
The above entanglement-energy analogy has lead to the extension of the โclassicalโ paradigm of LOCC operations by considering new class of entanglement processing called here entanglement enhanced LOCC operations (EELOCC). In particular, it suggested that entanglement can be pumped from one to other system producing different nonclasical chemical-like type processes. In fact it allowed to find a new quantum effect - activation of bound entanglement that corresponds to chemical activation process . Similarly, a recently discovered cathalysis of pure entanglement involves EELOCC operations . In result the second principle of entanglement processing (see Introduction) has been generalized to cover the EELOCC paradigm: By local action, classical communication and N qubits of quantum communication, entanglement cannot increase more than N e-bits.
Now, it is interesting in the above context to consider the problem of information balance in the cases where systems are in mixed states.
## VII Balance of information in distillation process
So far in our balance analysis the initial state of the QC system was pure. Let us consider the more general case. Suppose that initial state of the system $`S`$ is mixed. We have not generalized formalism to such case. We can however perform balance of information in the case of the distillation process (see in this context ). This task would be, in general, very difficult, because the almost all known distillation protocols are stochastic. As one knows, the distillation protocol aims at obtaining singlet pairs from a large amount of noisy pairs (in mixed state) by LOCC operations. A convenient form of such a process would be the following: Alice and Bob start with $`n`$ pairs, and after distillation protocol, end up with $`m`$ singlet pairs. Such a protocol we shall call deterministic. Unfortunately, in the stochastic protocols the situation is more complicated: Alice and Bob get with some probabilities different number of output distilled pairs:
$`\varrho _{in}=\underset{n}{\underset{}{\varrho \varrho \mathrm{}\varrho }}\{\begin{array}{c}p_0,\text{no output singlets}\hfill \\ p_1,\text{one output singlet}\hfill \\ p_2,\text{two output singlet}\hfill \\ \mathrm{}\hfill \end{array}`$
Since we must to describe the process in terms of closed system, we will not see the above probabilities, but only their amplitudes. As a result, we will have no clear distinction between the part of the system containing distilled singlet pairs and the part containing the remaining states of no useful entanglement.
Consider for example the first stage of the Bennett et al. recursive protocol. It involves the folowing steps
* take two two-spin 1/2 pairs, each in input state $`\varrho `$
* perform operation $`XORXOR`$
* measure locally the spins of the target pair, and:
+ if the spins agree (probability $`p_a`$), keep the source pair
+ if the spins disagree (probability $`p_d`$), discard both pairs
After this operation we have the following final โensembleโ
$`\{(p_a,\text{one pair in a new state }\stackrel{~}{\varrho }),(p_d,\text{no pairs})\}`$
If we include environment to the description, the events โno pairโ and โone pair in state $`\stackrel{~}{\varrho }`$โ will be entangled with states of measuring apparatuses (and environment) indicating these events. Then we see, that our total system becomes more and more entangled in a various possible ways, so that it is rather impossible to perform the balance of information.
Fortunately, in a recent work Rains showed that any distillation protocol can be replaced with a deterministic one, achieving the same distillation rate:
$`\varrho ^n\varrho _{out}|\psi _{distilled}\psi _{distilled}|\varrho _{rejected}`$
where $`\psi _{distilled}`$ is the state of $`m`$ distilled singlet pairs while $`\varrho _{rejected}`$ is the state of the rejected pairs. In this case the system can be divided into two parts
$$S=S_{distilled}+S_{rejected}$$
(22)
where $`S_{distilled}`$ is disentangled with the rest of universe $`S_{rejected}`$ is entangled with $`M`$, hence also with environment $`R`$.
This possibility of the clear partition into two systems is crucial for our purposes. Now the whole balance can be be preformed in this case as follows. As an input we have the state $`\varrho `$ with value of asymptotic entanglement of formation $`E=E_F(\varrho )`$. Because it is mixed we can take its purification adding come ancilla which would have the asymptotic entanglement $`E^{}=E+(E^{}E)`$. Now we can perform the distillation process, having no access to the ancilla. After the process the state of our whole system is still separated according to the formula (22) but now the state $`S_{rejected}`$ involves the degrees of freedom of the ancilla. The balance of the information can now be easily performed taking, in particular, into account that distillable entanglement $`E_D`$ can be interpreted as a useful work (15) $`W_u`$ (Alice can always teleport through state $`|\mathrm{\Psi }_{distilled}\mathrm{\Psi }_{distilled}|`$ if she wishes). To make the balance fully consistent one should substract from both input and output data the additional entanglement $`E^{}E`$ coming from extension of the system to the pure state. As the input physical work (connected with optimal distillation protocol) is the same regardless of the value $`E^{}E`$ and the kind of the ancilla itself, the whole balance is completely consistent. The input quantities of $`E`$, $`\mathrm{\Delta }=(E^{}E)`$ plus $`W_p`$ as well as the output ones $`E_D=W_u`$, $`\mathrm{\Delta }`$, $`E_{out}=E(\varrho _{rejected})=E_{bound}`$ are depicted on figure Fig. 2b. In particular if we deal with BE states then the corresponding diagram takes the form of Fig. 2c.
## VIII Objectivity of quantum information: information interpretation of quantum states
As we have dealt with balance of information in quantum composite systems it is natural to ask about objectivity of the entity we qualify. In this section we discuss that question and related ones in the context of quantum information theory and interpretational problems of quantum mechanics. As one knows the latter defens oneself wery well against commonly accepted interpretation. In result a number of different interpretations permanently grows while there is no operational criterions (exept, may be, Ockham reazor) to eliminate at least some of them.
It is characteristic that despite of dynamical development of interdisciplinary domain - quantum information there is no, to our knowledge, impact of the latter on interpretational problems. In this context a basic question arises: Does quantum information phenomena provide objective promises for existence of โnaturalโ ontology inherent in quantum formalism?
It is interesting that from among discovered recently quantum effects just quantum cryptography provides answer โyesโ. To see it clearly, consider quantum cryptographic protocol. A crucial observation is that the possibility of secret sharing key is due to the fact that we send quantum states themselves not merely the classical information about them! Clearly, the latter could be cloned by the eavesdropper and it is reason for which all classical cryptographic schemes are, in principle, not secure. Then the use of qubits is crucial if we would like to take any advantage of the novel possibilities offered quantum information theory.
Now, as there are experimental implementations of quantum information protocols , it follows that quantum information is objective and it can provide natural ontological basis for interpretation of quantum mechanics. Then we arrive at important conclusion: Quantum states carry two complementary kinds of information: the โclassicalโ information involving quantum measurements and โquantumโ information that can not be cloned \[komentarz\].
Note that it is consistent with proposed earlier information intepretation of the wave function in terms of objective information content . On the other hand it contradicts the Copenhagen interpretation according to which the wavefunctions have no objective meaning and only reality is the result of a measurement. It is remarkable that the above information interpretation of quantum states is compatible with the above mentioned unitary information field concept which rests in the assumption that information is physical and can be defined independently of probability itself. First axiomatic definition of classical information โwithout probabilitiesโ was considered by Ingarden and Urbanik . Quantum version of the definition was introduced by Ingarden and Kossakowski . On the other hand Kolmogorow , Solomonoff , Chaitin introduced the concept of classical alghoritmic information or complexity. Recently the classical alghoritmic information was incorporated to the definition of the so called physical entropy being a constant of โmotionโ under the โdemonic evolutionโ .
Quite recently alghoritmic information theory was extended in different ways to quantum states by Vitanyi and Berthiaume et al. . In fact one can convince oneself that the approaches and correspond to the above complementary kinds of information associated with quantum state. Indeed, Vitanyi alghoritmic complexity measures amount of โclassicalโ information in bits necessary to approximate the quantum state. Needless to say, form the point of view of quantum cryptography such information is useless. On the other hand the bounded fidelity version of quantum Kolmogorow complexity measures amount quantum information in a qubit string and it is closely related to quantum compression theory .
## IX Summary
In conclusion we have developed the entanglement โ energy analogy based on some natural postulates: (i) entanglement is a form of quantum information being counterpart of internal energy, (ii) the process of sending qubits as a counterpart of work. We also assume that the evolution of the quantum system is unitary.
Basing on the above postulates we have considered the balance of quantum information for bipartite quantum communication systems i. e. the systems composed of two spatially separated laboratories endowed with classical informational channel plus local quantum operations. Wa have introduced the notion of informational content of quantum state being a difference of maximal possible von Neumann entropy and the actual one. Then we have defined physical work as a number of qubits physically sent form Alice to Bob. We have also defined logical work as as increase of the informational content of Bob state. To have a proper description of quantum communication processes we have also introduced a notion of useful logical work as amount of qubits transmitted without decoherence.
Those tools have allowed us to perform the detailed balances of quantum information in two important processes of quantum communication: quantum teleportation and distillation of quantum noisy entanglement. In particular we have discuss the question of balance of quantum information for open systems. In the context of balance scheme and related notions we conclude that the irreversibility connected with existence of bound entanglement can be viewed as an analogue to irreversible thermodynamical processes where only the free energy (which is not equal to the total energy) can be converted to useful work. This allows us to interprete the equation for entanglement of formation as quantum information counterpart of the thermodynamical Gibbs-Hemholtz equation.
Finally we discuss the objectivity of quantum information in general context of some recent achievments of quantum information theory including quantum cryptography and recent propositions of classical and quantum alghorytmic information. This leads us to the conclusion that quantum states reflect properties of quantum information as objective entity involving โclassicalโ and โquantumโ components which correspond to recently introduced โclassicalโ and โquantumโ alghoritmic complexities. So the balance performed in the present paper concerns objective quantities rather than purely formal objects. We hope that the present informational approach to bipartite quantum communication systems, when suitably developed, may lead to deeper understanding of quantum information processing domain.
M. H. and P. H. thank Chris Fuchs and Paweล Masiak for discussions on quantum information. Part of this work was made during ESF-Newton workshop (Cambridge 1999). The work is supported by Polish Committee for Scientific Research, contract No. 2 P03B 103 16. |
warning/0002/astro-ph0002281.html | ar5iv | text | # Geometric Gaussianity and Non-Gaussianity in the Cosmic Microwave Background
## I introduction
In recent years, locally isortopic and homogeneous Friedmann-Robertson-Walker (FRW) models with non-trivial topology have attracted much attention. In the standard scenario, simply-connectivity of the spatial hypersurface is assumed for simplicity. However, the Einstein equations, being local equations, do not fix the global topology of the spacetime. In other words, a wide variety of topologically distinct spacetimes with the same local geometry described by a local metric element remain unspecified (see for review on the cosmological topology). The determination of the global topology of the universe is one of the most important problem of the modern observational cosmology.
For flat models without the cosmological constant, severest constraints have been obtained by using the COBE DMR data. The suppression of the fluctuations on scales beyond the topological identification scale $`L`$ leads to the decrease of the angular power spectra $`C_l`$ of the Cosmic Microwave Background (CMB) temperature fluctuations on large angular scales which puts a lower bound $`L2400h^1`$Mpc (with $`h=H_0/100\mathrm{kms}^1\mathrm{Mpc}^1`$) for a compact flat 3-torus model without the cosmological constant . Similar constraints have been obtained for other compact flat models . The maximum expected number of copies of the fundamental domain (cell) inside the last scattering surface is approximately 8 for the 3-torus model.
In contrast, for low density models, the constraint could be considerably milder than the locally isotropic and homogeneous flat (Einstein-de-Sitter) models since a bulk of large-angle CMB fluctuations can be produced by the so-called (late) integrated Sachs-Wolfe effect (ISW) which is the gravitational blueshift effect of the free streaming photons by the decay of the gravitational potential. As the gravitational potential decays in either $`\mathrm{\Lambda }`$-dominant epoch or curvature dominant epoch, the free streaming photons with large wavelength (the light travel time across the wavelength is greater than or comparable to the decay time) that climbed a potential well at the last scattering experience blueshifts due to the contraction of the comoving space along the trajectories of the photons. Because the angular sizes of the fluctuations produced at late time are large, the suppression of the fluctuations on scale larger than the topological identification scale does not lead to a significant suppression of the large-angle power if the ISW effect is dominant. Recent works have shown that the large-angle power ($`2l20`$) are completely consistent with the COBE DMR data for compact hyperbolic (CH) models which include a small CH orbifold, the Weeks and the Thurston manifolds with volume $`0.72,0.94`$ and $`0.98`$ in unit of the cube of the curvature radius, respectively. Note that the Weeks manifold is the smallest and the Thurston manifolds is the second smallest in the known CH manifolds. For instance, the number of copies of the fundamental domain inside the last scattering surface at present is approximately 190 for a Weeks model with $`\mathrm{\Omega }_0=0.3`$.
If the space is negatively curved, for a fixed number of the copies of the fundamental domain inside the present horizon, the large-angle fluctuations can be produced much effectively. In negatively curved spaces (hyperbolic spaces), trajectories of photons subtend a much smaller angle in the sky for a given scale. In other words, for a given angle of a pair of two photon trajectories, the physical distance of the trajectories is much greater than that in a flat space. Therefore, even if there is a number of copies of the fundamental domain which intersect the last scattering surface, the number of copies which intersect the wave front (a sphere with $`z=\text{const.}`$) of the free streaming photons is exponentially decreased at late time when the large-angle fluctuations are produced due to the ISW effect.
However, one may not be satisfied with the constraints using only the angular power spectrum $`C_l`$ since it contains only isotropic information of the ensemble averaged temperature fluctuations . If they have anisotropic structures, non-Gaussian signatures must be revealed. In fact, the global isotropy of the locally isotropic and homogeneous FRW models is generally broken. For instance, a flat 3-torus obtained by identifying the opposite faces of a cube is obviously anisotropic at any points. Thus the temperature fluctuations averaged over the initial conditions in these multiply-connected FRW models are no longer $`SO(3)`$ invariant at a certain point. The temperature fluctuations on the sky are written in terms of (real) spherical harmonics $`Q_{lm}(๐ง)`$ as
$$\frac{\mathrm{\Delta }T}{T}(๐ง)=\underset{l}{}\underset{m=l}{\overset{l}{}}b_{lm}Q_{lm}(๐ง).$$
(1)
If the distribution functions of the real expansion coefficients $`b_{lm}`$ are $`SO(3)`$ invariant, the temperature fluctuations must be Gaussian provided that $`b_{lm}`$โs are independent random numbers . Therefore, the temperature fluctuations at a certain point in the multiply-connected FRW models are not Gaussian if $`b_{lm}`$โs are independent.
For the simplest flat 3-torus models (without rotations in the identification maps) which are globally homogeneous, it is sufficient to choose one observing point and estimate how the power is distributed among the $`m`$โs for a given angular scale $`l`$ in order to see the effect of the global anisotropy. However, in general, one must consider an ensemble of fluctuations at different observing points because of the spatial (global) inhomogeneity. Previous analyses have not fully investigated the dependence of the temperature fluctuations on choice of the observing points.
Lack of analytical results on the eigenmodes makes it difficult to investigate the nature of the temperature fluctuations in CH models. However, we may expect a high degree of complexity in the eigenmodes since the corresponding classical systems (geodesic flows) are strongly chaotic. In fact, it has been numerically found that the expansion coefficients of the low-lying eigenmodes on the Thurston manifold at the point where the injectivity radius is maximal are Gaussian pseudo-random numbers which supports the previous analysis of the excited states (higher modes) of a two-dimensional asymmetrical CH model . We have put a prefix โpseudoโ since the eigenmodes are actually constrained by the periodic boundary conditions. These results imply that the statistical properties of the eigenmodes on CH spaces (orbifolds and manifolds) can be described by random-matrix theory (RMT). An investigation of the dependence of the property on the observing points is also important since CH spaces have symmetries (isometric groups) which may veil the random feature of the eigenmodes. In this paper, a detailed analysis on the statistical property of low-lying eigenmodes on the Weeks and the Thurston manifolds is conducted.
Assuming that the eigenmodes are Gaussian, one can expect that the anisotropic structure in the $`(l,m)`$ space is vastly erased when one averages the fluctuations over the space. This seems to be a paradox since the CH spaces are actually globally anisotropic. However, one should consider a spatial average of fluctuations with different initial conditions if one believes the Copernican principle that we are not in the center of the universe. Even if the space is anisotropic at a certain point, the averaged fluctuations may look isotropic by considering an ensemble of fluctuations at all the possible observing points. Note that the eigenmodes on CH spaces have no particular directions if they are Gaussian.
If the initial fluctuations are constant for each eigenmode, as we shall see, the Gaussian randomness of the temperature fluctuations can be solely attributed to the Gaussian pseudo-randomness of the eigenmodes. In this case, the Gaussian randomness of the temperature fluctuations has its origin in the geometrical property of the space (Geometric Gaussianity). Choosing an observing point is equivalent to fixing a certain initial condition. However, it is much natural to assume that the initial fluctuations are also random Gaussian as the standard inflationary scenarios predict. Then the temperature fluctuations may not obey the Gaussian statistics because they are written in terms of products of two different independent Gaussian numbers rather than sums while they remain almost spatially isotropic if averaged over the space.
In this paper, Gaussianity of eigenmodes and non-Gaussianity in the CMB for two smallest CH models (the Weeks and the Thurston models) are investigated. In Sec. II, numerical results on Gaussianity of eigenmodes are shown and we discuss to what extent the results are generic. In Sec. III, we study the non-Gaussian behavior of the temperature fluctuations in the ($`l,m`$) space. In Sec. IV, topological quantities (total length and genus) of isotemperature contours are numerically simulated for studying the non-Gaussian behavior in the real space. Finally, we summarize our conclusions in Sec. V.
## II GEOMETRIC GAUSSIANITY
In locally isotropic and homogeneous FRW background spaces, each type (scalar, vector and tensor) of first-order perturbations can be decomposed into a decoupled set of equations. In order to solve the decomposed linearly perturbed Einstein equations, it is useful to expand the perturbations in terms of eigenmodes of the Laplacian which satisfies the Helmholtz equation with certain boundary conditions,
$$(^2+k^2)u_k(x)=0,$$
(2)
since each eigenmode evolves independently in the linear approximation. Then one can easily see that the time evolution of the perturbations in the multiply-connected locally isotropic and homogeneous FRW spaces coincide with that in the FRW spaces while the global structure of the background space is described solely by these eigenmodes.
Unfortunately, no analytical expressions of eigenmodes on CH spaces have been known. Nevertheless, the correspondence between classical and quantum mechanics may provide us a clue for understanding the generic property of the eigenmodes. If one recognizes the Laplacian as the Hamiltonian in a quantum system, each eigenmode can be interpreted as a wavefunction in a stationary state. Because classical dynamical systems (=geodesic flows) on CH spaces are strongly chaotic (or more precisely they are K-systems with ergodicity, mixing and Bernoulli properties ), one can expect a high degree of complexity for each eigenstate. Imprint of the chaos in the classical systems may be hidden in the quantum counterparts. In fact, in many cases, the short-range correlations observed in the eigenvalues (energy states) have been found to be consistent with the universal prediction of RMT for three universality classes:the Gaussian orthogonal ensemble(GOE), the Gaussian unitary ensemble(GUE) and the Gaussian symplectic ensemble (GSE). In our case the statistical properties are described by GOE (which consist of real symmetric $`N\times N`$ matrices $`H`$ which obey the Gaussian distribution $`\mathrm{exp}(\text{Tr}H^2/(4a^2))`$ (where $`a`$ is a constant) as the systems possess a time-reversal symmetry. RMT also predicts that the squared expansion coefficients of an eigenstate with respect to a generic basis are distributed as Gaussian random numbers . Unfortunately, no analytic forms of generic bases(=eigenmodes) are known for CH spaces which seems to be an intractable problem. However, if the eigenmodes are continued onto the universal covering space by the periodic boundary conditions, they can be written in terms of a โgenericโ basis on the universal covering space (=3-hyperboloid $`H^3`$). In pseudospherical coordinates ($`R,\chi ,\theta ,\varphi `$), the eigenmodes are written in terms of complex expansion coefficients $`\xi _{\nu lm}`$ and eigenmodes on the universal covering space,
$$u_\nu =\underset{lm}{}\xi _{\nu lm}X_{\nu l}(\chi )Y_{lm}(\theta ,\varphi ),$$
(3)
where $`\nu =\sqrt{k^21}`$, $`X_{\nu l}`$ and $`Y_{lm}`$ denote the radial eigenfunction and (complex) spherical harmonic on the pseudosphere with radius $`R`$, respectively. Then the real expansion coefficients $`a_{\nu lm}`$ are given by
$`a_{\nu 00}`$ $`=`$ $`\text{Im}(\xi _{\nu 00}),a_{\nu l0}=\sqrt{c_{\nu l}}\text{Re}(\xi _{\nu l0}),`$ (4)
$`a_{\nu lm}`$ $`=`$ $`\sqrt{2}\text{Re}(\xi _{\nu lm}),m>0,`$ (5)
$`a_{\nu lm}`$ $`=`$ $`\sqrt{2}\text{Im}(\xi _{\nu lm}),m<0,`$ (6)
where
$`c_{\nu l}`$ $`=`$ $`{\displaystyle \frac{2}{(1+\text{Re}(F(\nu ,l)))}},`$ (7)
$`F(\nu ,l)`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Gamma }(l+\nu i+1)}{\mathrm{\Gamma }(\nu i)}}{\displaystyle \frac{\mathrm{\Gamma }(\nu i)}{\mathrm{\Gamma }(l\nu i+1)}}.`$ (8)
In this paper, the low-lying eigenmodes ($`k<13`$) on the Weeks and Thurston manifolds are numerically computed by the direct boundary element method. The identification matrices of the Dirichlet domains are obtained by a computer program โSnapPeaโ by J. R. Weeks . The computed eigenvalues are well consistent with that in the previous literature . The estimated errors in $`k`$ are within $`0.01`$. However, the last digits in $`k`$ may be incorrect. $`a_{\nu lm}`$ โs can be promptly obtained after the normalization and orthogonalization of these eigenmodes. The orthogonalization is achieved at the level of $`10^3`$ to $`10^4`$ (for the inner product of the normalized eigenmodes) which implies that each eigenmode is computed with relatively high accuracy.
In Fig.1 and Fig.2, one can see a high degree of complexity in the lowest eigenmodes on the Poincar$`\stackrel{ยด}{\text{e}}`$ ball which is isometric to the universal covering space $`H^3`$ whose coordinates are given by
$$x=R\mathrm{tanh}\frac{\chi }{2}\mathrm{sin}\theta \mathrm{cos}\varphi ,y=R\mathrm{tanh}\frac{\chi }{2}\mathrm{sin}\theta \mathrm{sin}\varphi ,z=R\mathrm{tanh}\frac{\chi }{2}\mathrm{cos}\theta .$$
(9)
Replacing $`\mathrm{tanh}\frac{\chi }{2}`$ by $`\mathrm{tanh}\chi `$ for each coordinate, one obtains the Klein (projective) coordinates. In the Poincar$`\stackrel{ยด}{\text{e}}`$ coordinates, angles of geodesics coincide with that of Euclidean ones. In the Klein coordinates, all geodesics are straight lines while angles does not coincide with that of Euclidean ones.
In what follows $`R`$ is normalized to 1 without loss of generality.
In Fig.3, one can see that the distribution of $`a_{\nu lm}`$โs which are ordered as $`l(l+1)+m+1`$ are qualitatively random. In order to estimate the randomness quantitatively, we consider a cumulative distribution of
$$b_{\nu lm}=\frac{|a_{\nu lm}\overline{a}_\nu |^2}{\sigma _\nu ^2}$$
(10)
where $`\overline{a}_\nu `$ is the mean of $`a_{\nu lm}`$โs and $`\sigma _\nu ^2`$ is the variance. If $`a_{\nu lm}`$โs are Gaussian then $`b_{\nu lm}`$ โs obey a $`\chi ^2`$ distribution $`P(x)=(1/2)^{1/2}\mathrm{\Gamma }(1/2)x^{1/2}e^{x/2}`$ with 1 degree of freedom. To test the goodness of fit between the the theoretical cumulative distribution $`I(x)`$ and the empirical cumulative distribution function $`I_N(x)`$, we use the Kolmogorov-Smirnov statistic $`D_N`$ which is the least upper bound of all pointwise differences $`|I_N(x)I(x)|`$ ,
$$D_N\underset{x}{sup}|I_N(x)I(x)|.$$
(11)
$`I_N(x)`$ is defined as
$`I_N(x)`$ $`=`$ $`\{\begin{array}{cc}0,\hfill & x<y_1,\hfill \\ j/N,\hfill & y_jx<y_{j+1},j=1,2,\mathrm{},N1,\hfill \\ 1,\hfill & y_Nx,\hfill \end{array}`$ (15)
where $`y_1<y_2<\mathrm{}<y_N`$ are the computed values of a random sample which consists of $`N`$ elements. For random variables $`D_N`$ for any $`z>0`$, it can be shown that the probability of $`D_N<d`$ is given by
$$\underset{N\mathrm{}}{lim}P(D_N<d=zN^{1/2})=L(z),$$
(16)
where
$$L(z)=12\underset{j=1}{\overset{\mathrm{}}{}}(1)^{j1}e^{2j^2z^2}.$$
(17)
From the observed maximum difference $`D_N=d`$, we obtain the significance level $`\alpha _D=1P`$ which is equal to the probability of $`D_N>d`$. If $`\alpha _D`$ is found to be large enough, the hypothesis $`I_N(x)I(x)`$ is not verified. The significance levels $`\alpha _N`$ for $`0l20`$ for eigenmodes $`k<13`$ on the Thurston manifold are shown in table 1.
The agreement with the RMT prediction is fairly good for most of eigenmodes which is consistent with the previous computation in . However, for five degenerated modes, the non-Gaussian signatures are prominent (in , two modes in ($`k<10`$) have been missed). Where does this non-Gaussianity come from?
First of all, we must pay attention to the fact that the expansion coefficients $`a_{\nu lm}`$ depend on the observing point. In mathematical literature the point is called the base point. For a given base point, it is possible to construct a particular class of fundamental domain called the Dirichlet (fundamental) domain which is a convex polyhedron. A Dirichlet domain $`\mathrm{\Omega }(x)`$ centered at a base point $`x`$ is defined as
$$\mathrm{\Omega }(x)=\underset{g}{}H(g,x),H(g,x)=\{z|d(z,x)<d(g(z),x)\},$$
(18)
where $`g`$ is an element of a Kleinian group $`\mathrm{\Gamma }`$(a discrete isometry group of $`PSL(2,)`$) and $`d(z,x)`$ is the proper distance between $`z`$ and $`x`$.
The shape of the Dirichlet domain depends on the base point but the volume is invariant. Although the base point can be chosen arbitrarily, it is a standard to choose a point $`Q`$ where the injective radius <sup>*</sup><sup>*</sup>*The injective radius of a point $`p`$ is equal to half the length of the shortest periodic geodesic at $`p`$. is locally maximal. More intuitively, $`Q`$ is a center where one can put a largest connected ball on the manifold. If one chooses other point as the base point, the nearest copy of the base point can be much nearer. The reason to choose $`Q`$ as a base point is that one can expect the corresponding Dirichlet domain to have many symmetries at $`Q`$ .
As shown in Fig.4, the Dirichlet domain at $`Q`$ has a $`Z2`$ symmetry (invariant by $`\pi `$-rotation) if all the congruent faces are identified. Generally, congruent faces are distinguished but it is found that these five modes have exactly the same values of eigenmodes on these congruent faces. Then one can no longer consider $`a_{\nu lm}`$โs as โindependentโ random numbers. Choosing the invariant axis by the $`\pi `$-rotation as the $`z`$-axis, $`a_{\nu lm}`$โs are zero for odd $`m`$โs which leads to the observed non-Gaussian behavior. It should be noted that the observed $`Z2`$ symmetry is not the subgroup of the isometry group (or symmetry group in mathematical literature) $`D2`$ (dihedral group with order $`2`$) of the Thurston manifold since the congruent faces must be actually distinguished in the manifoldThe observed $`Z2`$ symmetry is considered to be a โhidden symmetryโ which is a symmetry of the finite sheeted cover of the manifold (which tessellates the manifold as well as the universal covering space). For instance, the Dirichlet domain of the Thurston manifold can be tessellated by four pieces with three neighboring kite-like quadrilateral faces and one equilateral triangle on the boundary and seven faces which contain the center as a vertex. By identifying the four pieces (by a tetrahedral symmetry), one obtains an orbifold which has a $`Z2`$ symmetry.
Thus the observed non-Gaussianity is caused by a particular choice of the base point. However, in general, the chance that we actually observe any symmetries (elements of the isometry group of the manifold or the finite sheeted cover of the manifold) is expected to be very low. Because a fixed point by an element of the isometric group is either a part of 1-dimensional line (for instance, an axis of a rotation) or an isolated point (for instance, a center of an antipodal map).
In order to confirm that the chance is actually low, the KS statistics $`\alpha _D`$ of $`a_{\nu lm}`$โs are computed at 300 base points which are randomly chosen.
As shown in table 2, the averaged significance levels $`<\alpha _D>`$ are remarkably consistent with the Gaussian prediction. $`1\sigma `$ of $`\alpha _D`$ are found to be 0.26 to 0.30.
Next, we apply the run test for testing the randomness of $`a_{\nu lm}`$โs where each set of $`a_{\nu lm}`$ โs are ordered as $`l(l+1)+m+1`$ (see ). Suppose that we have $`n`$ observations of the random variable $`U`$ which falls above the median and n observations of the random variable $`L`$ which falls below the median. The combination of those variables into $`2n`$ observations placed in ascending order of magnitude yields
UUU LL UU LLL U L UU LL,
Each underlined group which consists of successive values of $`U`$ or $`L`$ is called run. The total number of run is called the run number. The run test is useful because the run number always obeys the Gaussian statistics in the limit $`n\mathrm{}`$ regardless of the type of the distribution function of the random variables.
As shown in table 3, averaged significance levels $`<\alpha _r>`$ are very high (1$`\sigma `$ is 0.25 to 0.31). Thus each set of $`a_{\nu lm}`$โs ordered as $`l(l+1)+m+1`$ can be interpreted as a set of Gaussian pseudo-random numbers except for limited choices of the base point where one can observe symmetries of eigenmodes.
Up to now, we have considered $`l`$ and $`m`$ as the index numbers of $`a_{\nu lm}`$ at a fixed base point. However, for a fixed $`(l,m)`$, the statistical property of a set of $`a_{\nu lm}`$โs at a number of different base points is also important since the temperature fluctuations must be averaged all over the places for spatially inhomogeneous models. From Fig.5, one can see the behavior of m-averaged significance levels
$$\alpha _D(\nu ,l)\underset{m=l}{\overset{l}{}}\frac{\alpha _D(a_{\nu lm})}{2l+1}$$
(19)
which are calculated based on 300 realizations of the base points. It should be noted that each $`a_{\nu lm}`$ at a particular base point is now considered to be โone realizationโ whereas a choice of $`l`$ and $`m`$ is considered to be โone realizationโ in the previous analysis (table 1). The agreement with the RMT prediction is considerably good for components $`l>1`$. For components $`l=1`$, the disagreement occurs for only several modes. However, the non-Gaussian behavior is distinct in $`l=0`$ components. What is the reason of the non-Gaussian behavior for $`l=0`$?
Let us estimate the values of the expansion coefficients for $`l=0`$. In general, the complex expansion coefficients $`\xi _{\nu lm}`$ can be written as,
$$\xi _{\nu lm}(\chi _0)=\frac{1}{X_{\nu l}(\chi _0)}u_\nu (\chi _0,\theta ,\varphi )Y_{lm}^{}(\theta ,\varphi )๐\mathrm{\Omega }.$$
(20)
For $`l=0`$, the equation becomes
$$\xi _{\nu 00}(\chi _0)=\frac{i}{2\sqrt{2}}\frac{\mathrm{sinh}\chi _0}{\mathrm{sin}\nu \chi _0}u_\nu (\chi _0,\theta ,\varphi )๐\mathrm{\Omega }.$$
(21)
Taking the limit $`\chi _00`$, one obtains,
$$\xi _{\nu 00}=\frac{2\pi u_\nu (0)i}{\nu }.$$
(22)
Thus $`a_{\nu 00}`$ can be written in terms of the value of the eigenmode at the base point. As shown in Fig.1, the lowest eigenmodes have only one โwaveโ on scale of the topological identification scale $`L`$ (which will be defined later on) inside a single Dirichlet domain which implies that the random behavior within the domain may be not present. Therefore, for low-lying eigenmodes, one would generally expect non-Gaussianity in a set of $`a_{\nu 00}`$ โs. However, for high-lying eigenmodes, this may not be the case since these modes have a number of โwavesโ on scale of $`L`$ and they may change their values locally in a almost random fashion.
The above argument cannot be applicable to $`a_{\nu lm}`$ โs for $`l0`$ where $`X_{\nu l}`$ approaches zero in the limit $`\chi _00`$ while the integral term
$$u_\nu (\chi _0,\theta ,\varphi )Y_{lm}^{}(\theta ,\varphi )๐\mathrm{\Omega }$$
(23)
also goes to zero because of the symmetric property of the spherical harmonics. Therefore $`a_{\nu lm}`$ โs cannot be written in terms of the local value of the eigenmode for $`l0`$. For these modes, it is better to consider the opposite limit $`\chi _0\mathrm{}`$. It is numerically found that the sphere with very large radius $`\chi _0`$ intersects each copy of the Dirichlet domain almost randomly (the pulled back surface into a single Dirichlet domain chaotically fills up the domain). Then the values of the eigenmodes on the sphere with very large radius vary in an almost random fashion. For large $`\chi _0`$, we have
$$X_{\nu l}(\chi _0)e^{2\chi _0+\varphi (\nu ,l)i},$$
(24)
where $`\varphi (\nu ,l)`$ describes the phase factor. Therefore, the order of the integrand in Eq. (20) is approximately $`e^{2\chi _0}`$ since Eq. (20) does not depend on the choice of $`\chi _0`$. As the spherical harmonics do not have correlation with the eigenmode $`u_\nu (\chi _0,\theta ,\varphi )`$, the integrand varies almost randomly for different choices of $`(l,m)`$ or base points. Thus we conjecture that Gaussianity of $`a_{\nu lm}`$โs have their origins in the chaotic property of the sphere with large radius in CH spaces. The property may be related to the classical chaos in geodesic flowsIf one considers a great circle on a sphere with large radius, the length of the circle is very long except for rare cases in which the circle โcomes backโ before it wraps around in the universal covering space. Because the long geodesics in CH spaces chaotically (with no particular direction and position) wrap through the manifold, it is natural to assume that the great circles also have this chaotic property. .
So far we have seen the Gaussian pseudo-randomness of the $`a_{\nu lm}`$โs. Let us now consider the statistical properties of the expansion coefficients. As the eigenmodes have oscillatory features, it is natural to expect that the averages are equal to zero. In fact, the averages of $`<a_{\nu lm}>`$ โs over $`0l20`$ and $`lml`$ and 300 realizations of base points for each $`\nu `$-mode are numerically found to be $`0.006\pm 0.040.02`$(1$`\sigma `$) for the Weeks manifold, and $`0.003\pm 0.040.02`$(1$`\sigma `$) for the Thurston manifold. Let us next consider the $`\nu `$-dependence ($`k`$-dependence) of the variances $`Var(a_{\nu lm})`$. In order to crudely estimate the $`\nu `$-dependence, we need the angular size $`\delta \theta `$ of the characteristic length of the eigenmode $`u_\nu `$ at $`\chi _0`$
$$\delta \theta ^2\frac{16\pi ^2Vol(M)}{k^2(\mathrm{sinh}(2(\chi _o+r_{ave}))\mathrm{sinh}(2(\chi _or_{ave}))4r_{ave})},$$
(25)
where $`Vol(M)`$ denotes the volume of a manifold $`M`$ and $`r_{ave}`$ is the averaged radius of the Dirichlet domain. There is an arbitrariness in the definition of $`r_{ave}`$. Here we define $`r_{ave}`$ as the radius of a sphere with volume equivalent to the volume of the manifold,
$$Vol(M)=\pi (\mathrm{sinh}(2r_{ave})2r_{ave}),$$
(26)
which does not depend on the choice of the base point. The topological identification length $`L`$ is defined as $`L=2r_{ave}`$. For the Weeks and the Thurston manifold, $`L=1.19`$ and $`L=1.20`$ respectively. From Eq. (25), for large $`\chi _0`$, one can approximate $`u_\nu (\chi _o)u_\nu ^{}(\chi _o^{})`$ by choosing an appropriate radius $`\chi _o^{}`$ which satisfies $`\nu ^2\mathrm{exp}(2\chi _o)=\nu ^2\mathrm{exp}(2\chi _o^{})`$. Averaging Eq. (20) over $`l`$โs and $`m`$โs or the base points, for large $`\chi _0`$, one obtains,
$$<|\xi _{\nu ^{}lm}|^2>\frac{\mathrm{exp}(2\chi _o)}{\mathrm{exp}(2\chi _o^{})}<|\xi _{\nu lm}|^2>,$$
(27)
which gives $`<|\xi _{\nu lm}|^2>\nu ^2`$. Thus the variance of $`a_{\nu lm}`$โs is proportional to $`\nu ^2`$. The numerical results for the two CH manifolds shown in Fig.6 clearly support the $`\nu ^2`$ dependence of the variance.
As we have seen, the property of eigenmodes on general CH manifolds is summarized in the following conjecture:
Conjecture: Except for the base points which are too close to any fixed points by symmetries, for a fixed $`\nu `$, a set of the expansion coefficients $`a_{\nu lm}`$ over $`(l,m)`$โs can be considered as Gaussian pseudo-random numbers. For a fixed $`(\nu lm)(l>0)`$, the expansion coefficients at different base points that are randomly chosen can also be considered as Gaussian pseudo-random numbers. In either case, the variance is proportional to $`\nu ^2`$ and the average is zero.
## III NON-GAUSSIANITY IN OBSERVABLE ANGULAR POWER SPECTRA
As mentioned in the last section, perturbations in CH models are written in terms of linear combinations of eigenmodes and the time evolution of the perturbations. Because the time evolution of the perturbations coincides with that in open models, once the expansion coefficients $`\xi _{\nu lm}`$ (or $`a_{\nu lm}`$) are given, the evolution of perturbations in CH models can be readily obtained.
If one assumes that the perturbation is a adiabatic scalar type without anisotropic pressure, and the subhorizon effects such as acoustic oscillations of the temperature and the velocity of the bulk fluid, and the effect of the radiation contribution at high z are negligible, the time evolution of the growing mode of the Newtonian curvature $`\mathrm{\Phi }`$ is analytically given as (see e.g. )
$$\mathrm{\Phi }(\eta )=\frac{5(\mathrm{sinh}^2\eta 3\eta \mathrm{sinh}\eta +4\mathrm{cosh}\eta 4)}{(\mathrm{cosh}\eta 1)^3},$$
(28)
where $`\eta `$ denotes the conformal time. In terms of $`\mathrm{\Phi }`$, the temperature fluctuation in the sky are written as
$`{\displaystyle \frac{\mathrm{\Delta }T(๐ง)}{T}}`$ $`=`$ $`{\displaystyle \underset{lm}{}}a_{lm}Y_{lm}(๐ง)`$ (29)
$`=`$ $`{\displaystyle \underset{\nu lm}{}}\mathrm{\Phi }_\nu (0)\xi _{\nu lm}F_{\nu l}(\eta _0)Y_{lm}(๐ง),`$ (30)
where
$$F_{\nu l}(\eta _0)\frac{1}{3}\mathrm{\Phi }(\eta _{})X_{\nu l}(\eta _0\eta _{})2_\eta _{}^{\eta _0}๐\eta \frac{d\mathrm{\Phi }}{d\eta }X_{\nu l}(\eta _0\eta ).$$
(31)
Here $`\mathrm{\Phi }_\nu (0)`$ is the initial value of the curvature perturbation and $`\eta _{}`$ and $`\eta _0`$ are the conformal time of the last scattering and the present conformal time, respectively. The angular power spectrum $`C_l`$ is defined as
$`(2l+1)C_l`$ $`=`$ $`{\displaystyle \underset{m=l}{\overset{l}{}}}|a_{lm}|^2`$ (32)
$`=`$ $`{\displaystyle \underset{\nu ,m}{}}{\displaystyle \frac{4\pi ^4๐ซ_\mathrm{\Phi }(\nu )}{\nu (\nu ^2+1)\text{Vol}(M)}}|\xi _{\nu lm}|^2|F_{\nu l}(\eta _0)|^2,`$ (33)
where $`๐ซ_\mathrm{\Phi }(\nu )`$ is the initial power spectrum. It should be noted that the above formula converges to that of open models in the short-wavelength limit (summation to integration) provided that $`<|\xi _{\nu lm}|^2>`$ is proportional to $`\nu ^2`$. The reason is as follows: Let us denote the number of eigenmodes with eigenvalues equal to or less than $`\nu `$ by $`N(\nu )`$. In the short-wavelength limit $`\nu >>1`$ one can use Weylโs asymptotic formula which leads to
$$\frac{dN(\nu )}{d\nu }=\frac{\text{Vol}(M)}{2\pi ^2}\nu ^2.$$
(34)
Thus the $`\nu ^2`$ dependence in Eq.(34) is exactly cancelled out by the $`\nu ^2`$ dependence of eigenmodes. In what follows we assume the extended Harrison-Zelโdovich spectrum, i.e. $`๐ซ_\mathrm{\Phi }(\nu )=Const.`$ (in the flat limit, it converges to the scale invariant Harrison-Zelโdovich spectrum) as the initial power spectrum.
In estimating the temperature correlations, the non-diagonal terms ($`ll^{}`$ or $`mm^{}`$) may not be negligible if the background spatial hypersurface is not isotropic, in other words, the angular power spectrum $`C_l`$ may not be sufficient in describing the temperature correlations since $`C_l`$ provides us with only an isotropic information of statistics of the correlations. However, this is not the case for CH models to which the conjecture proposed in Sec. II is applicable. Based on the Copernican principle, it is not likely that we are at the center of any symmetries. Therefore, in order to statistically estimate the temperature correlations in the globally inhomogeneous background space, one has to consider an ensemble of fluctuations with different initial conditions at different places (or base points) with different orientations. Almost all the anisotropic information is lost in the spatial averaging process since the eigenmodes are Gaussian.
As shown in Fig.7, for 300 realizations of observing points(left), the averaged absolute values of the off-diagonal elements in unit of diagonal elements are very small ($`0.016`$) whereas their contributions seem to be not negligible ($`0.25`$) at one particular observing point(right) where one can observe a symmetry of the Dirichlet domain. Thus the statistical property of the temperature correlation can be estimated by using $`C_l`$โs provided that the eigenmodes are Gaussian which validates the previous analyses using $`C_l`$โs for constraining the CH models . The spatial averaging process<sup>ยง</sup><sup>ยง</sup>ยงIn general, one should include an averaging process over different choices of orientation of coordinates as well as an averaging process over different choices of the observing point. Nevertheless, the Gaussian conjecture in Sec. II implies that the eigenmodes on CH spaces are โ$`SO(3)`$ invariantโ if averaged all over the space. Therefore, omission of the averaging procedure for different orientations of coordinates make no difference. must be taken into account since there is no reason to believe that we are in the center of any symmetries.
If the initial conditions satisfy $`(\mathrm{\Phi }_\nu (0))^2\nu (\nu ^2+1)`$ that corresponds to the extended Harrison-Zelโdovich spectrum , then Eq.(30) tells us that the temperature fluctuation is Gaussian since it is equal to a sum of Gaussian (pseudo-)random numbers at almost all the observing points. In this case, the Gaussian randomness of the temperature fluctuations in CH models can be solely attributed to the geometrical property of the space (geometric Gaussianity) which may be related to the deterministic chaos of the corresponding classical system. In other words, the Gaussian randomness can be explained in terms of the classical physical quantities without considering the initial quantum fluctuations provided that the above conditions are initially (deterministically) satisfied.
However, it is much natural to assume that $`\mathrm{\Phi }_\nu (0)`$โs are also random Gaussian as in the inflationary scenarios in which Gaussianity (on large scales) of the temperature fluctuations has its origin in Gaussianity of the initial quantum fluctuations because the angular powers are generally similar to the extended Harisson-Zelโdovich spectrum. Then the statistical properties of the temperature fluctuations are determined by the sum of the products of the two independent Gaussian random numbers (the initial fluctuations and the expansion coefficients of the eigenmodes).
Let us calculate the distribution function $`F(Z,\sigma _Z)`$ of a product of two independent random numbers $`X`$ and $`Y`$ that obey the Gaussian (normal) distributions $`N(X;0,\sigma _X)`$ and $`N(Y;0,\sigma _Y)`$, respectively.
$$N(X;\mu ,\sigma )\frac{1}{\sqrt{2\pi }\sigma }\text{e}^{(X\mu )^2/2\sigma ^2}.$$
(35)
Then $`F(Z=XY,\sigma _Z)`$ is readily given by
$`F(Z,\sigma _Z)`$ $`=`$ $`2{\displaystyle _0^{\mathrm{}}}N(Z/Y,0,\sigma _X)N(Y,0,\sigma _Y){\displaystyle \frac{dY}{Y}}`$ (36)
$`=`$ $`{\displaystyle \frac{1}{\pi \sigma _X\sigma _Y}}K_0\left({\displaystyle \frac{|Z|}{\sigma _X\sigma _Y}}\right),`$ (37)
where $`K_0(z)`$ is the modified Bessel function. The average of $`Z`$ is zero and the standard deviation satisfies $`\sigma _Z=\sigma _X\sigma _Y`$. As is well known, $`K_0(z)`$ is the Green function of the diffusion equation with sources distributed along an infinite line. Although $`K_0(z)`$ is diverged at $`z=0`$ its integration over $`(\mathrm{},\mathrm{})`$ is convergent.
From the asymptotic expansion of the modified Bessel function
$$K_0(z)\sqrt{\frac{\pi }{2z}}\text{e}^z\left[1\frac{1^2}{1!8z}+\frac{1^23^2}{2!(8z)^2}\frac{1^23^25^2}{3!(8z)^3}+\mathrm{}\right],z>>1,$$
(38)
one obtains in the lowest order approximation,
$$F(Z,\sigma )\frac{1}{\sqrt{2\pi \sigma |Z|}}\text{e}^{|Z|/\sigma },Z>>1.$$
(39)
Thus $`F(Z,\sigma )`$ is slowly decreased than the Gaussian distribution function with the same variance in the large limit. One can see the two non-Gaussian features in Fig.8(left):the divergence at $`Z0`$ and the slow convergence to zero at $`Z\mathrm{}`$. The slow convergence is an important feature, as we shall see, in distinguishing the non-Gaussian models with the Gaussian ones. In the modest region $`0.4<|Z|<2.4`$, $`F(Z,1)`$ is much less than $`N(Z,0,1)`$. Generally, the temperature fluctuation is written as a sum of the random variables $`Z_i`$ which obeys the distribution function $`F(Z_i,\sigma _{Z_i})`$ for a fixed set of cosmological parameters. For large-angle fluctuations, only the eigenmodes with large wavelength ($`2\pi /k`$)can contribute to the sum. Due to the finiteness of the space, the number of eigenmodes which dominantly contribute to the sum is finite. Therefore, the fluctuations are distinctively non-Gaussian. For small-angle fluctuations, the number of eigenmodes that contribute to the sum becomes so large that the distribution function converges to the Gaussian distribution as the central limit theorem implies. One can see from Fig.8 (right) that the distribution function $`G(W,1)`$ of $`W=Z1+Z2`$ where both $`Z1`$ and $`Z2`$ obey $`F(Z,\sqrt{2})`$ is much similar to the Gaussian distribution $`N(Z,0,1)`$ than $`F(W,1)`$ in the modest region.
Now let us see the non-Gaussian features of the observable angular power spectrum $`\widehat{C_l}`$ assuming that the initial fluctuations are Gaussian. First of all, we define a statistic $`\stackrel{~}{\chi }^2(2l+1)\widehat{C_l}/C_l`$ where
$$(2l+1)\widehat{C_l}=\underset{m=l}{\overset{l}{}}b_{lm}^2.$$
(40)
If the expansion coefficients $`b_{lm}`$ of the temperature fluctuation in the sky are Gaussian, $`\stackrel{~}{\chi }^2`$ must obey the $`\chi ^2`$ distribution with $`2m+1`$ degrees of freedom.
Fig.9 shows the two non-Gaussian features in the distribution of $`b_{lm}`$โs:a slight shift of the peak to the center(zero); slow convergence to zero for large $`\stackrel{~}{\chi }^2`$. As shown in Fig.10, the distribution of $`\stackrel{~}{\chi }^2`$ is approximately obtained by assuming that $`b_{lm}`$โs obey $`G(Z,1)`$ (actually, the distribution functions of $`b_{lm}`$โs are slightly much similar to the Gaussian distributions on large angular scales). The two non-Gaussian features are attributed to the nature of the distribution functions of each $`b_{lm}`$ which give large values at $`b_{lm}0`$ and decrease slowly at $`b_{lm}>>1`$ compared with the Gaussian distributions.
The slow decrease of the distribution of $`\stackrel{~}{\chi }^2`$ is important in discriminating the non-Gaussian models with the Gaussian models. As shown in Fig.11, observing $`\stackrel{~}{\chi }^250`$ are not improbable for the Weeks $`\mathrm{\Omega }_0`$ model ($`l=15`$) whereas it is almost unlikely for the Gaussian model. Because the distribution is slowly decreased for large $`\stackrel{~}{\chi }`$, the cosmic variances $`(\mathrm{\Delta }C_l)^2`$ are expected to be larger than that of the Gaussian models.
From Fig.12, on large angular scales($`2l15`$), one can see that the standard deviations $`\mathrm{\Delta }C_l`$ of $`\widehat{C_l}`$ in the two CH models are approximately 1-2 times of that for the Gaussian models.
## IV TOPOLOGICAL QUANTITIES
Topological measures:total area of the excursion regions, total length and the genus of the isotemperature contours have been used for testing Gaussianity of the temperature fluctuations in the COBE DMR data. Let us first summarize the known results for Gaussian fields (see ).
The genus $`G`$ of the excursion set for a random temperature field on a connected and simply-connected 2-surface can be loosely defined as
$`G`$ $`=`$ number of isolated high-temperature connected regions (41)
$``$ $`\text{number of isolated low-temperature connected regions}.`$ (42)
For instance, for a certain threshold, a hot spot will contribute $`+1`$ and a cold spot will contribute $`1`$ to the genus. If a hot spot contains a cold spot, the total contribution to the genus is zero. The genus which is the global property of the random field can be related to the integration of the local properties of the field. From the Gauss-Bonnet theorem, the genus of a closed curve $`C`$ being the boundary of a simply-connected region $`\mathrm{\Omega }_C`$ which consists of $`N`$ arcs with exterior angles $`\alpha _1,\alpha _2,\mathrm{}\alpha _N`$ can be written in terms of the geodesic curvature $`\kappa _s`$ and the Gaussian curvature $`K`$ as
$$G=\frac{1}{2\pi }\left[_C\kappa _g๐s+\underset{i=1}{\overset{N}{}}\alpha _i+_{\mathrm{\Omega }_C}K๐A\right].$$
(43)
For a random field on the 2-dimensional Euclidean space $`E^2`$ where the N arcs are all geodesic segments (straight line segments), $`K`$ and $`\kappa _g`$ vanish. Therefore, the genus is written as
$$G_{E^2}=\frac{1}{2\pi }\underset{i=1}{\overset{N}{}}\alpha _i.$$
(44)
The above formula is applicable to the locally flat spaces such as $`E^1\times S^1`$ and $`T^2`$ which have $`E^2`$ as the universal covering space since $`K`$ and $`\kappa _g`$ also vanish in these spaces. In these multiply-connected spaces, the naive definition Eq.(42) is not correct for excursion regions surrounded by a loop which cannot be contracted to a point.
In order to compute the genus for a random field on a sphere $`S^2`$ with radius equal to 1, it is convenient to use a map $`\psi `$:$`S^2\{p_1\}\{p_2\}S^1\times (0,\pi )`$ defined as
$$\psi :(\mathrm{sin}\theta \mathrm{cos}\varphi ,\mathrm{sin}\theta \mathrm{sin}\varphi ,\mathrm{cos}\theta )(\varphi ,\theta ),0\varphi <2\pi ,0<\theta <\pi ,$$
(45)
where $`p_1`$ and $`p_2`$ denote the north pole and the south pole, respectively. Because $`S^1\times (0,\pi )`$ can be considered as locally flat spaces ($`\varphi ,\theta `$) with metric $`ds^2=d\theta ^2+d\varphi ^2`$ which have boundaries $`\theta =0,\pi `$, the genus for excursion regions that do not contain the poles surrounded by straight segments in the locally flat ($`\varphi ,\theta `$) space is given by Eq.(44). It should be noted that the straight segments do not necessarily correspond to the geodesic segments in $`S^2`$. If a pole is inside an excursion region and the pole temperature is above the threshold then the genus is increased by one. If the pole temperature is below the threshold, it does not need any correction. Thus the genus for the excursions is
$$G_{S^2}=\frac{1}{2\pi }\underset{i}{}\alpha _i+N_p,$$
(46)
where $`\alpha _i`$ is the exterior angles at the intersection of two straight segments in the ($`\varphi ,\theta `$) space and $`N_p`$ is the number of poles above the threshold.
Now consider an isotropic and homogeneous Gaussian random temperature field on a sphere $`S^2`$ with radius 1. Let $`(x,y)`$ be the local Cartesian coordinates on $`S^2`$ and let the temperature correlation function be $`C(r)=<(\mathrm{\Delta }T/T)_0(\mathrm{\Delta }T/T)_r>`$ with $`r=x^2+y^2`$ and $`C_0=C(0)\sigma ^2`$, where $`\sigma `$ is the standard deviation and $`C_2=(d^2C/dr^2)_{r=0}`$. Then the expectation value of the genus for a threshold $`\mathrm{\Delta }T/T=\nu \sigma `$ is given as
$$<G_{S^2}>=\sqrt{\frac{2}{\pi }}\frac{C_2}{C_0}\nu \text{e}^{\nu ^2/2}+\text{erfc}\left(\frac{\nu }{\sqrt{2}}\right),$$
(47)
where erfc($`x`$) is the complementary error function. The first term in Eq.(47) is equal to the averaged contribution for the excursions which do not contain the poles while the second term in Eq.(47) is the expectation value of $`N_p`$.
The mean contour length per unit area for an isotropic homogeneous Gaussian random field is
$$<s>=\frac{1}{2}\left(\frac{C_2}{C_0}\right)^{\frac{1}{2}}\text{e}^{\nu ^2/2},$$
(48)
and the mean fractional area of excursion regions for the field is the cumulative probability of a threshold level,
$$<a>=\frac{1}{2}\text{erfc}\left(\frac{\nu }{\sqrt{2}}\right),$$
(49)
which gives the second term in Eq.(47).
As in Sec. III, the CMB anisotropy maps for the two CH adiabatic models are produced by using eigenmodes $`k<13`$ and angular components $`2l20`$ for $`\mathrm{\Omega }_0=0.2`$ and $`0.4`$. The contribution of higher modes are approximately 7 percent and 10 percent for $`\mathrm{\Omega }_0=0.2`$ and $`0.4`$, respectively. The initial power spectrum is assumed to be the extended Harrison-Zelโdovich spectrum. The beam-smoothing effect is not included. For comparison, sky maps for the Einstein-de-Sitter model with the Harrison-Zelโdovich spectrum $`C_l1/(l(l+1))`$ are also simulated.
In order to compute the genus and the contour length for each model, 10000 CMB sky maps on a 400$`\times `$200 grid in the $`(\varphi ,\theta )`$ space are produced. The contours are approximated by oriented straight segments. The genus comes from the sum of the exterior angles at the vertices of the contours and the number of poles at which the temperature is above the threshold. The total contour length is approximated by the sum of all the straight segments. Typical realizations of the sky map are shown in Fig.13.
Fig.14 and Fig.15 clearly show that the mean genuses and the mean total contours for the two CH models are well approximated by the theoretical values for the Gaussian models. This is a natural result since the distribution of the expansion coefficients $`b_{lm}`$ is very similar to the Gaussian distribution in the modest range. On the other hand, at high and low threshold levels, the variances of the total contour lengths and the genuses are larger than that for the Gaussian models that can be attributed to the nature of the distribution function of $`b_{lm}`$. One can easily notice the non-Gaussian signatures from Fig.16 and Fig.17. The excess variances for the Weeks model $`\mathrm{\Omega }_0=0.4`$ compared with the Gaussian flat Harrison-Zelโdovich model are observed at the absolute threshold level approximately $`|\nu |>1.4`$ for genus and $`|\nu |>0.6`$ for total contour length. If one assumes that the initial fluctuations are given by $`(\mathrm{\Phi }_\nu (0))^2\nu (\nu ^2+1)`$, the temperature fluctuations for CH models can be described as Gaussian pseudo-random fields. One can see from Fig.18 that the behavior of the variances of genus and total contour length for the Gaussian CH models is very similar to that for the flat Harrison-Zelโdovich model and the variances at high and low threshold levels are considerably smaller than that for the non-Gaussian models.
Because the mean behavior for the two non-Gaussian CH models is well described by the Gaussian models, the COBE DMR data which excludes grossly non-Gaussian models cannot constrain the two CH models by the topological measurements. However, one should take account of a fact that the signals in the $`10^o`$ smoothed COBE DMR 4-year sky maps are comparable to the noises that makes it hard to detect the non-Gaussian signals in the background fluctuations. In fact, some recent works using different statistical tools have shown that the COBE DMR 4-year sky maps are non-Gaussian although some authors cast doubts upon the cosmological origin of the observed non-Gaussian signals . Thus the evidence of Gaussianity in the CMB fluctuations is still not conclusive.
## V CONCLUSION
In this paper, Gaussianity of the eigenmodes and non-Gaussianity in the CMB temperature fluctuations in two smallest CH(Weeks and Thurston) models are investigated. As shown in Sec. II, it is numerically shown that the expansion coefficients of the two CH spaces behave as if they are random Gaussian numbers at almost all the places. If one recognizes the Laplacian as the Hamiltonian of a free particle, each eigenmode is interpreted as a wavefunction in a stationary state. The observed behavior is consistent with a prediction of RMT which has been considered to be a good empirical theory that describe the statistical properties of quantum mechanical systems whose classical counterparts are strongly chaotic. However, as we have seen, the global symmetries in the system can veil the generic properties. For instance, some eigenmodes on the Thurston manifold have a $`Z2`$ symmetry at a point where the injectivity radius is maximal. For these eigenmodes, the expansion coefficients are strongly correlated;hence they can no longer considered to be random Gaussian numbers.
Because the eigenmodes actually satisfy the periodic boundary conditions, there are points on a sphere $`S^2`$ which are identified with different points on $`S^2`$. These points form pairs of circles which are identified by the periodic boundary conditions . If one could identify all the circles on a sphere, one would be able to construct the corresponding CH space . Similarly, if one could identify all the fixed points and the corresponding symmetries, one would be able to construct a CH space which have these symmetries. The observed โrandomnessโ in the eigenmodes is actually determined by these simple structures.
In order to understand the symmetric structures of the CH spaces, it is useful to choose an observing point (base point) at which one enjoys symmetries as many as possible. However, in reality, there is no natural reason to consider fluctuations at only these particular points since the CH spaces are globally inhomogeneous.
Since the CMB fluctuations can be written in terms of a linear combination of eigenmodes, the fluctuations in CH models are almost spatially โisotropicโ if averaged all over the space except for very limited places at which the eigenmodes have certain symmetries provided that the eigenmodes are Gaussian. The spatial โisotropyโ implies that the contribution of non-diagonal terms in the two-point correlation functions are negligible. Thus the validity of the statistical tests using the angular power spectrum $`C_l`$ cannot be questioned on the ground that the background space is anisotropic at a certain point.
If one assumes that the initial fluctuations are Gaussian as in the standard inflationary scenarios, the temperature fluctuations are described by isotropic non-Gaussian random fields since they are written in terms of a sum of products of two independent random Gaussian variables, namely the initial perturbations and the expansion coefficients of the eigenmodes. The distribution functions of the expansion coefficients $`b_{lm}`$ for the sky maps at large values are slowly converged to zero than the Gaussian distribution with the same variance and the cosmic variances are found to be larger than that of the Gaussian models.
The increase in the variances are much conspicuous for topological quantities at large or small threshold levels. On the other hand, the mean behavior is well approximated by the Gaussian predictions. Therefore, the obtained results agree with the COBE DMR 4-year maps analyzed in. In real observations one has to tackle with what obscure the real signals such as pixel noises, galactic contaminations, beam-smoothing effect and systematic calibration errors which have not been considered in this paper. The absence of large deviations from the mean values at large or small threshold levels in the current data may be due to these effects, which will be much explored in the future work.
Although the recent observations seem to prefer the flat FRW models with the cosmological constant, the evidence is not perfectly conclusive. If one includes the cosmological constant for a fixed curvature radius, the radius of the last scattering surface (horizon) at present in unit of curvature radius becomes large. Therefore the observable imprints of the non-trivial topology of the background space become much prominent. For instance, the number $`N_f`$ of copies of the fundamental domains inside the last scattering at the present slice is approximately 27.9 for a Weeks model with $`\mathrm{\Omega }_\mathrm{\Lambda }=0.6`$ and $`\mathrm{\Omega }_m=0.2`$ whereas $`N_f=4.3`$ if $`\mathrm{\Omega }_\mathrm{\Lambda }=0`$ and $`\mathrm{\Omega }_m=0.8`$. Thus we have still great possibilities in detecting the non-trivial topology by the future satellite missions such as MAP and PLANCK which will provide us much better information on the statistical properties of the real signals. The large deviations of the topological quantities from the mean values would be the good signals that indicate the hyperbolicity (negative curvature) and the finiteness (smallness) of the universe in addition to the direct observation of the periodic structures peculiar to each non-trivial topology (see for recent developments).
## Acknowledgments
I would like to thank Jeff Weeks, Makoto Sakuma, Michihiko Fujii, and Craig Hodgson for answering many questions about symmetric structures of compact hyperbolic 3-spaces and topology of 3-manifolds. I would also like to thank N.J. Cornish, Naoshi Sugiyama and Kenji Tomita for their informative comments. The numerical computation in this work was carried out at the Data Processing Center in Kyoto University and Yukawa Institute Computer Facility. K.T. Inoue is supported by JSPS Research Fellowships for Young Scientists, and this work is supported partially by Grant-in-Aid for Scientific Research Fund (No.9809834). |
warning/0002/hep-th0002059.html | ar5iv | text | # References
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warning/0002/astro-ph0002423.html | ar5iv | text | # Photographic photometries and astrophysical parameters of the open clusters NGC1750 and NGC1758 This project is supported partly by the National Natural Science Foundation of China with Grant No. 19673012, 19603003 and 19733001, and in part by the astronomical foundation of Astronomical Committee of CAS.
## 1 Introduction
Open clusters, as systems of stars having a common origin, provide a very powerful tool on studying stellar evolution history. The homogeneity of photometric characteristics of stars in a cluster and their dynamics indicate that the cluster stars should have formed from one and the same primordial cloud within a relatively short time scale. Therefore, almost all member stars in a cluster should have roughly same age and chemical composition. Furthermore, open clusters can be used to understand both the formation and kinematics of the Galactic disk due to their wide age and mass distributions. For these reasons, open clusters constitute one of the most important research fields in observational and theoretical astronomy. Ages, distances, masses, luminosity functions and mass functions of open clusters are their basic researches.
The region of the open cluster NGC1750 has been paid more and more attention mainly because it is a complex area including two open clusters (NGC1750 and NGC1758) and the Taurus dark clouds, which located in the anticenter direction of the Galaxy. Cuffey(1937) was the first one to study this area systematically, obtaining extensive photographic photometry of stars in blue and red photometric bands to a limiting magnitude about $`R14`$mag. However he did not clearly point out the existence of two open clusters in this region. He considered the whole area as NGC1746 with distance about 590pc and BโR $``$0.30mag.
More recently, Galadรญ-Enrรญquez et al. (1998a, hereafter GJTR; 1998b; 1998c) did a series of studies of the open clusters NGC1750 and NGC1758, in which UBVRI CCD photometry of 3224 stars within the 45$`\times `$45 area was presented. At the same time, they combined different plates from several observatories to obtain proper motions for 45036 stars, and BVR photographic photometry of 39762 stars within 2.3$`\times `$2.3 in this area was completed down to $`V\mathrm{\hspace{0.17em}18.5}`$mag. Several physical parameters for NGC1750 and NGC1758 were discussed, including their positions, sizes, density profiles, extinctions, distances, ages, luminosity functions and masses, etc. Tian et al. (1998, hereafter TZSS) obtained high-precision proper motions and membership probabilities for 540 stars within 1.5$`\times `$1.5 in this area using 20 plates taken over a period of up to 68 years from Zว-Sรฉ 40cm Astrograph in Shanghai Astronomical Observatory, Chinese Academy of Sciences. In that work the two open clusters, NGC1750 and NGC1758, are successfully separated from each other. The core radii of the clusters have been estimated to be 17.20 and 2.25 respectively.
In this paper, we will present the results of photographic photometry of 789 stars to a limit of $`V16`$ mag within a 35$`\times `$35 area round the center of NGC1750. Combining the results of proper motions and membership probabilities for individual stars, we will more deeply study some astrophysical parameters of the two open clusters, including their H-R diagrams, distances, ages, masses, mass functions and kinematics. The program is as follows. In sec. 2, we will introduce the data reduction and compare with previous work. Detailed discussion on the basic astrophysical researches are presented in Sec.3. Summary is arranged in the final section.
## 2 Photographic Photometry
### 2.1 Observations, measurements and data reduction
The photographic photometries in the $`B`$ and $`V`$ bands were obtained from plates taken during 1992-1993 with the 1.56 meter reflecting telescope of Shanghai Astronomical Observatory, Chinese Academy of Sciences. The plates cover an area of about 35$`\times `$35, centered at $`\alpha _{2000}`$=5<sup>h</sup>3<sup>m</sup>30<sup>s</sup>, $`\delta _{2000}`$=2344, which includes the clusters of NGC1750 and NGC1758. The standard plate/filter combinations are IIaO+GG385 for the $`B`$ band and 156-01+GG495 for the $`V`$ band. The size of individual plates is 160mm$`\times `$160mm with a scale of 13<sup>โฒโฒ</sup>.25/mm, and the exposure time of individual plates is 30 minutes except for two plates, which have exposures of only 16 minutes. All plate materials are listed in Table 1, in which the first column denotes the plate number; the second column gives the epoch; filter and emulsion are shown respectively in Column 3 and 4; Column 5 is the exposure time; and the last two columns are the number of standard stars adopted and the reduced residual for each plate respectively (see below).
The density measurements of the whole areas of individual plates were done on the Photometric Data Systems (PDS) model 1010 automatic measuring machine at the Dominion Astrophysical Observatory (DAO) of Canada. The reductions of $`B`$,$`V`$ magnitudes were carried out according to the method presented by Stetson (1979). Because the range of candidate stars on each plate covers more than 7 magnitudes, the images of bright stars are generally saturated when the densities of faint stars can be measured. The fits of Gaussian point-spread-functions (PSF) to the unsaturated stars are good. However, the difference between the real PSF and Gaussian distribution of saturated stars are obvious, especially for stars in the cluster center region. Fortunately, Stetson (1979) has already concerned these two kinds of situations. The photographic magnitude index $`\mu `$ can be defined by the following equation
$$\mu =constant5\mathrm{Log}(D),$$
(1)
where $`D`$ is a measure of the total density of each star on a plate.
The magnitude indices for each band are averaged together among all plates, and the individual plates are transformed to this average plate by the least-squares method. At the same time, these fittings permit the standard error of a magnitude index on a typical plate to be estimated from the transformation residuals.
The Johnson-system BV photometric magnitudes are obtained by means of 141 secondary standard stars in this region with BV CCD photometry presented by GJTR before their paper was published. The magnitudes of these stars are all brighter than magnitude 14 in both the $`B`$ and $`V`$ bands. In addition, we choose about 40 standard stars fainter than magnitude 14 from Table 5 of GJTR, which have a homologous distribution on our plates. Our photographic magnitude transformation equations are cubic polynomials in $`B`$ and $`V`$, with a linear color term in each band. Generally, the candidate stars of the photometric standards should obey two principles (see Shu et al. 1998 for details): (1) the photometric standards should be distributed as homogeneously as possible along the magnitude interval covered on the plates; (2) they should be well isolated. The number of the standard stars used in the transformation of each plate and the residual rms are listed in Column 6 and 7 of Table 1 respectively.
The total number of stars we measured in the present study is 789. Their $`B`$ and $`V`$ magnitudes extend down to about 16.5 mag with accuracies estimated by both averaging the individual measurements of each star and considering the weighted internal error of the profile fitting. Because the centers of the available photographic plates are not exactly the same, some stars with large fitting errors must be discarded. At the same time, it is worthy noting that BV magnitudes of some stars are taken from only one or two plates. As a result, all 789 stars have $`V`$ magnitudes, with 540 stars among them measured on three or four plates, and only 653 stars have $`B`$ magnitudes, among which 481 were measured on three or four plates. The final accuracy for the stars which are measurable on at least three plates is given as a function of apparent visual magnitude in Table 2 and illustrated in Figure 1. One can see that average accuracies are almost the same for different magnitudes except for the faint end, which show larger scatters.
The accuracies of the derived magnitudes depend on the film of the plates and its homogeneity, the magnitude range of the standard stars used in the reductions, and the uncertainty resulting from transforming to the standard system. In general, it is difficult to obtain an accuracy of photographic photometry better than 0.1 mag. In the present study, it must be pointed out that the average accuracies of V magnitudes of the 540 stars and of B magnitudes of the 481 stars, which are measured at least three as mentioned above, are $`\pm 0.070`$mag and $`\pm 0.077`$mag, respectively. The reasons are: (1) most of our standard stars have B and V CCD data with very high accuracies; (2)the standard stars we adopted have a homogeneous distribution in both position and magnitude; (3)The method we chosen is reasonable.
The final reduced photometric results for individual stars in the region of NGC1750 and NGC1758 are given in Table 3, which is available only in electronic form. We present a small part of Table 3 here as an example. Column 1 in Table 3 is the ordinal star number in order of increasing right ascension; Column 2 and column 3 present the equatorial coordinates of J2000; Column 4-6 and 7-9 are $`V`$ and $`B`$ magnitudes with corresponding standard error and the number of measured plates respectively. The 10th column lists the identification of TZSS. The next two columns, 11 and 12, show the membership probabilities of individual stars in NGC1750 and NGC1758 taken from TZSS. The cross-identification with GJTR (their ordinal star number) is given in the last column.
### 2.2 Comparisons
Here we estimate the external error of our $`V`$ and $`B`$ magnitudes through comparison with GJTRโs CCD photometry. There are 448 and 347 common stars with $`V`$ and $`B`$ magnitudes respectively between our Table 3 and GJTRโs photometry catalogue. Their differences $`\mathrm{\Delta }V`$ and $`\mathrm{\Delta }B`$ in $`V`$ and $`B`$ bands as a function of magnitude for these stars are listed in Table 4 and shown in Figure 2. The mean differences in $`V`$ and $`B`$ magnitude are both $`\pm 0.045`$mag. It can be found clearly that this mean difference is small than the internal accuracy of stars available on at least three plates. It must be emphasized that this difference is not the external accuracy. It is because that our result is reduced from that of GJTR which has a very good accuracy. Furthermore, this also implies that the reduced method we adopted is reasonable. In fact, the mean difference of the 540 and 481 stars available on at least three plates in $`B`$ and $`V`$ band respectively are larger than above we estimated in Table 2. It is consistent with the normal principle that the external accuracy must be worse than the internal accuracy. Thus the difference seems better
In Figure 3 we compare the proper motions in the $`x`$ and $`y`$ directions for 517 stars common between our and GJTRโs proper motion catalogs. The two astrometric results are obtained from different plate sets, reference stars of proper motions and reduce method. It is found clearly that a fairly good linear relation exists in both components of the proper motions, but the slope is not unity. Although the proper motion in the present work is slightly larger than that of GJTR in both directions, which is due to the different reference frames adopted, it can be concluded that the two sets of results are consistent with each other, and there will be no significant difference in the results of membership determination because of the linear transformation.
## 3 Physical Parameters of NGC1750 and NGC1758
In order to investigate the basic astrophysical parameters of these two clusters, we must construct samples of members with positions, kinematics, membership probabilities and photometries available for individual stars. In the present work, there are 504 stars with both $`B`$, $`V`$ photometry and proper motion data, for which positions, proper motions and membership probabilities can be obtained from our previous work (TZSS). $`B`$ and $`V`$ magnitudes for 238 stars among them are taken from Table 3 and those for the remaining 266 are taken from the photometries done by GJTR. The sums of the membership probabilities for these 504 stars belong to NGC1750 and NGC1758 are 314 and 28 respectively. Meanwhile, the numbers of stars with membership probabilities higher than 0.7 for NGC1750 and NGC1758 are 311 and 23. We reasonably choose these 311 and 23 stars as our selected samples to analyze the CM diagram of the two clusters, in order to obtain the distances, ages and the kinematics of the individual clusters. On the other hand, all 504 stars are used to study the luminosity functions and mass functions (see below).
### 3.1 Color-magnitude diagrams, distances and ages
The CM diagram offers a powerful diagnostic of the evolutionary state of an open cluster. Because the locations of open clusters tend to be close to the plane of the Galaxy, their CM diagrams are liable to be heavily contaminated by unrelated field stars, and some caution must be taken into account to minimize this contamination by selecting stars on the basis of their kinematics, or by selecting only those stars with colors that are consistent with objects that have been reddened by the dust between us and the cluster. Figure 4 shows the observed CM diagrams of NGC1750 and NGC1758, respectively, based on the sample described above. The dots denote stars with $`B`$ and $`V`$ taken from the present work and the open circles denote the stars with photometries from GJTR. It can be seen that both observational color-magnitude diagrams show fairly clear main sequences. Toward the bottom of the diagram, the main sequence becomes boarder for NGC1750, with a width too large to be attributed to observational errors, which reflects the containination of field stars
In general, we do not know the cluster distances, so we cannot plot the CMDs on an absolute-magnitude scale. However, most cluster sizes are sufficiently small relative to their distance, so we can assume as usual that all stars belonging to a cluster lie at the same distance. To reduce contamination of field stars as much as possible, we trace the CM diagrams obtained by selecting stars with membership probabilities larger than 0.90, as shown in Fig. 5. It can be found that Fig.5 is much tighter than Fig.4. A careful discussion of the color excess due to dust absorption in front of these two clusters has been presented by GJRT. They found that the interstellar medium is relatively transparent toward the two clusters and the Johnson color excess for both of them is $`\overline{E(BV)}=(0.34\pm 0.07)`$mag, which corresponds to an extinction value of $`\overline{A_\mathrm{v}}=(1.1\pm 0.2)`$ mag. Based on these observational properties and the empirical ZAMS (Mermilliod 1981; Schaller et al. 1992), we can derive the distance modular of 8.60mag for NGC1750 and 9.50mag for NGC1758, which correspond to the distances of ($`525\pm 48`$pc) for NGC1750 and ($`794\pm 73`$pc) for NGC1758 with core radii of 2.6 pc and 0.5 pc respectively(TZSS).
The age distribution of open clusters plays an important role in many astrophysical researches, which can be used to estimate the lower limit for the age of the Galactic disk (Grenon, 1989), investigate the formation and evolution, especially the star formation history of our disk, as well as its dynamics (Janes & Phelps, 1994; Shu et al 1996). There are various methods to estimate ages of open clusters, which can lead to a significant scatter of the results for individual cluster. This is because of the differences among isochrone fitting, conversion from theoretical to observed stellar parameters, and so on. The most popular method adopted up to today is the fitting of theoretical isochrones to the observed CM diagram. The age determination of NGC1750 and NGC1758 in the present study is relatively difficult because of the relatively small number of photographic plates and some bright stars over-saturated, i.e., it is difficult to determine the turn-off points of these two clusters. Another reason is the relatively poor precision of the photographic photometries. Here, the same as GJTR did, we assume that the brightest stars on the main sequences of Fig. 5 denote the turn-off points of these two clusters, which are to be compared to the isochrones. After comparing the observed color-magnitude diagrams with those of theoretical results given by Schaller et al. (1992) for solar metallicity, we get the estimated ages of $`1.5\times 10^8`$yr for NGC1750 and $`6.3\times 10^8`$yr for NGC1758, which are shown in Fig. 5
The lifetimes of main sequence stars as a function of their absolute visual magnitudes $`M_V`$ are also presented by Meynet et al (1993). The brightest star, for NGC1750, which is assumed to be on the MS, is at $`V=8.42`$ mag, which corresponds to an age of about $`4.1\times 10^8\mathrm{yr}`$. Similarly, the fact that the brightest star on the MS for NGC1758 has $`V=10.77`$mag leads to its age of $`7.9\times 10^8`$yr. According to the relation among stellar mass, lifetime and its $`M_V`$ given by Miller & Scalo (1979), we can also obtain $`1.6\times 10^8`$yr and $`7.8\times 10^8`$yr for the ages of NGC1750 and NGC1758, respectively. If the relation between $`M_V`$ and lifetime is chosen as that presented by Mermilliod (1981), the ages of $`3.6\times 10^8\mathrm{yr}`$ and $`9.2\times 10^8\mathrm{yr}`$ for NGC1750 and NGC1758 are inferred respectively. All adopted relations in present work are the average results. The main reasons for these different results are: (1) the different evolution tracks for stars resulted from different stars evolutionary model; (2) the different weight of metalliaities. Combining all these results, we get the average age estimations for NGC1750 and NGC1758 should be $`(2.2\pm 1.0)\times 10^8`$yr and $`(7.8\pm 1.2)\times 10^8`$yr, respectively.
### 3.2 Luminosity functions and mass functions
It is important to study luminosity functions (LFs) and mass functions (MFs) of individual open clusters because they can provide information about both the initial mass function (IMF) and cluster dynamical evolution. Conceptually, the simplest estimation of a cluster luminosity function is to count stars within the cluster. In order to reduce the contamination of field stars, the sum of starsโ membership probabilities in different magnitude bins is one of the best to determine the luminosity functions $`\mathrm{\Phi }(\mathrm{V})`$ for individual clusters, i.e.,
$$\mathrm{\Phi }(V)=\frac{\mathrm{\Sigma }P(i)}{\mathrm{\Delta }V},$$
(2)
where $`P_c(i)`$ is the membership probability of star $`i`$ within the magnitude range of V to $`V+\mathrm{\Delta }V`$. Table 5 and Fig. 6 show the LFs for NGC1750 and NGC1758, respectively. One can see that there exists a peak for either clustersโ LFs, which to some extent reflects the complete magnitudes of the samples. The LF of the core region, which is within the center of $`2.6pc`$, of NGC1750 is also given in Fig. 6 as a dotted line. We did no do the same thing for the NGC1758 because of its small number of member stars within the core. It is clear that the profiles of luminosity functions in the central and whole observed region for NGC1750 are quite similar, i.e., there does not exist obvious mass segregation for NGC1750. The fact that the dynamical relaxation has not undergone thoroughly is consistent with its relatively young age (see last subsection). Combining the observed luminosity functions derived above and the mass-luminosity relations for main sequence stars given by Miller and Scalo (1979), we can infer the present-day mass functions of these two clusters, and the results are listed in Table 6 and also shown in Fig. 7 respectively. Here the average masses in individual mass bins are weighted by membership probability, and $`\mathrm{\Sigma }P`$ is summed over the stars in each bin as we did for their LFs, i.e.
$$\mathrm{\Psi }(M/M_{})=\frac{\mathrm{\Sigma }P_i}{\mathrm{\Delta }(M/M_{})_i},$$
(3)
with $`(M/M_{})_i=\frac{\mathrm{\Sigma }P_i(M_i/M_{})}{\mathrm{\Sigma }P_i}`$, here $`M_i`$ is the star mass with membership probabilities $`P_i`$ in the mass bin $`\mathrm{\Delta }(M/M_{})`$.
The slopes of the present-day mass functions of the two clusters are obtained by the least-squares linear regression. The results are shown by means of Log-Log plots in Figure 7. The slopes are $`(1.85\pm 0.19)`$ and $`(1.18\pm 0.33)`$ with the correlation coefficients of 0.83 and 0.66 for NGC1750 and NGC1758 respectively.Both clusters show the negative slopes. This also implies that they have not suffered the dynamical relaxation, which is consistent with the previous results.
Furthermore, based on the M/L relation given by Miller & Scalo (1979), the observed masses in the cluster region can be estimated to be about 390 $`\mathrm{M}_{}`$ and 40$`\mathrm{M}_{}`$ for NGC1750 and NGC1758, respectively. Here, binary stars are not considered, so these results are probably underestimated.
### 3.3 The kinematics
We might hope that direct studies of the kinematics of stars in NGC1750 and NGC1758 would reveal the effect of mass and space segregation. A reliable method for studying the kinematics of open clusters is based on proper motions of the member stars, which are comparatively easy to be obtained. In our sample, the average accuracy of proper motions is $`0.67mas\mathrm{yr}^1`$(TZSS), which corresponds to 1.7 km s<sup>-1</sup> for the distance of NGC1750 and 2.5 km s<sup>-1</sup> for the distance of NGC1758. Considering the stars with membership probabilities greater than 0.70, we estimate the intrinsic proper motion dispersions based on all the stars in the sample using the method outlined by Sagar & Bhatt (1988). The dependences of the intrinsic velocity dispersions on stellar masses and distances from the cluster centers are listed in Table 7 and Table 8 for NGC1750 and NGC1758 respectively, where the radial distances of each star is measured from the centers of the two clusters determined by TZSS, and N denotes the star number we used. One can see in Table 7 that there is no statistically significant radial dependence of the intrinsic proper motion dispersion $`\sigma _\mu `$. The values of $`\sigma _\mu `$ for different radial shells are almost the same within their uncertainties. It can also be seen from Table 7 that the intrinsic velocity dispersions $`\sigma _\mu `$ of different mass groups are almost the same. This means that the present data provide litter evidences of mass segregations in these two young open clusters, which is consistent with the results we obtained above. Even so, we can find from Table 7 that in the core region of NGC1750 $`(r<20arcmin)`$, the intrinsic velocity dispersions of the stars with larger mass are smaller than the intrinsic velocities of the stars with smaller mass, there exists some degree of both space and velocity mass segregation in the center region of NGC1750, where the dynamical relaxation is easy to undergo. But it is not clear for NGC1758 due to its small number of member stars.
To gain information about the isotropy or anisotropy of the velocity distribution, the radial and tangential components $`\sigma _{\mu r}`$, $`\sigma _{\mu t}`$ of the intrinsic dispersions of proper motions as a function radius are calculated and listed in Table 9, where the units of radial distance $`r`$ and of proper motion dispersions $`\sigma _\mu `$ have been converted into pc and $`\mathrm{kms}^1`$ respectively. The computation has been made for stars in the NGC1750 and NGC1758 regions with membership probabilities higher than 0.7. It can be found for both clusters that the ratios $`\sigma _{\mu r}/\sigma _{\mu t}`$ fluctuate around unity, which implies the absence of any significant evidence of velocity anisotropy. On the other hand, because the number of member stars is small, we can not get any certainly statistical results for NGC1758, but at least we can conclude that except the center region of NGC1750, there is no obvious velocity mass segregation or spatial mass segregation among the member stars of NGC1750, which suggests that this young open cluster has not reached energy equipartition.
## 4 Conclusion
In the present paper, based on the proper motions, photometries and membership probabilities of individual stars in the region of NGC1750 (TZSS, GJTR), we investigate basic astrophysical properties for two dynamically independent open clusters, NGC1750 and NGC1758. After detailed discussions on the photometric data and membership probabilities, the analysis samples with star number of 311 and 23 for NGC1750 and NGC1758 are constructed respectively. Comparing ZAMS (Mermilliod 1981; Schaller et al. 1992), we obtain the distances for these two clusters of $`(525\pm 48)`$pc for NGC1750 and $`(794\pm 73)`$pc for NGC1758 with their extinction being considered. Furthermore, many methods for the age determination are adopted to estimate the average ages of $`(2.2\pm 0.6)\times 10^8\mathrm{yr}`$ and $`(7.8\pm 1.6)\times 10^8\mathrm{yr}`$ with the observed masses 390 M and 40 M for NGC1750 and NGc1758 respectively.
According to the results of membership determination, luminosity functions and mass functions are given at the same time. It can be concluded that there exist no significant mass segregation effects for both clusters, which is consistent with the fact that their dynamical relaxation have not undergone thoroughly.
Finally, the velocity distributions of member stars for these two clusters are also discussed. It is found that both clusters seem to be isotropy in velocity space. Moreover, it is worth nothing that our statistical results could not be enough certain for NGC1758 because of its small number of member stars observed.
## Acknowledgements
The present work is part supported under the National Natural Science Fundation of China Grant No. 19673012 and 19733001 and by the astronomical fundation of Astronomical Committee of CAS. This work is also supported in part under Joint Laboratory for Optional Astronomy of CAS. K.P.Tian and J.L.Zhao are grateful to The National Research Council of Canada, which supported the living expenses while they visited the Dominion Astrophysical Observatory. |
warning/0002/cond-mat0002220.html | ar5iv | text | # Thermodynamic and diamagnetic properties of weakly doped antiferromagnets
## I Introduction
Anomalous normal-state properties of superconducting cuprates imad have stimulated intense theoretical investigations of models of strongly correlated electrons describing the interplay between antiferromagnetic (AFM) ordering of reference (undoped) insulating substances and the itinerant character of charge carriers introduced by doping. For the understanding of superconductivity the most challenging regime is that of intermediate (optimum) doping. However, even the apparently simplest region of weak doping is not fully understood theoretically.
Recently, the attention in experimental and theoretical investigations of cuprates has been given to characterization and understanding of different doping regimes batl . In a simple picture, weak doping should correspond to the regime where properties vary linearly with the concentration of holes, i.e. one can deal with a semiconductor-like model where charge carriers (holes) are independent and well defined quasiparticles. This requires a nonsingular variation of thermodynamic quantities with doping. However, this scenario has been questioned near the metalโinsulator transition based also on numerical solutions for some model systems imad , e.g. the Hubbard model. Alternative possibilities include phase separation emer , quantum critical behavior soko or other instabilities at low doping. Still, singular behavior in a planar (2D) system is expected only at $`T=0`$, while $`T>0`$ should lead to a regular variation with doping.
Among the least understood properties of charge carriers in cuprates and correlated systems in general are those related to the coupling of their orbital motion to an external magnetic field. Evidently anomalous and not understood is the Hall constant in cuprates which reveals unusual temperature and doping dependence ong . Another quantity is the diamagnetic (orbital) susceptibility $`\chi _d`$, which for noninteracting electrons corresponds to Landau diamagnetism land and seems to be connected to the Hall response rojo . Anomalous paramagnetic-like variation with magnetic field has been noticed within the ground state of the $`t`$-$`J`$ model bera at low doping. Recent $`T>0`$ studies of a single hole within the $`t`$-$`J`$ model vebe confirm the existence of a paramagnetic regime at intermediate $`T`$, though the systems studied were quite small. Conclusive experimental results on diamagnetic susceptibility are lacking wals , since the orbital part appears quite hidden by other contributions, although it could be distinguished via the anisotropy.
The aim of this paper is to study the thermodynamic properties and orbital response of correlated electrons at finite temperature in the low-doping regime. Most numerical studies of the $`t`$-$`J`$ model have so far focused on the ground-state properties dago , employing exact diagonalization of small systems, projector Monte Carlo, and density matrix renormalization group whit (DMRG). Recently, the finite-temperature Lanczos method (FTLM) has been introduced, which allows insight into the statics and dynamics at $`T>0`$. In previous applications certain thermodynamic quantities have also been investigated as a function of doping. In this paper we focus on the low doping regime, where the method can be compared with the alternative approach, a novel adaptation of the worldline quantum Monte Carlo (QMC) cluster method ever which allows for the study of much larger systems at least for temperatures $`T>T_{}`$ below which the minus-sign problem sets in. Large systems are particularly important for the study of diamagnetic response which appears to be quite sensitive to finite size effects. In both cases, new ways of dealing with the magnetic field are introduced. Related QMC methods have been used to study nonmagnetic properties of the $`t`$-$`J`$ model, in an exploratory calculation for doped chains and for ladders with 1 and 2 holes ammo , in two dimensions at $`J0`$ with 1 or 2 holes brun , and for chains at finite $`J`$ in a background of no holes assa .
In the following, the planar $`t`$-$`J`$ model as a representative model for strongly correlated electrons and electronic properties of cuprates is studied,
$$H=t\underset{ij\sigma }{}(\stackrel{~}{c}_{j\sigma }^{}\stackrel{~}{c}_{i\sigma }^{}+\text{H.c.})+J\underset{ij}{}\left(\stackrel{}{S}_i\stackrel{}{S}_j\frac{1}{4}n_in_j\right),$$
(1)
where $`\stackrel{~}{c}_{i\sigma }^{}`$, $`\stackrel{~}{c}_{i\sigma }^{}`$ are fermionic operators, projecting out sites with double occupancy. To approach the regime of strong correlations close to the real situation in cuprates, $`J/t=0.4`$ is used in most numerical calculations. We also use $`k_B=\mathrm{}=1`$.
The paper is organized as follows. Section II of the paper is devoted to a brief introduction of both numerical techniques employed, QMC and FTLM. In Sec. III results for several thermodynamic properties in the low-doping regime are presented and discussed. Sec. IV is devoted to the discussion of the orbital susceptibility of the system.
## II Numerical Methods
Results are obtained independently by the worldline QMC method and the FTLM. Wherever possible, results of both methods for doped systems are compared and presented relative to the undoped Heisenberg AFM. For large enough systems we expect to reach a typical behavior in the low doping regime.
### II.1 Worldline quantum Monte Carlo method
The loop cluster algorithm (LCA) for the world-line QMC has been introduced by one of the present authors ever and recently adapted also to the $`t`$-$`J`$ model ammo ; brun .
We briefly describe the worldline representation of the quantum QMC. The Hamiltonian, Eq. 1, on a 2D square lattice can be split within the standard Trotter-Suzuki decomposition trot ; suzu into four parts $`H=H_1+H_2+H_3+H_4`$ consisting of mutually commuting terms. This is equivalent to the well known checkerboard decomposition of Hamiltonians in 1D. The partition function is
$`Z`$ $`=`$ $`\mathrm{Tr}e^{\beta H}=\underset{M\mathrm{}}{lim}\mathrm{Tr}[e^{\stackrel{~}{\beta }(H_1+H_2+H_3+H_4)}]^M=`$ (2)
$`=`$ $`\mathrm{Tr}[e^{\stackrel{~}{\beta }H_1}e^{\stackrel{~}{\beta }H_2}e^{\stackrel{~}{\beta }H_3}e^{\stackrel{~}{\beta }H_4}]^M+O(\stackrel{~}{\beta }^2)`$
$``$ $`{\displaystyle \underset{\varphi _1\mathrm{}\varphi _{4M}}{}}\varphi _{4M}\left|e^{\stackrel{~}{\beta }H_1}\right|\varphi _1\varphi _1\left|e^{\stackrel{~}{\beta }H_2}\right|\varphi _2\mathrm{}`$
$`\mathrm{}`$ $`\varphi _{4M1}\left|e^{\stackrel{~}{\beta }H_4}\right|\varphi _{4M},`$
where $`\stackrel{~}{\beta }=\beta /M`$ and $`\beta =1/T`$. The summation is taken over the complete orthonormal set of states $`|\varphi _i`$. Within each imaginary time step $`\stackrel{~}{\beta }`$ the time evolution operator is applied. Since the Hamiltonian is total spin conserving, we can track time evolution of a particular spin along its so called worldline (WL). Because of the cyclic property of the trace, the WLs are periodic in the imaginary time interval $`[0,\beta ]`$. The time evolution operator acts only on $`2\times 2`$ plaquettes and the weight of the configuration $`W(๐)`$ factorizes into a product of plaquette weights. The partition function
$$Z=\underset{๐}{}W(๐)=\underset{๐}{}\underset{p๐}{}W(p)$$
(3)
is formally that of a ($`2+1`$)-dimensional classical system. The thermal average of an observable $`๐ช`$ can be obtained by
$$๐ช=\frac{1}{Z}\underset{๐}{}W(๐)๐ช(๐).$$
(4)
Such thermal expectation values are calculated by means of Monte Carlo (MC) importance sampling, where a sequence of configurations $`๐_i`$ (Markov chain) is constructed, which obeys detailed balance and reproduces the correct Boltzmann distribution $`W(๐)/Z`$. Thermal expectation values now become simple averages
$$๐ช=\underset{K\mathrm{}}{lim}\frac{1}{K}\underset{i}{\overset{K}{}}๐ช(๐_i).$$
(5)
In practice, Monte Carlo runs are finite, $`K<\mathrm{}`$, leading to statistical errors which can be calculated from the standard deviation of partial data sets ever .
In standard local algorithms an update from one configuration $`๐`$ to another $`๐^{}`$ in the Markov chain represents a small local change of the WLs. Therefore, consecutive configurations are highly correlated, which drastically increases the necessary number of Monte Carlo steps. Such difficulties are overcome in the LCA ever which introduces global (nonlocal) stochastic updates that effectively reduce the correlations. In the LCA formulation also the continuous time limit $`\stackrel{~}{\beta }0`$ can be taken bear avoiding the second order systematic error of Eq. 2. For certain observables improved estimators can be easily constructed allowing a potential reduction of statistical errors. For more details we refer to the introductory paper ever .
The LCA has recently been adapted to the $`t`$-$`J`$ model. The update procedure is split into three substeps, allowing the application of the standard LCA for the $`S=1/2`$ antiferromagnetic Heisenberg model or for free fermions in all three cases. Within each substep, only updates between two of the possible three states ($``$, $``$, and hole $``$) are performed. For the weights of particular plaquettes and other technical details we refer to ammo .
In case of negative weights $`W(๐)<0`$, their magnitude $`|W(๐)|`$ is taken for construction of the MC procedure, since the negative $`W(๐)`$ cannot be taken as a probability. Eq. (4) becomes
$$๐ช=\frac{\text{sign}๐ช_{|W|}}{\mathrm{sign}_{|W|}},$$
(6)
where $`\mathrm{}_{|W|}`$ denotes the expectation value with respect to the absolute value of the weight. In systems with such a โsign problemโ, the average sign $`\mathrm{sign}_{|W|}`$ often becomes exponentially small with increasing system size and decreasing temperature $`T`$, leading to a blow up of statistical errors suzu .
Let us briefly comment on the origin of negative signs in the WL formulation of the $`t`$-$`J`$ model. In the system with no doped holes the only source of negative weights are plaquettes, where two opposite spins exchange their positions representing a spin flip. Because of the periodicity of WLs in time direction and the absence of holes spin flips always occur in even numbers, producing no net negative sign. In the pure Heisenberg model, this sign can also be transformed away by rotation of spins on one sublattice, resulting in all-positive plaquette weights ever .
For one hole doped into the AFM one would naively not expect a sign problem, e.g., in this case there is no sign in the exact diagonalization approach. Examining the particle WLs surrounding the hole WL one finds, however, that an exchange of two fermions can occur when $`t0`$ and $`J0`$, producing an odd number of spin flips, i.e. a negative sign, as can be seen schematically in a small $`2\times 2`$ system,
$$\left[\begin{array}{ccc}& & \overline{}\\ & & \\ & & \end{array}\right]\left[\begin{array}{ccc}& & \\ & & \\ & & \overline{}\end{array}\right]\left[\begin{array}{ccc}& & \\ & & \\ \overline{}& & \end{array}\right],$$
where a loop motion of the hole $``$ around the system and a consecutive spin flip $``$ reproduce the original configuration with two fermions $``$ and $`\overline{}`$ interchanged. Measuring the sign here reduces to spin flip counting.
For higher concentration of holes a more general expression for the sign of the configuration can be obtained. It links fermion WL ($`\text{perm}_f`$) and hole WL permutation ($`\text{perm}_h`$)
$$\mathrm{sign}(๐)=(1)^{\text{perm}_f}=(1)^{\text{perm}_h}\mathrm{sign}(W(๐)),$$
(7)
so that for low doping it is preferable to measure $`\text{perm}_h`$ rather than $`\text{perm}_f`$. The sign problem also complicates the use of the improved estimators since for every observable a separate algorithm must be devised.
To follow the emergence of the sign problem as well as the development of diamagnetic properties it is convenient to generalize the isotropic spin interaction term of the model Eq. (1) to an anisotropic one with general anisotropy parameter $`\gamma `$,
$$H_J=J\underset{ij}{}[\frac{\gamma }{2}(S_i^+S_j^{}+S_j^+S_i^{})+S_i^zS_j^z].$$
(8)
This modifies the pure spin substep ($``$, $``$) of the $`t`$-$`J`$ LCA so that otherwise independent loops are โfrozenโ together ever into clusters and updated stohasticaly.
The results for $`\mathrm{sign}`$ as a function of inverse temperature $`\beta t=1/k_BT`$ are presented in Fig. 1. Note that for a single hole in the system, $`N_h=1`$, the relevant temperature scale in the anisotropic case is $`T_{}\gamma J`$, i.e. there is no sign problem for $`\gamma =0`$. At $`TT_{}`$ the sign starts to deteriorate rapidly, as can be seen in Fig. 1, preventing the investigation of low temperature properties. As expected $`\mathrm{sign}`$ decreases by adding additional holes $`N_h>1`$. For this reason, within the doped $`t`$-$`J`$ model only chains and coupled chains ammo have been investigated by LCA so far, and in two dimensions the limit $`J0`$ with $`N_h=1,2`$ brun . Recently, though, a way around the sign problem has been proposed by calculating fermion propagators for a background of no holes assa .
For a fixed number of holes, the average sign will converge as the system size increases. This convergence could be taken as a another criterion that the limit of a dilute system has been reached.
### II.2 Finite-temperature Lanczos method
In the analysis of the $`t`$-$`J`$ model the exact diagonalization of small systems using the Lanczos algorithm has been extensively employed dago , predominantly in the investigation of the static and dynamic properties of the ground state. More recently a FTLM combining the Lanczos procedure and random sampling was introduced jakl1 ; jakl2 , allowing the calculation of $`T>0`$ static and dynamic properties of correlated systems. The application is particularly simple for an arbitrary function of conserved quantities, e.g. $`f(H,S_z)`$,
$`Z`$ $``$ $`{\displaystyle \frac{N_{\text{st}}}{K}}{\displaystyle \underset{n=1}{\overset{K}{}}}{\displaystyle \underset{i=0}{\overset{M1}{}}}e^{\beta E_i^n}|n|\psi _i^n|^2,`$ (9)
$`f`$ $``$ $`{\displaystyle \frac{N_{\text{st}}}{KZ}}{\displaystyle \underset{n=1}{\overset{K}{}}}{\displaystyle \underset{i=0}{\overset{M1}{}}}f(E_i^n,S_z^n)e^{\beta E_i^n}|n|\psi _i^n|^2,`$ (10)
where $`|\psi _i^n`$, $`E_i^n`$ are (approximate) eigenfunctions and energies, respectively, obtained by diagonalization within the reduced orthonormal set, generated from the initial functions $`|n`$ in $`M`$ Lanczos steps. $`N_{\text{st}}`$ is the dimension of the complete basis. $`K`$ initial functions $`|n`$ are chosen at random but with good quantum number $`S_z`$. Usually it is enough to choose $`M,KN_{st}`$. For a more detailed discussion of the method and results we refer to jakl2 .
It is expected that $`T>0`$ reduces the finite-size effects of the measured quantities. It is however important to realize that for a particular system, finite-size effects start to be pronounced at $`T<T_{\text{fs}}`$ where, e.g., some characteristic length-scale becomes larger than the system size. In our case of low doping, $`N20`$ and $`J=0.4t`$, we find $`T_{\text{fs}}0.4J`$. All our results are presented for $`T>T_{\text{fs}}`$ where $`Z(T_{\text{fs}})Z^{}`$. In the present study $`Z^{}=30`$ so that at least 30 states are sampled in the thermal averages jakl2 . It should be stressed that the FTLM gives also the correct ground state within the chosen small system.
## III Thermodynamic Properties
Thermodynamic properties $`๐ช(c_h)`$ depend on the hole concentration $`c_h=N_h/N`$. For the weak doping limit, one would expect a linear dependence for most quantities. For a finite system size, the relevant parameter is thus the number of holes $`N_h`$ doped into the AFM. In the low-doping regime it makes sense to represent the results as a difference
$$\mathrm{\Delta }๐ช_i=๐ช(N_h=i)๐ช(N_h=i1).$$
(11)
To distinguish the change in a particular quantity with doping, this notation is used in the following. If, e.g., $`\mathrm{\Delta }๐ช_2`$ behaves quantitatively as $`\mathrm{\Delta }๐ช_1`$ one can conclude that the quantity changes linearly with the number of added holes $`N_h`$, i.e. the holes behave as independent entities, and the system sizes are large enough so that the low-doping regime has indeed been reached. Such behavior is however not the only possibility at low doping, since one can expect, e.g., even-odd effects in the case of pairing of holes.
### III.1 Internal energy, specific heat and entropy
The internal energy, defined by
$$E=\frac{\beta F}{\beta }=\frac{1}{Z}\frac{Z}{\beta },$$
(12)
is calculated within FTLM as $`E`$ in Eq. (10) and in QMC as an expectation value of the corresponding operator, $`E=_๐E(๐)W(๐)`$. Results of both methods are presented in Fig. 2.
From Fig. 2 we first conclude that for $`\mathrm{\Delta }E_1`$ the results at $`N=20`$ (which overlap also with the results at $`N=18`$) obtained via the FTLM are essentially equivalent with QMC results for much larger lattices, at least for $`T0.2t`$ reached by the QMC method. We note also a close agreement between QMC $`\mathrm{\Delta }E_2`$ and $`\mathrm{\Delta }E_1`$, confirming the assumption of the low-doping regime and holes as independent quasiparticles in the $`T`$ window presented. In FTLM results, on the other hand, the difference between $`\mathrm{\Delta }E_2`$ and $`\mathrm{\Delta }E_1`$ is already visible, since $`N_h=2`$ here means already an appreciable doping $`c_h=0.1`$. For $`T=0`$ the difference $`\mathrm{\Delta }E_2\mathrm{\Delta }E_1`$ equals to the binding energy dago and in the continuum corresponds to the second derivative of the ground state energy with respect to the doping. In the chosen parameter regime ($`J=0.4t`$) the binding energy is negative and thus pointing to the attractive interaction between the holes. With the increase of the temperature the bound state disintegrates and the difference $`\mathrm{\Delta }E_2\mathrm{\Delta }E_1`$ approaches zero. In the case of a small system the difference can even become positive but vanishes with increasing system size.
It is also evident from Fig. 2 that $`\mathrm{\Delta }E(T)`$ is not a monotonous function. The ground state of a single hole introduced into the AFM is quite well understood via analytical approaches schm and numerical calculations dago . For $`J/t=0.4`$, the zero temperature result $`\mathrm{\Delta }E(0)1.44t`$ can be explained well by the interplay between the gain of the kinetic energy represented by the hopping term $`H_t`$ and the loss of local AFM correlation energy around the hole.
$`\mathrm{\Delta }E(T)`$ has not been considered so far. An interpretation of its behavior can be given as follows. Introducing a single hole into an AFM destroys the local AFM spin order and thus increases the exchange energy. The increase is however expected to disappear at $`T>J`$ where the spin system becomes disordered. On the other hand the ground-state kinetic energy in a disordered spin system is quite similar to the one in an AFM, hence the decrease of the internal energy $`\mathrm{\Delta }E`$ for $`T>J`$. This remains valid for $`T<t`$ where also higher hopping-related states become populated and finally $`\mathrm{\Delta }E(T\mathrm{})0`$, explaining turn back up for $`T0.7t`$.
The specific heat defined by
$$C=\frac{E}{T}=\beta ^2\left[\frac{1}{Z}\frac{^2Z}{\beta ^2}\left(\frac{1}{Z}\frac{Z}{\beta }\right)^2\right]$$
(13)
is obtained as $`\beta ^2[E^2E^2]`$ within the FTLM and the QMC method. The results are presented in Fig. 3.
The main effect of introducing holes into the AFM insulator on $`C`$ is to decrease the peak at $`TJ`$. This appears in $`\mathrm{\Delta }C`$ as a pronounced dip which slightly weakens and shifts its energy scale $`J`$ to lower values with doping, as can be seen from the line-shape in Fig. 3.
The entropy is
$$S=\beta (EF)=\beta E+\mathrm{ln}Z$$
(14)
We reconstruct it from the specific heat
$$C=T\frac{S}{T}$$
(15)
by numerical integration from high temperatures $`T\mathrm{}`$
$$\mathrm{\Delta }S(T)\mathrm{\Delta }S(\mathrm{})=_{\mathrm{}}^T\frac{\mathrm{\Delta }C}{T}\text{d}T$$
(16)
The high-temperature integration constants are chosen so that $`\mathrm{\Delta }S(\mathrm{})=\mathrm{\Delta }\mathrm{ln}N_{\text{st}}`$.
In contrast to $`\mathrm{\Delta }E`$ and $`\mathrm{\Delta }C`$ discussed previously, $`\mathrm{\Delta }S`$ is not linear in $`N_h`$ in the low doping limit. In analogy with low-concentration systems, like the dilute classical gas, entropy is expected to scale as $`\mathrm{\Delta }S|\mathrm{\Delta }(Nc_h\mathrm{ln}c_h)|`$, i.e. $`\mathrm{\Delta }S_1\mathrm{ln}N`$, so the change still depends explicitly on the system size. This is also realized in Fig.4, where $`\mathrm{\Delta }S_1`$ still varies with $`N`$, yet curves for different $`N`$ appear parallel down to the lowest reachable temperatures.
It is particularly remarkable how large the entropy increase $`\mathrm{\Delta }S_11`$ is even at the lowest $`T<J`$. This is indeed consistent with $`\mathrm{\Delta }S`$ measured in cuprates lora . While this could be attributed partly to the logarithmic dependence on $`c_h`$, at the same it is evident that the behavior is much closer to a system of classical particles than to a degenerate electron gas.
### III.2 Spin susceptibility
The uniform spin susceptibility can be evaluated as a thermodynamic quantity from
$$\chi _s=\frac{\beta S_z^2}{N},$$
(17)
where $`S_z=_iS_i^z`$ is the conserved total spin. In the FTLM then Eq. (10) can be applied, while within the QMC method $`\chi _s`$ is related to the number of spin up and down WLs.
It is instructive to present results both for $`\mathrm{\Delta }\chi _s`$ with respect to the undoped AFM, Fig. 5, as well as for the effective Curie constant (difference of the square moment) per hole $`\mathrm{\Delta }S_z^2=N\mathrm{\Delta }\chi _s/\beta `$ in Fig. 6.
The results in Figs. 5, 6 are easy to interpret for high $`T>t`$. Each hole introduced into the system reduces the effective Curie constant by one spin, i.e., $`\mathrm{\Delta }S_z^2=1/4`$. On the other hand, at low $`T<J`$ the situation is reversed since $`\mathrm{\Delta }\chi _s>0`$. This increase can be attributed to the relaxation of the AFM order by the hole doping. Note that in an AFM, $`\chi _s`$ achieves a maximum at $`TJ`$ while below that temperature it is reduced due to the longer range AFM order. It is interesting to note that at the lowest reachable temperature $`TJ/2`$ each hole effectively adds just one spin, i.e. $`\mathrm{\Delta }S_z^21/4`$.
## IV Orbital Susceptibility
In order to investigate the orbital response of the system, a homogeneous magnetic field $`B`$ perpendicular to the plane has to be introduced. When we discuss the orbital magnetization and susceptibility, $`B`$ enters only in the kinetic term of Eq. (1), via the Peierls construction
$$H_t=t\underset{ij\sigma }{}(e^{i\theta _{ij}}\stackrel{~}{c}_{j\sigma }^{}\stackrel{~}{c}_{i\sigma }^{}+\text{H.c.}),$$
(18)
where the phases are given within Landau gauge as
$$\theta _{ij}=\frac{e}{\mathrm{}}๐(๐ซ_i)๐_{ij},๐=B(0,x,0),$$
(19)
with $`๐_{ij}=๐ซ_j๐ซ_i`$. The relevant parameter for the strength of $`B`$ is the dimensionless flux per plaquette $`\alpha =2\pi Ba^2/\varphi _0`$, where $`a`$ is the lattice spacing and $`\varphi _0=h/e`$ is the unit quantum flux.
The dc orbital susceptibility of the system in the external magnetic field is
$`\chi _d`$ $`=`$ $`\mu _0{\displaystyle \frac{^2F}{B^2}}={\displaystyle \frac{\mu _0e^2a^4}{\mathrm{}^2}}{\displaystyle \frac{^2F}{\alpha ^2}}=`$ (20)
$`=`$ $`{\displaystyle \frac{\chi _0}{\beta }}\left[{\displaystyle \frac{1}{Z}}{\displaystyle \frac{^2Z}{\alpha ^2}}\left({\displaystyle \frac{1}{Z}}{\displaystyle \frac{Z}{\alpha }}\right)^2\right],`$
where $`\chi _0=\mu _0e^2a^4/\mathrm{}^2`$. So far $`\chi _d`$ has been investigated only for a single hole $`N_h=1`$ by high-temperature expansion at $`J=0`$, and by the FTLM for $`J>0`$ vebe . In the latter study it was realized that results are quite sensitive to finite-size effects, so it is desirable to get corresponding results also via the QMC method, where much larger lattices can be studied.
Let us first derive the expression for the orbital susceptibility within the QMC method. As seen from Eq. (18), the magnetic field affects only the hopping of the electrons. In the WL representation this concerns matrix elements $`\varphi \left|e^{\stackrel{~}{\beta }H}\right|\varphi ^{}`$ in Eq. (2). For the plaquette representing the hopping between sites $`i`$ and $`j`$ the matrix elements become
$$H_{ij}^{\text{hop}}=\left(\begin{array}{cc}0& te^{i\theta _{ij}}\\ te^{i\theta _{ij}}& 0\end{array}\right),$$
(21)
written in the $`|`$, $`|`$ base. The imaginary time propagator in the same base is
$$e^{\stackrel{~}{\beta }H_{ij}^{\text{hop}}}=\left(\begin{array}{cc}\mathrm{ch}\stackrel{~}{\beta }t& e^{i\theta _{ij}}\mathrm{sh}\stackrel{~}{\beta }t\\ e^{i\theta _{ij}}\mathrm{sh}\stackrel{~}{\beta }t& \mathrm{ch}\stackrel{~}{\beta }t\end{array}\right).$$
(22)
Thus the plaquette weights along the hole WL obtain an additional phase factor $`W(p,\alpha )=W(p)e^{i\theta (p)}`$. The weight of the whole configuration is a product, Eq. (3),
$$W(๐,\alpha )=\underset{p๐}{}W(p,\alpha )=W(๐)\underset{p๐}{}\mathrm{exp}\left(i\theta (p)\right)$$
(23)
where the phases sum up
$`\mathrm{exp}\left[i{\displaystyle \underset{p๐}{}}\theta (p)\right]=\mathrm{exp}\left[i{\displaystyle \theta (๐ซ)\text{d}๐ซ}\right]=`$
$`=\mathrm{exp}\left[i{\displaystyle \frac{e}{\mathrm{}}}{\displaystyle ๐(๐ซ)\text{d}๐ซ}\right]=e^{i\alpha ๐ฎ}.`$ (24)
The integral runs along the hole WL. Here $`๐ฎ`$ is defined as the oriented area of the hole WL projected onto the plane in units of the lattice plaquette area $`a^2`$. For more holes, $`๐ฎ`$ generalizes similarly to the sum of all hole WL areas.
Now we can write the partition function in the magnetic field as
$$Z=\underset{๐}{}W(๐,\alpha )=\underset{๐}{}e^{i\alpha ๐ฎ(๐)}W(๐).$$
(25)
For a given configuration $`๐`$ there always exist $`๐^{}`$ (imaginary time inversion) with the same weight but $`๐ฎ(๐^{})=๐ฎ(๐)`$, therefore the exponential in Eq. (25) can be replaced by a $`\mathrm{cos}`$ function. For $`B=0`$ we have $`๐ฎ=0`$ and obtain the zero-field susceptibility from Eq. (20),
$$\chi _d=\chi _0\frac{๐ฎ^2}{\beta }.$$
(26)
$`\chi _d`$ can be thus measured without the presence of a magnetic field. This is just another consequence of the more general fluctuationโdissipation theorem. Even though $`๐ฎ^2`$ is strictly positive, the thermal average in Eq. (26) can become negative because of correlations between the sign of the weight and the area $`๐ฎ`$ of the hole WL. From Eq. (6) we can deduce that $`๐ฎ^2<0`$ when the configurations with negative sign tend to have larger $`๐ฎ^2(๐)|W(๐)|`$ than configurations with positive ones.
The hole WL can obtain nonzero spatial winding number due to the periodic boundary conditions and small system size. E.g., the WL can run along the imaginary time, cross the system boundary, and complete time periodicity reconnecting with its spatially periodic image. In that case the area $`๐ฎ`$ as defined by Eq. (24) has no physical meaning. Therefore we restrict our simulation so that only the zero spatial winding loop updates are generated (a discussion on fixing the winding numbers can be found in hene ). The hole is allowed to cross the system boundary as long as it does not increase the winding number. The effect of the restriction is analogous to the movement of the hole doped into an infinite periodic spin background with a unit cell equal to the size of the system. The results for thermodynamic quantities presented in the previous section agree within the errorbars with the unrestricted case. This restriction is weaker than closed boundary conditions, resulting in smaller finite size effects.
The introduction of finite $`B>0`$ into the model, Eq. (18), reduces the translational symmetry and thus for a given system size increases the required minimal base set used in FTLM. In the present study a few mobile holes on a system of tilted squares with $`N`$ up to 20 sites and periodic boundary conditions are considered. It is nontrivial to incorporate phases $`\theta _{ij}`$ corresponding to a homogeneous $`B`$, being at the same time compatible with periodic boundary conditions frad ; vebe . This is possible only for quantized magnetic fields $`B=mB_0`$, where $`B_0=\varphi _0/N`$.
$`\chi _d`$ from Eq. (26) can be calculated in FTLM only by taking a numerical derivative of the free energy $`F=T\mathrm{ln}Z`$, with $`Z`$ from Eq. (9). Finite systems provide $`F(\alpha )`$ only for discrete values of $`\alpha `$. $`\chi _d`$ can be obtained by fitting the $`F(\alpha )`$ dependence to the parabolic form $`F(\alpha )=F(0)+c\alpha ^2`$ in two points $`\alpha =i2\pi /N`$ and $`\alpha =j2\pi /N`$. The corresponding results for the susceptibility are denoted by $`\chi _{ij}`$. In small systems ($`N<20`$) the introduction of $`B>0`$ can lift some zero-field degeneracies. The $`\chi _{01}`$ values are therefore systematically affected by larger finite-size effects. For the system with $`N=20`$ both values $`\chi _{01}`$ and $`\chi _{12}`$ agree quite reasonably with the QMC data.
Let us first comment on the validity of results for $`\chi _d`$. Since $`\chi _d`$ deals with an orbital current represented by loop motion of charge carriers (holes), it is much more sensitive to finite-size effects vebe than most other correlation functions. This was also the main motivation to employ the QMC method, where much larger systems can be reached. In Fig. 7, finite-size scaling is performed for the QMC data for $`\chi _d(T)`$ for the case of $`N_h=1`$. We can see that the different $`T`$ points do not cross upon changing the system size $`N`$. Thus, at least qualitatively, results do not depend on the system size. Therefore, the sizes of choice for the QMC systems will be $`6\times 6`$ and $`8\times 8`$.
In Fig. 8, $`\chi _d`$ obtained via both methods is presented. For $`Tt`$, the response is diamagnetic and proportional to $`T^3`$ as well as essentially $`J`$-independent vebe . The most striking effect is that the orbital response below some temperature $`T_p`$ turns from diamagnetic to paramagnetic, consistent with the preliminary results obtained via the FTLM vebe . In order to locate the origin of this phenomenon, results for different $`J`$ and anisotropies $`\gamma `$ are shown in Fig. 9. It appears that $`T_p`$ scales with $`\gamma J`$, i.e. at $`J=0`$ the response is clearly diamagnetic at all $`T`$, and for $`\gamma =0,J>0`$ no crossing is observed with either method.
At lower temperatures $`T<T_dT_p`$, the diamagnetic behavior is expected to be restored. This follows from the argument that at $`T0`$ a hole in an AFM should behave as a quasiparticle with a finite effective mass, exhibiting a cyclotron motion in $`B0`$. The latter behavior should lead to $`\chi _d(T0)\mathrm{}`$ vebe . Numerically it is easiest to test this conjecture for a single hole and $`\gamma =0`$ (also true for $`J=0`$). Namely, the QMC has no sign problem at $`\gamma =0`$, so that error bars are only due to the finite MC sampling.
Results in Fig. 9 are quite interesting even for $`\gamma =0`$. At $`J=0`$ a monotonous increase of $`|\chi _d|`$ is observed, diverging as $`T0`$ vebe . It can be explained as a gradual transition from a hole in a random spin background to a well defined quasiparticle, i.e. the ferromagnetic polaron, at $`T0`$. The situation is more complicated for $`J>0`$. It is plausible that the difference to the $`J=0`$ case shows up at $`T<J`$, where the AFM short-range correlations appear. Relative to $`J=0`$, a spin ordered state blocks the loop motion of holes, necessary for finite diamagnetic $`\chi _d`$. The effect is thus first a decrease of $`|\chi _d(T)|`$ with decreasing $`T`$, as seen in Fig. 9. Turnover to the diverging diamagnetic $`\chi _d`$ should happen only when the coherent quasiparticle is formed at $`T_dJ`$. $`T_d`$ should scale with the inverse of the quasiparticle effective mass $`1/m^{}`$. It is known that $`m^{}`$ can become very large for the extreme $`\gamma =0`$ case, in particular for larger $`J`$. This explains why we cannot reach the coherent regime for $`J=0.4t`$ even at $`T=0.2t`$ (Fig. 8), while the downturn is indeed observed for $`J/t=0.1`$, $`0.2`$, $`\gamma =0`$ (Fig. 9).
At $`\gamma >0`$, the results are qualitatively different. The most pronounced effect is the change into a paramagnetic $`\chi _d`$ for $`T<T_p`$. The width of this $`T`$ window is quite large. In fact within the FTLM and QMC data we are unable to locate the reentrance temperature $`T_d`$ into the diamagnetic response, although the latter is expected vebe . An argument for the low value of $`T_d`$ can be given in terms of a very shallow energy minimum which defines the quasiparticle dispersion near the ground-state of a hole in an AFM within the $`t`$-$`J`$ model dago , hence the quasiparticle looses its character already at very low excitation energies. Still the paramagnetic response in the window $`T_d<T<T_p`$ remains to be explained.
Let us finally discuss also results for $`\chi _d`$ for finite doping $`c_h>0`$, as presented in Fig. 10. The easiest regime to interpret is that of a nearly empty band, i.e. $`c_h>0.7`$, where $`\chi _d`$ is diamagnetic and nearly independent of $`T`$. In this regime the electron system is dilute and strong correlations are unimportant, hence Landau diamagnetism is expected. At moderate temperatures $`T>J`$ and for an intermediate-doping regime, $`0.2<c_h<0.7`$, $`\chi _d`$ is dominated by a paramagnetic response with a peak at approximately $`c_h=1/2`$. As consistent with results at low doping, there is a weak diamagnetism at $`c_h<0.2`$ and $`T>T_p`$, while the paramagnetic regime extends to $`c_h=0`$ for $`T<T_p`$. For low temperatures $`TJ`$ quite pronounced oscillations in $`\chi _d(c_h)`$ appear and can be partly attributed to finite-system effects.
Certain aspects of the above results for $`\chi _d(c_h)`$ can be understood using the HTE. One is the asymmetry between $`c_h0`$ and $`c_h1`$ at higher $`T>t`$. In the lowest order of the HTE, only hopping of the electrons around a basic plaquette loop has to be considered. The signal on the $`c_h0`$ side of $`\chi _d(c_h)`$ is thus reduced by a factor of $`2^3=8`$ against the $`c_h1`$ side since in the first case only the plaquettes with ferromagnetically aligned spins contribute.
## V Conclusions
Let us first compare both numerical methods used in the analysis of the $`t`$-$`J`$ model at low doping. The QMC method allows the studies of larger systems and the loop algorithm solves some serious drawbacks of the MC methods. It is indeed very efficient in cases where there is no sign problem, e.g. the anisotropic model $`\gamma =0`$ at $`N_h=1`$. Within the WL approach it is also very easy to formulate and measure certain responses like the orbital susceptibility $`\chi _d`$. Still the method suffers from a sign problem (for $`\gamma >0`$) even for a single hole $`N_h=1`$ in an AFM (though not in a background of no holes assa ). Results are thus in practice limited to $`TJ/3`$ for the isotropic case $`\gamma =1`$. On the other hand the FTLM has no minus-sign problem but rather limitations due to small systems which can be studied. These are even more pronounced in cases with $`B>0`$ where the translational symmetry is lost. It is an interesting observation that within the FTLM, the limiting temperature $`T_{fs}`$ for the $`t`$-$`J`$ model is in most cases quite close to the lowest $`T`$ reached by the QMC method.
We have presented several results for thermodynamic quantities, i.e. energy $`E`$, specific heat $`C`$, entropy $`S`$ and spin susceptibility $`\chi _s`$, as a function of $`T`$ at low doping. At $`T>T_{fs}`$ all results are consistent with the picture where holes introduced into the AFM behave as independent (nondegenerate) particles. The perturbation introduced into the AFM by holes is quite large even at lowest $`T<J`$, in particular visible from $`\mathrm{\Delta }C`$ and $`\mathrm{\Delta }S`$, consistent with experiments in cuprates lora .
Results for the orbital susceptibility $`\chi _d`$, now obtained also for much larger systems using the QMC method, confirm the preliminary FTLM results vebe , indicating an anomalous paramagnetic response at low doping in an intermediate window of temperatures $`T_d<T<T_pJ`$ (for the isotropic model $`\gamma =1`$). In fact within our numerical studies it is hard to reach the lower end of this window, meaning that $`T_d<J/3`$. Still, the reentrance into a diamagnetic behavior is expected from theoretical arguments on the existence of a well defined quasiparticle at $`T0`$, as well as from more reliable QMC results for the $`\gamma =0`$ case vebe . The paramagnetic response at intermediate $`T`$ can be viewed also as an extension of a more pronounced $`\chi _d>0`$ regime observed at finite doping $`0.2<c_h<0.5`$ at all $`T`$. The explanation can thus go in the direction proposed by Laughlin laug ; bera , that at low doping $`c_h0`$ we are dealing with quasiparticles (with a diamagnetic response), being a bound composite of charge (holon) and spin (spinon) elementary excitations. The binding appears however to be quite weak and thus easily destroyed by finite $`T`$ or $`c_h`$, enabling the independent response of constituents, which apparently is paramagnetic.
###### Acknowledgements.
The authors wish to thank I. Sega for helpful suggestions. This work was supported by the Ministry of Science and Technology of Slovenia under Project No. J1-0231. |
warning/0002/hep-ph0002023.html | ar5iv | text | # Strong and electromagnetic decays of p-wave heavy baryons ฮ_{๐โข1},ฮ^โ_{๐โข1}
## I Introduction
Now most of the ground state charm baryons have been found experimentally . Important progress has been made in the search of orbitally excited heavy baryons. The ARGUS , E687 and CLEO ) collaborations have observed a pair of states in the channel $`\mathrm{\Lambda }_c^+\pi ^+\pi ^{}`$, which were interpreted as the lowest lying orbitally excited states: $`\mathrm{\Lambda }_{c1}(2593)`$ with $`J^P=\frac{1}{2}^{}`$ and $`\mathrm{\Lambda }_{c1}^{}(2625)`$ with $`J^P=\frac{3}{2}^{}`$. The total decay width of the $`\mathrm{\Lambda }_{c1}(2593)`$ is $`3.6_{1.3}^{+2.0}`$ MeV while only an upper limit of $`<1.9`$ MeV has been set for $`\mathrm{\Lambda }_c^{}(2625)`$ up to now . Recently there emerges evidence for the $`\mathrm{\Xi }_{c1}^+`$ with $`J^P=\frac{3}{2}^{}`$, the strange partner of the $`\mathrm{\Lambda }_{c1}^{}(2625)`$. Its width is less than $`2.4`$ MeV. In the near future much more data will be expected. We will focus on the strong and electromagnetic decays of the $`\mathrm{\Lambda }_{c1}`$ doublet since they are the only well established states .
There exist many theoretical discussions on this topic. In the single pion and two pion strong decays and radiative decays of the $`\mathrm{\Lambda }_{c1}`$ doublet were discussed within the framework of heavy baryon chiral perturbation theory. Due to unknown couplings constants in the chiral Lagrangian, no actual decay widths were given. Within the same framework the pionic decay widths were calculated assuming the heavy quark effective theory is still valid for the strange quark . The coupling constants in the chiral Lagrangian were fixed using the p-wave strange baryon decay widths, which were later used to predict the strong decays of the p-wave charm baryons . The two pion width of $`\mathrm{\Lambda }_{c1}`$ was estimated to be around $`2.5`$ MeV, which was comparable to the total one pion width $`3.0`$ MeV. And the decays of $`\mathrm{\Lambda }_{c1}^{}`$ was suppressed by more than an order . In the p-wave doublet was treated as the bound state of the nucleon and heavy meson. It was found that the decays $`\mathrm{\Lambda }_{c1},\mathrm{\Lambda }_{c1}^{}\mathrm{\Lambda }_c\gamma `$ were suppressed due to the kinematic suppression of the electric dipole moment . In the constituent quark model was employed to study the orbitally excited heavy baryons. Sum rules were derived to constrain the coupling constants. The light front quark model, together with underlying $`SU(2N_f)\times O(3)`$ symmetry for the light diquark system, was used to relate and analyse the pionic coupling . However, the results have strong dependence on the constituent quark mass $`m_q`$. Varying $`m_q`$ from $`220`$ MeV to $`340`$ MeV, the decay widths increase by more than a factor of two . Within the same framework the electromagnetic decays of the p-wave baryons were calculated in . In both strong and radiative decays were calculated using a relativistic three-quark model. After this paper was submitted there appears an interesting paper discussing the radiative decays of the ground state heavy baryon multiplets in the framework of heavy baryon chiral perturbation theory. In some cases the loop corrections yield sizeable enhancement of the deca widths .
It will be helpful to extract these pionic and photonic coupling constants at the quark gluon level using QCD Lagrangian. We will treat this problem using QCD sum rules (QSR) , which are successful to extract the low-lying hadron masses and couplings. In this approach the nonperturbative effects are introduced via various condensates in the vacuum. The light cone QCD sum rule (LCQSR) differs from the conventional short-distance QSR in that it is based on the expansion over the twists of the operators. The main contribution comes from the lowest twist operators. Matrix elements of nonlocal operators sandwiched between a hadronic state and the vacuum define the hadron wave functions. In the present case our sum rules involve with the pion and photon wave function. When the LCQSR is used to calculate the coupling constant, the double Borel transformation is always invoked so that the excited states and the continuum contribution can be subtracted quite cleanly. We have calculated the pionic and electromagnetic coupling constants and decay widths of the ground state heavy hadrons and possible hybrid heavy mesons . In this work we extend the same framework to study the strong and radiative decays of lowest p-wave heavy baryons, i.e., $`\mathrm{\Lambda }_{c1}`$ doublet.
Our paper is organized as follows: Section I is an introduction. In the next section we derive the mass sum rule. The light cone sum rules for the pionic coupling constants are derived in Section III. Numerical analysis is presented. In Section IV we extend the same framework to analyse the electromagnetic processes $`\mathrm{\Lambda }_{c1}\mathrm{\Sigma }_c\gamma `$ etc. In Section V we discuss the processes $`\mathrm{\Lambda }_{c1},\mathrm{\Lambda }_{c1}^{}\mathrm{\Lambda }_c\gamma `$ and compare our results with other theoretical approaches. The last section is a summary.
## II The mass sum rules for the heavy hybrid mesons in HQET
### A Heavy quark effective theory
The effective Lagrangian of the HQET, up to order $`1/m_Q`$, is
$$_{\mathrm{eff}}=\overline{h}_vivDh_v+\frac{1}{2m_Q}๐ฆ+\frac{1}{2m_Q}๐ฎ+๐ช(1/m_Q^2),$$
(1)
where $`h_v(x)`$ is the velocity-dependent field related to the original heavy-quark field $`Q(x)`$ by
$$h_v(x)=e^{im_Qvx}\frac{1+\text{/}v}{2}Q(x),$$
(2)
$`v_\mu `$ is the heavy hadron velocity. $`๐ฆ`$ is the kinetic operator defined as
$$๐ฆ=\overline{h}_v(iD_t)^2h_v,$$
(3)
where $`D_t^\mu =D^\mu (vD)v^\mu `$, with $`D^\mu =^\mu igA^\mu `$ is the gauge-covariant derivative, and $`๐ฎ`$ is the chromomagnetic operator
$$๐ฎ=\frac{g}{2}C_{mag}(m_Q/\mu )\overline{h}_v\sigma _{\mu \nu }G^{\mu \nu }h_v,$$
(4)
where $`C_{mag}=\left({\displaystyle \frac{\alpha _s(m_Q)}{\alpha _s(\mu )}}\right)^{3/\beta _0}`$, $`\beta _0=112n_f/3`$. Note the heavy quark propogator has a simple form in coordinate space.
$$<0|T\{h_v(x),\overline{h}_v(0)\}|0>=_0^{\mathrm{}}๐t\delta (xvt)\frac{1+\widehat{v}}{2}.$$
(5)
### B The interpolating currents
We introduce the interpolating currents for the relevant heavy baryons:
$$\eta _{\mathrm{\Lambda }_c}(x)=ฯต_{abc}[u_{}^{a}{}_{}{}^{T}(x)C\gamma _5d^b(x)]h_v^c(x),$$
(6)
$$\eta _{\mathrm{\Sigma }_c^+}(x)=ฯต_{abc}[u_{}^{a}{}_{}{}^{T}(x)C\gamma _\mu d^b(x)]\gamma _t^\mu \gamma _5h_v^c(x),$$
(7)
$$\eta _{\mathrm{\Sigma }_{c}^{++}{}_{}{}^{}}^\mu (x)=ฯต_{abc}[u_{}^{a}{}_{}{}^{T}(x)C\gamma _\nu u^b(x)]\mathrm{\Gamma }_t^{\mu \nu }h_v^c(x),$$
(8)
$$\eta _{\mathrm{\Lambda }_{c1}}(x)=ฯต_{abc}[u_{}^{a}{}_{}{}^{T}(x)C\gamma _5d^b(x)]\gamma _t^\mu \gamma _5D_\mu ^th_v^c(x),$$
(9)
$$\eta _{\mathrm{\Lambda }_{c1}}^\mu (x)=ฯต_{abc}[u_{}^{a}{}_{}{}^{T}(x)C\gamma _5d^b(x)]\mathrm{\Gamma }_t^{\mu \nu }D_\nu ^th_v^c(x),$$
(10)
where $`a`$, $`b`$, $`c`$ is the color index, $`u(x)`$, $`d(x)`$, $`h_v(x)`$ is the up, down and heavy quark fields, $`T`$ denotes the transpose, $`C`$ is the charge conjugate matrix, $`\mathrm{\Gamma }_t^{\mu \nu }=g_t^{\mu \nu }+\frac{1}{3}\gamma _t^\mu \gamma _t^\nu `$, $`g_t^{\mu \nu }=g^{\mu \nu }v^\mu v^\nu `$, $`\gamma _t^\mu =\gamma _\mu \widehat{v}v^\mu `$, and $`v^\mu `$ is the velocity of the heavy hadron.
The overlap amplititudes of the interpolating currents with the heavy baryons are defined as:
$$0|\eta _{\mathrm{\Lambda }_c}|\mathrm{\Lambda }_c=f_{\mathrm{\Lambda }_c}u_{\mathrm{\Lambda }_c},$$
(11)
$$0|\eta _{\mathrm{\Lambda }_{c1}}|\mathrm{\Lambda }_{c1}=f_{\mathrm{\Lambda }_{c1}}u_{\mathrm{\Lambda }_{c1}},$$
(12)
$$0|\eta _{\mathrm{\Lambda }_{c1}^{}}^\mu |\mathrm{\Lambda }_{c1}^{}=\frac{f_{\mathrm{\Lambda }_{c1}^{}}}{\sqrt{3}}u_{\mathrm{\Lambda }_{c1}^{}}^\mu ,$$
(13)
$$0|\eta _{\mathrm{\Sigma }_c}|\mathrm{\Sigma }_c=f_{\mathrm{\Sigma }_c}u_{\mathrm{\Sigma }_c},$$
(14)
$$0|\eta _{\mathrm{\Sigma }_c^{}}^\mu |\mathrm{\Sigma }_c^{}=\frac{f_{\mathrm{\Sigma }_c^{}}}{\sqrt{3}}u_{\mathrm{\Sigma }_c^{}}^\mu ,$$
(15)
where $`u_{\mathrm{\Lambda }_{c1}^{}}^\mu `$, $`u_{\mathrm{\Sigma }_c^{}}^\mu `$ are the Rarita-Schwinger spinors in HQET. In the leading order of HQET, $`f_{\mathrm{\Sigma }_c}=f_{\mathrm{\Sigma }_c^{}}`$ and $`f_{\mathrm{\Lambda }_{c1}}=f_{\mathrm{\Lambda }_{c1}^{}}`$ due to heavy quark symmetry.
### C The $`\mathrm{\Lambda }_{Q1}`$ mass sum rules
In order to extract the binding energy of the p-wave heavy baryons in the leading order of HQET, we consider the correlators
$$id^4xe^{ikx}0|T\{\eta _{\mathrm{\Lambda }_{c1}}(x),\overline{\eta }_{\mathrm{\Lambda }_{c1}}(0)\}|0=\mathrm{\Pi }(\omega )\frac{1+\widehat{v}}{2},$$
(16)
with $`\omega =kv`$.
The dispersion relation for $`\mathrm{\Pi }(\omega )`$ reads
$$\mathrm{\Pi }(\omega )=\frac{\rho (s)}{s\omega iฯต}๐s,$$
(17)
where $`\rho (s)`$ is the spectral density in the limit $`m_Q\mathrm{}`$.
At the phenomenological side
$$\mathrm{\Pi }(\omega )=\frac{f_{\mathrm{\Lambda }_{c1}}^2}{\mathrm{\Lambda }_{c1}\omega }+\text{continuum}.$$
(18)
In order to suppress the continuum and higher excited states contribution we make Borel transformation with the variable $`\omega `$ to (17). We have
$$f_{\mathrm{\Lambda }_{c1}}^2e^{\frac{\overline{\mathrm{\Lambda }}_{\mathrm{\Lambda }_{c1}}}{T}}=_0^{s_0}\rho (s)e^{\frac{s}{T}}๐s,$$
(19)
where $`\overline{\mathrm{\Lambda }}_{\mathrm{\Lambda }_{c1}}`$ is the $`\mathrm{\Lambda }_{c1}`$ binding energy of in the leading order and $`s_0`$ is the continuum threshold. Starting from $`s_0`$ we have modeled the phenomenological spectral density with the parton-like one including both the perturbative term and various condensates.
The spectral density $`\rho (s)`$ at the quark level reads,
$$\rho (s)=\frac{3}{140\pi ^4}s^7\frac{1}{384\pi ^4}g_s^2G^2s^3+\frac{m_0^2a^2}{128\pi ^4}\delta (s)$$
(20)
where $`a=4\pi ^2\overline{q}q=0.55`$GeV<sup>3</sup>, $`g_s^2G^2=0.48`$GeV<sup>4</sup>, $`\overline{q}g_s\sigma Gq=m_0^2\overline{q}q`$, and $`m_0^2=0.8`$ GeV<sup>2</sup>. An interesting feature of (20) is that the gluon condensate is of the opposite sign as the leading perturbative term, in contrast with the ground state baryon mass sum rules. This may be interpreted as some kind of gluon excitation since we are considering p-wave baryons. In the present case the gluon in the covariant derivative also contributes to various condensates.
Two common approaches exist to extract the masses. One is the derivative method.
$$\overline{\mathrm{\Lambda }}_{\mathrm{\Lambda }_{c1}}=\frac{_0^{s_0}s\rho (s)e^{\frac{s}{T}}๐s}{_0^{s_0}\rho (s)e^{\frac{s}{T}}๐s}.$$
(21)
The other is the fitting method, which involves with fitting the left hand side and right hand side of Eq. (19) with the most suitable parameters $`\overline{\mathrm{\Lambda }}_{\mathrm{\Lambda }_{c1}},f_{\mathrm{\Lambda }_{c1}},s^0`$ in the working region of the Borel parameter. With both methods we get consistent results,
$`\overline{\mathrm{\Lambda }}_{\mathrm{\Lambda }_{c1}}`$ $`=`$ $`(1.1\pm 0.2)\text{GeV},`$ (22)
$`f_{\mathrm{\Lambda }_{c1}}`$ $`=`$ $`(0.025\pm 0.005)\text{GeV}^4,`$ (23)
$`s_{\mathrm{\Lambda }_{c1}}^0`$ $`=`$ $`(1.45\pm 0.2)\text{GeV}`$ (24)
in the working region $`0.51.3`$ GeV for the Borel parameter $`T`$. For later use we also need the mass and overlapping amplitude of the $`\mathrm{\Sigma }`$, $`\mathrm{\Lambda }`$ heavy baryon doublet, $`\overline{\mathrm{\Lambda }}_{\mathrm{\Sigma }_c}`$, $`\overline{\mathrm{\Lambda }}_{\mathrm{\Lambda }_c}`$, $`f_{\mathrm{\Sigma }_c}`$, $`f_{\mathrm{\Lambda }_c}`$ in the leading order of $`\alpha _s`$ .
$`\overline{\mathrm{\Lambda }}_{\mathrm{\Sigma }_c}`$ $`=`$ $`(1.0\pm 0.1)\text{GeV},`$ (25)
$`f_{\mathrm{\Sigma }_c}`$ $`=`$ $`(0.04\pm 0.004)\text{GeV}^3,`$ (26)
$`s_{\mathrm{\Sigma }_c}^0`$ $`=`$ $`(1.25\pm 0.15)\text{GeV}`$ (27)
$`\overline{\mathrm{\Lambda }}_{\mathrm{\Lambda }_c}`$ $`=`$ $`(0.8\pm 0.1)\text{GeV},`$ (28)
$`f_{\mathrm{\Lambda }_c}`$ $`=`$ $`(0.018\pm 0.004)\text{GeV}^3,`$ (29)
$`s_{\mathrm{\Lambda }_c}^0`$ $`=`$ $`(1.2\pm 0.15)\text{GeV}`$ (30)
## III LCQSR for the pionic couplings
### A The correlator for pionic couplings
We introduce the following amplitudes
$$M(\mathrm{\Lambda }_{c1}\mathrm{\Sigma }_c\pi )=g_s\overline{u}_{\mathrm{\Sigma }_c}u_{\mathrm{\Lambda }_{c1}},$$
(31)
$$M(\mathrm{\Lambda }_{c1}^{}\mathrm{\Sigma }_c\pi )=\sqrt{3}g_d\overline{u}_{\mathrm{\Sigma }_c}\gamma _5q_\mu ^t\widehat{q}u_{\mathrm{\Lambda }_{c1}^{}}^\mu ,$$
(32)
$$M(\mathrm{\Lambda }_{c1}\mathrm{\Sigma }_c^{}\pi )=\sqrt{3}g_d^1\overline{u}_{\mathrm{\Sigma }_c}^\mu \gamma _5q_\mu ^t\widehat{q}u_{\mathrm{\Lambda }_{c1}},$$
(33)
$$M(\mathrm{\Lambda }_{c1}^{}\mathrm{\Sigma }_c^{}\pi )=\overline{u}_{\mathrm{\Sigma }_c}^\mu [g_s^{}g_{\mu \nu }^t+3g_d^2(q_\mu ^tq_\nu ^t\frac{1}{3}g_{\mu \nu }^tq_t^2)]u_{\mathrm{\Lambda }_{c1}^{}}^\nu ,$$
(34)
where $`\widehat{q}=q_\mu \gamma ^\mu `$, $`q_\mu `$ is the pion momentum. Only the first two decay processes are kinematically allowed. Due to heavy quark symmetry, $`g_s^{}=g_s`$, $`g_d^1=g_d^2=g_d`$ in the limit of $`m_Q\mathrm{}`$. In other words there are two independent coupling constants correpsonding to s-wave and d-wave decays. Note we are unable to determine the sign of $`g_s`$ and $`g_d`$. And we are mainly interested in the decay widths of the p-wave heavy baryons. In the following our convention is to let both couplings be positive.
We consider the following correlators
$$d^4xe^{ikx}0|T\left(\eta _{\mathrm{\Lambda }_{c1}}(x)\overline{\eta }_{\mathrm{\Sigma }_c}(0)\right)|\pi (q)=\frac{1+\widehat{v}}{2}G_s(\omega ,\omega ^{}),$$
(35)
$$d^4xe^{ikx}0|T\left(\eta _{\mathrm{\Lambda }_{c1}^{}}^\mu (x)\overline{\eta }_{\mathrm{\Sigma }_c}(0)\right)|\pi (q)=\frac{1+\widehat{v}}{2}q_\alpha q_\nu \mathrm{\Gamma }_t^{\mu \alpha }\gamma _t^\nu \gamma _5G_d(\omega ,\omega ^{}),$$
(36)
where $`k^{}=kq`$, $`\omega =vk`$, $`\omega ^{}=vk^{}`$ and $`q^2=m_\pi ^2=0`$.
The function $`G_{s,d}(\omega ,\omega ^{})`$ has the following pole terms from double dispersion relation. For $`G_s`$ we have
$`{\displaystyle \frac{f_{\mathrm{\Lambda }_{c1}}f_\mathrm{\Sigma }g_s}{(\overline{\mathrm{\Lambda }}_{\mathrm{\Lambda }_{c1}}\omega ^{})(\overline{\mathrm{\Lambda }}_{\mathrm{\Sigma }_c}\omega )}}+{\displaystyle \frac{c}{\overline{\mathrm{\Lambda }}_{\mathrm{\Lambda }_{c1}}\omega ^{}}}+{\displaystyle \frac{c^{}}{\overline{\mathrm{\Lambda }}_{\mathrm{\Sigma }_c}\omega }}.`$ (37)
### B Pion light cone wave functions
To go futher we need the two- and three-particle pion light cone wave functions :
$`<\pi (q)|\overline{d}(x)\gamma _\mu \gamma _5u(0)|0>`$ $`=`$ $`if_\pi q_\mu {\displaystyle _0^1}๐ue^{iuqx}(\phi _\pi (u)+x^2g_1(u)+๐ช(x^4))`$ (38)
$`+`$ $`f_\pi \left(x_\mu {\displaystyle \frac{x^2q_\mu }{qx}}\right){\displaystyle _0^1}๐ue^{iuqx}g_2(u),`$ (39)
$`<\pi (q)|\overline{d}(x)i\gamma _5u(0)|0>`$ $`=`$ $`{\displaystyle \frac{f_\pi m_\pi ^2}{m_u+m_d}}{\displaystyle _0^1}๐ue^{iuqx}\phi _P(u),`$ (40)
$`<\pi (q)|\overline{d}(x)\sigma _{\mu \nu }\gamma _5u(0)|0>`$ $`=`$ $`i(q_\mu x_\nu q_\nu x_\mu ){\displaystyle \frac{f_\pi m_\pi ^2}{6(m_u+m_d)}}{\displaystyle _0^1}๐ue^{iuqx}\phi _\sigma (u).`$ (41)
$`<\pi (q)|\overline{d}(x)\sigma _{\alpha \beta }\gamma _5g_sG_{\mu \nu }(ux)u(0)|0>=`$ (42)
$`if_{3\pi }[(q_\mu q_\alpha g_{\nu \beta }q_\nu q_\alpha g_{\mu \beta })(q_\mu q_\beta g_{\nu \alpha }q_\nu q_\beta g_{\mu \alpha })]{\displaystyle ๐\alpha _i\phi _{3\pi }(\alpha _i)e^{iqx(\alpha _1+v\alpha _3)}},`$ (43)
$`<\pi (q)|\overline{d}(x)\gamma _\mu \gamma _5g_sG_{\alpha \beta }(vx)u(0)|0>=`$ (44)
$`f_\pi \left[q_\beta \left(g_{\alpha \mu }{\displaystyle \frac{x_\alpha q_\mu }{qx}}\right)q_\alpha \left(g_{\beta \mu }{\displaystyle \frac{x_\beta q_\mu }{qx}}\right)\right]{\displaystyle ๐\alpha _i\phi _{}(\alpha _i)e^{iqx(\alpha _1+v\alpha _3)}}`$ (45)
$`+f_\pi {\displaystyle \frac{q_\mu }{qx}}(q_\alpha x_\beta q_\beta x_\alpha ){\displaystyle ๐\alpha _i\phi _{}(\alpha _i)e^{iqx(\alpha _1+v\alpha _3)}}`$ (46)
and
$`<\pi (q)|\overline{d}(x)\gamma _\mu g_s\stackrel{~}{G}_{\alpha \beta }(vx)u(0)|0>=`$ (47)
$`if_\pi \left[q_\beta \left(g_{\alpha \mu }{\displaystyle \frac{x_\alpha q_\mu }{qx}}\right)q_\alpha \left(g_{\beta \mu }{\displaystyle \frac{x_\beta q_\mu }{qx}}\right)\right]{\displaystyle ๐\alpha _i\stackrel{~}{\phi }_{}(\alpha _i)e^{iqx(\alpha _1+v\alpha _3)}}`$ (48)
$`+if_\pi {\displaystyle \frac{q_\mu }{qx}}(q_\alpha x_\beta q_\beta x_\alpha ){\displaystyle ๐\alpha _i\stackrel{~}{\phi }_{}(\alpha _i)e^{iqx(\alpha _1+v\alpha _3)}}.`$ (49)
The operator $`\stackrel{~}{G}_{\alpha \beta }`$ is the dual of $`G_{\alpha \beta }`$: $`\stackrel{~}{G}_{\alpha \beta }=\frac{1}{2}ฯต_{\alpha \beta \delta \rho }G^{\delta \rho }`$; $`๐\alpha _i`$ is defined as $`๐\alpha _i=d\alpha _1d\alpha _2d\alpha _3\delta (1\alpha _1\alpha _2\alpha _3)`$. Due to the choice of the gauge $`x^\mu A_\mu (x)=0`$, the path-ordered gauge factor $`P\mathrm{exp}\left(ig_s_0^1๐ux^\mu A_\mu (ux)\right)`$ has been omitted.
The wave function $`\phi _\pi (u)`$ is associated with the leading twist 2 operator, $`g_1(u)`$ and $`g_2(u)`$ correspond to twist 4 operators, and $`\phi _P(u)`$ and $`\phi _\sigma (u)`$ to twist 3 ones. The function $`\phi _{3\pi }`$ is of twist three, while all the wave functions appearing in eqs.(46), (49) are of twist four. The wave functions $`\phi (x_i,\mu )`$ ($`\mu `$ is the renormalization point) describe the distribution in longitudinal momenta inside the pion, the parameters $`x_i`$ ($`_ix_i=1`$) representing the fractions of the longitudinal momentum carried by the quark, the antiquark and gluon.
The wave function normalizations immediately follow from the definitions (38)-(49): $`_0^1๐u\phi _\pi (u)=_0^1๐u\phi _\sigma (u)=1`$, $`_0^1๐ug_1(u)=\delta ^2/12`$, $`๐\alpha _i\phi _{}(\alpha _i)=๐\alpha _i\phi _{}(\alpha _i)=0`$, $`๐\alpha _i\stackrel{~}{\phi }_{}(\alpha _i)=๐\alpha _i\stackrel{~}{\phi }_{}(\alpha _i)=\delta ^2/3`$, with the parameter $`\delta `$ defined by the matrix element: $`<\pi (q)|\overline{d}g_s\stackrel{~}{G}_{\alpha \mu }\gamma ^\alpha u|0>=i\delta ^2f_\pi q_\mu `$.
### C The pionic sum rules
Now the expressions of $`G_s`$, $`G_d`$ at the quark level read,
$`G_s(\omega ,\omega ^{})=if_\pi {\displaystyle _0^{\mathrm{}}}dt{\displaystyle _0^1}due^{i(1u)\omega t}e^{iu\omega ^{}t}\{{\displaystyle \frac{6\mu _\pi }{\pi ^2t^4}}\phi _P(u)+{\displaystyle \frac{\mu _\pi }{3\pi ^2t^4}}(3\phi _\sigma ^{}(u)+[u\phi _\sigma (u)]^{^{\prime \prime }})`$ (50)
$`+(\overline{q}q+{\displaystyle \frac{t^2}{16}}\overline{q}g_s\sigma Gq)g_2(u)+(\overline{q}q+{\displaystyle \frac{t^2}{16}}\overline{q}g_s\sigma Gq){\displaystyle \frac{[u\phi _\pi (u)]^{^{\prime \prime }}+t^2[uG_2(u)+ug_1(u)]^{^{\prime \prime }}}{3t^2}}\}`$ (51)
$`+{\displaystyle \frac{i}{\pi ^2}}f_{3\pi }{\displaystyle \frac{dt}{t^2}_0^1๐u(1u)๐\alpha _ie^{i\omega t[1(\alpha _1+u\alpha _3)]}e^{i\omega ^{}t(\alpha _1+u\alpha _3)}[(qv)^2it(qv)^3(\alpha _1+u\alpha _3)]\phi _{3\pi }(\alpha _i)}`$ (52)
$`{\displaystyle \frac{2i}{\pi ^2}}f_{3\pi }{\displaystyle \frac{dt}{t^2}_0^1๐uu๐\alpha _ie^{i\omega t[1(\alpha _1+u\alpha _3)]}e^{i\omega ^{}t(\alpha _1+u\alpha _3)}(qv)^2\phi _{3\pi }(\alpha _i)}`$ $`,`$ (53)
$`G_d(\omega ,\omega ^{})=if_\pi {\displaystyle _0^{\mathrm{}}}dt{\displaystyle _0^1}duue^{i(1u)\omega t}e^{iu\omega ^{}t}\{{\displaystyle \frac{\mu _\pi }{3\pi ^2t^2}}\phi _\sigma (u)`$ (54)
$`+{\displaystyle \frac{1}{3}}(\overline{q}q+{\displaystyle \frac{t^2}{16}}\overline{q}g_s\sigma Gq)(\phi _\pi (u)+t^2[G_2(u)+g_1(u)])\}`$ (55)
$`{\displaystyle \frac{i}{\pi ^2}}f_{3\pi }{\displaystyle \frac{dt}{t^2}_0^1๐u(1u)๐\alpha _ie^{i\omega t[1(\alpha _1+u\alpha _3)]}e^{i\omega ^{}t(\alpha _1+u\alpha _3)}[1it(qv)(\alpha _1+u\alpha _3)]\phi _{3\pi }(\alpha _i)}`$ (56)
$`{\displaystyle \frac{i}{\pi ^2}}f_{3\pi }{\displaystyle \frac{dt}{t^2}_0^1๐uu๐\alpha _ie^{i\omega t[1(\alpha _1+u\alpha _3)]}e^{i\omega ^{}t(\alpha _1+u\alpha _3)}\phi _{3\pi }(\alpha _i)}`$ $`,`$ (57)
where $`\mu _\pi =1.65`$GeV, $`f_\pi =132`$MeV, $`F^{}(u)=\frac{dF(u)}{du}`$ and $`F^{\prime \prime }(u)=\frac{d^2F(u)}{du^2}`$. There are two three particle terms in the form of $`\phi _{3\pi }`$ in (50), (54). The gluon arises from the light quark propagator in the first term and from the covariant derivative in the second term. For large euclidean values of $`\omega `$ and $`\omega ^{}`$ this integral is dominated by the region of small $`t`$, therefore it can be approximated by the first few terms with lowest twists.
After Wick rotations and making double Borel transformation with the variables $`\omega `$ and $`\omega ^{}`$ the single-pole terms in (37) are eliminated. Subtracting the continuum contribution which is modeled by the dispersion integral in the region $`s,s^{}s_0`$, we arrive at:
$`g_sf_{\mathrm{\Lambda }_{c1}}f_{\mathrm{\Sigma }_c}=`$ $`{\displaystyle \frac{f_\pi }{\pi ^2}}e^{\frac{\overline{\mathrm{\Lambda }}_{\mathrm{\Lambda }_{c1}}+\overline{\mathrm{\Lambda }}_{\mathrm{\Sigma }_c}}{2T}}\{6\mu _\pi \phi _P(u_0)T^5f_4({\displaystyle \frac{s_0}{T}})+{\displaystyle \frac{f_{3\pi }}{f_\pi }}(2I_3I_4I_6)T^5f_4({\displaystyle \frac{s_0}{T}})`$ (58)
$`+{\displaystyle \frac{\mu _\pi }{3}}\left(3\phi _\sigma ^{}(u_0)+[u\phi _\sigma (u)]_{u_0}^{^{\prime \prime }}\right)T^5f_4({\displaystyle \frac{s_0}{T}})+{\displaystyle \frac{a}{12}}[u\phi _\pi (u)]_{u_0}^{^{\prime \prime }}(1{\displaystyle \frac{m_0^2}{16T^2}})T^3f_2({\displaystyle \frac{s_0}{T}})`$ (59)
$`{\displaystyle \frac{a}{4}}(g_2(u_0)+{\displaystyle \frac{1}{3}}[uG_2(u)+ug_1(u)]_{u_0}^{^{\prime \prime }})(1{\displaystyle \frac{m_0^2}{16T^2}})Tf_0({\displaystyle \frac{s_0}{T}})\},`$ (60)
where $`f_n(x)=1e^x\underset{k=0}{\overset{n}{}}\frac{x^k}{k!}`$ is the factor used to subtract the continuum, $`s_0`$ is the continuum threshold. $`u_0=\frac{T_1}{T_1+T_2}`$, $`T\frac{T_1T_2}{T_1+T_2}`$, $`T_1`$, $`T_2`$ are the Borel parameters. The functions $`I_i`$ are defined below. In obtaining (58) we have used the Borel transformation formula: $`\widehat{}_\omega ^Te^{\alpha \omega }=\delta (\alpha \frac{1}{T})`$ and integration by parts to absorb the factors $`(qv)`$ and $`1/(qv)`$. In this way we arrive at the simple form after double Borel transformation.
Similarly we have:
$`g_df_{\mathrm{\Lambda }_{c1}^{}}f_{\mathrm{\Sigma }_c}=`$ $`{\displaystyle \frac{f_\pi }{\pi ^2}}e^{\frac{\overline{\mathrm{\Lambda }}_{\mathrm{\Lambda }_1^{}}+\overline{\mathrm{\Lambda }}_{\mathrm{\Sigma }_c}}{2T}}\{{\displaystyle \frac{\mu _\pi }{3}}u_0\phi _\sigma (u_0)T^3f_2({\displaystyle \frac{s_0}{T}}){\displaystyle \frac{f_{3\pi }}{f_\pi }}(I_1+I_2+I_5)T^3f_2({\displaystyle \frac{s_0}{T}})`$ (61)
$`+{\displaystyle \frac{a}{12}}u_0\phi _\pi (u_0)(1{\displaystyle \frac{m_0^2}{16T^2}})Tf_0({\displaystyle \frac{s_0}{T}}){\displaystyle \frac{a}{12T}}u_0[G_2(u_0)+g_1(u_0)](1{\displaystyle \frac{m_0^2}{16T^2}})\}.`$ (62)
The functions $`G_2(u_0)`$, $`I_i`$ are defined as:
$$G_2(u_0)=_0^{u_0}g_2(u)๐u,$$
(63)
$$I_1=_0^{u_0}๐\alpha _1_0^{1u_0}๐\alpha _2\frac{u_0\alpha _1}{\alpha _3^2}\phi _{3\pi }(\alpha _i),$$
(64)
$$I_2=_0^{u_0}๐\alpha _1_0^{1u_0}๐\alpha _2\frac{1u_0\alpha _2}{\alpha _3^2}\phi _{3\pi }(\alpha _i),$$
(65)
$`I_3={\displaystyle _0^{u_0}}๐\alpha _1{\displaystyle \frac{d}{d\alpha _3}}[{\displaystyle \frac{\phi _{3\pi }(\alpha _1,1\alpha _1\alpha _3,\alpha _3)}{\alpha _3}}]|_{\alpha _3=u_0\alpha _1}`$ (66)
$`{\displaystyle _0^{u_0}}๐\alpha _1{\displaystyle \frac{\phi _{3\pi }(\alpha _1,1u_0,u_0\alpha _1)}{(u_0\alpha _1)^2}}+{\displaystyle _0^{1u_0}}๐\alpha _2{\displaystyle \frac{\phi _{3\pi }(u_0,\alpha _2,1u_0\alpha _2)}{(1u_0\alpha _2)^2}}`$ $`,`$ (67)
$`I_4={\displaystyle _0^{1u_0}}{\displaystyle \frac{d\alpha _3}{\alpha _3}}[{\displaystyle \frac{d\phi _{3\pi }(\alpha _1,1\alpha _1\alpha _3,\alpha _3)}{d\alpha _1}}]|_{\alpha _1=u_0}`$ (68)
$`+{\displaystyle _0^{u_0}}๐\alpha _1{\displaystyle \frac{\phi _{3\pi }(\alpha _1,1u_0,u_0\alpha _1)}{(u_0\alpha _1)^2}}{\displaystyle _0^{1u_0}}๐\alpha _2{\displaystyle \frac{\phi _{3\pi }(u_0,\alpha _2,1u_0\alpha _2)}{(1u_0\alpha _2)^2}}`$ $`,`$ (69)
$$I_5=_0^{1u_0}u_0\frac{d\alpha _3}{\alpha _3}\phi _{3\pi }(u_0,1u_0\alpha _3,\alpha _3)+_0^{u_0}๐\alpha _1_0^{1u_0}๐\alpha _2\frac{2u_01+\alpha _2}{\alpha _3^2}\phi _{3\pi }(\alpha _i),$$
(70)
$`I_6={\displaystyle \frac{d[\alpha _1\phi _{3\pi }(\alpha _1,1\alpha _1\alpha _3,\alpha _3)]}{d\alpha _1}}|_{\alpha _3=1u_0}^{\alpha _1=u_0}`$ (71)
$`{\displaystyle _0^{1u_0}}๐\alpha _3{\displaystyle \frac{d^2}{d\alpha _{1}^{}{}_{}{}^{2}}}[\phi _{3\pi }(\alpha _1,1\alpha _1\alpha _3,\alpha _3){\displaystyle \frac{\alpha _1}{\alpha _3}}]|_{\alpha _1=u_0}`$ (72)
$`+[\phi _{3\pi }(\alpha _1,1\alpha _1\alpha _3,\alpha _3){\displaystyle \frac{\alpha _3\alpha _1}{\alpha _3^2}}]|_{\alpha _3=u_0}^{\alpha _1=0}`$ (73)
$`+{\displaystyle _0^{u_0}}๐\alpha _3{\displaystyle _0^{u_0\alpha _3}}๐\alpha _1{\displaystyle \frac{d}{d\alpha _1}}[\phi _{3\pi }(\alpha _1,1\alpha _1\alpha _3,\alpha _3){\displaystyle \frac{\alpha _3\alpha _1}{\alpha _3^2}}]`$ (74)
$`+[\phi _{3\pi }(\alpha _1,1\alpha _1\alpha _3,\alpha _3){\displaystyle \frac{\alpha _3\alpha _1}{\alpha _3^2}}]|_{\alpha _3=1u_0}^{\alpha _1=u_0}`$ (75)
$`{\displaystyle _0^{1u_0}}๐\alpha _3{\displaystyle \frac{d}{d\alpha _1}}[\phi _{3\pi }(\alpha _1,1\alpha _1\alpha _3,\alpha _3){\displaystyle \frac{\alpha _3\alpha _1}{\alpha _3^2}}]|_{\alpha _1=u_0}`$ (76)
$`2[{\displaystyle \frac{\phi _{3\pi }(\alpha _i)}{\alpha _3}}]|_{\alpha _3=u_0}^{\alpha _1=0}2{\displaystyle _0^{u_0}}๐\alpha _3{\displaystyle _0^{u_0\alpha _3}}๐\alpha _1{\displaystyle \frac{d}{d\alpha _1}}[{\displaystyle \frac{\phi _{3\pi }(\alpha _1,1\alpha _1\alpha _3,\alpha _3)}{\alpha _3}}]`$ (77)
$`2{\displaystyle _0^{u_0}}๐\alpha _3{\displaystyle _0^{u_0\alpha _3}}๐\alpha _1{\displaystyle \frac{\phi _{3\pi }(\alpha _1,1\alpha _1\alpha _3,\alpha _3)}{\alpha _3^2}}`$ (78)
$`+2{\displaystyle _0^{1u_0}}๐\alpha _3{\displaystyle _0^{1u_0\alpha _3}}๐\alpha _2{\displaystyle \frac{\phi _{3\pi }(1\alpha _2\alpha _3,\alpha _2,\alpha _3)}{\alpha _3^2}}`$ $`,`$ (79)
where $`\alpha _3,\alpha _1`$ are the longitudinal momentum fraction of gluon and down quark inside the pion respectively.
### D Determination of the parameters for pionic sum rules
The mass difference between $`\mathrm{\Lambda }_{c1}`$ and $`\mathrm{\Sigma }_c`$ is only about $`0.1`$GeV in the leading order of HQET. And the values of the Borel parameter $`T_1,T_2`$ is around $`2`$ GeV in the working region. So we choose to work at the symmetric point $`T_1=T_2=2T`$, i.e., $`u_0=\frac{1}{2}`$, which diminishes the uncertainty arising from the pion wave functions and enables a rather clean subtraction of the continuum contribution.
The pion wave functions and their values at the middle point are discussed in . At the scale $`\mu =1.0`$GeV the values of the various functions appearing in (58)-(61) at $`u_0=\frac{1}{2}`$ are: $`\phi _\pi (u_0)=(1.5\pm 0.2)`$ , $`\phi _P(u_0)=1.142`$, $`\phi _\sigma (u_0)=1.463`$, $`g_1(u_0)=0.034`$GeV<sup>2</sup>, $`G_2(u_0)=0.02`$GeV<sup>2</sup> , $`\phi _\sigma ^{}(u_0)=0`$, $`g_2(u_0)=0`$, $`[u\phi _\pi (u)]_{u=u_0}^{^{\prime \prime }}=[u\phi _\sigma (u)]_{u=u_0}^{^{\prime \prime }}=6`$, $`[ug_1(u)+uG_2(u)]_{u=u_0}^{^{\prime \prime }}=0.29`$, $`I_1=1.17`$, $`I_2=1.17`$, $`I_3=31.9`$, $`I_4=31.9`$, $`I_5=1.64`$, $`I_6=247.5`$, $`f_{3\pi }=0.0035`$GeV<sup>2</sup>. We have used the forms in for $`\phi _{3\pi }(\alpha _i)`$ to calculate integrals $`I_i`$. The three particle wave functions are known to next order in the conformal spin expansion up to now. The second derivatives need knowledge of the detailed shape of the pion wave functions at the middle point. Various sources indicate $`\phi _\pi (u)`$ is very close to the asymptotic form , which is exactly known. Based on these considerations we have employed the asymptotic forms to extract the second derivatives for $`\phi _\sigma (u)`$ and $`\phi _\pi (u)`$.
### E Numerical analysis of pionic sum rules
Note the spectral density of the sum rule (58)-(61) is either proptional to $`s^2`$ or $`s^4`$, the continuum has to be subtracted carefully. We use $`s_0=(1.3\pm 0.15)`$ GeV, which is the average of the thresholds of the $`\mathrm{\Lambda }_{c1}`$ and $`\mathrm{\Sigma }_c`$ mass sum rules. The variation of $`g_{s,d}`$ with the Borel parameter $`T`$ and $`s_0`$ is presented in Fig. 1 and Fig. 2. The curves correspond to $`s_0=1.2,1.3,1.4`$GeV from bottom to top respectively. Stability develops for these sum rules in the region $`0.5`$ GeV $`<`$$`T`$$`<`$$`1.5`$ GeV, we get:
$`g_sf_{\mathrm{\Lambda }_{c1}}f_\mathrm{\Sigma }=(0.5\pm 0.3)\times 10^3\text{GeV}^7,`$ (80)
$`g_df_{\mathrm{\Lambda }_1^{}}f_\mathrm{\Sigma }=(2.8\pm 0.6)\times 10^3\text{GeV}^5,`$ (81)
where the errors refers to the variations with $`T`$ and $`s_0`$ in this region. And the central value corresponds to $`T=1`$GeV and $`s_0=1.3`$GeV.
Combining (22), (25) we get
$`g_s=(0.5\pm 0.3),`$ (82)
$`g_d=(2.8\pm 0.6)\text{GeV}^2.`$ (83)
We collect the values of the pionic couplings from various approaches TABLE I. Note in our notation $`3g_d`$ corresponds to those in .
We use the following formulas to calculate the pionic decay widths of p-wave heavy baryons.
$$\mathrm{\Gamma }(\mathrm{\Lambda }_{c1}\mathrm{\Sigma }_c\pi )=\frac{g_s^2}{2\pi }\frac{m_{\mathrm{\Sigma }_c}}{m_{\mathrm{\Lambda }_{c1}}}|q|,$$
(84)
$$\mathrm{\Gamma }(\mathrm{\Lambda }_{c1}^{}\mathrm{\Sigma }_c\pi )=\frac{g_d^2}{2\pi }\frac{m_{\mathrm{\Sigma }_c}}{m_{\mathrm{\Lambda }_{c1}^{}}}|q|^5,$$
(85)
where $`|q|`$ is the pion decay momentum. We use the values $`m_{\mathrm{\Lambda }_{c1}}=2.593`$ GeV, $`m_{\mathrm{\Lambda }_{c1}^{}}=2.625`$ GeV, $`m_{\mathrm{\Sigma }_c}=2.452`$ GeV . In the $`\mathrm{\Lambda }_{c1}`$ decays due to isospin symmetry violations of the pion and $`\mathrm{\Sigma }_c`$ multiplet masses, the pion decay momentum is $`17,23,32`$ MeV for the final states $`\mathrm{\Sigma }_c^{++}\pi ^{},\mathrm{\Sigma }_c^0\pi ^+,\mathrm{\Sigma }_c^+\pi ^0`$ respectively. This effect causes significant difference in the decay widths, which are collected in TABLE II. Summing all the three isospin channels we get $`\mathrm{\Gamma }(\mathrm{\Lambda }_{c1}\mathrm{\Sigma }_c\pi )=2.7`$ MeV and $`\mathrm{\Gamma }(\mathrm{\Lambda }_{c1}^{}\mathrm{\Sigma }_c\pi )=33`$ keV. The later is nearly suppressed by two oders of magnitude due to d-wave decays.
ยฟFrom TABLE II we see that our results are numerically close to those from fixing the unknown coupling constants from the p-wave strange baryon strong decay widths assuming heavy quark effective theory could be extended to the strange quark case . The values of d-wave decay widths from the above approach and ours are much smaller than those from the quark models . As for the s-wave decays various approaches yield consistent results.
## IV Radiative decays of p-wave heavy baryons
### A The correlator
The light cone photon wave functions have been used to discuss radiative decay processes in in the framework of QCD sum rules. We extend the same formalism to extract the electromagnetic coupling consants for the $`\mathrm{\Lambda }_{Q1}`$ doublet decays.
The radiative coupling constants are defined through the following amplitudes:
$$M(\mathrm{\Lambda }_{c1}\mathrm{\Sigma }_c\gamma )=eฯต_{\beta \nu \rho \mu }q^\beta e_{}^{\nu }{}_{}{}^{}\overline{u}_{\mathrm{\Sigma }_c}[f_sg_t^{\rho \alpha }+f_dq^\alpha v^\rho ]\gamma _\alpha ^t\gamma _t^\mu u_{\mathrm{\Lambda }_{c1}},$$
(86)
$$M(\mathrm{\Lambda }_{c1}\mathrm{\Sigma }_c^{}\gamma )=\sqrt{3}eฯต_{\beta \nu \rho \mu }q^\beta e_{}^{\nu }{}_{}{}^{}\overline{u}_{\mathrm{\Sigma }_c^{}}^\alpha [f_s^1g_t^{\rho \alpha }+f_d^1q^\alpha v^\rho ]\gamma _5\gamma _t^\mu u_{\mathrm{\Lambda }_{c1}},$$
(87)
$$M(\mathrm{\Lambda }_{c1}^{}\mathrm{\Sigma }_c\gamma )=\sqrt{3}eฯต_{\beta \nu \rho \mu }q^\beta e_{}^{\nu }{}_{}{}^{}\overline{u}_{\mathrm{\Sigma }_c}[f_s^2g_t^{\rho \alpha }+f_d^2q^\alpha v^\rho ]\gamma _5\gamma _\alpha ^tu_{\mathrm{\Lambda }_{c1}}^\mu ,$$
(88)
$$M(\mathrm{\Lambda }_{c1}^{}\mathrm{\Sigma }_c^{}\gamma )=3eฯต_{\beta \nu \rho \mu }q^\beta e_{}^{\nu }{}_{}{}^{}\overline{u}_{\mathrm{\Sigma }_c^{}}^\alpha [f_s^3g_t^{\rho \alpha }+f_d^3q^\alpha v^\rho ]u_{\mathrm{\Lambda }_{c1}}^\mu ,$$
(89)
where $`e_\mu (\lambda )`$ and $`q_\mu `$ are the photon polarization vector and momentum respectively, $`e`$ is the charge unit. Due to heavy quark symmetry, we have $`f_s^1=f_s^2=f_s^3=f_s`$, $`f_d^1=f_d^2=f_d^3=f_d`$. As in the case of pionic couplings there are only two independent coupling constants associated with the E1 and M2 decays.
We consider the correlator
$$d^4xe^{ikx}\gamma (q)|T\left(\eta _{\mathrm{\Lambda }_{c1}}(0)\overline{\eta }_{\mathrm{\Sigma }_c}(x)\right)|0=e\frac{1+\widehat{v}}{2}\gamma _\alpha ^t\gamma _t^\mu ฯต_{\beta \nu \rho \mu }v^\beta e_{}^{\nu }{}_{}{}^{}\{F_s(\omega ,\omega ^{})g_t^{\rho \alpha }+F_d(\omega ,\omega ^{})q^\alpha v^\rho \},$$
(90)
$`F_{s,d}(\omega ,\omega ^{})`$ has the same pole structures as $`G_{s,d}(\omega ,\omega ^{})`$.
The light cone two-particle photon wave functions are :
$`<\gamma (q)|\overline{q}(x)\sigma _{\mu \nu }q(0)|0>=ie_qe\overline{q}q{\displaystyle _0^1}due^{iuqx}\{(e_\mu q_\nu e_\nu q_\mu )[\chi \varphi (u)+x^2h_1(u)]`$ (91)
$`+[(qx)(e_\mu x_\nu e_\nu x_\mu )+(ex)(x_\mu q_\nu x_\nu q_\mu )x^2(e_\mu q_\nu e_\nu q_\mu )]h_2(u)\}`$ $`,`$ (92)
$$<\gamma (q)|\overline{q}(x)\gamma _\mu \gamma _5q(0)|0>=\frac{f}{4}e_qeฯต_{\mu \nu \rho \sigma }e^\nu q^\rho x^\sigma _0^1๐ue^{iuqx}\psi (u).$$
(93)
The $`\varphi (u),\psi (u)`$ is associated with the leading twist two photon wave function, while $`g_1(u)`$ and $`g_2(u)`$ are twist-4 PWFs. All these PWFs are normalized to unity, $`_0^1๐uf(u)=1`$.
We want to emphasize that the photon light cone wave functions include the complete perturbative and non-perturbative electromagnetic interactions for the light quarks in principle. Yet the interaction of the photon with the heavy quark is not parametrized and constrained by the photon light cone wave functions. It seems possible that the photon couples directly to the heavy quark line. This is different from the QCD sum rules for the pionic couplings since pions can not couple directly to the heavy quark. However the real photon coupling to heavy quark involves a spin-flip transition, which is suppressed by a factor of $`1/m_Q`$ . So it vanishs in the leading order of $`1/m_Q`$ expansion. Since we are interested in the leading order couplings $`f_{s,d}`$, itโs enough to keep the photon light cone wave functions for the light quarks only.
Expressing (90) with the photon wave functions, we arrive at:
$`F_s(\omega ,\omega ^{})={\displaystyle \frac{1}{\pi ^2}}(e_ue_d)\overline{q}q{\displaystyle _0^{\mathrm{}}}dt{\displaystyle _0^1}due^{i(1u)\omega t}e^{iu\omega ^{}t}\{[{\displaystyle \frac{1}{t^4}}\chi \varphi (u)`$ (94)
$`+{\displaystyle \frac{1}{t^2}}(h_1(u)h_2(u))]+{\displaystyle \frac{\pi ^2}{24}}f\psi (u)t\}+\mathrm{}`$ $`.`$ (95)
$`F_d(\omega ,\omega ^{})={\displaystyle \frac{i}{\pi ^2}}(e_ue_d)\overline{q}q{\displaystyle _0^{\mathrm{}}}dt{\displaystyle _0^1}due^{i(1u)\omega t}e^{iu\omega ^{}t}u\{[{\displaystyle \frac{1}{t^3}}\chi \varphi (u)`$ (96)
$`+{\displaystyle \frac{1}{t}}(h_1(u)h_2(u))]+{\displaystyle \frac{\pi ^2}{24}}f\psi (u)t\}+\mathrm{}`$ $`.`$ (97)
The final sum rules are:
$`f_sf_{\mathrm{\Lambda }_{c1}}f_{\mathrm{\Sigma }_c}={\displaystyle \frac{a}{4\pi ^4}}(e_ue_d)e^{\frac{\overline{\mathrm{\Lambda }}_{\mathrm{\Lambda }_{c1}}+\overline{\mathrm{\Lambda }}_{\mathrm{\Sigma }_c}}{2T}}\{\chi \varphi (u_0)T^5f_4({\displaystyle \frac{s_0}{T}})`$ (98)
$`[h_1(u_0)h_2(u_0)]T^3f_2({\displaystyle \frac{s_0}{T}})+{\displaystyle \frac{\pi ^2}{24}}f\psi (u_0)T^1f_0({\displaystyle \frac{s_0}{T}})\}`$ $`,`$ (99)
$`f_df_{\mathrm{\Lambda }_{c1}}f_{\mathrm{\Sigma }_c}={\displaystyle \frac{a}{4\pi ^4}}(e_ue_d)e^{\frac{\overline{\mathrm{\Lambda }}_{\mathrm{\Lambda }_{c1}}+\overline{\mathrm{\Lambda }}_{\mathrm{\Sigma }_c}}{2T}}u_0\{\chi \varphi (u_0)T^4f_3({\displaystyle \frac{s_0}{T}})`$ (100)
$`[h_1(u_0)h_2(u_0)]T^2f_1({\displaystyle \frac{s_0}{T}})+{\displaystyle \frac{\pi ^2}{24}}f\psi (u_0)\}`$ $`.`$ (101)
### B Numerical analysis of the photonic sum rules
The leading photon wave functions receive only small corrections from the higher conformal spins so they do not deviate much from the asymptotic form. We shall use
$$\varphi (u)=6u(1u),$$
(102)
$$\psi (u)=1,$$
(103)
$$h_1(u)=\frac{1}{8}(1u)(3u),$$
(104)
$$h_2(u)=\frac{1}{4}(1u)^2.$$
(105)
with $`f=0.028`$GeV<sup>2</sup> and $`\chi =4.4`$GeV<sup>2</sup> at the scale $`\mu =1`$GeV.
The variation of $`f_{s,d}`$ with the Borel parameter $`T`$ and $`s_0`$ is presented in FIG. 3 and FIG. 4. Stability develops for the sum rules (98), (100) in the region $`0.5`$ GeV $`<`$$`T`$$`<`$$`1.5`$ GeV, we get:
$`f_sf_{\mathrm{\Lambda }_{c1}}f_\mathrm{\Sigma }=(2.0\pm 0.8)\times 10^4\text{GeV}^6,`$ (106)
$`f_df_{\mathrm{\Lambda }_{c1}}f_\mathrm{\Sigma }=(4.8\pm 1.2)\times 10^4\text{GeV}^5,`$ (107)
where the errors refers to the variations with $`T`$ and $`s_0`$ in this region. And the central value corresponds to $`T=1.0`$GeV and $`s_0=1.3`$GeV. Our final result is
$`f_s=(0.20\pm 0.08)\text{GeV}^1,`$ (108)
$`f_d=(0.48\pm 0.12)\text{GeV}^2.`$ (109)
The decay width formulas in the leading order of HQET are
$`\mathrm{\Gamma }(\mathrm{\Lambda }_{c1}\mathrm{\Sigma }_c\gamma )=16\alpha |\stackrel{}{q}|^3{\displaystyle \frac{m_{\mathrm{\Sigma }_c}}{m_{\mathrm{\Lambda }_{c1}}}}[f_s^2+{\displaystyle \frac{1}{2}}f_d^2|\stackrel{}{q}|^2],`$ (110)
$`\mathrm{\Gamma }(\mathrm{\Lambda }_{c1}\mathrm{\Sigma }_c^{}\gamma )=8\alpha |\stackrel{}{q}|^3{\displaystyle \frac{m_{\mathrm{\Sigma }_c^{}}}{m_{\mathrm{\Lambda }_{c1}}}}[f_s^2+{\displaystyle \frac{1}{2}}f_d^2|\stackrel{}{q}|^2],`$ (111)
$`\mathrm{\Gamma }(\mathrm{\Lambda }_{c1}^{}\mathrm{\Sigma }_c\gamma )=4\alpha |\stackrel{}{q}|^3{\displaystyle \frac{m_{\mathrm{\Sigma }_c}}{m_{\mathrm{\Lambda }_{c1}^{}}}}[f_s^2+{\displaystyle \frac{1}{2}}f_d^2|\stackrel{}{q}|^2],`$ (112)
$`\mathrm{\Gamma }(\mathrm{\Lambda }_{c1}^{}\mathrm{\Sigma }_c^{}\gamma )=20\alpha |\stackrel{}{q}|^3{\displaystyle \frac{m_{\mathrm{\Sigma }_c^{}}}{m_{\mathrm{\Lambda }_{c1}^{}}}}[f_s^2+{\displaystyle \frac{1}{2}}f_d^2|\stackrel{}{q}|^2],`$ (113)
where $`|\stackrel{}{q}|=134,72,164,103`$ MeV is the photon decay momentum for the above four processes. The d-wave decay is negligible. The decay width values are collected in TABLE III. The uncertainty is typically about $`50\%`$.
The decays $`\mathrm{\Lambda }_{Q1}\mathrm{\Sigma }_Q\gamma `$ do not occur in the leading order in the bound state picture . Due to the unknown coupling constant $`c_{RS}`$ in the chiral lagrangian for the heavy quark electromagnetic interactions, no numerical values are available . However the decay width ratios of the four final states are exactly the same as ours if we ignore the isospin violations of the heavy multiplet masses in the heavy quark limit. Our results are much smaller than those from various versions of quark models , which may indicate that the $`1/m_c`$ correction is important.
## V The process $`\mathrm{\Lambda }_{c1}\mathrm{\Lambda }_c\gamma `$ etc
As can be seen later the radiative decay processes of p-wave $`\mathrm{\Lambda }_{c1}`$ doublet to $`\mathrm{\Lambda }_c`$ is quite different from those in the previous section. We present more details here. The possible E1 decay amplitudes are
$$M(\mathrm{\Lambda }_{c1}\mathrm{\Lambda }_c\gamma )=eh_pe_\mu ^{}\overline{u}_{\mathrm{\Lambda }_c}[g_t^{\mu \nu }vqv^\mu q^\nu ]\gamma _\mu \gamma _5u_{\mathrm{\Lambda }_{c1}},$$
(114)
$$M(\mathrm{\Lambda }_{c1}^{}\mathrm{\Lambda }_c\gamma )=\sqrt{3}eh_p^{}e_\mu ^{}\overline{u}_{\mathrm{\Lambda }_c}[g_t^{\mu \nu }vqv^\mu q^\nu ]u_{\mathrm{\Lambda }_{c1}^\nu }.$$
(115)
Due to heavy quark symmetry $`h_p=h_p^{}`$.
We consider the correlator
$$\mathrm{\Pi }=id^4xe^{ikx}\gamma (q)|T\left(\eta _{\mathrm{\Lambda }_{c1}}(x)\overline{\eta }_{\mathrm{\Lambda }_c}(0)\right)|0=\frac{1+\widehat{v}}{2}e_\mu ^{}[g_t^{\mu \nu }vqv^\mu q^\nu ]\gamma _\mu \gamma _5H_p(\omega ,\omega ^{}).$$
(116)
We first calculate the part solely involved with the light quark, which can be expressed with the photon wave functions. We get
$$\mathrm{\Pi }=2i_0^{\mathrm{}}dtd^4xe^{ikx}\widehat{D}_t\delta (xvt)\gamma _5\frac{1+\widehat{v}}{2}\{Tr[\gamma _5CiS_u^T(x)C\gamma _5<\gamma (q)|d(x)\overline{d}(0)|0>]+(ud)\},$$
(117)
where summation over color has been performed. There are two types of terms with even $`\gamma `$ matrices in the trace. The first one is connected with $`\psi (u)`$ and the trace looks like $`Tr[\gamma _5C\widehat{x}^TC\gamma _5\gamma _\mu \gamma _5]`$. The second is involved with $`\varphi (u),h_1(u),h_2(u)`$ and the trace looks like $`Tr[\gamma _5C1C\gamma _5\sigma _{\mu \nu }]`$. In both $`\mathrm{\Lambda }_{Q1},\mathrm{\Lambda }_Q`$ states the up and down quarks are in the $`0^+`$ state, which leads to the presence of $`\gamma _5C`$ and $`C\gamma _5`$ in both traces. Clearly both traces vanish. This property results from the underlying flavor and spin structure of the light quark sector. In other words the light quark contribution is zero to all orders of the heavy quark expansion in the framework of LCQSR with the commonly used interpolating currents (6) and (9) for $`\mathrm{\Lambda }_Q`$ and $`\mathrm{\Lambda }_{Q1}`$ respectively. The decays $`\mathrm{\Lambda }_{Q1}\mathrm{\Lambda }_Q\gamma `$ and $`\mathrm{\Lambda }_{Q1}^{}\mathrm{\Lambda }_Q\gamma `$ happens only when the photon couples directly to the heavy quark line.
Now letโs move to the part involved with the heavy quark. At first sight there are two types of terms in the leading order of heavy quark expansion. The first one comes from the insertion of the operator $`i\overline{[}h_v(y)ivDh_v(y)]d^4y`$ in (116), which contributes a factor $`ve^{}(\lambda )`$ to the decay amplitude. For the real photon $`ve^{}(\lambda )=0`$ so it drops out. The other possible term arises from the covariant derivative in $`\eta _{\mathrm{\Lambda }_{c1}}`$, which leads to a nonzero correlator. For the tensor structure $`i\widehat{e}^{}\gamma _5\frac{1+\widehat{v}}{2}`$ we have
$$\mathrm{\Pi }(\omega ,\omega ^{})=\frac{e}{\pi ^4}_0^{\mathrm{}}๐te^{i\omega ^{}t}\{\frac{6}{t^6}+\frac{<g_s^2G^2>}{64t^2}\frac{a^2}{96}\},$$
(118)
where the photon field has contributed a factor $`e^{iqx}`$. Itโs important to note only the variable $`\omega ^{}`$ appears in (118). Itโs a single pole term which must vanish after we make double Borel transformation to the variables $`\omega ,\omega ^{}`$ simultaneously. We have shown there is no leading order E1 transition in (114) arising from the photon couplings to the heavy quark line in the leading order of heavy quark expansion. Based on the same spin and flavor consideration we know that radiative decay processes like $`\mathrm{\Sigma }_{Q1}\mathrm{\Lambda }_{Q1}\gamma `$, $`\mathrm{\Lambda }_{Q1}\mathrm{\Lambda }_{Q1}\gamma `$, $`\mathrm{\Sigma }_{Q1}\mathrm{\Sigma }_{Q1}\gamma `$ are also forbidden in the leading order of $`1/m_Q`$ expansion, where we have used notations in .
We may rewrite the decay amplitudes as
$$M(\mathrm{\Lambda }_{c1}\mathrm{\Lambda }_c\gamma )=ef_pF_{\mu \nu }\overline{u}_{\mathrm{\Lambda }_c}\sigma ^{\mu \nu }\gamma _5u_{\mathrm{\Lambda }_{c1}},$$
(119)
$$M(\mathrm{\Lambda }_{c1}^{}\mathrm{\Lambda }_c\gamma )=2\sqrt{3}ef_p^1F_{\mu \nu }v^\mu \overline{u}_{\mathrm{\Lambda }_c}u_{\mathrm{\Lambda }_{c1}^\nu },$$
(120)
$$M(\mathrm{\Lambda }_{c1}^{}\mathrm{\Lambda }_{c1}\gamma )=2\sqrt{3}ef_p^2F_{\mu \nu }\overline{u}_{\mathrm{\Lambda }_{c1}}\gamma _t^\mu \gamma _5u_{\mathrm{\Lambda }_{c1}^\nu }.$$
(121)
Due to heavy quark symmetry we have
$$f_p=f_p^1=f_p^2.$$
(122)
Note $`f_p=\frac{1}{4}h_p`$.
In these decays we know the light quarks do not contribute. However, the $`J^P`$ of the light diquark changes from $`1^{}`$ to $`0^+`$ which ensures the decay $`\mathrm{\Lambda }_{c1}\mathrm{\Lambda }_c\gamma `$ is an E1 transition. The angular momentum and parity $`J^P=\frac{1}{2}^+`$ of the heavy quark does not change so the coupling constant $`f_p`$ is the same as that for the heavy quark $`M1`$ transition, which is induced by the magnetic moment operator
$$f_p=\frac{\mu _c}{2}=\frac{e_c}{4m_c}.$$
(123)
Another approach is to consider the three point correlation function for the tensor structure $`\widehat{e}_t\gamma _5\frac{1+\widehat{v}}{2}`$
$$id^4xd^4ze^{ikxik^{}z}0|T\{\eta _{\mathrm{\Lambda }_{c1}}(x),\frac{๐ฆ(0)+๐ฎ(0)}{2m_c},\overline{\eta }_{\mathrm{\Lambda }_{c1}}(z)\}|0=\mathrm{\Pi }_3(\omega ,\omega ^{})\frac{1+\widehat{v}}{2},$$
(124)
with $`\omega =kv,\omega ^{}=k^{}v`$.
$$\mathrm{\Pi }_3(\omega ,\omega ^{})=\frac{2e_c}{m_c}\frac{1}{\pi ^4}_0^{\mathrm{}}๐t_1๐t_2e^{i\omega t_1+i\omega ^{}t_2}\{\frac{18}{(t_1+t_2)^8}+\frac{<g_s^2G^2>}{64(t_1+t_2)^4}\},$$
(125)
After the double Borel transformation and continuum subtraction we get the sum rule for $`h_p`$
$$h_p(\overline{\mathrm{\Lambda }}_{\mathrm{\Lambda }_{c1}}\overline{\mathrm{\Lambda }}_{\mathrm{\Lambda }_c})f_{\mathrm{\Lambda }_{c1}}f_{\mathrm{\Lambda }_c}e^{\frac{\overline{\mathrm{\Lambda }}_{\mathrm{\Lambda }_{c1}}+\overline{\mathrm{\Lambda }}_{\mathrm{\Lambda }_c}}{2T}}=\frac{1}{\pi ^4}\frac{e_c}{m_c}\{36T^8f_7(\frac{s_0}{T})+\frac{<g_s^2G^2>}{32}T^4f_3(\frac{s_0}{T})\}.$$
(126)
Dividing (126) by (19) we get
$$h_p=\frac{e_c}{m_c}e^{\frac{\overline{\mathrm{\Lambda }}_{\mathrm{\Lambda }_{c1}}\overline{\mathrm{\Lambda }}_{\mathrm{\Lambda }_c}}{2T}}\frac{f_{\mathrm{\Lambda }_{c1}}}{3(\overline{\mathrm{\Lambda }}_{\mathrm{\Lambda }_{c1}}\overline{\mathrm{\Lambda }}_{\mathrm{\Lambda }_c})f_{\mathrm{\Lambda }_c}}\frac{T^8f_7(\frac{s_0}{T})+\frac{<g_s^2G^2>}{1152}T^4f_3(\frac{s_0}{T})}{T^8f_7(\frac{s_0}{T})\frac{<g_s^2G^2>}{6912}T^4f_3(\frac{s_0}{T})+\frac{m_0^2a^2}{8192}}.$$
(127)
Numerically we have $`h_p\frac{e_c}{m_c}`$, which is consistent with (123). The decay widths formulas are
$`\mathrm{\Gamma }(\mathrm{\Lambda }_{c1}\mathrm{\Lambda }_c\gamma )=e_c^2\alpha |\stackrel{}{q}|^3{\displaystyle \frac{m_{\mathrm{\Lambda }_c}}{m_{\mathrm{\Lambda }_{c1}}m_c^2}},`$ (128)
$`\mathrm{\Gamma }(\mathrm{\Lambda }_{c1}^{}\mathrm{\Lambda }_c\gamma )=e_c^2\alpha |\stackrel{}{q}|^3{\displaystyle \frac{m_{\mathrm{\Lambda }_c}}{m_{\mathrm{\Lambda }_{c1}^{}}m_c^2}},`$ (129)
$`\mathrm{\Gamma }(\mathrm{\Lambda }_{c1}^{}\mathrm{\Lambda }_{c1}\gamma )=e_c^2\alpha |\stackrel{}{q}|^3{\displaystyle \frac{m_{\mathrm{\Lambda }_{c1}}}{m_{\mathrm{\Lambda }_{c1}^{}}m_c^2}}.`$ (130)
The decay momentum is $`290,320,32`$ MeV respectively. We take $`m_c=1.4`$ GeV. The numerical values are collected in TABLE III. These widths comes solely from the $`๐ช(1/m_Q)`$ correction. But their numerical values are greater than those leading order widths for the channels $`\mathrm{\Sigma }_c\gamma ,\mathrm{\Sigma }_c^{}\gamma `$. The reason is purely kinematical. The decay momentum for the final state $`\mathrm{\Lambda }_c\gamma `$ is three times larger. For the p-wave decay there appears an enhancement factor of 27.
These widths in (128) are propotional to $`\frac{e_c^2}{m_c^2}`$. Therefore the corresponding radiative decays $`\mathrm{\Lambda }_{b1}\mathrm{\Lambda }_b\gamma `$, $`\mathrm{\Lambda }_{b1}^{}\mathrm{\Lambda }_b\gamma `$, $`\mathrm{\Lambda }_{b1}^{}\mathrm{\Lambda }_{b1}\gamma `$ are further suppressed by a factor $`(\frac{e_b}{e_c}\frac{m_c}{m_b})^240`$. The widths of the first two decays are around $`1`$ keV.
If we use naive dimensional analysis to let $`c_{RT}`$ in TABLE III be of the order of unity or simply assume that the E1 transition coupling constant $`h_p`$ in (114) is of the same order of M1 transition one , we would get a width $`๐ช(100)`$ keV. Our result is in strong contrast with those from the bound state picture , where $`\mathrm{\Gamma }(\mathrm{\Lambda }_{c1},\mathrm{\Lambda }_{c1}^{}\mathrm{\Lambda }_c\gamma )=16,21`$ keV and $`\mathrm{\Gamma }(\mathrm{\Lambda }_{b1},\mathrm{\Lambda }_{b1}^{}\mathrm{\Lambda }_b\gamma )=90,119`$ keV. Future experiments should be able to judge which mechanism is correct.
It was noted that the radiative decays $`\mathrm{\Lambda }_{Q1}\mathrm{\Lambda }_Q`$ was forbidden in the leading order of heavy quark symmetry assuming one-body transition operators, which arises from a complete cancellation due to the specific spins of light constituent quarks in the antisymmetric initial and final state . The point is consistent with our observation of the vanishing contribution of the light quark sector to this radiative process.
ยฟFrom our calculation we know the d-wave single pion width of $`\mathrm{\Lambda }_{c1}^{}`$ is $`33`$ keV and the estimate in yielded $`35`$ keV for the two pion decay width. Itโs interesting to notice that the radiative decay widths are $`48,5,6,0.014`$ keV for the final states $`\mathrm{\Lambda }_c\gamma ,\mathrm{\Sigma }_c\gamma ,\mathrm{\Sigma }_c^{}\gamma ,\mathrm{\Lambda }_{c1}\gamma `$ respectively. The width of the decay channel $`\mathrm{\Lambda }_c\gamma `$ is bigger than either of that of the strong decay modes. The $`\mathrm{\Lambda }_{c1}^{}`$ should be a narrow state with a total width about $`130`$ keV.
The two pion width of $`\mathrm{\Lambda }_{c1}`$ is about $`2.5`$ MeV . ยฟFrom TABLE II and III the one pion and electromagnetic widths are $`\mathrm{\Gamma }(\mathrm{\Sigma }_c\pi ,\mathrm{\Lambda }_c\gamma ,\mathrm{\Sigma }_c\gamma ,\mathrm{\Sigma }^{}\gamma )=2.7,0.048,0.011,0.001`$ MeV. Its total width is about $`5.4`$ MeV.
Itโs believed that $`\mathrm{\Lambda }_{b1}`$ lies below $`\mathrm{\Sigma }_b\pi ,\mathrm{\Lambda }_b\pi \pi `$ threshold. If so its dominant decays are electromagnetic. From our calculation we see $`\mathrm{\Gamma }(\mathrm{\Lambda }_{b1}\mathrm{\Lambda }_b\gamma ,\mathrm{\Sigma }_b\gamma ,\mathrm{\Sigma }_b^{}\gamma )=1,11,1`$ keV if we assume the same decay momentum as in the $`\mathrm{\Lambda }_{c1}`$ decays. Its total width is about $`13`$ keV. It will be a very narrow state. Clearly the radiative channels $`\mathrm{\Sigma }_b\gamma `$ will be very useful to find them experimentally.
The major decay modes of $`\mathrm{\Lambda }_{b1}^{}`$ might be d-wave one pion decay and electromagnetic decays to $`\mathrm{\Sigma }_b`$ doublet if the two pion mode is not allowed. Their widths are $`\mathrm{\Gamma }(\mathrm{\Lambda }_{b1}^{}\mathrm{\Sigma }_b\pi ,\mathrm{\Lambda }_b\gamma ,\mathrm{\Sigma }_b\gamma ,\mathrm{\Sigma }_b^{}\gamma )=33,1,5,6`$ keV if we assume the same decay momentum as in the $`\mathrm{\Lambda }_{c1}^{}`$ case. Itโs also a very narrow state with a width of $`45`$ keV.
Before ending this section we want to improve our previous calculation of radiative decays of excited heavy mesons . (1) First the s-wave terms involved with $`g_s`$ should not appear in $`(1^+,2^+)(0^{},1^{})\gamma `$ processes. All decays are M2 transitions. The $`g_s^2`$ in the decay width formulas should be replaced by $`\frac{1}{9}g_d^2|\stackrel{}{q}|^4`$. The last eight widths in Eq. (94) should read $`2,8,3,11,6,23,7,27`$ keV respectively, which is much smaller than original wrong ones. (2) The E1 transition $`(0^+,1^+)(0^{},1^{})\gamma `$ decays was identified as s-wave decays. This was misleading. The factor $`(qv)`$ should be in the tensor structure to ensure the E1 transition structure in Eq. (47) in . We present the correct sum rules for $`g_1`$ below.
$$g_1f_{,1/2}f_{+,1/2}=\frac{a}{4\pi ^2}e^{\frac{\mathrm{\Lambda }_{,1/2}+\mathrm{\Lambda }_{+,1/2}}{2T}}\{\chi \varphi (u_0)Tf_0(\frac{s_0}{T})g_1(u_0)\frac{1}{T}\},$$
(131)
where $`s_0=\omega _c/2=(1.5\pm 0.2)`$ GeV. Numerically we have $`g_1=(1.6\pm 0.2)`$ GeV<sup>-1</sup>.
## VI Discussions
In our calculation only the errors due to the variations of $`T`$ and $`s_0`$ are included in the final results for $`g_{s,d}`$, $`f_{s,d}`$. The various input parameters like quark condensate, gluon condensate, $`\chi `$, $`f`$ etc also have some uncertainty. Among these the values of the pion and photon wave functions introduce largest uncertainty. Although their values are constrained by either experimental data or other QCD sum rule analysis, they may still lead to $`25\%`$ unceritainty. Keeping the light cone wave functions up to twist four also leads to some errors. However the light cone sum rules are dominated by the lowest twist wave functions. Take the sum rule (100) for $`f_d`$ for an example. At $`T=1`$ GeV, the twist-four term involved with $`h_1,h_2`$ is only $`9\%`$ of the leading twist term after the continuum subtraction. In other words the light cone expansion converges quickly. So we expect the contribution of higher twist terms to be small. There are other sources of uncertainty which is difficult to estimate. One is the QCD radiative correction, which is not small in both the mass sum rule and LCQSRs for the pionic coupling constants of the ground state heavy hadrons in HQET. But their ratio depends only weakly on these corrections because of large cancellation . Numerically the radiative corrections are around $`10\%`$ of the tree level result.
Another possible source is the $`1/m_Q`$ correction for the charmed p-wave baryons. The leading order coupling constants $`g_{s,d}`$ etc will be corrected by terms like $`g_{s,d}^{}/m_Q`$, which will affect decay widths. For the charmed hadrons $`1/m_Q`$ corrections are sizable and may reach $`30\%`$ while such corrections are generally less than $`10\%`$ of the leading order term for the bottom system . Especially for the E1 transition coupling constant $`f_s`$, the correction is of the order $`\frac{e_c}{4m_c}`$, which may be comparable with the leading order one for the charm system. One is justified to use these coupling constants to calculate the decay widths of the p-wave bottom baryons. Unfortunately data is still not available for the p-wave bottom baryons. So we have calculated the p-wave $`\mathrm{\Lambda }_{c1}`$ doublet decay widths with some reservation.
In summary we have calculated the pionic and electromagnetic coupling constants and decay widths of the lowest p-wave heavy baryon doublet. We compare our calculation with different approaches in literature. We hope these results will be useful in the future experimental search of $`\mathrm{\Lambda }_{b1},\mathrm{\Lambda }_{b1}^{}`$ baryons.
Acknowledgements: S.-L.Z. is grateful to Prof. C.-S. Huang for bringing the topic of excited heavy baryon to his attention.
Figure Captions
Fig. 1. Dependence of $`g_sf_{\mathrm{\Lambda }_{c1}}f_{\mathrm{\Sigma }_c}`$ on the Borel parameter $`T`$ for different values of the continuum threshold $`s_0`$. ยฟFrom top to bottom the curves correspond to $`s_0=1.4,1.3,1.2`$ GeV.
Fig. 2. Dependence of $`g_df_{\mathrm{\Lambda }_1^{}}f_{\mathrm{\Sigma }_c}`$ on $`T`$, $`s_0`$.
Fig. 3. Dependence of $`f_sf_{\mathrm{\Lambda }_{c1}}f_{\mathrm{\Sigma }_c}`$ on $`T`$, $`s_0`$.
Fig. 4. Dependence of $`f_df_{\mathrm{\Lambda }_{c1}}f_{\mathrm{\Sigma }_c}`$ on $`T`$, $`s_0`$. |
warning/0002/hep-ph0002212.html | ar5iv | text | # Dynamical fermion masses under the influence of Kaluza-Klein fermions in extra dimensions
## I Introduction
It is an interesting idea to assume an existence of the extra-dimensional space which eventually compactifies leaving our 4-dimensional space-time as a real world. The recent proposal for the mass scale of the compactified space to be much smaller than the Planck scale gave a strong impact on the onset of studying phenomenological evidences of extra-dimensional effects. In recent approaches with extra dimensions it is usually assumed that the standard model particles reside in the 4-dimensional brane while the graviton may move around the bulk, the space-time with extra dimensions. In our present analysis we introduce bulk fermions in addition to the graviton and see what effects could be observed on the standard model particles. The bulk fermions interact with themselves as well as with fermions in the 4-dimensional brane through the exchange of the graviton and its Kaluza-Klein excited modes, or through the exchange of gauge bosons which may be assumed to exist in the bulk. The interactions among fermions generated as a result of the exchange of all the Kaluza-Klein excited modes of the graviton or gauge bosons may be expressed as effective four-fermion interactions. According to the four-fermion interactions we expect that the dynamical generation of fermion masses will take place.
In the present communication we look for a possibility of the dynamical fermion mass generation under the influence of the bulk fermions through the effective four-fermion interactions. Although our argument is applicable to any higher dimensional models, we confine ourselves to the 5-dimensional space-time for the convenience of explanations. In 5 dimensions fermion mass terms are forbidden if we require the symmetry under the chiral projection. The possible source of fermion masses in 4 dimensions is two-fold, i. e. the dynamically generated fermion masses and masses of the Kaluza-Klein excited modes of the bulk fermions. The mass of the Kaluza-Klein excited modes is known to be of order $`1/R`$ where $`R`$ is the radius of the compactified fifth dimension.
We show in Sec. III that the mixing between the brane fermion and bulk fermion does not lead to the large mass of order $`1/R`$ for the fermion in 4 dimensions. We then consider in Secs. IV and V whether the dynamically generated fermion mass can be made small compared with the mass of the Kaluza-Klein excited modes. We calculate an effective potential for a composite operator composed of a fermion and an anti-fermion in the leading order of the $`1/N_f`$ expansion with $`N_f`$ the number of fermion species. We find that the mass of fermions in the 4-dimensional brane is generated dynamically if the compactification radius $`R`$ passes its critical value $`R_C`$ and the phase transition associated with the mass generation is of second order. This means that the fermion masses in the 4-dimensional brane is small as far as the radius of the compactified fifth dimension is close to its critical value.
## II 5-Dimensional Fermion Theory and Torus Compactification
We assume an existence of 5-dimensional bulk fermions $`\psi `$ in interaction with fermions $`L`$ on the 4-dimentional brane. Effective interactions among these fermions can be given in the form of the four-fermion interaction. We imagine that such effective interactions originate from the exchange of the bulk gravitons between fermions. In fact it is known that the exchange of the Kaluza-Klein excited modes of the bulk graviton results in effective four-fermion interactions. After the Fierz transformation on the four-fermion interactions we generate the transition-type interactions. Accordingly we start with the following Lagrangian for our model
$$^{(5)}=\overline{\psi }i\gamma ^M_M\psi +[\overline{L}i\gamma ^\mu _\mu L+g^2\overline{\psi }\gamma ^ML\overline{L}\gamma _M\psi ]\delta (x^4),$$
(1)
where $`g`$ is the coupling constant with mass dimension -3/2 and index $`M`$ runs from 0 to 4 while index $`\mu `$ runs from 0 to 3. Fermions $`\psi `$ and $`L`$ are assumed to be of $`N_f`$ components.
It is easy to see that the Lagrangian is symmetric under the chiral projection $`x^4x^4`$, $`\psi (x^4)i\gamma ^4\psi (x^4)`$ and $`L(x^4)i\gamma ^4L(x^4)`$. Thus, if we impose this symmetry on a Lagrangian describing our fermion system, any fermion mass terms are forbidden thus resulting in the above Lagrangian. It should be noted here that, as is easily understood by referring to the Clifford algebra, an irreducible representation of the 5-dimensional fermion field is given by a 4-component field just as in the case of 4 dimensions. Hence we can use the same field both for the 4 and 5 dimensions. In odd dimensions it is well-known that there exists no object like $`\gamma _5`$ in 4 dimensions which commutes with all the $`\gamma `$ matrices. Hence we do not have such object in 5 dimensions while the fifth component of the $`\gamma `$ matrices in 5 dimensions, $`\gamma ^4`$, turns out to be $`i\gamma _5`$ in 4 dimensions.
For the later convenience we rewrite the above Lagrangian by using auxiliary field $`\sigma _M`$ in the following equivalent form,
$`^{(5)}=\overline{\psi }i\gamma ^M_M\psi +[\overline{L}i/L\sigma ^M\sigma _M^{}+(g\sigma _M\overline{\psi }\gamma ^ML+h.c.)]\delta (x^4).`$ (2)
Since we are mainly interested in the dynamical fermion mass generation in the leading order of the $`1/N_f`$ expansion, we neglect the irrelevant terms in the Lagrangian by assuming $`\sigma _\mu =0`$ where $`\mu `$ runs from 0 to 3. After chiral rotation $`\psi e^{i\frac{\pi }{4}\gamma _5}\psi `$ and $`Le^{i\frac{\pi }{4}\gamma _5}L`$ we have
$`^{(5)}=\overline{\psi }i/\psi i\overline{\psi }_4\psi +[\overline{L}i/L|\sigma |^2+(g\sigma \overline{\psi }L+h.c.)]\delta (x^4),`$ (3)
where $`\sigma =\sigma _4`$. The Lagrangian (3) is considered to be a composite-Higgs version of the mixing interaction adopted in Ref. .
We now consider that the space of the fifth dimension is compactified on a circle with radius $`R`$. By adopting the periodic boundary condition at $`x_4=0`$ and $`x_4=2\pi R`$ the bulk fermion field is expressed as a Fourier series,
$$\psi (x^\mu ,x^4)=N\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}\psi _n(x^\mu )e^{i\frac{n}{R}x^4},$$
(4)
where $`N`$ is the normalization constant.
We define the 4-dimensional Lagrangian in the following way,
$`^{(4)}`$ $``$ $`{\displaystyle _0^{2\pi R}}๐x^4^{(5)}`$ (5)
$`=`$ $`{\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}\overline{\psi }_ni/\psi _n+{\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}{\displaystyle \frac{n}{R}}\overline{\psi }_n\psi _n+\overline{L}i/L`$ (7)
$`|\sigma |^2+(m{\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}\overline{\psi }_nL+h.c.),`$
where $`mNg\sigma `$. Note that we have normalized the kinetic terms in Eq. (7) by choosing $`N1/\sqrt{2\pi R}`$.
## III Mass Spectrum in 4-Dimensional Space-time
In the following arguments we employ the matrix expressions
$`\mathrm{\Psi }^t(L,\psi _0,\psi _1,\psi _1,\psi _2,\psi _2\mathrm{}),`$ (8)
and
$`M\left(\begin{array}{ccccccc}0& m^{}& m^{}& m^{}& m^{}& m^{}& \mathrm{}\\ m& 0& 0& 0& 0& 0& \mathrm{}\\ m& 0& \frac{1}{R}& 0& 0& 0& \mathrm{}\\ m& 0& 0& \frac{1}{R}& 0& 0& \mathrm{}\\ m& 0& 0& 0& \frac{2}{R}& 0& \mathrm{}\\ m& 0& 0& 0& 0& \frac{2}{R}& \mathrm{}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\end{array}\right).`$ (16)
By using Eqs. (8) and (16) we rewrite the mass term and mixing term in the 4-dimensional Lagrangian such that
$`_{\mathrm{mixing}}^{(4)}={\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}{\displaystyle \frac{n}{R}}\overline{\psi }_n\psi _n+(m{\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}\overline{\psi }_nL+h.c.)=\overline{\mathrm{\Psi }}M\mathrm{\Psi }.`$ (17)
If auxiliary field $`\sigma `$ acquires a non-vanishing vcuum expectation value, we replace $`\sigma `$ in $`m`$ by its vacuum expectation value $`\sigma `$, i. e. $`m=Ng\sigma `$. The eigenvalues of matrix $`M`$ determine the masses of 4-dimensional fermions. The eigenvalue equation is given by
$`det(M\lambda I)`$ (18)
$`=`$ $`\left[{\displaystyle \underset{j=1}{\overset{\mathrm{}}{}}}\left(\lambda ^2\left({\displaystyle \frac{j}{R}}\right)^2\right)\right]\left[\lambda ^2|m|^22|m|^2\lambda ^2{\displaystyle \underset{l=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{\lambda ^2\left(l/R\right)^2}}\right]=0.`$ (19)
It should be noted that the solutions $`\lambda =j/R`$ obtained by setting the first factor in the middle of Eq. (19) to vanish is not the eigenvalues of Eq. (19) since they are canceled by the same factor in the denominator in the second factor. The real eigenvalues are obtained from the second factor in the middle of Eq. (19). The summation in the third term of the second factor in Eq. (19) can be performed as it is the Fourier series expansion of the cotangent function and thus the eigenvalue equation reduces to the simpler form,
$`\lambda R=\pi |mR|^2\mathrm{cot}(\pi \lambda R).`$ (20)
We would like to see whether we have a possibility of getting the light solution in the eigenvalue equation (20). We confine ourselves to the case $`|m|1/R`$. By expanding the solution of Eq. (20) in powers of $`|mR|`$ we obtain
$`\lambda _{\pm 0}R`$ $`=`$ $`\pm |mR|\left(1{\displaystyle \frac{\pi ^2}{6}}|mR|^2+๐ช(|mR|^4)\right),`$ (21)
$`\lambda _{\pm n}R`$ $`=`$ $`\pm n\left(1+{\displaystyle \frac{|mR|^2}{n^2}}+๐ช(|mR|^4)\right)(n0).`$ (22)
Obviously we find from Eqs. (21) and (22) that the lightest eigenvalue is given by
$`\lambda _{\pm 0}=\pm |m|\text{for}|m|1/R.`$ (23)
Thus we conclude that within our scheme there is a possibility of having the light fermion masses which is much smaller than the mass of the Kaluza-Klein modes of the bulk fermion. It should be noted here that in Eq. (23) we have two kinds of fermions with the positive and negative mass. These fermions are, however, indistinguishable since they have exactly the same properties. The next step that we have to proceed is to show that this light fermion mass is obtained as a result of the dynamical mass generation mechanism and can really be small.
## IV Effective Potential for Composite Fields
We are now interested in the actual value of the lightest fermion mass in the 4-dimensional brane as is given in Eq. (23). Thus we have to study the dynamical mechanism of generating the non-vanishing vacuum expectation value of composite field $`\sigma `$ which is involved in the expression $`m=Ng\sigma `$. For this purpose we would like to calculate the effective potential for composite field $`\sigma `$.
Our 4-dimensional Lagrangian (7) is rewritten with the matrix representation introduced in the last section as follows,
$`^{(4)}=\overline{\mathrm{\Psi }}(M+Ii/)\mathrm{\Psi }\left|\sigma \right|^2.`$ (24)
The generating functional $`Z`$ for our system is given by
$`Z={\displaystyle [๐\overline{\mathrm{\Psi }}][๐\mathrm{\Psi }][๐\sigma ][๐\sigma ^{}]e^{i{\scriptscriptstyle d^4x^{\left(4\right)}}}}.`$ (25)
By performing the path-integration for fermion field $`\mathrm{\Psi }`$ in Eq. (25) we find in the leading order of the $`1/N_f`$ expansion,
$`Z`$ $`=`$ $`{\displaystyle [๐\sigma ][๐\sigma ^{}]e^{i{\scriptscriptstyle d^4xV(\sigma )}}},`$ (26)
$`V(\sigma )`$ $``$ $`|\sigma |^2{\displaystyle \frac{d^4k}{i(2\pi )^4}\mathrm{ln}det(M+I/k)},`$ (27)
where $`V`$ is the effective potential for the auxiliary field $`\sigma `$. By performing the Wick rotation in the momentum integration in Eq. (27) we rewrite the effective potential such that
$`V(\sigma )=|\sigma |^2{\displaystyle \frac{1}{8\pi ^2}}{\displaystyle _0^{\mathrm{\Lambda }^2}}d(k_E^2)k_E^2\mathrm{ln}det(M^2+Ik_E^2),`$ (28)
where $`k_E`$ stands for the Euclidean momentum and $`\mathrm{\Lambda }`$ is the momentum cut-off. Note that in addition to the divergence in the momentum integration the divergence in the Kaluza-Klein mode sum shows up in general. In our 5-dimensional model this divergence is not present while, if we start from the space-time dimensions higher than six, the divergence does exist and we have to rely on the regularization method developed recently in Ref. . The same argument applies to the case where we derived Eq. (20).
After some calculations we find that
$`V(\sigma )=|\sigma |^2`$ $``$ $`{\displaystyle \frac{1}{2\pi ^2}}{\displaystyle _0^\mathrm{\Lambda }}๐xx^3\mathrm{ln}\left[x^2+|m|^2(\pi xR)\mathrm{coth}(\pi xR)\right]`$ (29)
$``$ $`{\displaystyle \frac{1}{2\pi ^2}}{\displaystyle \underset{j=1}{\overset{\mathrm{}}{}}}{\displaystyle _0^\mathrm{\Lambda }}๐xx^3\mathrm{ln}\left[x^2+\left({\displaystyle \frac{j}{R}}\right)^2\right].`$ (30)
The gap equation to determine the vacuum expectation value $`<|\sigma |>`$ of $`|\sigma |`$ reads
$`{\displaystyle \frac{V(\sigma )}{|\sigma |}}=2|\sigma |\left\{1{\displaystyle \frac{g^2}{2\pi ^2}}{\displaystyle _0^\mathrm{\Lambda }}๐x{\displaystyle \frac{x^3}{2x\mathrm{tanh}(\pi xR)+g^2|\sigma |^2}}\right\}=0.`$ (31)
By numerical observation of Eq. (31) we find that there exists a non-trivial solution for $`|\sigma |`$ for a suitable range of parameters $`g`$ and $`R`$ and the solution corresponds to the true minimum of the effective potential. Accordingly the fermion mass is generated dynamically. Here the auxiliary field $`\sigma `$ (or the composite field $`\overline{L}\psi `$) acquires a vacuum expectation value. The continuous symmetry that is broken by this dynamical process is the U(1) symmetry existed in the original Lagrangian (1):
$$\psi e^{i\beta }\psi ,Le^{i\beta }L.$$
(32)
Moreover the phase transition associated with this symmetry breaking is of second order as is seen in Fig. 1.
In the case of the second order phase transition the critical radius may be obtained by solving the equation derived from Eq. (31) by setting $`|\sigma |=0`$,
$$1\frac{g^2}{4\pi ^2}_0^\mathrm{\Lambda }๐x\frac{x^2}{\mathrm{tanh}(\pi xR)}=0.$$
(33)
Eq. (33) determines the relation between $`R`$ and $`g`$ with the suitable choice of the cut-off parameter $`\mathrm{\Lambda }`$.
It is interesting to note here that, if we started from 3 dimensions instead of 5 dimensions and regarded the 2-dimensional world as a physical world, we would have obtained an equation for the critical radius corresponding to Eq. (33) as follows,
$$1\frac{g^2}{2\pi }_\lambda ^\mathrm{\Lambda }๐x\frac{1}{\mathrm{tanh}(\pi xR)}=0,$$
(34)
where parameter $`\lambda `$ is the cut-off of the low momentum integration which is needed to remove an infrared divergence created by letting $`\sigma 0`$. The integration in Eq. (34) is easily performed and we find that
$$\frac{g^2}{2\pi ^2}=R/\mathrm{ln}\frac{\mathrm{sinh}(\pi \mathrm{\Lambda }R)}{\mathrm{sinh}(\pi \lambda R)}.$$
(35)
Eq. (35), if it is inverted, gives us a formula for the critical radius as a function of the four-fermion coupling constant. Similarly Eq. (33) gives us a relation between $`R`$ and $`g`$ if we solve the equation numerically.
## V Generation of Dynamical Fermion Masses
As was shown in Eq. (23) the lowest fermion mass on the 4-dimensional brane is $`m=Ng<|\sigma |>`$ where $`<|\sigma |>`$ is determined by solving Eq. (31). By using the numerical estimation we calculate the fermion mass $`m`$ as a function of radius $`R`$ for fixed coupling constant $`g`$. The result is shown in Fig. 2 where $`m`$, $`R`$ and $`g`$ are normalized by the cut-off parameter $`\mathrm{\Lambda }`$. It is immediately recognized that the fermion mass generation takes place below the critical radius and so is kept small near the critical radius. The critical curve which represents the critical radius as a function of the coupling constant is shown in Fig. 3 as the curve for $`m=0`$. It is easy to see that $`g^2\mathrm{\Lambda }^3=12\pi ^2`$ as $`R\mathrm{}`$ on the critical curve.
In our model the bulk fermions reside in the 5-dimensional space-time and their Kaluza-Klein modes show up as fermions with masses $`n/R`$ in the 4-dimensional space-time. Since those fermions are not observed in the present experimental situation except for the zero mode which mixes with the fermion on the brane, their masses have to be very high and hence $`1/R`$ should be much higher than several TeV. Thus the compactification scale $`R`$ for the fifth dimension is considered to be very small. On the other hand the bulk graviton lives in the higher dimensional space-time and the compactification scale $`1/R_G`$ for the extra-dimensional space associated with the bulk graviton could be smaller than $`1/R`$ as is suggested in Ref. 2. Our cut-off scale $`\mathrm{\Lambda }`$ is introduced to suppress the divergence appearing in the integration in the effective potential. Hence it is to be determined by the region of the validity of our effective four-fermion theory. We suppose that $`\mathrm{\Lambda }>1/R`$ since the fundamental theory which derives our effective theory is realized at much higher scale than the compactification scale for the bulk fermions. In Fig. 3 we recognize that, if $`g^2\mathrm{\Lambda }^312\pi ^2`$, the fermion mass is kept small for a wide range of radius $`R`$ except for the small $`R`$ region.
## VI Conclusions
We have found within our model that in spite of the presence of the large mass scale $`1/R`$ in the theory the fermion masses on the 4-dimensional brane can be made small as a consequence of the interaction among the bulk and brane fermions: the mixing of the brane fermions with the bulk fermions does not lead to the lightest fermion masses of order $`1/R`$ and also the dynamically generated fermion masses are not of order $`1/R`$. This result is obtained because the dynamical fermion masses generated under the second-order phase transition are small irrespectively of $`1/R`$ near the critical radius. In our model the possibility of having low mass fermions resulted from the dynamical origin. This mechanism is quite different from the ones in other approaches in which low mass fermions are expected to show up as a result of the kinematical origins.
It is tempting to make a final comment on the conjecture that the nature chooses the compactification scale near the critical radius so as to have low mass fermions. In this connection it may be interesting to suppose that the existence of the critical radius is a result of the physical process that the compactification of the fifth dimensional space is driven by the force unknown to us and this force is balanced by the pressure coming from the fermionic Casimir energy.
###### Acknowledgements.
The authors would like to thank Masahiro Yamaguchi (Tohoku U.) and Koichi Yoshioka (Kyoto U.) for fruitful discussions and correspondences. They are also indebted to Tak Morozumi for useful comments and to Keith Dienes, Emilian Dudas and Tony Gherghetta for a useful correspondence. The present work is supported financially by the Monbusho Grant, Grant-in-Aid for Scientific Research (C) with contract number 11640280. |
warning/0002/astro-ph0002205.html | ar5iv | text | # Photometric catalog of nearby globular clusters (I) Based on data collected at the European Southern Observatory, La Silla, Chile.
## 1 Introduction
There are two main properties which make the study of the Galactic globular clusters (GGC) particularly interesting: 1) each cluster (with possible rare exceptions) is made up by a single population of stars, all born at the same time, in the same place, and out of the same material; 2) GGC stars have the oldest measurable age in the Universe, and therefore we believe they are the oldest fossil records of the formation history of our Galaxy.
Among the many tools we have to investigate the properties of a stellar population, the color-magnitude diagrams (CMD) are the most powerful ones, as they allow to recover for each individual star its evolutionary phase, giving precious information on the age of the entire stellar system, its chemical content, and its distance. This information allows us to locate the system in the space, giving a base for the distance scale, study the formation histories of the Galaxy, and test our knowledge of stellar evolution models.
In particular, the study of a large sample of simple stellar systems, as the GGCs, provides important clues to the Milky Way formation history. Recently, many studies on the relative ages of the GGCs have been presented with results at least controversial: while some authors find a notable age spread ($`5`$ Gyrs) among the clusters, others find that the bulk of GGCs is coeval. This controversy is surely mainly due to the heterogeneity of the data used in each study, where the combination of photographic and/or CCD data from the early epochs of solid state detectors has been frequently used. For this reason, a survey of both southern and northern GGCs has been started two years ago by means of 1-m class telescopes, i.e. the 91cm European Southern Observatory (ESO) / Dutch telescope and the 1m Isaac Newton Group (ING) / Jacobus Kapteyn telescope (JKT). We were able to collect the data for 52 of the 69 known GGCs with $`(mM)_V16.15`$. Thirty-nine have been observed with the Dutch telescope (data that are presented in this paper, hereafter Paper I), and the remaining ones with the JKT (the corresponding CMDs will be presented in a companion paper, Rosenberg et al. rosenberg00 (2000), hereafter Paper II).
As a first exploitation of this new data base, we have conducted a GGC relative age investigation based on the best 34 CMDs of our catalog (Rosenberg et al. rosenberg99 (1999), hereafter Paper III), showing that most of the GGCs have the same age. We have also used our data base to obtain a photometric metallicity ranking scale (Saviane et al. saviane00 (2000), hereafter Paper IV), based on the red giant branch (RGB) morphology. We measured a complete set of metallicity indices, based on the morphology and position of the RGB. Using a grid of selected RGB fiducial points, we defined a function in the $`(VI)_0`$, $`M_\mathrm{I}`$, \[Fe/H\] space which is able to reproduce the whole set of GGC RGBs in terms of a single parameter (the metallicity). The use of this function will improve the current determinations of metallicity and distances within the Local Group.
There are many other parameters that can be measured from a homogeneous, well calibrated CMD data base: the horizontal branch (HB) level, homogeneous reddening and distances, etc. We are presently working on these problems. However, we believe it is now the time to present to the community this data base to give to anyone interested the opportunity to take advantage of it.
In the next section, we will describe the observations collected at the ESO/Dutch telescope during two runs in 1997. The data reduction and calibration is presented in Sect. 3, while in Sect. 4 a cross check of the calibration between the two runs is given. In order to facilitate the readerโs work, we have included the main parameters characterizing our clusters in Sect. 5. Finally, the observed fields for each cluster, and the obtained CMDs are presented and briefly discussed in Sect. 6.
## 2 Observations
The data were collected during two runs in 1997: the first in April ($`11^{\mathrm{th}}15^{\mathrm{th}}`$) and the second in December ($`23^{\mathrm{rd}},24^{\mathrm{th}}`$ and $`26^{\mathrm{th}}`$). All nights of the first run and the first two of the second run were photometric and had a stable seeing.
Observations were done with the ESO 91cm DUTCH telescope, at La Silla (Chile). The same same CCD$`\mathrm{\#}33`$ was used in both runs, a thinned CCD with $`512\times 512`$ pixels, each projecting $`0.\mathrm{}442`$ on the sky, with a total field of view of $`3.77\times 3.77(\mathrm{})^2`$, and the same set of $`V`$ Johnson and $`i`$ Gunn filters.
Two short ($`1045`$s), one medium ($`90120`$s) and one long ($`6001800`$s) exposures were taken in each band (depending on the cluster distance modulus) for one to three fields (in order to ensure a statistically significant sample of stars) for each of the proposed objects. Also a large number of Landolt (landolt92 (1992)) standard stars were measured during each night.
In Table 1 the 39 observed GGCs are presented. Column 1 gives an identification number adopted in this paper; cluster NGC numbers and alternative names are given in columns 2 and 3. The observing dates are in column 4, the number of covered fields in column 5, the mean seeing for each filter in column 6, and the integration time for the long exposures in column 7. In Fig. 1 we show the heliocentric distribution of the clusters of our entire catalog.
## 3 Data reduction and calibration
The images were corrected for a constant bias, dark current, and for spatial sensitivity variations using the respective master flats, computed as the median of all available sky flats of the specific run. Afterwards, photometry was performed using the DAOPHOT/ALLSTAR/ALLFRAME software, made available to us by Dr. Stetson (see Stetson stetson87 (1987), stetson94 (1994)). A preliminary photometry was carried out in order to construct a short list of stars for each single frame. This list was used to accurately match the different frames. With the correct coordinate transformations among the frames, we obtained a single image, median of all the frames, regardless of the filter. In this way we could eliminate all the cosmic rays and obtain the highest signal/noise image for star finding. We ran the DAOPHOT/FIND routine on the median image and performed PSF fitting photometry in order to obtain the deepest list of stellar objects free from spurious detections. Finally, this list was given as input to ALLFRAME, for the simultaneous profile fitting photometry of all the individual frames. We constructed the model PSF for each image using typically from 60 to 120 stars.
The absolute calibration of the observations to the V-Johnson and I-Cousins systems is based on a set of standard stars from the catalog of Landolt (landolt92 (1992)). Specifically, the observed standard stars were in the fields: PG0231, SA95 (41, 43, 96, 97, 98, 100, 101, 102, 112, 115), SA98 (556, 557, 563, 580, 581, L1, 614, 618, 626, 627, 634, 642), RUBIN 149, RUBIN 152, PG0918, PG0942, PG1047, PG1323, PG1525, PG1530, PG1633, and Mark A. At least 3 exposures were taken for each standard field, with a total of $`100`$ standard star measurements per night and per filter.
The reduction and aperture photometry of standard star fields were performed in the same way as for the cluster images. The aperture magnitudes were corrected for atmospheric extinction, assuming $`A_V=0.14`$ and $`A_I=0.08`$ as extinction coefficients for the $`V`$ and $`i`$ filters, respectively.
As shown in Fig. 2, a straight line well reproduces the calibration equations. As the seeing and the overall observing conditions were stable during the run, the slopes of the calibration equations for each observing run and for each filter have been computed using the data from all the nights. As it can be seen in Table 2, the standard deviations of the calibration constants for each run and filter is $`0.015`$mag, corroborating our assumption that all nights were photometric, and that we can assume a constant slope for each filter and run.
Standard stars for which previous problems were reported (PG 1047C, RU149A, RU149G, PG1323A; see Johnson & Bolte johnsonbolte98 (1998)) were excluded, as well as saturated stars, those close to a cosmic ray, etcโฆ After this cleaning, the mean slope was computed, and finally the different night constants were found using this slope to fit the individual data, night by night. The adopted values are presented in Table 2. The typical errors ($`rms`$) are also given.
The calibration curves are shown in Fig. 2 for both runs. In this figure, the dotted line represents the best fitting equation, while the continuous line is obtained by best fitting the data imposing the adopted mean slope. The two lines are almost overlapping. The mean number of standard star measures used for computing the curve per night and filter is $`75`$. Notice the wide color coverage for the standard stars.
The last step on the calibration is the aperture correction. As no available bright and isolated stars exist on the cluster images, we used DAOPHOT to subtract from the image the stars in the neighborhood of the brightest ones, in order to compute the difference between the aperture and the PSF-fitting magnitudes. In view of the stable seeing conditions, we used the same aperture for calculating the aperture photometry of the standard and cluster stars.
## 4 Photometric homogeneity of the two runs
In order to check the photometric homogeneity of the data and of the calibration to the standard photometric system, one cluster (NGC 3201) was observed in both runs. Having one common field, it is possible to analyze the individual star photometry, and test if any additional zero point difference and/or color term exist. The latter check is crucial when measures of the relative position of CMD features are going to be done. The comparison between the two runs is presented in Fig. 3, where 456 common stars with internal photometric errors (as given by ALLFRAME) smaller than 0.02 mag are used. Fig. 3 shows that there are no systematic differences between the two runs.
The slope of the straight lines best fitting all the points in both the (V,$`\mathrm{\Delta }V_{\mathrm{dec}}^{\mathrm{apr}}`$) plane and the (V,$`\mathrm{\Delta }I_{\mathrm{dec}}^{\mathrm{apr}}`$) plane is $`0.001\pm 0.002`$, and $`0.002\pm 0.003`$ in the (V-I,$`\mathrm{\Delta }(VI)_{\mathrm{dec}}^{\mathrm{apr}}`$) plane. The zero point differences are always $`0.01`$ mag. This ensures the homogeneity of our database, particularly for relative measurements within the CMDs.
## 5 Parameters for the GGC sample
In order to facilitate the readers work, we present in Tables 3, 4 and 5 the basic parameters available for our GGCs sample<sup>1</sup><sup>1</sup>1Unless otherwise stated, the data presented in these tables are taken from the McMaster catalog described by Harris (harris96 (1996))..
In Table 3 we give the coordinates, the position, and the metallicity of the clusters: right ascension and declination (epoch J2000, columns 3 and 4); Galactic longitude and latitude (columns 5 and 6); Heliocentric (column 7) and Galactocentric (column 8) distances (assuming $`R_{\mathrm{}}`$=8.0 kpc); spatial components (X,Y,Z) (columns 9, 10 and 11) in the Sun-centered coordinate system (X pointing toward the Galactic center, Y in direction of Galactic rotation, Z toward North Galactic Pole) and, finally, the metallicity given in Rutledge et al. (rutledge97 (1997)), on both the Zinn & West (zinnwest84 (1984)) and Carretta & Gratton (carretagratton97 (1997)) scales.
In Table 4, the photometric parameters are given. Column 3 lists the foreground reddening; column 4, the $`V`$ magnitude level of the horizontal branch; column 5, the apparent visual distance modulus; integrated $`V`$ magnitudes of the clusters are given in column 6; column 7 gives the absolute visual magnitude. Columns 8 to 11 give the integrated color indices (uncorrected for reddening). Column 12 gives the specific frequency of RR Lyrae variables, while column 13 list the horizontal-branch morphological parameter (Lee lee90 (1990)).
In Table 5, we present the kinematical and structural parameters for the observed clusters. Column 3 gives the heliocentric radial velocity (km/s) with the observational (internal) uncertainty; column 4, the radial velocity relative to the local standard of rest; column 5, the concentration parameter ($`c=\mathrm{log}(r_\mathrm{t}/r_\mathrm{c})`$); a โcโ denotes a core-collapsed cluster; columns 6 and 7, the core and the half mass radii in arcmin; column 8, the logarithm of the core relaxation time, in years; and column 9 the logarithm of the relaxation time at the half mass radius. Column 10, the central surface brightness in $`V`$; and column 11, the logarithm of central luminosity density (Solar luminosities per cubic parsec).
## 6 The Color-Magnitude Diagrams
In this section the $`V`$ vs. $`(VI)`$ CMDs for the 39 GGCs and the covered fields are presented.
The same color and magnitude scale has been used in plotting the CMDs, so that differential measures can be done directly using the plots. Two dot sizes have been used, with the bigger ones corresponding to the better measured stars, normally selected on the basis of error ($`0.1`$) and sharpness parameter (Stetson 1987). In some exceptional cases, a selection on radius is also done in order to make evident the cluster stars over the field stars, or to show differential reddening effects. The smaller size dots show all the measured stars with errors (as calculated by DAOPHOT) smaller than 0.15 mag.
The images of the fields are oriented with the North at the top and East on the left side. As explained in Sect. 2, each field covers $`3.77\times 3.77(\mathrm{})^2`$, and the overlaps between fields of the same object are about $`2025\%`$ of the area. For some clusters, only short exposures were obtained for the central fields.
In the next subsections, we present the single CMDs and clusters, and give some references to the best existing CMDs. This is by no means a complete bibliographical catalog: a large number of CMDs are available in the literature for many of the clusters of this survey, but we will concentrate just on the best CCD photometric works. The tables with the position and photometry of the measured stars will be available via a web interface at IAC and Padova in the near future.
### NGC 104 (47 Tucanae).
(Fig. 4)
The cluster 47 Tucanae is (after $`\omega `$ Centauri) the second brightest globular cluster in the sky, and consequently a lot of work has been done on this object. 47 Tucanae has been often indicated as the prototype of the metal-rich GGCs, characterized by a well populated red HB (RHB) clump and an extended RGB that also in our CMD spans $``$2 mag in color from the RHB to the reddest stars at the tip.
A classical CMD of 47 Tucanae is that presented by Hesser et al. (hesser87 (1987)) where a composite CMD was obtained from the superposition of $`B`$ and $`V`$ CCD photometry for the main sequence (MS) and photographic data for the evolved part of the diagram. The same year, Alcaino & Liller (1987a ) published a BVI CCD photometry. One year later, Armandroff (armandroff88 (1988)) presented the RGB $`V`$ and $`I`$ bands photometry for this cluster (together with other five). In 1994, Sarajedini & Norris (sarajedininorris94 (1994)) presented a study of the RGB and HB stars in the $`B`$ and $`I`$ bands. Sosin et al. (1997a ) and Rich et al. (rich97 (1997)) have published a $`B`$,$`V`$ photometry based on HST data.
A recent work in the $`V`$ and $`I`$ bands has been presented by Kaluzny et al. (kaluzny98 (1998)), who focussed their study on the variable stars. They do not find any RR-Lyrae, but many other variables (mostly located in the BSS region), identified as binary stars. As already stated by these authors, a small difference is found between their and our photometry. Indeed, their magnitudes coincide with ours at $`12.5`$ mag in both bands, but there is a small deviation from linearity of $`0.015`$ magnitudes per magnitude (with the Kaluzny et al. stars brighter than ours), in both bands (computed from 90 common stars with small photometric errors) within a magnitude range of $`3`$ mag. We are confident that our calibration, within the quoted errors, is correct, as further confirmed by the comparison with other authors for other objects, as discussed below. Although small, these differences could be important in relative measures, if they appear randomly in different CMDs. For example, in this case, the $`\mathrm{\Delta }V_{\mathrm{TO}}^{\mathrm{HB}}`$ parameter is $`0.05`$ mag smaller in Kaluzny et al. (kaluzny98 (1998)) CMD than in our one, implying, for the 47 Tuc metallicity, an age difference of $`0.8`$ Gyrs. We want to stress the importance of a homogeneous database for a reliable measurement of differential parameters on the CMDs.
### NGC 288 and NGC 362.
(Figs. 5 and 6)
The diagram of NGC 288 is well defined and presents an extended blue horizontal branch (EBHB) which extends from the blue side of the RR-Lyrae region, to just above the TO. Conversely, NGC 362 has a populated RHB with just a few blue HB stars.
These two clusters define one of the most studied second parameter couple: despite their similar metallicities, their HB morphologies are different. Much work have been done on both clusters in order to try to understand the origin of such differences: Bolte (bolte89 (1989)) and Sarajedini & Demarque (sarajedinidemarque90 (1990)) in the $`B`$ and $`V`$ bands, and Green & Norris (greennorris90 (1990)) in the $`B`$ and $`R`$ bands, based on homogeneous CCD photometry, obtain an age difference of $`3`$ Gyrs, NGC 288 being older than NGC 362. A similar conclusion is obtained in our study (Paper III), where NGC 362 is found $`20\%`$ younger than NGC 288. It has also been proposed (e.g. Green & Norris greennorris90 (1990)) that these age differences might be responsible of the HB differences between the two clusters. On the other side, Buonanno et al. (buonanno98 (1998)) and Salaris & Weiss (salarisweiss98 (1998)) do not find significant age differences. Another $`B`$,$`V`$ photometry of NGC 362 based on HST data is in Sosin et al. (1997a ).
It might be worth to remark here that, as it will be discussed in Paper II, there are clusters with different HB morphologies, though with the same metallicities and ages (within errors). This means that the analysis of a single couple of GGCs can not be considered conclusive for understanding the second parameter problem, while a large scale study (as that feasible with this catalog) can be of more help.
### NGC 1261.
(Fig. 7)
This cluster is the object with the largest distance in our southern hemisphere sample. It is located at $`16`$ kpc from the Sun.
Three major CCD CMDs have been published for NGC 1261: Bolte & Marleau (boltemarleau89 (1989)) in $`B,V`$, Alcaino et al. (alcaino92 (1992)) in $`B,V,R,I`$, and Ferraro et al. (1993b ) in the $`B`$ and $`V`$ bands.
The CMD is characterized by an HB which is similar to the HB of NGC 1851. From here on, clusters with an HB well populated both on the red and blue side of the RR-Lyrae gap will be named bimodal HB clusters, though a more objective classification would require taking into account the color distribution of stars along the HB including the RR Lyrae (Catelan et al. catelan98 (1998)). NGC 1261 has a metallicity very close to that of the previous couple; Chaboyer et al. (chaboyer96 (1996)), Richer et al. (richer96 (1996)) and Rosenberg et al. (Paper III) find that it is younger (similar in age to NGC 362) than the bulk of GGCs. A blue straggler (BS) is clearly visible in Fig. 7.
### NGC 1851.
(Fig. 8)
This cluster has a bimodal HB, with very well defined RHB and blue HB (BHB). Also in this case, a BS sequence is visible in Fig. 8. It is curious that, again, a bimodal cluster results to be younger than the GGCs bulk. From the 34 clusters in the present catalog, only 4 result to be surely younger, i.e. the already described NGC 362 and NGC 1261, this cluster, and NGC 2808: three of them have a bimodal HB (cf. Rosenberg et al. rosenberg99 (1999) for a detailed discussion). There exist other two recent $`(V,I)`$ CCD photometries of NGC 1851 by Walker (walker98 (1998)) and Saviane et al. (saviane98 (1998)). The three photometries are all in agreement within the errors, confirming our calibration to the standard system. A CMD of NGC 1851 in the $`B`$,$`V`$ bands from HST is in Sosin et al. (1997a ). Older CCD photometries are found in Alcaino et al (1990a ) ($`B,V,I`$ bands) and Walker (walker92 (1992)) ($`B,V`$ bands)
### NGC 1904 (M 79).
(Fig. 9)
M 79 is the farthest cluster ($`R_{\mathrm{GC}}=18.5`$ kpc) from the Galactic center in our sample. The main feature in the CMD of Fig. 9 is the EBHB. Previous CMDs from CCD photometry are in Heasley et al. (heasley83 (1983)) ($`U,B,V`$ bands), Gratton & Ortolani (grattonortolani86 (1986)) ($`B,V`$ bands), Ferraro et al. (1993a ) ($`B,V`$ bands), Alcaino et al. (alcaino94 (1994)) ($`B,V,R,I`$ bands), and Kravtsov et al. (kravstov97 (1997)) ($`U,B,V`$ bands), and the $`B`$,$`V`$ photometry from HST in Sosin et al. (1997a ).
### NGC 2298.
(Fig. 11)
This cluster is poorly sampled, particularly for the bright part of the diagram (due to problems with a short exposure). Only four BHB stars are present in the HB region. Recent photometric works on this object are in Gratton & Ortolani (grattonortolani86 (1986)) ($`B,V`$ bands), Alcaino & Liller (1986a ) ($`B,V,R,I`$ bands), Janes & Heasley (janesheasley88 (1988)) ($`U,B,V`$ bands), and Alcaino et al. (1990b ) ($`B,V,R,I`$ bands).
### NGC 2808.
(Fig. 11)
This cluster has some differential reddening (Walker 1999), as it can be inferred also from the broadening of the sequences in the CMD of Fig. 11, and a moderate field contamination. The most interesting features of the CMD are the bimodal HB and the EBHB tail with other two gaps, as extensively discussed in Sosin et al. (1997b ). As previously discussed, NGC 2808 is another bimodal HB cluster at intermediate metallicity with a younger age (Rosenberg et al. rosenberg99 (1999)). Apart from the already quoted $`B`$, and $`V`$ band photometry from HST data by Sosin et al. (1997b ), there are many other CCD photometries: Gratton & Ortolani (grattonortolani86 (1986)) ($`B,V`$ bands), Buonanno et al. (buonanno89 (1989)) ($`B,V`$ bands), Ferraro et al. (ferraro90 (1990)) ($`B,V`$ bands), Alcaino et al. (1990c ) ($`B,V,R,I`$ bands), Byun & Lee (byunlee93 (1993)), Ferraro et al. (ferraro97 (1997)) ($`V,I`$ bands), and more recently Walker (walker99 (1999)) ($`B,V`$ bands).
### E3.
(Fig. 12)
This cluster is one of the less populated clusters in our Galaxy, resembling some Palomar-like globular as Pal 1 (Rosenberg et al. rosenberg98 (1998)). As in Pal 1, there are no HB stars in the CMD, and the entire population of observed stars is smaller than 1000 objects. E3 is suspected to have a metallicity close to that of Pal 1. From the $`\delta (VI)_{\mathrm{@}2.5}`$ (Paper III) measured on Fig. 12, E3 is coeval with the other GGCs of similar metallicity, though the result is necessarily very uncertain, due to the high contamination and the small number of RGB stars. E3 is the cluster with the better defined MS binary sequence (Veronesi et al. veronesi96 (1996)), which can be also seen in Fig. 12. Previous CCD CMDs are in McClure et al. (mcclure85 (1985)) ($`B,V`$ bands), Gratton & Ortolani (grattonortolani87 (1987)) ($`B,V`$ bands), and Veronesi et al. (veronesi96 (1996)) ($`B,V,R,I`$ bands).
### NGC 3201.
(Fig. 13)
The two lateral fields presented in Fig. 13 were observed in both runs, in order to test the homogeneity of the data and instrumentation (see Sect. 4). The HB of NGC 3201 has a bimodal appearance, though it is not younger than the bulk of GGCs of the same metallicity group, at variance with the previously discussed cases. It has a small differential reddening. A blue straggler (BS) sequence is visible in Fig. 13. Previous CCD studies of this cluster include Penny (penny84 (1984)) ($`B,V,I`$ bands), Alcaino et al. (alcaino89 (1989)) ($`B,V,R,I`$ bands), Brewer et al. (brewer93 (1993)) ($`U,B,V,I`$ bands) and Covino & Ortolani (covinoortolani97 (1997)) ($`B,V`$ bands).
### NGC 4372.
(Fig. 14)
The principal characteristic of the CMD of this cluster is the broadening of all the sequences, consequence of the high differential reddening, probably due to the Coal-sack Nebulae. In the CMD of Fig. 14 the darker dots are from the stars in the lowest reddening region (south east) of the observed fields. We have computed the reddening field for this cluster from the shift of the CMDs obtained in different positions, finding that it is homogeneously distributed in space and quite easy to correct by a second order polynomial surface. Two previous CCD photometries can be found in Alcaino et al. (alcaino91 (1991)) ($`B,V,R,I`$ bands) and Brocato et al. (brocato96 (1996)) ($`B,V`$ bands).
### NGC 4590 (M 68).
(Fig. 15)
This cluster is probably the lowest metallicity cluster of the present sample. It has a well defined CMD, with an HB populated on both sides of the instability strip, and including some RR-Lyrae stars. It has sometimes been classified as one of the oldest GGCs (Salaris et al salaris97 (1997)), and, in fact, we find that M68 is old, though coeval with the rest of the metal poor clusters (Paper III). Other CCD CMDs for this cluster are in McClure et al. (mcclure87 (1987)) ($`B,V`$ bands), Alcaino et al. (1990d ) ($`B,V,R,I`$ bands) and Walker (walker94 (1994)) ($`B,V,I`$ bands).
### NGC 4833.
(Fig. 16)
NGC 4833 is another metal-poor cluster, with an extended BHB, likely with gaps, for which we have not found any previous CCD photometry.
### NGC 5139 ($`\omega `$ Centauri).
(Fig. 17)
NGC 5139 is the intrinsically brightest cluster in our Galaxy. Apart from this, there are many other properties of $`\omega `$ Centauri which make it a very particular object. Its stellar population shows metallicity variations as large as $`1.5`$ dex from star to star (Norris et al. norris96 (1996)). Its overall properties suggest that this clusters could have a different origin from the bulk of GGCs. It has an extended BHB and probably numerous BSS. The broad sequences in the CMD are mainly due to the metallicity variations though likely there is some differential reddening in the field of $`\omega `$ Centauri. Due to its peculiarities, $`\omega `$ Centauri has been (and is!) extensively studied; there is a large number of photometries, and we cannot cite all of them. The most recent and interesting CCD CMDs are in: Alcaino & Liller (1987b ), who present a multi-band ($`B,V,R,I`$) photometry, but poorly sampled, specially for the evolved part of the diagram; Noble et al. (noble91 (1991)) present a deep $`B,V`$ diagram, where the MS is well sampled, but the RGB is not so clear and only 3-5 stars are present in the HB; Elson et al. (elson95 (1995)) present a HST $`V,I`$ photometry of the MS; Lynga (lynga96 (1996)) presents a $`BVRI`$ study of the evolved part of the diagram ($`2`$ mag below the HB); Kaluzny et al. (kaluzny96 (1996), kaluzny97a (1997)) present a $`V,I`$ CMD covering more than $`10^5`$ stars.
### NGC 5897.
(Fig. 18)
NGC 5897 is a metal poor cluster with a blue, not extended HB, typical for its metallicity. All the sequences of Fig. 18 are well defined and populated, including a BS sequence. Two CCD photometric studies exist for this cluster: Sarajedini (sarajedini92 (1992)) ($`B,V`$ bands) and Ferraro et al. (ferraro92 (1992)) ($`U,B,V,I`$ bands).
### NGC 5927.
(Fig. 19)
NGC 5927 has the highest metallicity among the objects of our catalog. It has, as most of the GGCs with \[Fe/H\]$`>0.8`$, a well populated red horizontal branch (RHB), and an extended RGB, which, in our CMD, covers more than $``$2.5 mag in ($`VI`$), from the RHB (partially overlapped with the RGB) to the reddest stars of the RGB tip. It has a high reddening, possibly differential, judging from the broadening of the RGB, and, due to its location (projected towards the Galactic center), the field object contamination (disk and bulge stars) is very high. Previous CCD photometries are in Friel & Geisler (frielgeisler91 (1991)) (Washington photometry), Sarajedini & Norris (sarajedininorris94 (1994)) ($`B,V`$ bands), Samus et al. (samus96 (1996)) ($`B,V,I`$ bands), Sosin et al. (1997a ), and Rich et al. (rich97 (1997)) (HST $`B,V`$ bands).
### NGC 5986.
(Fig. 20)
To our knowledge, this is the first CCD photometry for this cluster. NGC 5986 is an intermediate metallicity cluster, but with a metal-poor like HB. The broadening of the CMD suggests some differential reddening. Contamination by field stars is clearly visible, as expected on the basis of the position within the Galaxy of this cluster.
### NGC 6093 (M 80).
(Fig. 21)
NGC 6093 is a bright and moderately metal poor cluster, and one of the densest globular clusters in the Galaxy. It has an EBHB, which extends well below the TO as clearly visible also in the CMD of Fig. 21, with gaps (Ferraro et al. 1998). Three recent CCD photometries that cover the entire object, with CMD from the brightest stars to above the TO exist for this cluster: Brocato et al. (brocato98 (1998)) ($`B,V`$ bands) and Ferraro et al. (ferraro98 (1998)) (HST $`U,V`$, and far-UV (F160BW) bands). A ground-based multicolor $`U,B,V,I`$ CCD CMD has been published also by Alcaino et al. (alcaino98 (1998)).
### NGC 6101.
(Fig. 22)
NGC 6101 was observed under not very good seeing conditions, and this is the reason for the brighter limiting magnitude. Its CMD has the morphology expected for a metal-poor cluster: the HB is predominantly blue, and the giant branch is steep. In Fig. 22 we note that, starting from the BSS sequence, there is a sequence of stars parallel to the RGB on its blue side. In view of the position of the cluster ($`l,b`$)=(318,$`16`$) these can unlikely be bulge stars; it is possible that on the same line of sight there is an open cluster, though the slope of the two RGBs are quite similar, implying an unlikely similar metallicity. A larger field coverage of NGC 6101 is desirable. The only previous CCD photometry that exists for this cluster is the $`B`$ and $`V`$ study by Sarajedini & Da Costa (sarajedinidacosta91 (1991)), which shows these stars in the same CMD location. However, being the background-foreground stellar contamination heavier, the sequences we discussed can hardly be seen.
### NGC 6121 (M 4).
(Fig. 23)
This cluster is the closest GGC, located approximately at $`2.2`$ kpc from the Sun, though, due to the large reddening caused by the nebulosity in Scorpio-Ophiuchus, it has an apparent visual distance modulus larger than NGC 6397. The reddening is differential, though (as in the case of NGC 4372) it is homogeneously distributed in space. The mean regions of the CMD can be improved using an appropriate second order polynomial fit to the reddening distribution, at least on the two fields shown in Fig. 23. The stars from the southern field have been plotted as darker dots; they are located on the redder (more reddened) part of the CMD. The two most recent CMD of M4 are in Ibata et al. (ibata99 (1999)) ($`V,I,U`$ filters) and Pulone et al. (pulone99 (1999)), who present (near IR) HST studies of the faint part of the MS and of the WD sequence. Other recent CMDs from the RGB tip to below the MSTO are in Alcaino et al. (1997a ), who presented an $`UBVI`$ CCD photometry, and Kanatas et al. (kanatas95 (1995)), who obtained a composite ($`B,V`$) CMD from $`V12`$ to $`V25`$.
### NGC 6171 (M 107).
(Fig. 24)
Previous CCD studies of NGC 6171 are the $`(J,K)`$ and $`(B,V)`$ photometry by Ferraro et al. (ferraro95 (1995) and ferraro91 (1991), respectively). This cluster is affected by a moderate reddening, which could be slightly differential. It has a RHB, with a few stars bluer than the instability strip blue edge.
### NGC 6266 (M 62).
(Fig. 25)
This cluster is located very close to the Galactic center, and it has a high differential reddening. It seems to have both a RHB and a BHB resembling the HB of NGC 1851. Previous $`B,V`$ bands CCD works are in Caloi et al. (caloi87 (1987)), and Brocato et al. (brocato96 (1996)). A de-reddened CMD and RR-Lyrae stars are also studied in Malakhova et al. (malakhova97 (1997)).
### NGC 6304.
(Fig. 26)
NGC 6304 is a high metallicity cluster very close to the Galactic center, and has one of the highest reddenings in our sample. It has some disk and bulge star contamination. There is a second RGB fainter and redder than the main RGB (bulge star contamination or a more absorbing patch?), but the most noticeable feature is the extremely long RGB. The reddest star of its RGB is located $``$3.7 mag redward from the RHB! To our knowledge, this is the most extended RGB known for a GGC. The most recent CCD CMD for this cluster comes from the $`V`$ and $`K`$ photometry by Davidge et al. (davidge92 (1992)) which covers the hottest RGB stars and the HB.
### NGC 6352.
(Fig. 27)
NGC 6352 is another high metallicity bulge GGC, with a CMD typical of a cluster with this metal content. The most recent CCD study on this cluster is in Fullton et al. (fullton95 (1995)), where a $`VI_\mathrm{c}`$ CMD from HST data combined with ground-based observations is presented. Another study of the RGB and HB regions of this cluster is presented by Sarajedini & Norris (sarajedininorris94 (1994)) in the $`B,V`$ bands.
### NGC 6362.
(Fig. 28)
NGC 6362 presents a well defined CMD with a bimodal HB. The most recent CMD on this cluster is given by Piotto et al. (piotto99 (1999)), who present observations of the center of the cluster obtained with the HST/WFPC2 camera in the $`B`$ and $`V`$ bands. The only previous ground-based CCD photometry is in Alcaino & Liller (1986b ). Our field has been also observed in the same filters by Walker (priv. comm.), who made available to us his data for a cross-check of the photometric calibration. We find that the two photometries agree within the errors. In particular, we found a zero point difference of 0.02 mag for the $`V`$ band and 0.01 mag for the $`I`$ band, with a negligible -0.001 color term difference between Walker and our data. These discrepancies are well within the uncertainties, and allow to further confirm our calibration to the standard (Landolt landolt92 (1992)) system.
### NGC 6397.
(Fig. 29)
This cluster is the GGC with the smallest apparent distance modulus. Cool et al. (cool96 (1996)) and King et al. (king98 (1998)) present an extremely well defined CMD of the main sequence of this cluster, from HST data, from just below the TO down to $`I=24.5`$, which correspond to a mass of less than $`0.1M_{}`$. Other HST studies on this cluster have been presented by Burgarella et al. (burgarella94 (1994)), De Marchi & Paresce (demarchiparesce94 (1994)), Cool et al. (cool95 (1995)) and King et al. (king95 (1995)). Many ground-based CCD data have also been published: Auriere et al. (auriere90 (1990)), Anthony-Twarog et al. (anthonytwarog92 (1992)) (Stromgren photometry), Lauzeral et al. (lauzeral92 (1992), lauzeral93 (1993)), Kaluzny (kaluzny97b (1997)) ($`B,V`$ bands) and Alcaino et al. (alcaino87 (1987): $`B,V`$ bands; 1997b : $`U,B,V,I`$ bands).
### NGC 6496.
(Fig. 30)
NGC 6496 is another metal rich GGC which presents an extended RGB. In this case, the reddest stars are $``$ 2 mag redder than the RHB. It has also a remarkably tilted RHB, already noted by Richtler et al. (richtler94 (1994)), who present a CCD $`(B,V)`$ photometry of this cluster; Armandroff (armandroff88 (1988)) gives $`(V,I)`$ CCD photometry. A tilted RHB can be noted not only in this CMD, but also in the CMDs of most of the very metal-rich clusters of our sample. Such a feature is usually not present in the canonical models. The RHB is well populated, and there are two stars located on the BHB region. This is quite unusual considering the metallicity of NGC 6496, and it would be interesting to study the membership and to obtain a CMD on a larger field. Another CCD photometry of this cluster is in Friel & Geisler (frielgeisler91 (1991)) in the Washington system. Sarajedini & Norris (sarajedininorris94 (1994)) present a $`B`$ and $`V`$ photometry for the RGB and HB region.
### NGC 6541.
(Fig. 31)
NGC 6541 is located rather close to the Galactic center, and this explains the high field star contamination of the CMD. It has a BHB, as expected from its metal content. The only previous CCD study of this cluster is the multicolor photometry by Alcaino et al. (1997c ).
### NGC 6544.
(Fig. 32)
This is an example of a terrible โspottyโ field with a high (the highest in our sample) and highly differential reddening, due to the location of NGC 6544, which is very close to the Galactic plane and projected towards the Galactic center. Interestingly enough, despite its intermediate metallicity, there are only BHB stars. Probably, the use of the HST in this case is almost inevitable if we want to estimate the age of this kind of clusters. We have not found any previous CCD photometry of this cluster.
### NGC 6624.
(Fig. 33)
Another member of the metal-rich group is presented in Fig. 33. Despite of being the cluster closest to the Galactic center, NGC 6624 has a moderate field star contamination, and a very well defined RGB and RHB. The reddest stars of the RGB are in this case $`2.2`$ mag redder than the RHB.
Richtler et al. (richtler94 (1994)) present a $`B`$ and $`V`$ CCD CMD of this cluster extending well below the TO, while Sarajedini & Norris (sarajedininorris94 (1994)) present a photometric study of the RGB and HB in the same bands. A $`B,V`$ CMD from HST data is in Sosin & King (sosinking95 (1995)) and Sosin et al. (1997a ).
### NGC 6626 (M 28).
(Fig. 34)
Again a high differential reddening is present in the field of NGC 6626, which is located close to the Galactic center. NGC 6626 seems to have an extended BHB, and maybe a few RHB stars, though the field star contamination makes it rather difficult to see them. Previous CCD photometry is given by Davidge et al. (davidge96 (1996)), who present a deep near infrared photometry.
### NGC 6637 (M 69).
(Fig. 35)
The CMD of NGC 6637 presents the typical distribution in color for the RGB stars discussed for other metal rich clusters, with the reddest stars $`2.4`$ mag redder than the RHB. Previous $`B`$ and $`V`$ CCD photometry is presented by Richtler et al. (richtler94 (1994)), and the RGB-HB region is also studied by Sarajedini & Norris (sarajedininorris94 (1994)) in the same bands.
### NGC 6638.
Fig. 36
Affected by high differential reddening, the CMD this cluster is not very well defined. However, the HB is clearly populated on both sides of the instability strip, and probably there are many RR-Lyrae. We have not found any previous CCD photometries of this cluster.
### NGC 6656 (M 22).
(Fig. 37)
A possible internal dispersion in metallicity has been proposed for M22. It presents an EBHB with some HB stars as faint as the TO, and several possible RR-Lyrae stars. It is close to the Galactic center and to the Galactic plane, with a high reddening.
Piotto & Zoccali (piottozoccali99 (1999)) published the most recent study of this cluster. From a combination of HST data and ground based CCD photometry, they produced a CMD extending from the tip of the RGB to below $`0.2M_{}`$. Anthony-Twarog et al. (anthonytwarog95 (1995)) present $`uvbyCa`$ data for over 300 giant and HB stars, while in Davidge & Harris (davidgeharris96 (1996)) there is a deep near infrared study.
### NGC 6681 (M 70).
(Fig. 38)
NGC 6681 has a predominantly blue HB with a few HB stars on the red side of the instability strip. Brocato et al. (brocato96 (1996)) present the only other available CCD photometry for this cluster in the $`B`$ and $`V`$ bands.
### NGC 6717 (Palomar 9).
(Fig. 39)
NGC 6717 is a poorly populated cluster (as most of the โPalomar-likeโ objects), and the CMD is contaminated by bulge stars. The RGB is difficult to identify, and its HB is blue, resembling that of NGC 288. Notice that there is a very bright field star close to the cluster, located at the north side of it. Brocato et al. (brocato96 (1996)) present the first CCD photometry for this cluster; their $`B`$ and $`V`$ CMD resembles that of Fig. 39. Recently, Ortolani et al. (ortolani99 (1999)) presented a new CMD, in the same bands, but the CMD branches are more poorly defined.
### NGC 6723.
(Fig. 40)
NGC 6723 has both a red and blue HB, and the overall morphology is typical of a cluster of intermediate metallicity. Alcaino et al. (alcaino99 (1999)) present the most recent CCD study (multicolor photometry), with a CMD extending down to $`V21`$. Fullton & Carney (fulltoncarney96 (1996)) have obtained a deep $`B`$ and $`V`$ photometry, extending to $`V24`$, though the results of this study have not been completely published, yet.
### NGC 6752.
(Fig. 41)
NGC 6752 has been largely studied in the past. It has a very well defined EBHB. Penny & Dickens (pennydickens86 (1986)) presented a $`B`$ and $`V`$ CCD study from a combination of data from two telescopes, and published a CMD from the RGB tip to $`V24`$ mag, though with a small number of measured stars. In the same year, Buonanno et al. (buonanno86 (1986)) present a CMD in the same bands for stars from $`1`$ mag above the TO to $`5`$ mag below it. More recently, Renzini et al. (renzini96 (1996)) and Rubenstein & Baylin (rubensteinbaylin97 (1997)) published a CMD from HST data.
### NGC 6809 (M 55).
(Fig. 42)
Also the CMD of NGC 6809 is typical for its (low) metallicity. A very well defined BS sequence is visible in Fig. 42. The most recent CCD study is in Piotto & Zoccali (piottozoccali99 (1999)), who study the cluster luminosity function based on deep HST data combined with ground-based CCD data for the evolved part of the CMD. Zaggia et al. (zaggia97 (1997)) present $`V`$ and $`I`$ CCD photometry of $`34000`$ stars covering an entire quadrant of the cluster (out to $`1.5`$ times the tidal radius) down to $`V21`$. Mateo et al. (mateo96 (1996)) and Fahlman et al. (fahlman96 (1996)) presented photometric datasets of M 55 that have been mainly used to study the age and the tidal extension of the Sagittarius dwarf galaxy. Mandushev et al. (mandushev96 (1996)) published the first deep (down to $`V24.5`$) photometry of the cluster. |
warning/0002/cond-mat0002050.html | ar5iv | text | # Fermiโs golden rule in a mesoscopic metal ring
## I Introduction
Consider a mesoscopic metal ring threaded by a time-dependent magnetic flux $`\varphi (t)`$ that has a static component $`\varphi `$ and a part that oscillates with frequency $`\omega `$,
$$\varphi (t)=\varphi +\varphi _\omega \mathrm{sin}(\omega t).$$
(1)
By Faradayโs law of induction, the oscillating part generates a time-dependent electric field directed along of circumference of the ring, $`E(t)=E_\omega \mathrm{cos}(\omega t)`$, with amplitude
$$eLE_\omega =2\pi \omega \frac{\varphi _\omega }{\varphi _0}.$$
(2)
Here $`L`$ is the circumference of the ring, $`e`$ is the charge of the electron, and $`\varphi _0`$ is the flux quantum. We would like to know the induced current around the ring. In the limit $`\omega 0`$ this is just the usual persistent current. But what happens for frequencies in the range between $`10^8`$ and $`10^{13}`$ Hz, which for experimentally relevant rings corresponds to $`\mathrm{\Delta }\omega \tau ^1`$? Here $`\mathrm{\Delta }`$ is the average level spacing at the Fermi energy, and $`\tau `$ is the elastic lifetime. We use units where $`\mathrm{}`$ is set equal to unity. This problem has first been studied by Kravtsov and Yudson (KY), who found that in quadratic order the time-dependent field induces (among other terms that oscillate) a time-independent non-equilibrium current $`I_0^{(2)}`$. Calculating the disorder average of this current perturbatively, KY found that it has the peculiar property that for frequencies exceeding the Thouless energy $`E_c=\mathrm{}๐/L^2`$ (where $`๐`$ is the diffusion coefficient) the average of $`I_0^{(2)}`$ does not vanish exponentially, but only as $`\omega ^2`$. This is in disagreement with the intuitive expectation that the external frequency $`\omega `$ leads to a similar exponential suppression of this mesoscopic non-equilibrium current as a dephasing rate in the case of the equilibrium persistent current. The perturbative calculation of KY is based on the assumption of a continuous energy-spectrum, which means that the level-broadening due to dephasing, $`1/\tau _\phi `$, must exceed the average level spacing at the Fermi energy, $`\mathrm{\Delta }`$. If we assume that for low temperature $`T`$ the dominant dephasing effect comes from electron-electron interactions, a simple estimate shows that $`1/\tau _\phi (\omega )<\mathrm{\Delta }`$ for $`|\omega |E_c`$ in the limit $`T0`$. Hence, for frequencies smaller than the Thouless energy the spectrum is discrete and the perturbative analysis breaks down. In this work we shall show that in this case the term considered by KY is not constant, but grows linearly in time, a result which can be understood simply in terms of Fermiโs golden rule of time-dependent perturbation theory.
It is important to point out the difference between the current considered here and the direct current due to the usual photovoltaic effect. It is well known that irradiation of a medium without an inversion center by an alternating electric field can give rise to a direct current (photovoltaic effect). The lack of inversion symmetry can be due to impurities and defects in a finite sample. For mesoscopic junctions the photovoltaic direct current has been studied in Ref.. In this case the average current vanishes, because disorder averaging restores the inversion symmetry. In our case, however, we calculate the direct current induced in a mesoscopic ring threaded by a magnetic flux. Because the magnetic flux breaks the time-reversal symmetry, the direct current considered here has a finite disorder average. Thus, the physical origin of a mesoscopic non-equilibrium current discussed in this work is quite different from Ref..
## II The quadratic response function: What is wrong with the Greenโs function approach?
We consider non-interacting disordered electrons of mass $`m`$ on a mesoscopic metal ring threaded by the time-dependent magnetic flux given in Eq.(1). Suppose that we have diagonalized the Hamiltonian in the absence of the oscillating flux (i.e. for $`\varphi _\omega =0`$ in Eq.(1)) for the given realization of the disorder. The time-independent part of the Hamiltonian is then $`\widehat{H}_0=_\alpha \epsilon _\alpha c_\alpha ^{}c_\alpha `$, where $`\epsilon _\alpha `$ are the exact electronic eigen-energies for fixed disorder, which are labeled by appropriate quantum numbers $`\alpha `$. The operators $`c_\alpha ^{}`$ create electrons in the corresponding eigenstates $`|\alpha `$. If we now switch on the time-dependent part of the field, the Hamiltonian becomes $`\widehat{H}=\widehat{H}_0+\widehat{V}(t)`$, with
$`\widehat{V}(t)`$ $`=`$ $`{\displaystyle \frac{2\pi }{mL}}\delta \phi (t){\displaystyle \underset{\alpha ,\beta }{}}\alpha |\widehat{P}_x|\beta c_\alpha ^{}c_\beta `$ (3)
$`+`$ $`{\displaystyle \frac{1}{2m}}\left({\displaystyle \frac{2\pi }{L}}\delta \phi (t)\right)^2{\displaystyle \underset{\alpha }{}}c_\alpha ^{}c_\alpha .`$ (4)
Here $`\delta \phi (t)=(\varphi _\omega /\varphi _0)\mathrm{sin}(\omega t)`$, and $`\widehat{P}_x=id/dx+(2\pi /L)(\varphi /\varphi _0)`$ is the $`x`$-component of the one particle momentum operator. As usual, the coordinate along the circumference is called the $`x`$-direction, and we impose periodic boundary conditions. Using standard non-equilibrium Greenโs function methods, the contribution to the non-equilibrium current that is quadratic in the external field is easily obtained:
$`I^{(2)}(t)`$ $`=`$ $`{\displaystyle \frac{(e)(2\pi )^2}{(mL)^3}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}๐\omega _1๐\omega _2\delta \phi _{\omega _1}\delta \phi _{\omega _2}`$ (5)
$`\times `$ $`K^{(2)}(\omega _1,\omega _2)e^{i(\omega _1+\omega _2)t},`$ (6)
where $`\phi _\omega `$ is the Fourier transform of the time-dependent part of the flux (1) in units of the flux quantum (i.e. $`\varphi (t)\varphi =\varphi _0๐\omega ^{}\delta \phi _\omega ^{}e^{i\omega ^{}t}`$) and the response function $`K^{(2)}(\omega _1,\omega _2)`$ is given by
$`K^{(2)}(\omega _1,\omega _2)`$ $`=`$ $`{\displaystyle \underset{\alpha \beta \gamma }{}}{\displaystyle \frac{P_{\alpha \beta \gamma }}{\epsilon _\gamma \epsilon _\alpha +\omega _1+\omega _2+i0}}`$ (8)
$`\times \left[{\displaystyle \frac{f(\epsilon _\gamma )f(\epsilon _\beta )}{\epsilon _\gamma \epsilon _\beta +\omega _2+i0}}{\displaystyle \frac{f(\epsilon _\beta )f(\epsilon _\alpha )}{\epsilon _\beta \epsilon _\alpha +\omega _1+i0}}\right],`$
with
$$P_{\alpha \beta \gamma }=\alpha |\widehat{P}_x|\beta \beta |\widehat{P}_x|\gamma \gamma |\widehat{P}_x|\alpha .$$
(9)
Here $`f(\epsilon _\alpha )=c_\alpha ^{}c_\alpha `$ is the occupation number, which in a grand-canonical ensemble is the Fermi function. Keeping in mind that the time-dependent part of the flux (1) corresponds to
$$\delta \phi _\omega ^{}=\frac{\varphi _\omega }{2i\varphi _0}[\delta (\omega ^{}+\omega )\delta (\omega ^{}\omega )],$$
(10)
it is clear that in this case Eq.(6) contains not only oscillating terms, but also a time-independent contribution,
$$I_0^{(2)}=A_\omega \left[K^{(2)}(\omega ,\omega )+K^{(2)}(\omega ,\omega )\right],$$
(11)
where
$$A_\omega =\frac{(e)(2\pi \varphi _\omega )^2}{4(Lm)^3\varphi _0^2},$$
(12)
and
$`K^{(2)}(\omega ,\omega )`$ $`=`$ $`{\displaystyle \underset{\alpha \beta \gamma }{}}{\displaystyle \frac{P_{\alpha \beta \gamma }}{\epsilon _\gamma \epsilon _\alpha +i0}}`$ (14)
$`\times \left[{\displaystyle \frac{f(\epsilon _\gamma )f(\epsilon _\beta )}{\epsilon _\gamma \epsilon _\beta \omega +i0}}{\displaystyle \frac{f(\epsilon _\beta )f(\epsilon _\alpha )}{\epsilon _\beta \epsilon _\alpha +\omega +i0}}\right].`$
Defining retarded and advanced Greenโs functions,
$$G_\alpha ^R(\epsilon )=\frac{1}{\epsilon \epsilon _\alpha +i0},G_\alpha ^A(\epsilon )=\frac{1}{\epsilon \epsilon _\alpha i0},$$
(15)
Eq.(14) can also be written as
$`K^{(2)}(\omega ,\omega )`$ $`=`$ $`{\displaystyle \frac{1}{2\pi i}}{\displaystyle \underset{\alpha \beta \gamma }{}}P_{\alpha \beta \gamma }`$ (19)
$`\times \{{\displaystyle _{\mathrm{}}^{\mathrm{}}}d\epsilon f(\epsilon +\omega )[G_\alpha ^R(\epsilon +\omega )G_\beta ^R(\epsilon )G_\gamma ^R(\epsilon +\omega )`$
$`G_\alpha ^A(\epsilon +\omega )G_\beta ^A(\epsilon )G_\gamma ^A(\epsilon +\omega )]`$
$`{\displaystyle _{\mathrm{}}^{\mathrm{}}}๐\epsilon [f(\epsilon +\omega )f(\epsilon )]`$
$`\times `$ $`[G_\alpha ^R(\epsilon +\omega )G_\beta ^A(\epsilon )G_\gamma ^A(\epsilon +\omega )`$ (21)
$`G_\alpha ^R(\epsilon +\omega )G_\beta ^R(\epsilon )G_\gamma ^A(\epsilon +\omega )]\}.`$
The structure of the Greenโs functions agrees with the one given by KY in Ref.. Note, however, that these authors work in a different gauge: they represent the electric field by a scalar potential, so that their expressions contain only a single current vertex. The introduction of Greenโs function is useful for calculating disorder averages. It is common wisdom that for the calculation of the disorder average of Eq.(21) the terms involving products of only retarded or only advanced Greenโs functions can be neglected. In this approximation a perturbative calculation of the disorder average of Eq.(21) has been given by KY, with the result that the associated time-independent part of the non-equilibrium current is proportional to $`\omega ^2`$ for frequencies larger than the Thouless energy. As explained in Sec.I, for frequencies $`\omega <E_c`$ the perturbative expansion is not controlled anymore since the energy spectrum becomes discrete. In fact, it will turn out, that the physical behavior is completely different in this regime.
To demonstrate the break down of the diagrammatic perturbation theory for systems with a discrete spectrum, we now show that an exact evaluation of the disorder average of Eq.(21) should actually yield an infinite result. Let us therefore go back to the exact spectral representation (14) of the response function. Using the formal identity
$$\frac{1}{x+i0}=\mathrm{}\frac{1}{x}i\pi \delta (x),$$
(22)
where $`\mathrm{}`$ denotes the Cauchy principal part, we can rewrite Eq.(14) as
$$K^{(2)}(\omega ,\omega )=K_{\mathrm{}}^{(2)}(\omega ,\omega )+K_{\delta \delta }^{(2)}(\omega ,\omega ),$$
(23)
with
$`K_{\mathrm{}}^{(2)}(\omega ,\omega )`$ $`=`$ $`2{\displaystyle \underset{\alpha \beta \gamma }{}}{\displaystyle \frac{\mathrm{Re}P_{\alpha \beta \gamma }}{\epsilon _\gamma \epsilon _\alpha }}\mathrm{}{\displaystyle \frac{f(\epsilon _\gamma )f(\epsilon _\beta )}{\epsilon _\gamma \epsilon _\beta \omega }},`$ (24)
$`K_{\delta \delta }^{(2)}(\omega ,\omega )`$ $`=`$ $`2\pi ^2{\displaystyle \underset{\alpha \beta \gamma }{}}\mathrm{Re}P_{\alpha \beta \gamma }[f(\epsilon _\gamma )f(\epsilon _\beta )]`$ (25)
$`\times `$ $`\delta (\epsilon _\gamma \epsilon _\alpha )\delta (\epsilon _\gamma \epsilon _\beta \omega ).`$ (26)
The terms with $`\alpha =\gamma `$ in Eqs. (24) and (26) yield the following contributions,
$`K_{\mathrm{},\mathrm{diag}}^{(2)}(\omega ,\omega )`$ $`=`$ $`\mathrm{}{\displaystyle \underset{\alpha \beta }{}}P_{\alpha \beta \alpha }{\displaystyle \frac{}{\epsilon _\alpha }}{\displaystyle \frac{f(\epsilon _\alpha )f(\epsilon _\beta )}{\epsilon _\alpha \epsilon _\beta \omega }}`$ (28)
$`=\mathrm{}{\displaystyle \underset{\alpha \beta }{}}P_{\alpha \beta \alpha }\left[{\displaystyle \frac{\frac{}{\epsilon _\alpha }f(\epsilon _\alpha )}{\epsilon _\alpha \epsilon _\beta \omega }}{\displaystyle \frac{f(\epsilon _\alpha )f(\epsilon _\beta )}{(\epsilon _\alpha \epsilon _\beta \omega )^2}}\right],`$
$`K_{\delta \delta ,\mathrm{diag}}^{(2)}(\omega ,\omega )`$ $`=`$ $`2\pi ^2\delta (0){\displaystyle \underset{\alpha \beta }{}}\mathrm{Re}P_{\alpha \beta \alpha }[f(\epsilon _\alpha )f(\epsilon _\beta )]`$ (30)
$`\times \delta (\epsilon _\alpha \epsilon _\beta \omega ).`$
The right-hand side of Eq.(30) is proportional to the infinite factor $`\delta (0)`$. Hence, the term $`K_{\delta \delta }^{(2)}(\omega ,\omega )`$ must also be infinite. Because the singular prefactor $`\delta (0)`$ in Eq.(30) does not depend on the disorder, this singularity survives disorder averaging. Keeping in mind that Eq.(21) is mathematically equivalent with Eq.(14), we conclude that a correct evaluation of the disorder average $`\overline{K^{(2)}(\omega ,\omega )}`$ must yield an infinite result. Unfortunately, in an approximate evaluation of Eq.(21) by means of the usual diagrammatic methods this $`\delta `$-function singularity is artificially smoothed out, and one obtains a finite result.
## III Adiabatic switching on
The infinite term (30) is clearly unphysical. This term is closely related to the infinitesimal imaginary parts $`i0`$ that have been added to the real frequencies in the spectral representation (14) for the response function $`K^{(2)}(\omega ,\omega )`$. As emphasized by KY, the infinitesimal imaginary parts are a consequence of the fact that the response function must be causal when the time-dependent part of the Hamiltonian is adiabatically switched on. Let us examine the โadiabatic switching onโ of the time-dependent perturbation more carefully. Following the usual recipe, we replace the Hamiltonian $`\widehat{H}_0+\widehat{V}(t)`$ by $`\widehat{H}_0+\widehat{V}_\eta (t)`$, where $`\widehat{V}_\eta (t)=\mathrm{exp}(\eta t)\widehat{V}(t)`$. The limit $`\eta 0`$ is then taken at the end of the calculation of physical quantities. For large enough times $`t`$ the physical result should be independent of the switching on procedure. Indeed, in the appendix we show by explicit calculation that sudden switching on produces the same result for the long-time response as adiabatic switching on. However, in the latter case one still has to be careful to take the limit $`\eta 0`$ only after the physical quantity of interest has been calculated. We now show that the singularity in Eq.(30) has been artificially generated by taking the limit $`\eta 0`$ at an intermediate step of the calculation.
By direct expansion of the time evolution operator in the interaction representation to second order in the time-dependent perturbation, we obtain the current for adiabatic switching on with finite $`\eta `$
$`I_\eta ^{(2)}(t)`$ $`=`$ $`{\displaystyle \frac{(e)(2\pi )^2}{(mL)^3}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}๐\omega _1๐\omega _2\delta \phi _{\omega _1}\delta \phi _{\omega _2}`$ (31)
$`\times `$ $`K_{\eta t}^{(2)}(\omega _1,\omega _2)e^{i(\omega _1+\omega _2)t},`$ (32)
with
$`K_{\eta t}^{(2)}(\omega _1,\omega _2)`$ $`=`$ $`e^{2\eta t}{\displaystyle \underset{\alpha \beta \gamma }{}}{\displaystyle \frac{P_{\alpha \beta \gamma }}{\epsilon _\gamma \epsilon _\alpha +\omega _1+\omega _2+2i\eta }}`$ (34)
$`\times \left[{\displaystyle \frac{f(\epsilon _\gamma )f(\epsilon _\beta )}{\epsilon _\gamma \epsilon _\beta +\omega _2+i\eta }}{\displaystyle \frac{f(\epsilon _\beta )f(\epsilon _\alpha )}{\epsilon _\beta \epsilon _\alpha +\omega _1+i\eta }}\right].`$
Comparing Eq.(34) with Eq.(8), we see that the former is multiplied by an extra factor of $`e^{2\eta t}`$. If we directly take the limit $`\eta 0`$, this factor is replaced by unity. This is the limiting procedure adopted in the usual Greenโs function approach, where one takes first the limit $`\eta 0`$ in Eq.(34) and then inserts the resulting expression into Eq.(32). In this case we recover Eqs.(11) and (14), which lead to the divergence in Eq.(30). We now show that this unphysical divergence does not appear if the limit $`\eta 0`$ is taken after the physical current has been calculated. Substituting Eq.(34) into Eq.(32) we obtain
$`I_\eta ^{(2)}(t)`$ $`=`$ $`A_\omega [K_{\eta t}^{(2)}(\omega ,\omega )+K_{\eta t}^{(2)}(\omega ,\omega )`$ (36)
$`+K_{\eta t}^{(2)}(\omega ,\omega )e^{2i\omega t}+K_{\eta t}^{(2)}(\omega ,\omega )e^{2i\omega t}].`$
In analogy with Eq.(23), we express $`K_{\eta t}^{(2)}(\omega ,\omega )`$ in terms of products of real and imaginary parts
$$K_{\eta t}^{(2)}(\omega ,\omega )=K_{\eta t,\mathrm{}}^{(2)}(\omega ,\omega )+K_{\eta t,\delta \delta }^{(2)}(\omega ,\omega ),$$
(37)
with
$`K_{\eta t,\mathrm{}}^{(2)}(\omega ,\omega )`$ $`=`$ $`2e^{2\eta t}{\displaystyle \underset{\alpha \beta \gamma }{}}\mathrm{Re}P_{\alpha \beta \gamma }[f(\epsilon _\gamma )f(\epsilon _\beta )]`$ (39)
$`\times \left[{\displaystyle \frac{\epsilon _\gamma \epsilon _\alpha }{(\epsilon _\gamma \epsilon _\alpha )^2+(2\eta )^2}}{\displaystyle \frac{\epsilon _\gamma \epsilon _\beta \omega }{(\epsilon _\gamma \epsilon _\beta \omega )^2+\eta ^2}}\right],`$
$`K_{\eta t,\delta \delta }^{(2)}(\omega ,\omega )`$ $`=`$ $`2e^{2\eta t}{\displaystyle \underset{\alpha \beta \gamma }{}}\mathrm{Re}P_{\alpha \beta \gamma }[f(\epsilon _\gamma )f(\epsilon _\beta )]`$ (41)
$`\times \left[{\displaystyle \frac{2\eta }{(\epsilon _\gamma \epsilon _\alpha )^2+(2\eta )^2}}{\displaystyle \frac{\eta }{(\epsilon _\gamma \epsilon _\beta \omega )^2+\eta ^2}}\right].`$
From Eq.(39) it is now obvious that $`K_{\eta t,\mathrm{}}^{(2)}`$ does not have any contributions from the terms $`\alpha =\gamma `$. The finite contribution in Eq.(28) is thus an artifact of taking the limit $`\eta 0`$ before calculating any physical quantities. Let us now focus on the term (41). If we directly take the limit $`\eta 0`$ using
$$\underset{\eta 0}{lim}\frac{\eta }{ฯต^2+\eta ^2}=\pi \delta (ฯต),$$
(42)
we recover the infinite result (30). However, the structure of the $`\eta `$-dependent part of Eq.(41) is familiar from the derivation of Fermiโs golden rule of elementary quantum mechanics. As discussed for example in the classic textbook by Baym, terms with this structure should be interpreted as a rate, i.e. as a contribution to the current that grows linearly in time. It is therefore clear that after taking the derivative of Eq.(41) with respect to $`t`$ we obtain a finite result if we then let $`\eta 0`$. A simple calculation yields
$`\underset{\eta 0}{lim}{\displaystyle \frac{d}{dt}}K_{\eta t,\delta \delta }^{(2)}(\omega ,\omega )`$ (45)
$`=2\underset{\eta 0}{lim}{\displaystyle \underset{\alpha \beta }{}}P_{\alpha \beta \alpha }{\displaystyle \frac{[f(\epsilon _\alpha )f(\epsilon _\beta )]\eta }{(\epsilon _\alpha \epsilon _\beta \omega )^2+\eta ^2}}`$
$`=2\pi {\displaystyle \underset{\alpha \beta }{}}P_{\alpha \beta \alpha }[f(\epsilon _\alpha )f(\epsilon _\beta )]\delta (\epsilon _\alpha \epsilon _\beta \omega ).`$
Because this expression contains only a single $`\delta `$-function, after averaging over disorder it becomes a smooth function of $`\omega `$. We conclude that to quadratic order in the field the non-equilibrium current induced by the time-dependent flux (1) has the following three contributions,
$$I^{(2)}(t)\underset{\eta 0}{lim}I_\eta ^{(2)}(t)=I_{\mathrm{th}}^{(2)}+t\frac{dI_{\mathrm{kin}}^{(2)}}{dt}+I_{\mathrm{osc}}^{(2)}(t),$$
(46)
where the time-independent part is given by
$`I_{\mathrm{th}}^{(2)}`$ $`=`$ $`A_\omega \underset{\eta 0}{lim}\left[K_{\eta t,\mathrm{}}^{(2)}(\omega ,\omega )+K_{\eta t,\mathrm{}}^{(2)}(\omega ,\omega )\right]`$ (47)
$`=`$ $`2A_\omega {\displaystyle \underset{\alpha \beta \gamma ,\alpha \gamma }{}}{\displaystyle \frac{\mathrm{Re}P_{\alpha \beta \gamma }}{\epsilon _\gamma \epsilon _\alpha }}`$ (48)
$`\times `$ $`\mathrm{}\left[{\displaystyle \frac{f(\epsilon _\gamma )f(\epsilon _\beta )}{\epsilon _\gamma \epsilon _\beta \omega }}+(\omega \omega )\right].`$ (49)
The coefficient of the term linear in time is
$`{\displaystyle \frac{dI_{\mathrm{kin}}^{(2)}}{dt}}`$ $`=`$ $`A_\omega \underset{\eta 0}{lim}\left[{\displaystyle \frac{d}{dt}}K_{\eta t,\delta \delta }^{(2)}(\omega ,\omega )+{\displaystyle \frac{d}{dt}}K_{\eta t,\delta \delta }^{(2)}(\omega ,\omega )\right]`$ (50)
$`=`$ $`2\pi A_\omega {\displaystyle \underset{\alpha \beta }{}}P_{\alpha \beta \alpha }[f(\epsilon _\alpha )f(\epsilon _\beta )]`$ (52)
$`\times \left[\delta (\epsilon _\alpha \epsilon _\beta \omega )+(\omega \omega )\right],`$
and the oscillating part is
$$I_{\mathrm{osc}}^{(2)}(t)=A_\omega \underset{\eta 0}{lim}\left[K_{\eta t}^{(2)}(\omega ,\omega )e^{2i\omega t}+K_{\eta t}^{(2)}(\omega ,\omega )e^{2i\omega t}\right].$$
(53)
Thus, a time-dependent electric field with frequency $`\omega `$ induces in quadratic order three fundamentally different currents. (a) A time-independent contribution $`I_{\mathrm{th}}^{(2)}`$; as shown in the next section, this contribution can be derived from a thermodynamic calculation. (b) A contribution $`tdI_{\mathrm{kin}}^{(2)}/dt`$ which increases linearly in time; this term can be understood in terms of the usual golden rule of time-dependent perturbation theory. (c) Finally, there is also a time-dependent contribution $`I_{\mathrm{osc}}^{(2)}`$ oscillating with frequency $`2\omega `$. When this term is averaged over a time-interval larger than $`\omega ^1`$, its contribution to the current is negligibly small.
From the above analysis it is clear that the contribution that is proportional to $`t`$ cannot be calculated within the usual Greenโs function machinery, because in this approach the limit $`\eta 0`$ is taken at an intermediate step of the calculation, causing an unphysical divergence. To further support the correctness of the limiting procedure adopted here we show in the appendix that Eqs.(4652) can also be re-derived if the perturbation is suddenly (instead of adiabatically) switched on.
## IV The thermodynamic origin of the time-independent part of the current
The time-independent part $`I_{\mathrm{th}}^{(2)}`$ of the non-equilibrium current in Eq.(46) has been discussed by us in Ref.. This contribution can be obtained from a thermodynamic calculation. In Ref. we have assumed (without further justification) the existence of such a relation. Let us now put this assumption on a more solid basis. Within the Matsubara (imaginary time) formalism one can directly calculate the imaginary frequency version of the response function $`K^{(2)}(\omega _1,\omega _2)`$ given in Eq.(8), i.e.
$`K^{(2)}(i\omega _1,i\omega _2)`$ $`=`$ $`{\displaystyle \underset{\alpha \beta \gamma }{}}{\displaystyle \frac{P_{\alpha \beta \gamma }}{\epsilon _\gamma \epsilon _\alpha +i\omega _1+i\omega _2}}`$ (54)
$`\times `$ $`\left[{\displaystyle \frac{f(\epsilon _\gamma )f(\epsilon _\beta )}{\epsilon _\gamma \epsilon _\beta +i\omega _2}}{\displaystyle \frac{f(\epsilon _\beta )f(\epsilon _\alpha )}{\epsilon _\beta \epsilon _\alpha +i\omega _1}}\right].`$ (55)
As pointed out by KY, in order to obtain the causal response function, one should first continue both frequencies to the real axis with positive imaginary part ($`i\omega _1\omega _1+i0`$, $`i\omega _2\omega _2+i0`$), and then set $`\omega _1=\omega _2`$ to obtain the constant part of the physical current. On the other hand, if one performs these steps in opposite order (i.e. first sets $`i\omega _1=i\omega _2`$ and then continues $`i\omega _1\omega +i0`$) one obtains for the current response function
$`K_{\mathrm{th}}^{(2)}(\omega ,\omega )`$ $`=`$ $`\mathrm{Re}{\displaystyle \underset{\alpha \beta \gamma }{}}{\displaystyle \frac{P_{\alpha \beta \gamma }}{\epsilon _\gamma \epsilon _\alpha }}`$ (57)
$`\times \left[{\displaystyle \frac{f(\epsilon _\gamma )f(\epsilon _\beta )}{\epsilon _\gamma \epsilon _\beta \omega i0}}{\displaystyle \frac{f(\epsilon _\beta )f(\epsilon _\alpha )}{\epsilon _\beta \epsilon _\alpha +\omega +i0}}\right].`$
Comparing this expression with Eqs.(2326), it is easy to see that
$$K_{\mathrm{th}}^{(2)}(\omega ,\omega )=K_{\mathrm{}}^{(2)}(\omega ,\omega ).$$
(58)
Hence, the time-independent part $`I_{\mathrm{th}}^{(2)}`$ of the current can indeed be obtained from a thermodynamic calculation. Note, however, that our analysis of Sec.III (see also the appendix) implies that the terms with $`\alpha =\gamma `$ in Eq.(57) should be omitted from the sum, i.e. the physical current is given by
$$I_{\mathrm{th}}^{(2)}=A_\omega \left[\stackrel{~}{K}_{\mathrm{th}}^{(2)}(\omega ,\omega )+\stackrel{~}{K}_{\mathrm{th}}^{(2)}(\omega ,\omega )\right].$$
(59)
where
$`\stackrel{~}{K}_{\mathrm{th}}^{(2)}(\omega ,\omega )`$ $`=`$ $`K_{\mathrm{}}^{(2)}(\omega ,\omega )K_{\mathrm{},\mathrm{diag}}^{(2)}(\omega ,\omega )`$ (62)
$`={\displaystyle \underset{\alpha \beta \gamma ,\alpha \gamma }{}}{\displaystyle \frac{P_{\alpha \beta \gamma }}{\epsilon _\gamma \epsilon _\alpha }}`$
$`\times \mathrm{}\left[{\displaystyle \frac{f(\epsilon _\gamma )f(\epsilon _\beta )}{\epsilon _\gamma \epsilon _\beta \omega }}{\displaystyle \frac{f(\epsilon _\beta )f(\epsilon _\alpha )}{\epsilon _\beta \epsilon _\alpha +\omega }}\right],`$
see Eq.(28). The direct diagrammatic calculation of the disorder average of $`I_{\mathrm{th}}^{(2)}`$ is difficult, because the restriction $`\alpha \gamma `$ in Eq.(62) is not so easy to implement. In Ref. the following limiting procedure was adopted: Instead of directly calculating $`\overline{\stackrel{~}{K}_{\mathrm{th}}^{(2)}(\omega ,\omega )}`$, consider the generalization of the imaginary frequency response function (55) for electric fields with finite wave-vector $`๐ช`$, which we denote by $`\overline{K^{(2)}(i\omega ,i\omega ,๐ช)}`$. The limit $`๐ช0`$ is taken after the disorder averaged current has been calculated. As shown in Ref., in the diffusive regime the function $`\overline{K^{(2)}(i\omega ,i\omega ,๐ช)}`$ is a smooth function of $`๐ช`$, so that the limit $`๐ช0`$ is well defined. The so-defined averaged response function vanishes for frequencies exceeding the Thouless energy as $`\mathrm{exp}(\sqrt{|\omega |/2E_c})`$. On the other hand, perturbative averaging of the contribution from the (unwanted) diagonal term (28) shows that this term vanishes as $`\omega ^2`$ for large frequencies. This indicates that the above limiting procedure indeed eliminates the contribution of the unphysical diagonal term (28) to the time-independent part of the non-equilibrium current.
## V Conclusion
In this work we have shown that a time-dependent flux oscillating with frequency $`\omega `$ that pierces the center of a mesoscopic metal ring generates to quadratic order three fundamentally different contributions to the current: a constant non-equilibrium current $`I_{\mathrm{th}}^{(2)}`$, a current $`tdI_{\mathrm{kin}}^{(2)}/dt`$ that grows linearly in time, and a current oscillating with frequency $`2\omega `$. As shown in Ref., the disorder average of the constant term $`\overline{I_{\mathrm{th}}^{(2)}}`$ vanishes for frequencies exceeding the Thouless energy as $`\mathrm{exp}(\sqrt{|\omega |/2E_c})`$. The calculation of the disorder average of the contribution $`tdI_{\mathrm{kin}}^{(2)}/dt`$ remains an open problem. A direct perturbative calculation by means of the impurity diagram technique is not straightforward, because Eq.(52) involves three matrix elements but only one energy denominator. Therefore this expression cannot be simply written in terms of Greenโs functions.
The main result of this work is the prediction of a current $`tdI_{\mathrm{kin}}^{(2)}/dt`$ increasing linearly with time. From the well-known derivation of Fermiโs golden rule it is clear that this result is only valid in an intermediate time interval. In particular, the calculation of the long-time behavior of the non-equilibrium current requires non-perturbative methods.
One should keep in mind that our calculation has been performed for non-interacting electrons in a random potential, so that our results are valid as long as the spectrum of the system is discrete. We have argued in Sec.I that at low enough temperatures this should be the case for small external frequencies, $`|\omega |<E_c`$. On the other hand, for frequencies exceeding $`E_c`$ the spectrum is effectively continuous. In this regime the conventional Greenโs function methods can be used to calculate the direct current, so that the results of KY should be valid.
Let us also point out that the linear time-dependence of the current is a consequence of the discrete spectrum, and is not related to the adiabatic switching on procedure in Eq.(34). In the appendix we show that sudden switching on yields the same linear time-dependence of the current. It seems reasonable to expect that for sufficiently short times the constant part $`I_{\mathrm{th}}^{(2)}`$ of the current is dominant. We would like to encourage experimentalists to measure the non-equilibrium response of mesoscopic metal rings to a time-dependent flux in the frequency range $`10^8\mathrm{Hz}\omega 10^{13}\mathrm{Hz}`$.
## Acknowledgement
This work was supported by the Deutsche Forschungsgemeinschaft (SFB 345). We thank V. E. Kravtsov for his comments.
## Sudden switching on
To confirm that the โswitching on procedureโ outlined in Sec.III yields the correct physical results, let us consider a harmonic perturbation that is suddenly turned on at time $`t=0`$,
$$\varphi (t)=\varphi +\varphi _\omega \mathrm{\Theta }(t)\mathrm{sin}(\omega t),$$
(63)
where $`\mathrm{\Theta }(t)`$ is the step function. To second order in $`\varphi _\omega `$ the induced current is
$`I^{(2)}(t)`$ $`=`$ $`2A_\omega \mathrm{Re}{\displaystyle \underset{\alpha \beta \gamma }{}}P_{\alpha \beta \gamma }[f(\epsilon _\beta )f(\epsilon _\alpha )]`$ (68)
$`\times [{\displaystyle \frac{e^{2i\omega t}e^{i(\epsilon _\gamma \epsilon _\alpha )t}}{(\epsilon _\alpha \epsilon _\gamma +2\omega )(\epsilon _\alpha \epsilon _\beta +\omega )}}`$
$`{\displaystyle \frac{1e^{i(\epsilon _\gamma \epsilon _\alpha )t}}{(\epsilon _\alpha \epsilon _\gamma )(\epsilon _\alpha \epsilon _\beta \omega )}}`$
$`+{\displaystyle \frac{2\omega }{(\epsilon _\alpha \epsilon _\beta )^2\omega ^2}}\left[{\displaystyle \frac{e^{i(\epsilon _\beta \epsilon _\alpha +\omega )t}e^{i(\epsilon _\gamma \epsilon _\alpha )t}}{\epsilon _\beta \epsilon _\gamma +\omega }}\right]`$
$`+(\omega \omega )].`$
The diagonal term $`\alpha =\gamma `$ is
$`I_{\mathrm{diag}}^{(2)}(t)`$ $`=`$ $`4A_\omega {\displaystyle \underset{\alpha \beta }{}}P_{\alpha \beta \alpha }[f(\epsilon _\alpha )f(\epsilon _\beta )]`$ (69)
$`\times `$ $`[{\displaystyle \frac{\mathrm{sin}^2(\omega t)}{(\epsilon _\alpha \epsilon _\beta )^2\omega ^2}}`$ (72)
$`{\displaystyle \frac{\mathrm{sin}^2\left(\frac{\epsilon _\beta \epsilon _\alpha +\omega }{2}t\right)+\mathrm{sin}^2\left(\frac{\epsilon _\beta \epsilon _\alpha \omega }{2}t\right)}{(\epsilon _\alpha \epsilon _\beta )^2\omega ^2}}`$
$`+\left[{\displaystyle \frac{\mathrm{sin}\left(\frac{\epsilon _\beta \epsilon _\alpha +\omega }{2}t\right)}{\epsilon _\beta \epsilon _\alpha +\omega }}\right]^2+\left[{\displaystyle \frac{\mathrm{sin}\left(\frac{\epsilon _\beta \epsilon _\alpha \omega }{2}t\right)}{\epsilon _\beta \epsilon _\alpha \omega }}\right]^2].`$
The terms in the last line can be interpreted in the same way as is done in Fermiโs golden rule by using the identity
$$\left[\frac{\mathrm{sin}\left(\frac{\mathrm{\Delta }\epsilon }{2}t\right)}{\mathrm{\Delta }\epsilon }\right]^2\frac{\pi }{2}t\delta (\mathrm{\Delta }\epsilon )\text{for}t\mathrm{}.$$
(73)
It is now easy to see that for large times $`I_{\mathrm{diag}}^{(2)}(t)`$ yields exactly the same linear in time contribution as given in Eq.(52). The terms with no explicit time dependence in Eq.(68) can be identified with $`I_{\mathrm{th}}^{(2)}`$ in Eq.(49). |
warning/0002/astro-ph0002006.html | ar5iv | text | # ULTRA HIGH ENERGY COSMIC RAYS: the theoretical challenge
## 1 Introduction
The detection of cosmic rays with energies above $`10^{20}`$ eV has triggered considerable interest on the origin and nature of these particles. As reviewed by Watson in this volume, many hundreds of events with energies above $`10^{19}`$ eV and about 20 events above $`10^{20}`$ eV have now been observed by a number of experiments such as AGASA , Flyโs Eye , Haverah Park , Yakutsk , and most recently the High Resolution Flyโs Eye .
Most unexpected is the significant flux of events observed above $`7\times 10^{19}`$ eV with no sign of the Greisen-Zatsepin-Kuzmin (GZK) cutoff . A cutoff should be present if the ultra-high energy particles are protons, nuclei, or photons from extragalactic sources. Cosmic ray protons of energies above a few $`10^{19}`$ eV lose energy to photopion production off the cosmic microwave background (CMB) and cannot originate further than about $`50`$Mpc away from Earth. Nuclei are photodisintegrated on shorter distances due to the infrared background while the radio background constrains photons to originate from even closer systems .
In addition to the presence of events past the GZK cutoff, there has been no clear counterparts identified in the arrival direction of the highest energy events. If these events are protons, cosmic ray observations should finally become astronomy! At these high energies the Galactic and extragalactic magnetic fields do not affect their orbits significantly so that they should point back to their sources within a few degrees. Protons at $`10^{20}`$ eV propagate mainly in straight lines as they traverse the Galaxy since their gyroradii are $``$ 100 kpc in $`\mu `$G fields which is typical in the Galactic disk. Extragalactic fields are expected to be $`\mu `$G , and induce at most $``$ 1<sup>o</sup> deviation from the source. Even if the Local Supercluster has relatively strong fields, the highest energy events may deviate at most $``$ 10<sup>o</sup> . At present, no correlations between arrival directions and plausible optical counterparts such as sources in the Galactic plane, the Local Group, or the Local Supercluster have been clearly identified. Ultra high energy cosmic ray (UHECR) data are consistent with an isotropic distribution of sources in sharp contrast to the anisotropic distribution of light within 50 Mpc from Earth.
The absence of a GZK cutoff and the isotropy of arrival directions are two of the many challenges that models for the origin of UHECRs face. This is an exciting open field, with many scenarios being proposed but no clear front runner. Not only the origin of these particles may be due to physics beyond the standard model of particle physics, but their existence can be used to constrain extensions of the standard model such as violations of Lorentz invariance (see, e.g., ).
In the next section, a brief summary of the challenges faced by all theoretical models is given. In ยง3, astrophysical accelerators or โbottom-upโ scenarios are reviewed, hybrid models are discussed in ยง4, and top-down scenarios in ยง5. To conclude, future observational tests of UHECR models and their implications are discussed in ยง6. For previous reviews of UHECR models, the reader is encouraged to consult .
## 2 The Challenge
In attempting to explain the origin of UHECRs, models confront a number of challenges. The extreme energy is the greatest challenge that models of astrophysical acceleration face while for top-down models the observed flux represents the highest hurdle. To complete the puzzle, models have to match the spectral shape, the primary composition, and the arrival direction distribution of the observed events.
2.1 Energy
The observed highest energy event at $`3.2\times 10^{20}`$ eV argues for the existence of Zevatrons in nature , accelerators that reach as high as one ZeV (ZeV=10<sup>21</sup> eV) which is a billion times the energy limit of current terrestrial accelerators. The energetic requirements at the source may be even more stringent if the distance traveled by the UHE primaries from source to Earth is larger than typical interaction lengths. As can be seen from Figure 1 of (from ), if $`3\times 10^{20}`$ eV is taken as a typical energy for protons travelling in straight lines, accelerators located further than 30 Mpc need to reach above 1 ZeV while those located further than 60 Mpc require over 10 ZeV energies. Depending on the strength and structure of the magnetic field along the primaryโs path, the distance traveled can be significantly larger than the distance to the source. As magnetic fields above $`10^8`$ G may thread extragalactic space , protons travel in curved paths and sources need to be either more energetic or located closer to Earth .
There are great difficulties with finding plausible accelerators for such extremely energetic particles . As discussed in ยง3, even the most powerful astrophysical objects such as radio galaxies and active galactic nuclei can barely accelerate charged particles to energies as high as $`10^{20}`$ eV. If the origin of these events date back to the early universe, then the energy is not as challenging since typical symmetry breaking scales that give rise to early universe relics can be well above the ZeV scale (ยง5).
2.2 Flux
At 10<sup>20</sup> eV, the observed flux of UHECRs is about $``$ 1 event/km<sup>2</sup>/century which has strongly limited our ability to gather more than 20 events after decades of observations . Although challenging to observers, the flux is not particularly constraining in terms of general requirements on astrophysical sources. In fact, this flux equals the flux of gamma-rays in one gamma-ray burst that may have taken place in a 50 Mpc radius volume around us . In terms of an average energy density, UHECRs correspond to $`10^{21}`$ erg/cm<sup>3</sup>, about 8 orders of magnitude less than the cosmic background radiation.
Although less constraining to astrophysical accelerators, flux requirements are very challenging for top-down scenarios. The dynamics of topological defect generation and evolution generally selects the present horizon scale as the typical distance between defects which implies a very low flux. Some scenarios such as monopolia, cosmic necklaces, and vortons have additional scales and may avoid this problem. The possibility of a long lived relic particle that cluster as dark matter can also more easily meet the flux requirements than general top-down models (ยง5).
2.3 Spectrum
The energy spectrum of cosmic rays below the expected GZK cutoff (i.e., between $`10^8`$eV and $`<10^{19}`$eV) is well established to have a steep energy dependence: $`N(E)E^\gamma `$, with $`\gamma 2.7`$ up to the โkneeโ at $`E10^{15}`$ eV and $`\gamma 3.1`$, for $`10^{15}`$ eV $`<E<10^{19}`$ eV. Cosmic rays of energy below the knee are widely accepted to originate in shocks associated with galactic supernova remnants (see, e.g., ), but this mechanism has difficulties producing particles of higher energies . Larger shocks, such as those associated with galactic winds, could reach energies close to the knee and supernova explosions into stellar winds may explain cosmic rays beyond the knee . Although the source of cosmic rays above the knee is not clear, the steepening of the spectrum argues for a similar origin with an increase in losses or decrease in confinement time above the knee. However, the events with energy above $`10^{19.5}`$eV show a much flatter spectrum with $`1<\gamma <2`$. The drastic change in slope suggests the emergence of a new component of cosmic rays at ultra-high energies. This new component is generally thought to be extragalactic , although, depending on its composition, it may also originate in the Galaxy , in an extended halo , or in the dark matter halo . Galactic and halo origins for UHECRs ease the difficulties with the lack of a GZK cutoff but represent an even greater challenge to acceleration mechanisms.
2.4 Propagation - Losses and Magnetic Fields
In order to contrast plausible candidates for UHECR sources with the observed spectrum and arrival direction distribution, the propagation from source to Earth needs to be taken into account. Propagation studies involve both the study of losses along the primariesโ path as well as the structure and magnitude of cosmic magnetic fields that determine the trajectories of charged primaries and influence the development of the electromagnetic cascade (see, e.g., ).
For primary protons the main loss processes are pair production and photopion production off the CMB that gives rise to the GZK cutoff . For straight line propagation, loss processes limit sources of 10<sup>20</sup> eV to be within $``$ 50 Mpc from us and a clear cutoff should be present at $`7\times 10^{19}`$ eV. Even with the small number of accumulated events at the highest energies, the AGASA spectrum seems incompatible with a GZK cutoff for a homogeneous extragalactic source distribution . The shape of the cutoff can be modified if the distribution of sources is not homogeneous and if the particle trajectories are not rectilinear (e.g., the case of sizeable intergalactic magnetic fields) . In fact, if the observed distribution of galaxies in the local universe is used to simulate the range of possible cutoff shapes, the AGASA spectrum is still consistent with sources distributed with the luminous matter given the poor statistics . The need for a new component should become apparent with the increased statistics of future observatories .
Charged particles of energies up to $`10^{20}`$eV can be deflected significantly in cosmic magnetic fields. In a constant magnetic field of strength $`B=B_6\mu `$G, particles of energy $`E=E_{20}10^{20}\mathrm{eV}`$ and charge $`Ze`$ have Larmor radii of $`r_L110`$ kpc $`(E_{20}/B_6Z)`$. If the UHECR primaries are protons, only large scale intergalactic magnetic fields affect their propagation significantly unless the Galactic halo has extended fields . For higher $`Z`$, the Galactic magnetic field can strongly affect the trajectories of primaries .
Whereas Galactic magnetic fields are reasonably well studied, extragalactic fields are still very ill understood . Faraday rotation measures indicate large magnitude fields ($`\mu `$G) in the central regions of clusters of galaxies. In regions between clusters, the presence of magnetic fields is evidenced by synchrotron emission but the strength and structure are yet to be determined. On the largest scales, limits can be imposed by the observed isotropy of the CMB and by a statistical interpretation of Faraday rotation measures of light from distant quasars. The isotropy of the CMB can constrain the present horizon scale fields $`B_{H_0^1}<3\times 10^9`$ G . Although the distribution of Faraday rotation measures have large non-gaussian tails, a reasonable limit can be derived using the median of the distribution in an inhomogeneous universe: for fields assumed to be constant on the present horizon scale, $`B_{H_0^1}<10^9`$ G; for fields with 50 Mpc coherence length, $`B_{50\mathrm{M}\mathrm{p}\mathrm{c}}<6\times 10^9`$ G; while for 1 Mpc coherence length, $`B_{Mpc}<10^8`$ G . These limits apply to a $`\mathrm{\Omega }_bh^2=0.02`$ universe and use quasars up to redshift $`z=2.5`$. Local structures can have fields above these upper limits as long as they are not common along random lines of site between $`z`$ = 0 and 2.5 .
Of particular interest is the field in the local 10 to 20 Mpc volume around us. If the Local Supercluster has fields of about $`10^8`$ G or larger, the propagation of ultra high energy protons becomes diffusive and the spectrum and angular distribution at the highest energies are significantly modified . As shown in Figure 1 (from ), a source with spectral index $`\gamma >2`$ that can reach $`E_{max}>10^{20}`$ eV is constrained by the overproduction of lower energy events around 1 to 10 EeV (EeV $`10^{18}`$ eV). Furthermore, the structure and magnitude of magnetic fields in the Galactic halo or in a possible Galactic wind can also affect the observed UHECRs. In particular, if our Galaxy has a strong magnetized wind, what appears to be an isotropic distribution in arrival directions may have originated on a small region of the sky such as the Virgo cluster . In the future, as sources of UHECRs are identified, large scale magnetic fields will be better constrained .
If cosmic rays are heavier nuclei, the attenuation length is shorter than that for protons due to photodisintegration on the infrared background . However, UHE nuclei may be of Galactic origin. For large enough charge, the trajectories of UHE nuclei are significantly affected by the Galactic magnetic field such that a Galactic origin can appear isotropic . The magnetically induced distortion of the flux map of UHE events can give rise to some higher flux regions where caustics form and some much lower flux regions (blind spots) even for an originally isotropic distribution of sources . Such propagation effects are one of the reasons why full-sky coverage is necessary for resolving the UHECR puzzle.
The trajectories of neutral primaries are not affected by magnetic fields. If associated with luminous systems, sources of UHE neutral primaries should point back to their nearby sources. The lack of counterpart identifications suggests that if the primaries are neutral, their origin involves physics beyond the standard model (ยง4 & ยง5).
2.5 Cosmography
The distribution of arrival directions of UHECRs can in principle hold the key to solving the UHECR puzzle. Within a 50 Mpc radius volume around us, the most well-known luminous structures are the Galactic plane, the Local Group and the large-scale galaxy distribution with a relative overdensity around the Local Supercluster. The Galactic halo is another noteworthy structure that is expected to be a spheroidal overdensity of dark matter centered at the Galactic disk while the dark matter distribution on larger scales correlates with the luminous matter distribution. For the few highest energy events, there is presently no strong evidence of correlations between the eventsโ arrival direction and any of the known nearby luminous structures: the distribution is consistent with isotropy . For slightly lower energies, some correlations may have been detected. For events around 40 EeV, a positive correlation with the Supergalactic plane is found but only at the 1 $`\sigma `$ level . For even lower energies, a more significant correlation was recently announced by AGASA: the arrival direction distribution of EeV events shows a correlation with the Galactic center and the nearby Galactic spiral arms . If confirmed, this correlation would be strong evidence for a Galactic origin of EeV cosmic rays.
2.6 Composition
An excellent discriminator between proposed models is the composition of the primaries. In general, Galactic disk models have to invoke heavier nuclei such as iron to be consistent with the isotropic distribution, while extragalactic astrophysical models tend to favor proton primaries. Photon primaries are more common among top-down scenarios although nucleons can reach comparable fluxes for some models . Experimentally, the composition can be determined by the muon content of the shower in ground arrays and the depth of shower maximum in fluorescence detectors . Unfortunately, the muon content analysis is not very effective at the highest energies. Data from the largest air shower array, AGASA, disfavor photon primaries and indicate a fixed composition across the EeV to 100 EeV range but does not distinguish nuclei from proton primaries . The shower development of the highest energy event ever detected, the 320 EeV Flyโs Eye event, is consistent with either proton or iron and also disfavors a photon primary . This event constrains hypothetical hadronic primaries to have masses below $``$ 50 GeV . Since fluctuations in shower development are usually large, strong composition constraints await larger statistics of future experiments.
2.7 Clusters of Events
A final challenge for models of UHECRs is the possible small scale clustering of arrival directions . AGASA reported that their 47 events above 40 EeV show three double coincidences (doublets) and one triple coincidence (triplet) in arrival directions, a $`<1`$% chance probability . Adding to the AGASA data that of Haverah Park, Volcano Ranch, and Yakutsk, the 51 events above 50 EeV show one doublet and two triplets . Although these could be due to a statistical fluctuation since the chance probability for the combined set is $`10\%`$ , they may indicate the position of the sources. (When limited to $`\pm 10^o`$ around the Supergalactic plane the chance probability decreases to $`1\%`$.) If these clusters indicate the position of sources, the arrival times and energies of some of the events are inconsistent with a burst and require long lived sources. Furthermore, if the clustering is confirmed by larger data sets and their distribution correlates with some known matter distributions in the nearby universe, the composition of the primaries as well as the magnitude of extragalactic magnetic fields would be strongly constrained . Alternative explanations for the clustering involve either the effect of caustics in the propagation due to magnetic fields or the clustering of dark matter in the halo of the Galaxy.
## 3 Facing the Challenge with Zevatrons
The challenge put forth by these observations has generated two different approaches to reaching a solution: a โbottom-upโ and a โtop-downโ. A bottom-up approach involves looking for Zevatrons , possible acceleration sites in known astrophysical objects that can reach ZeV energies, while a top-down approach involves the decay of very high mass relics from the early universe and physics beyond the standard model of particle physics. Bottom-up models are discussed first and top-down models in the next section.
Acceleration of UHECRs in astrophysical plasmas occurs when large-scale macroscopic motion, such as shocks and turbulent flows, is transferred to individual particles. The maximum energy of accelerated particles, $`E_{\mathrm{max}}`$, can be estimated by requiring that the gyroradius of the particle be contained in the acceleration region. Therefore, for a given strength, $`B`$, and coherence length, $`L`$, of the magnetic field embedded in an astrophysical plasma, $`E_{\mathrm{max}}=ZeBL`$, where $`Ze`$ is the charge of the particle. The โHillas plotโ in Figure 2 shows that, for $`E_{max}>10^{20}`$ eV and $`Z1`$, the only known astrophysical sources with reasonable $`BL`$ products are neutron stars ($`B10^{13}`$ G, $`L10`$ km), active galactic nuclei (AGNs) ($`B10^4`$ G, $`L10`$ AU), radio lobes of AGNs ($`B0.1\mu `$G, $`L10`$ kpc), and clusters of galaxies ($`B\mu `$G, $`L100`$ kpc).
In general, when these sites are considered more carefully, one finds great difficulties due to either energy losses in the acceleration region or the great distances of known sources from our Galaxy . In many of these objects shock acceleration is invoked as the primary acceleration mechanism. Although effective in the acceleration of lower energy cosmic rays, shock acceleration is unable to reach ZeV energies for most plausible acceleration sites with the possible exception of shocks in radio lobes. Unipolar inductors are often invoked as plausible alternative to shocks .
3.1 Cluster Shocks
Moving from right to left on Figure 2, cluster shocks are reasonable sites to consider for UHECR acceleration, since $`E_{max}`$ particles can be contained by cluster fields. However, the propagation of these high energy particles inside the cluster medium is such that they do not escape without significant energy losses. In fact, efficient losses occur on the scales of clusters of galaxies for the same reason that a GZK cutoff is expected, namely, the photopion production off the CMB. Losses limit UHECRs in cluster shocks to reach at most $``$ 10 EeV .
3.2 AGN - Jets and Radio Lobes
Extremely powerful radio galaxies are likely astrophysical UHECR accelerators (for a recent review see ). Jets from the central black-hole of the active galaxy end at a termination shock where the interaction of the jet with the intergalactic medium forms radio lobes and โhot spotsโ. Of special interest are the most powerful AGNs such as Fanaroff-Riley class II objects . Particles accelerated in hot spots of FR-II sources via first-order Fermi acceleration may reach energies well above an EeV and may explain the spectrum up to the GZK cutoff . A nearby specially powerful source may be able to reach energies past the cutoff . Alternatively, the crossing of the tangential discontinuity between the relativistic jet and the surrounding medium may also be able to make protons reach the necessary energies . The spectrum of UHECR primaries formed by the latter proposal is flatter than the Fermi acceleration at the hot spots scenario. Improved statistics of events past the GZK cutoff by future experiments should better determine the spectral index, and therefore, discriminate between plausible sites for UHECR acceleration in radio sources.
Both hot spots and tangential jet discontinuity models avoid the efficient loss processes faced by acceleration models in AGN central regions (ยง3.3). However, the location of possible sources is problematic for both types of mechanisms. Extremely powerful AGNs with radio lobes and hot spots are rare and far apart. The closest known object is M87 in the Virgo cluster ($``$ 18 Mpc away) and could be a main source of UHECRs. Although a single nearby source may be able to fit the spectrum for a given strength and structure of the intergalactic magnetic field , it is unlikely to match the observed arrival direction distribution. After M87, the next known nearby source is NGC315 which is already too far at a distance of $``$ 80 Mpc.
A recent proposal gets around this challenge by invoking a Galactic wind with a strongly magnetized azimuthal component . Such a wind can significantly alter the paths of UHECRs such that all the observed arrival directions of events above 10<sup>20</sup> eV trace back to the Virgo cluster close to M87 . If our Galaxy has a wind with the required characteristics to allow for this magnetic focusing is yet to be determined. Future observations of UHECRs from the Southern Hemisphere (e.g., the Southern Auger Site ) will provide data on previously unobserved parts of the sky and help distinguish plausible proposals for the effect of local magnetic fields on arrival directions. Once again full sky coverage is a key discriminator of such proposals.
3.3 AGN - Central Regions
The powerful engines that give rise to the observed jets and radio lobes are located in the central regions of active galaxies and are powered by the accretion of matter onto supermassive black holes. It is reasonable to consider the central engines themselves as the likely accelerators . In principle, the nuclei of generic active galaxies (not only the ones with hot spots) can accelerate particles via a unipolar inductor not unlike the one operating in pulsars . In the case of AGNs, the magnetic field is provided by the infalling matter and the spinning black hole horizon provides the imperfect conductor for the unipolar induction. Close to the horizon of a black hole ($`RGM/c^2`$) with a mass M $`=10^9M_9\mathrm{M}_{}`$, the electromotive force is : $`emfcBR4.4\times 10^{20}B_4M_9\mathrm{Volts}`$ for a magnetic field $`B=10^4B_4`$ G. It is reasonable to expect that such fields are reached in some nearby AGNs. In addition, the arrival direction of events above $`5\times 10^{19}`$ eV correlate qualitatively well with active galaxies within 100 Mpc . Although it is not clear how statistically significant the correlation is, the clustering of UHECR events in the same regions of the sky where clusters of AGNs reside is certainly tantalizing.
The problem with AGNs as UHECR sources is two-fold: first, UHE particles face debilitating losses in the acceleration region due to the intense radiation field present in AGNs, and second, the spatial distribution of objects should give rise to a GZK cutoff of the observed spectrum. In the central regions of AGNs, loss processes are expected to downgrade particle energies well below the maximum achievable energy. This limitation has led to the proposal that quasar remnants, supermassive black holes in centers of inactive galaxies, are more effective UHECR accelerators . In this case, losses are not as significant. In addition, the problem with the rarity of very luminous radio sources (ยง3.2) is also avoided since any galaxy with a supermassive quiescent black hole could host a UHECR accelerator.
Quasar remnants are manifestly underluminous such that losses in the acceleration region are kept at a reasonably low level . Although presently underluminous, the underlying supermassime black holes are likely to be sufficiently spun-up for individual particles to be accelerated. An incomplete sample of 32 massive dark objects (MDOs) in the nearby universe (of which 8 are within 50 Mpc) finds about 14 MDOs which could have fields strong enough for an $`emf>10^{20}`$ Volts . From the number density and accretion evolution of quasars, more than 40 quasar remnants are expected to have $`>4\times 10^8`$ M within a 50 Mpc volume while more than a dozen would have $`>10^9`$ M .
The second difficulty with AGNs mentioned above, namely the spatial distribution and the GZK cutoff induced by the more distance galaxies, is not avoided by the quasar remnants proposal unless the spectrum is fairly hard. However, it is still within the errors of the current UHECR spectrum the possibility that a GZK cutoff is presently hidden due to the effect of the local clustering of galaxies . This ambiguity should be lifted and a GZK cutoff made apparent by future experiments.
3.4 Neutron Stars
From Figure 2, the last astrophysical objects capable of accelerating UHECRs are neutron stars (see, e.g., ). With the recent identification of โmagnetarsโ (neutron stars with fields of $`>10^{14}`$ G) as the sources of soft gamma ray repeaters , neutron stars have strong enough fields to reach past the required $`E_{max}`$ as in Figure 2. Acceleration processes inside the neutron star light cylinder are bound to fail much like the AGN central region case: ambient magnetic and radiation fields induce significant losses . However, the plasma that expands beyond the light cylinder is freer from the main loss processes and may be accelerated to ultra high energies. One possible solution to the UHECR puzzle is the proposal that the early evolution of neutron stars may be responsible for the flux of cosmic rays beyond the GZK cutoff . In this case, UHECRs originate mostly in the Galaxy and the arrival directions require that the primaries have large $`Z`$ (i.e., primaries are heavier nuclei).
Newly formed, rapidly rotating neutron stars may accelerate iron nuclei to UHEs through relativistic MHD winds beyond their light cylinders . The nature of the relativistic wind is not yet clear, but observations of the Crab Nebula indicate that most of the rotational energy emitted by the pulsar is converted into the flow kinetic energy of the particles in the wind (see, e.g., ). Recent observations of the Crab Nebula by the Chandra satellite indicate both a complex disk and jet structure that is probably associated with the magnetic wind as well as the presence of iron in the expanding shell. Understanding the structure of observable pulsar winds such as the Crab nebula will help determine if during their first years pulsars were efficient Zevatrons.
If most of the magnetic energy in the wind zone is converted into particle kinetic energy and the rest mass density of the wind is not dominated by electron-positron pairs, particles in the wind can reach a maximum energy of $`E_{max}8\times 10^{20}Z_{26}B_{13}\mathrm{\Omega }_{3k}^2\mathrm{eV},`$ for iron nuclei ($`Z_{26}Z/26=1)`$, neutron star surface fields $`B=10^{13}B_{13}`$ G, and initial rotation frequency $`\mathrm{\Omega }=3000\mathrm{\Omega }_{3k}`$ s<sup>-1</sup>. In the rest frame of the wind, the plasma is relatively cold while in the starโs rest frame the plasma moves with Lorentz factors $`\gamma 10^910^{10}`$.
Iron nuclei can escape the remnant of the supernova without suffering significant spallation about a year after the explosion. As the ejected envelope of the pre-supernova star expands, the young neutron star spins down and $`E_{max}`$ decreases. Thus, a requirement for relativistic winds to supply UHECRs is that the column density of the envelope becomes transparent to UHECR iron before the spin rate of the neutron star decreases significantly. The allowed parameter space for this model is shown in Figure 3. Magnetars with the largest surface fields spin down too quickly for iron nuclei to escape unless the remnant is asymmetric with lower density โholes.โ The spectrum of UHECRs accelerated by young neutron star winds is determined by the evolution of the rotational frequency which gives $`\gamma 1`$, at the hard end of the allowed $`\gamma `$ range (ยง2.3).
Depending on the structure of Galactic magnetic fields, the trajectories of iron nuclei from Galactic neutron stars may be consistent with the observed arrival directions of the highest energy events . Moreover, if cosmic rays of a few times $`10^{18}`$ eV are protons of Galactic origin, the isotropic distribution observed at these energies is indicative of the diffusive effect of the Galactic magnetic fields on iron at $`10^{20}`$ eV.
Another recent proposal involving neutron stars suggests that relativistic winds formed around neutron star binaries may generate high energy cosmic rays in a single shot $`\mathrm{\Gamma }^2`$ acceleration , where $`\mathrm{\Gamma }`$ is the bulk Lorentz factor. However, the $`\mathrm{\Gamma }^2`$ acceleration process is likely to be very inefficient which renders the proposal insufficient for explaining UHECRs .
In general, there is an added bonus to considering the existence of Zevatrons in Galactic systems: one may find Pevatrons or Evatrons instead. These may explain the origin of cosmic rays from the knee at 10<sup>15</sup> eV up to the โankleโ at 10<sup>18</sup> eV that remain largely unidentified.
3.5 Gamma-Ray Bursts
Before moving on to more exotic explanations for the origin of UHECRs, one should consider astrophysical phenomena that may act as Zevatrons not included in Figure 2. In effect, transient high energy phenomena such as gamma-ray bursts (bursts of $`0.11`$ MeV photons that last up to a few seconds) may accelerate protons to ultra-high energies . The systems that generate gamma-ray bursts (GRBs) remain unknown but evidence that GRBs are of cosmological origin and involve a relativistic fireball has been mounting with the recent discovery of X-ray, optical, and radio afterglows and the subsequent identification of host galaxies and their redshifts.
Aside from both having unknown origins, GRBs and UHECRs have some similarities that argue for a common origin. Like UHECRs, GRBs are distributed isotropically in the sky , and the average rate of $`\gamma `$-ray energy emitted by GRBs is comparable to the energy generation rate of UHECRs of energy $`>10^{19}`$ eV in a redshift independent cosmological distribution of sources , both have $`10^{44}\mathrm{erg}\mathrm{Mpc}^3\mathrm{yr}^1.`$
Although the systems that generate GRBs have not been identified, they are likely to involve a relativistic fireball (see, e.g., ). Cosmological fireballs may generate UHECRs through Fermi acceleration by internal shocks . In this model the generation spectrum is estimated to be $`dN/dEE^2`$ which is consistent with observations provided the efficiency with which the wind kinetic energy is converted to $`\gamma `$-rays is similar to the efficiency with which it is converted to UHECRs . Acceleration to $`>10^{20}`$ eV is possible provided that $`\mathrm{\Gamma }`$ of the fireball shocks are large enough and that the magnetic field is close to equipartition.
There are a few problems with the GRBโUHECR common origin proposal. First, events past the GZK cutoff require that only GRBs from $`<50`$ Mpc contribute. However, only one burst is expected to have occurred within this region over a period of 100 yr. Therefore, a very large dispersion of $`>`$$``$ 100 yr in the arrival time of protons produced in a single burst is a necessary condition. The deflection by random magnetic fields combined with the energy spread of the particles is usually invoked to reach the required dispersion . If the dispersion in time is achieved, the energy spectrum for the nearby source(s) is expected to be very narrowly peaked $`\mathrm{\Delta }E/E1`$ . Second, the fireball shocks may not be able to reach the required $`\mathrm{\Gamma }`$ factors for UHECR shock acceleration . Third, UHE protons are likely to loose most of their energy as they expand adiabatically with the fireball . However, if acceleration happens by internal shocks in regions where the expansion becomes self-similar, protons may escape without significant losses . Fourth, the observed arrival times of different energy events in some of the UHE clusters argues for long lived sources not bursts (ยง2.7). These clusters can still be due to fluctuations but should become clear in future experiments . Finally, the present flux of UHE protons from GRBs is reduced to $`<10^{42}\mathrm{erg}\mathrm{Mpc}^3\mathrm{yr}^1`$, if a redshift dependent source distribution that fits the GRB data is considered (see also ).
## 4 Hybrid Models
The UHECR puzzle has inspired proposals that use Zevatrons to generate UHE particles other than protons, nuclei, and photons. These use physics beyond the standard model in a bottom-up approach, thus, named hybrid models.
The most economical among such proposals involves a familiar extension of the standard model, namely, neutrino masses. The most common solution to the atmospheric or the solar neutrino problems entails neutrino oscillations, and hence, neutrino masses (see, e.g., ). Recently, the announcement by SuperKamiokande on atmospheric neutrinos has strengthened the evidence for neutrino oscillations and the possibility that neutrinos have a small mass . If some flavor of neutrinos have masses $`1`$ eV, the relic neutrino background will cluster in halos of galaxies and clusters of galaxies. High energy neutrinos ($`10^{21}`$ eV) accelerated in Zevatrons can annihilate on the neutrino background and form UHECRs through the hadronic Z-boson decay .
This proposal is aimed at generating UHECRs nearby (in the Galactic halo and Local Group halos) while using Zevatrons that can be much further than the GZK limited volume, since neutrinos do not suffer the GZK losses. It is not clear if the goal is actually achieved since the production in the uniform non-clustered neutrino background may be comparable to the local production depending on the neutrino masses . In addition, the Zevatron needed to accelerate protons above ZeVs that can produce ZeV neutrinos as secondaries is quite spectacular and presently unknown, requiring an energy generation in excess of $`10^{48}\mathrm{erg}\mathrm{Mpc}^3\mathrm{yr}^1`$ .
Another suggestion is that the UHECR primary is a new particle. For instance, a stable or very long lived supersymmetric neutral hadron of a few GeV, named uhecron, could explain the UHECR events and evade the present laboratory bounds . (Note that the mass of a hypothetical hadronic UHECR primary can be limited by the shower development of the Flyโs Eye highest energy event to be below $`<50`$ GeV .) Both the long lived new particle and the neutrino Z-pole proposals involve neutral particles which are usually harder to accelerate (they are created as secondaries of even higher energy charged primariess) but can traverse large distances without being affected by the cosmic magnetic fields. Thus, a signature of such hybrid models for future experiments is a clear correlation between the position of powerful Zevatrons in the sky such as distant compact radio quasars and the arrival direction of UHE events .
Topological defects have also been suggested as possible UHE primaries . Monopoles of masses between $`10^910^{10}`$ GeV have relic densities below the Parker limit and can be easily accelerated to ultra high energies by the Galactic magnetic field . The main challenges to this proposal are the observed shower development for the Flyโs Eye event that seems to be inconsistent with a monopole primary and the arrival directions not showing a preference for the local Galactic magnetic field .
Another exotic primary that can use a Zevatron to reach ultra high energies is the vorton. Vortons are small loops of superconducting cosmic string stabilized by the angular momentum of charge carriers . Vortons can be a component of the dark matter in galactic halos and be accelerated in astrophysical Zevatrons . Although not yet clearly demonstrated, the shower development profile is also the likely liability of this model.
## 5 Top-Down Models
It is possible that none of the astrophysical scenarios are able to meet the challenge posed by the UHECR data as more observations are accumulated. In that case, the alternative is to consider top-down models. For example, if the primaries are not iron, the distribution in the sky remains isotropic with better statistics, and the spectrum does not show a GZK cutoff, UHECRs are likely to be due to the decay of very massive relics from the early universe.
This possibility was the most attractive to my dear colleague and friend, David N. Schramm, to whom this volume is dedicated. After learning with the work of Hill that high energy particles would be produced by the decay of supermassive Grand Unified Theory (GUT) scale particles (named X-particles) in monopole-antimonopole annihilation, Schramm joined Hill in proposing that such processes would be observed as the highest energy cosmic rays . Schramm realized the potential for explaining UHECRs with physics at very high energies well beyond those presently available at terrestrial accelerators. One winter in Aspen, CO, he remarked pointing to the ski lift โwhy walk up if we can start at the topโ. His enthusiasm for this problem only grew after his pioneering work . In the last conference he attended, an OWL workshop at the University of Maryland , he summarized the meeting by reminding us that in this exciting field the most conventional proposal involves supermassive black holes and that the best fit models involve physics at the GUT scale and beyond. In this field our imagination is the limit (as well as the low number of observed events).
The lack of a clear astrophysical solution for the UHECR puzzle has encouraged a number of interesting proposals based on physics beyond the standard model such as monopolia annihilation, the decay of ordinary and superconducting cosmic strings, cosmic necklaces, vortons, and superheavy long-lived relic particles, to name a few. Due to the lack of space and a number of recent thorough reviews, only a brief summary of the general features of these proposals will be given here. The interested reader is encouraged to consult the following reviews by long-time collaborators of David Schramm and references therein.
The idea behind top-down models is that relics of the very early universe, topological defects (TDs) or superheavy relic (SHR) particles, produced after or at the end of inflation, can decay today and generate UHECRs. Defects, such as cosmic strings, domain walls, and magnetic monopoles, can be generated through the Kibble mechanism as symmetries are broken with the expansion and cooling of the universe (see, e.g., ). Topologically stable defects can survive to the present and decompose into their constituent fields as they collapse, annihilate, or reach critical current in the case of superconducting cosmic strings. The decay products, superheavy gauge and higgs bosons, decay into jets of hadrons, mostly pions. Pions in the jets subsequently decay into $`\gamma `$-rays, electrons, and neutrinos. Only a few percent of the hadrons are expected to be nucleons . Typical features of these scenarios are a predominant release of $`\gamma `$-rays and neutrinos and a QCD fragmentation spectrum which is considerably harder than the case of shock acceleration.
ZeV energies are not a challenge for top-down models since symmetry breaking scales at the end of inflation typically are $`10^{21}`$ eV (typical X-particle masses vary between $`10^{22}10^{25}`$ eV) . Fitting the observed flux of UHECRs is the real challenge since the typical distances between TDs is the Horizon scale, $`H_0^13h^1`$ Gpc. The low flux hurts proposals based on ordinary and superconducting cosmic strings . Monopoles usually suffer the opposite problem, they would in general be too numerous. Inflation succeeds in diluting the number density of monopoles usually making them too rare for UHECR production. To reach the observed UHECR flux, monopole models usually involve some degree of fine tuning. If enough monopoles and antimonopoles survive from the early universe, they can form a bound state, named monopolium, that decay generating UHECRs through monopole-antimonopole annihilation . The lifetime of monopolia may be too short for this csenario to succeed unless they are connected by strings .
Once two symmetry breaking scales are invoked, a combination of horizon scales gives room to reasonable number densities. This can be arranged for cosmic strings that end in monopoles making a monopole string network or even more clearly for cosmic necklaces . Cosmic necklaces are hybrid defects where each monopole is connected to two strings resembling beads on a cosmic string necklace. Necklace networks may evolve to configurations that can fit the UHECR flux which is ultimately generated by the annihilation of monopoles with antimonopoles trapped in the string .
In addition to fitting the UHECR flux, topological defect models are constrained by limits on the flux of high energy photons observed by EGRET (10 MeV to 100 GeV). The energy density of lower energy cascade photons generated by UHE photons and electrons off the CMB and radio background is limited to $`<10^6`$ eV/cm<sup>3</sup>. Figure 4 shows the predicted flux for necklace models given different radio backgrounds and different masses for the X-particle (from ). As can be seen from the Figure, protons dominate the flux at lower energies while photons tend to dominate at higher energies depending on the radio background. If future data can settle the composition of UHECRs from 0.01 to 1 ZeV, these models will be well constrained.
Another interesting possibility is the recent proposal that UHECRs are produced by the decay of unstable superheavy relics that live much longer than the age of the universe . SHRs may be produced at the end of inflation by non-thermal effects such as a varying gravitational field, parametric resonances during preheating, instant preheating, or the decay of topological defects (see, e.g., ). SHRs have unusually long lifetimes insured by discrete gauge symmetries and a sufficiently small percentage decays today producing UHECRs . As in the topological defects case, the decay of these relics also generate jets of hadrons. These particles behave like cold dark matter and could constitute a fair fraction of the halo of our Galaxy. Therefore, their halo decay products would not be limited by the GZK cutoff allowing for a large flux at UHEs. The flux of UHECRs predicted by SHRs clustered in our halo is plotted in Figure 5 (from ). It is clear that the spectrum is not power law (unlike the case of shock acceleration) and that photon fluxes dominate.
From Figures 4 and 5 it is clear that future experiments should be able to probe these hypotheses. For instance, in the case of SHR and monopolium decays, the arrival direction distribution should be close to isotropic but show an asymmetry due to the position of the Earth in the Galactic Halo . Studying plausible halo models and the expected asymmetry will help constrain halo distributions especially when larger data sets are available from future experiments. High energy gamma ray experiments such as GLAST will also help constrain the SHR models due to the products of the electromagnetic cascade .
## 6 Conclusion
Next generation experiments such as the High Resolution Flyโs Eye which recently started operating, the Pierre Auger Project which is now under construction, the proposed Telescope Array , and the OWL-Airwatch satellite will significantly improve the data at the extremely-high end of the cosmic ray spectrum . With these observatories a clear determination of the spectrum and spatial distribution of UHECR sources is within reach. The lack of a GZK cutoff should become apparent with Auger and most extragalactic Zevatrons may be ruled out. The observed spectrum will distinguish Zevatrons from top-down models by testing power laws versus QCD fragmentation fits. The cosmography of sources should also become clear and able to discriminate between plausible populations for UHECR sources. The correlation of arrival directions for events with energies above $`10^{20}`$ eV with some known structure such as the Galaxy, the Galactic halo, the Local Group or the Local Supercluster would be key in differentiating between different models. For instance, a correlation with the Galactic center and disk should become apparent at extremely high energies for the case of young neutron star winds , while a correlation with the large scale galaxy distribution should become clear for the case of quasar remnants. If SHRs or monopolia are responsible for UHECR production, the arrival directions should correlate with the dark matter distribution and show the halo asymmetry. For these signatures to be tested, full sky coverage is essential. Finally, an excellent discriminator would be an unambiguous composition determination of the primaries. In general, Galactic disk models invoke iron nuclei to be consistent with the isotropic distribution, extragalactic Zevatrons tend to favor proton primaries, while photon primaries are more common for early universe relics. The hybrid detector of the Auger Project should help determine the composition by measuring simultaneously the depth of shower maximum and the muon content of the same shower.
In addition to explaining the origin of UHECRs, GUT to Planck scale physics can potentially be probed by the existence of UHECRs. For instance, the breaking of Lorentz invariance can change the threshold for photopion production significantly in such a way as to be constrained by a clear observation of the GZK cutoff . There are great gains to be made if the data at the highest energies is improved by a few orders of magnitude. The prospect of testing extremely high energy physics as well as solving the UHECR puzzle given all the presently proposed models sends a strong message that the challenge is back in the observational arena. Fortunately, observers have accepted the challenge and are building and planning experiments large enough to resolve these open questions .
## Acknowledgment
It has been a great pleasure to have known David N. Schramm and to be able to contribute to this volume in his honor. Dave was a kind mentor and friend who is missed with saudades. Many thanks to my โultra-high energyโ collaborators, P. Blasi and R. Epstein, for the careful reading of the manuscript, many ongoing discussions, and for providing most of the figures. I am also very grateful to I. Albuquerque, V. Berezinsky, P. Biermann, J. Cronin, G. Farrar, T. Gaisser, M. Lemoine, G. Sigl, T. Stanev, A. Watson, and T. Weiler for their comments on the manuscript. This work was supported by NSF through grant AST 94-20759 and DOE grant DE-FG0291 ER40606. |
warning/0002/math0002003.html | ar5iv | text | # Generalized Hopfian property, minimal Haken manifold, and J. Simonโs conjecture for 3-manifold groups
We address a conjecture that $`\pi _1`$-surjective maps between closed aspherical 3-manifolds having the same rank on $`\pi _1`$ must be of non-zero degree. The conjecture is proved for Seifert manifolds, which is used in constructing the first known example of minimum Haken manifold. Another motivation is to study epimorphisms of 3-manifold groups via maps of non-zero degree between 3-manifolds.
support: The first author was supported by the Royal Society, the NSF and The Alfred P. Sloan Foundation. The second and third author are supported by Outstanding Youth Fellowship of NSFC.
Section 1. Introduction and some examples.
Let $`M`$ and $`N`$ be closed 3-manifolds and $`f:MN`$ a map of non-zero degree, then the image of $`f_{}`$ is a subgroup of finite index in $`\pi _1(N)`$. If $`M`$ and $`N`$ are aspherical, any homomorphism $`\varphi :\pi _1(M)\pi _1(N)`$ determines a unique map $`f:MN`$ up to homotopy such that $`f_{}=\varphi `$. It seems natural to ask when there exists $`f:MN`$ of non-zero degree given a homomorphism $`\varphi `$ surjecting $`\pi _1(M)`$ on a subgroup of finite index in $`\pi _1(N)`$? There are elementary constructions of examples (see below) that show in general that the answer is no. Before discussing some examples we make the following definition.
Definition 1.1. A map $`f:MN`$ between 3-manifolds is $`\pi _1`$-surjective (resp. $`\pi _1`$-finite-index) if $`f_{}:\pi _1(M)\pi _1(N)`$ is surjective (resp. the image of $`f_{}`$ is a subgroup of finite index).
Recall that if $`M`$ is an $`n`$-manifold, the rank of $`\pi _1(M)`$ (or simply by abuse just $`M`$) is the minimal cardinality of a generating system for $`\pi _1(M)`$.
Let us first have a look of the situation in dimension 2 which is quite simple. Throughout the paper $`\mathrm{\Sigma }_k`$ will denote a closed orientable surface of genus $`k`$.
Example 1.1. It is not difficult to see that there is a $`\pi _1`$-surjective map $`f:\mathrm{\Sigma }_l\mathrm{\Sigma }_k`$ which is of degree zero when $`l2k`$.
On the otherhand, if $`f:\mathrm{\Sigma }_l\mathrm{\Sigma }_k`$ is a $`\pi _1`$-surjective map with $`0<l<2k`$, then we claim that $`f`$ is of non-zero degree. The proof of this result is direct. Choose a 1-skeleton of $`\mathrm{\Sigma }_k`$ to be a one point wedge of $`2k`$ circles $`๐ฑ=C_i`$. Fix a point $`x_i`$ on $`C_i`$. If $`f`$ is of degree zero, then the image of $`f`$ can be deformed into $`๐ฑ`$. We assume therefore that this is the case. Since $`f_{}`$ is $`\pi _1`$-surjective, $`f:\mathrm{\Sigma }_l๐ฑ`$ must be surjective. We may also assume that $`f`$ is transverse to each $`x_i`$, $`i=1,2,\mathrm{},2k`$. So $`f^1(x_i)`$ is a set of essential circles. Partition $`f^1(x_i)`$ into sets $`G_1,\mathrm{},G_h`$ such that two components are in the same set if and only if they are parallel. For each $`G_j`$, find an annulus $`A_j`$ containing $`G_j`$. Then squeeze each $`A_j`$ to an arc $`a_j`$ and the part $`\mathrm{\Sigma }_kA_j`$ to a point. The quotient $`Q`$ will be a bouquet of $`h`$ circles. Since $`๐ฑ\{x_i,i=1,2,\mathrm{},2k\}`$ is contractible, the map $`f:\mathrm{\Sigma }_l๐ฑ`$ factors through $`q:Q๐ฑ`$ which is still $`\pi _1`$-surjective. It follows that $`h2k`$. In particular, there are at least $`h2k`$ disjoint essential non-separating non-parallel circles. By a well known argument in surface topology, we must have that the $`l`$, the genus of $`\mathrm{\Sigma }_l`$, is at least $`2k`$. $`\mathrm{}`$
Let us come back to dimension 3. The first example illustrates the aspherical assumption.
Example 1.2:
Let $`f=ep:S^2\times S^1S^2\times S^1`$, where $`p`$ is a map which pinches $`S^2\times S^1`$ to $`S^1`$, and $`e`$ identifies $`S^1`$ to a fiber $`\times S^1S^2\times S^1`$. Clearly $`f`$ is of zero degree but $`\pi _1`$-surjective. $`\mathrm{}`$
The second example shows that if we do not require that the manifolds have the same rank, then the answer to the question is no.
Example 1.3:
We construct a map $`f:\mathrm{\Sigma }_{g+1}\times S^1\mathrm{\Sigma }_g\times S^1`$ of zero degree which is $`\pi _1`$-surjective. The map $`f`$ is the composition of the following four geometric operations.
Project $`\mathrm{\Sigma }_{g+1}\times S^1`$ to $`\mathrm{\Sigma }_{g+1}`$.
Squeeze a suitable separating circle on $`\mathrm{\Sigma }_{g+1}`$ to a point in such a way that the quotient space is a one point union of a torus and $`\mathrm{\Sigma }_g`$.
Squeeze the torus to a circle in such a way that the quotient space is a one point union of the circle and $`\mathrm{\Sigma }_g`$.
Send $`\mathrm{\Sigma }_g`$ and the circle to a section $`\mathrm{\Sigma }_g\times `$ and the circle fiber of $`\mathrm{\Sigma }_g\times S^1`$ respectively.$`\mathrm{}`$
The third example has the same purpose as the second one, but the manifolds in this case are hyperbolic.
Example 1.4:
Let $`M`$ be a closed hyperbolic 3-manifold whose fundamental group surjects the free group of rank $`2`$. Such examples are easily constructed by doing hyperbolic surgery on a null-homotopic hyperbolic knot in $`\mathrm{\Sigma }_2\times S^1`$ \[Section 3, BW\]. Let $`\varphi _1:\pi _1(M)`$ denote such a map. Let $`N`$ be any hyperbolic 3-manifold such that $`\pi _1(N)`$ has two generators, then there is an epimorphism $`\varphi _2:\pi _1(N)`$. If we choose $`N`$ such that the volume of $`N`$ is larger than the volume of $`M`$, then the map realizing the epimorphism $`\varphi =\varphi _2\varphi _1`$ must be zero degree by the work of Gromov and Thurston, \[T\]. We remark that the volumes of hyperbolic 3-manifolds of rank $`2`$ are unbounded. Briefly, it follows from work of Adams that the volumes of hyperbolic 2-bridge knot complements are unbounded. Doing large enough hyperbolic Dehn surgeries on these gives the required family, see \[CR\].
In fact it can be seen directly that the map realizing $`\varphi `$ must be of zero degree since such a map factors through a 1-dimensional complex. $`\mathrm{}`$
As a consequence of these examples, we state the following more refined version of the question posed above.
###### \bfQuestion 1.5
Let $`M`$ and $`N`$ be closed aspherical 3-manifolds such that the rank of $`\pi _1(M)`$ equals the rank of $`\pi _1(N)`$. Assume, that $`\varphi :\pi _1(M)\pi _1(N)`$ is surjective or whose image is a subgroup of finite index. Does $`\varphi `$ determine a map $`f:MN`$ of non-zero degree?
Note if $`M`$ and $`N`$ are homeomorphic and satisfy Thurstonโs geometrization conjecture, then a $`\pi _1`$-surjective map $`f:MM`$ must be degree one. For since $`\pi _1(M)`$ is hopfian, $`f_{}`$ is surjective implies $`f_{}`$ is an isomorphism. Since $`M`$ is aspherical $`f`$ must be a homotopy equivalence, and so in particular, $`f`$ is of degree one. Thus the question above is a kind of generalization of the Hopfian property: the condition โhomeomorphic manifoldsโ is replaced by โmanifolds of the same rankโ, the condition โ$`\pi _1`$-surjectiveโ is replaced by โ$`\pi `$-surjectiveโ or โ$`\pi _1`$-finite-indexโ, and the conclusion replace โdegree oneโ by โnon-zero degreeโ. It is easy to construct examples to show that โnon-zero degreeโ cannot be sharpened to โdegree oneโ, see the examples in Section 3.
One of the main results of this paper is to prove that for Seifert fibered 3-manifolds Question 1.5 has a positive answer (see Theorem 2.1 and Remark 2.4). In ยง4 we use this result to construct the first known example of a Haken 3-manifold which is minimal with respect to degree 1 mappings in Thurstonโs picture of 3-manifolds (Theorem 4.1). The manifold is a graph manifold built from the union of two trefoil knot complements. An orientable 3-manifold $`M`$ is minimal if there is a degree one map $`f:MN`$ implies either $`N=S^3`$ or $`M=N`$. Usually it is difficult tell if a 3-manifold is minimal. We remark that all minimal Seifert manifolds are non-Haken \[LWZ\], and that the known minimal hyperbolic 3-manifolds are also non-Haken \[RW\], see \[BW\], \[RW\] and \[LWZ\] for a further discussion of such matters.
We were also motivated by the following posed by J. Simon.
###### \[K. Problem 1.12\]
Let $`G_K=\pi _1(S^3K)`$ for a knot $`K`$ in $`S^3`$. Conjectures: if there is an epimorphism $`\varphi :G_KG_L`$, then
(A) rank $`G_L>`$ rank $`G_K`$.
(B) genus$`(L)`$ genus$`(K)`$.
(C) Given $`K`$, there is a number $`N_K`$ such that any sequence of epimorphisms of knot groups $`G_KG_{L_1}\mathrm{}.G_{L_n}`$ with $`nN_K`$ contains an isomorphism.
(D) Given $`K`$, there are are only finitely many knot groups $`G`$ for which there is an epimorphism $`G_KG`$.
These conjectures have seen little progress. On the otherhand, more recently, questions similar to (C) and (D) have been raised for degree one maps and there are already several substantial results in this setting.
###### \[K Problem 3.100\]
Let $`M`$ be a closed orientable 3-manifold.
(A) Are there only finitely many irreducible 3-manifolds $`N`$ such that there exists a degree one map $`MN`$?
(B) Does there exist an integer $`N_M`$ such that if a sequence of degree one map $`MM_1\mathrm{}.M_k`$ with $`k>N_{M_0}`$, the sequence contains an homotopy equivalence
If one assumes Thurstonโs Geometrization conjectural picture of 3-manifolds, the answer for (B) is Yes if $`k=\mathrm{}`$ \[Ro2\]; the answer for (A) is Yes if the targets are hyperbolic \[So\], or the domain is non-Haken \[RW\], or the targets have finite $`\pi _1`$ \[LMWZ\].
Thus it seems natural to study the conjectures of J. Simon for closed orientable 3-manifolds (Question 1.6 in ยง3). We find that the positive answer for Question 1.5 are important for studying the conjectures. This will be addressed in ยง3.
Section 2. $`\pi _1`$-surjective maps between aspherical Seifert manifolds.
###### Theorem 2.1
Let $`M_1`$ and $`M_2`$ be closed orientable aspherical Seifert fiber spaces with the same rank and whose base orbifolds are orientable. Then any $`\pi _1`$-surjective map $`f:M_1M_2`$ is of non-zero degree.
To prove Theorem 2.1, we will make use of \[Ro1\], in particular we refer the reader to \[Ro1\] for the definition of a vertical pinch and a squeeze, and squeeze torus. Call a squeeze is vertical, if in the squeeze torus, the squeezing circle meets the regular fiber exactly one point.
Also remember that any orientable Seifert manifold $`M`$ with orientable base orbifold of genus $`g`$ and with $`n`$ singular fibers has a unique normal form $`(g;b;\alpha _1,\beta _1;\mathrm{};\alpha _n,\beta _n)`$, where $`0\beta _i\alpha _i`$, $`i=1,\mathrm{},n`$. The orbifold $`O_1`$ of $`M_1`$ will denoted by $`(g;\alpha _1,\mathrm{},\alpha _n)`$. $`g`$ is often omitted if $`g=0`$.
We first need two lemmas.
###### Lemma 2.2 \[Ro1, Lemma 3.5\]
Let $`f:MN`$ be a map between aspherical Seifert manifolds and $`1f_{}(h)h^{}`$, where $`M`$ is closed and $`N\mathrm{}`$, $`h`$ and $`h^{}`$ are regular fibers of $`M`$ and $`N`$ respectively. Then either $`f`$ admits a vertical squeeze, or $`f`$ can be homotoped so that the image of $`f`$ lies in a fiber of $`N`$.
###### Lemma 2.3
Suppose $`f:FO`$ is an orbifold branch covering, where $`F`$ is a surface of genus $`g`$, and $`O`$ is a orbifold, both are orientable and have non-positive Euler characteristic, then $`rank(\pi _1(F))rank(\pi _1(O))1`$ if $`f`$ is a double branched cover over 2-sphere and $`rank(\pi _1(F))rank(\pi _1(O))`$ otherwise.
###### Demonstration Proof
The proof is based on the results about the ranks of Fuchsian groups \[ZVC, Theorem 4.16.1\] and the Riemann-Hurwitz formula.
Suppose $`O`$ has $`k`$ singular points of index $`v_i`$, $`i=1,\mathrm{},k`$, with the underlying space of genus $`g^{}`$ and the degree of $`f`$ is $`n`$. Then we have
$$22g=n(22g^{}\mathrm{\Sigma }_{i=1}^k(1\frac{1}{v_i}))$$
For the case $`g=1`$, the verification is direct, so we assume below that $`g>1`$.
If $`n=2`$ then all $`v_i=2`$, $`k=2m`$ and we have $`22g=2(22g^{}m)`$, i.e., $`g=2g^{}+m1`$. Now $`rank(\pi _1(F))=2g=4g^{}+2m2`$ and the $`rank(\pi _1(O))`$ is at most $`2g^{}+2m1`$ if $`g^{}>0`$ and is $`2m1`$ if $`g^{}=0`$ by \[ZVC, Theorem 4.16.1\]. In any case the lemma follows.
If $`n3`$, then
$$22g3(22g^{}\mathrm{\Sigma }_{i=1}^k(1\frac{1}{v_i}))3(22g^{}\frac{k}{2})$$
i.e., $`2g6g^{}4+\frac{3k}{2}`$. If $`g^{}>0`$, $`2g2g^{}+\frac{3k}{2}`$. But the rank of $`\pi _1(O)`$ is at most $`2g^{}+k1`$. If $`g^{}=0`$, then we have $`2g4+\frac{3k}{2}`$ if $`k`$ is even and $`g4+\frac{3k}{2}+\frac{1}{2}`$ if $`k`$ is odd. The rank of $`\pi _1(O)`$ is at most $`k1`$. It follows that if $`k5`$, then $`2gk1`$. If $`k4`$ we still have $`2gk1`$ since we assume that $`g>1`$. โ
###### Demonstration Proof of Theorem 2.1
Suppose $`f`$ is of zero degree. For clarity, the proof is divided up three steps.
Step (1) We prove the following:
Claim : $`f(h)`$ is homotopically non-trivial, where $`h`$ is the regular fiber of $`M_1`$.
Proof of Claim:
Let $`M_1=(g;b;a_1,b_1;\mathrm{};a_k,b_k)`$ and $`G_1=\pi _1(M_1)/<h>`$, where $`<h>`$ is the cyclic group generated by regular fiber of $`M_1`$. Let $`r=rank(\pi _1(M_1))=rank(\pi _2(M_2))`$.
By \[BZ, Theorem 1.1\] and \[ZVC, Theorem 4.16.1\], one of the following cases holds.
(1) $`rank(\pi _1(M_1))>rank(G_1)`$,
(2) $`rank(\pi _1(M_1))=rank(G_1)`$ there is a set of generators of $`G_1`$ which realizes the rank and contains at least one torsion element,
(3) $`rank(\pi _1(M_1))=rank(G_1)=rank(G_1/T)`$, where $`T`$ is the normal subgroup normally generated by the torsion elements and $`G_1/T`$ is a surface group.
If $`f(h)`$ is homotopically trivial, then $`f_{}:\pi _1(M_1)\pi _1(M_2)`$ induces an epimorphism $`\varphi :G_1\pi _1(M_2)`$.
In Case (1) the Claim is clearly true.
In Case (2) the Claim is also true since $`\pi _1(M_2)`$ is torsion free.
In Case (3), $`f_{}`$ induces an epimorphism $`\varphi ^{}:G_1/T\pi _1(M_2)`$. Let $`f^{}:FM_2`$ be the map which realizes $`\varphi ^{}`$. Since $`\varphi ^{}`$ is not injective, (otherwise $`\varphi ^{}`$ would be an isomorphism and $`\pi _1(M_2)`$ would be surface group), by the simple loop theorem for maps from a surface to a Seifert manifold \[H\], there are essential simple loops in the kernel of $`\varphi ^{}`$. Assume first there is an essential non-separating simple loop, which we $`\alpha `$, in the kernel. Then the map $`f^{}`$ induces a map $`f^{\prime \prime }:F^{}M_2`$, where $`F^{}`$ is a complex obtained by squeezing $`F`$ along $`\alpha `$. It is easy to see that the rank of $`\pi _1(F^{})`$ is $`r1`$. We reach a contradiction. If all essential simple loops in kernel of $`\varphi ^{}`$ are separating, let $`\alpha `$ be a maximal family of non-parallel separating essential simple closed curves in kernel $`\varphi ^{}`$. Again $`f^{}`$ can factors through $`f^{\prime \prime }:F^{}M_2`$, where $`F^{}`$ is a complex of obtained by squeezing $`F`$ along $`\alpha `$, which is union of closed surfaces connected by arcs. Let $`S`$ be a surface in $`F^{}`$. Due to the maximality of $`\alpha `$, the restriction $`f^{\prime \prime }|_S`$ $`\pi _1`$-injective, which must be either horizontal or vertical by \[H\]. If $`f^{\prime \prime }|_S`$ is horizontal, than $`p_2f^{\prime \prime }|:SO(M_2)`$ is an orbifold branched covering, where $`p_2:M_2O(M_2)`$ is the fiber map. But the rank of $`\pi _1(S)`$ is at most $`r2`$. This is also ruled out by Lemma 2.3. If $`f^{\prime \prime }|_S`$ is vertical for each surface $`S`$ of $`F^{}`$, then $`F^{}`$ contains at most $`g`$ such surfaces and and each of them is a torus. Clearly the rank of $`f_{}^{\prime \prime }\pi _1(F^{})`$ is at most $`g+1`$, which is at most $`r1`$ (since $`g>1`$ and $`r2g`$). Again we reach a contradiction.
Step (2)
We will factor $`f:M_1XM_2`$, where the 2-dimensional complex $`X`$ is a quotient of $`M`$ with rank $`r_X`$.
Since $`f(h)`$ is homotopically non-trivial, a standard argument in 3-manifold topology shows that $`f:M_1M_2`$ can be deformed to be a fiber preserving map (see \[J\] for example). Since a vertical pinch reduces the rank of $`\pi _1`$, $`f:M_1M_2`$ admits no vertical pinch. Suppose the mapping degree is zero. We can further deform the map so that the image $`f(M_1)`$ misses a regular fiber $`h^{}`$ of $`M_2`$. To see this, $`f:M_1M_2`$ is fiber preserving. We can further deform $`f`$ so that for each singular fiber of $`M_2`$, its preimage consists of finitely many fibers of $`M_1`$. Now remove all singular fibers of $`M_1`$ and their f-images, and remove all singular fibers of $`M_2`$ and their f-preimages. The restriction of $`f`$ gives a proper map $`f^{}:M_1^{}M_2^{}`$, which is fiber preserving map between circle bundles. Since $`f`$ is assumed to be degree zero, $`f^{}`$ is of zero degree. Since $`f(h)`$ is non-trivial, the induced proper map $`\overline{f}^{}:F_1^{}F_2^{}`$ between base surfaces must be degree zero. Hence by $`\overline{f}^{}`$ can be deformed so that its image misses a point of $`F_2^{}`$. This deformation can be lifted to the bundle map $`f^{}`$ whose image then misses a circle fiber in $`M_2^{}`$. With this we reach the situation claimed above.
Now remove an open fibered neighborhood of $`h^{}`$, and denote the resulting manifold by $`N`$. Then we have a fiber preserving map $`f:M_1N`$, where $`N\mathrm{}`$.
According Lemma 2.2 either $`f:M_1N`$ admits a fiber squeeze along an incompressible vertical torus, or $`f(M_1)`$ a fiber of $`N`$. Using this we can reformulate the above so that either $`f:M_1M_2`$ admits a vertical squeeze along an incompressible vertical torus, or $`f(M_1)`$ a fiber of $`M_2`$.
Since $`f`$ is $`\pi _1`$-surjective, and $`M_2`$ is an closed aspherical Seifert fiber space, the situation that $`f(M_1)`$ a fiber of $`M_2`$ cannot happen. Let $`๐ฏ`$ be a maximal family of disjoint non-parallel incompressible tori along which $`f`$ admits vertical squeeze. Let $`X=๐ฌ๐`$ be the space obtained after the squeezing, where $`๐ฌ`$ is a union of Seifert fiber spaces with the induced Seifert fibration, $`๐`$ is a union of annuli and $`๐`$ is a union of regular fibers of $`๐ฌ`$. Then $`f`$ induces a $`\pi _1`$-surjective map $`XM_2`$, which we continue to denote by $`f`$.
Now all components of $`๐ฌ`$ are Seifert fibered spaces with the induced Seifert fibrations, so we may assume that $`Q_1,\mathrm{}.,Q_{k_1}`$ are Seifert manifolds of $`๐ฌ`$ which are not the trivial circle bundle over $`S^2`$ and $`Q_{k_1+1},\mathrm{}.,Q_{k_1+k_2}`$ are trivial circle bundle over $`S^2`$. Clearly each $`Q_j`$, $`j>k_1`$, is $`S^2\times S^1`$.
For each $`j>k_1`$, there is an annulus $`A`$ in $`๐`$, with two components $`C_1`$ and $`C_2`$ of $`A`$ such that $`C_1`$ belongs to $`Q_j`$, Since $`C_1`$ is a regular fiber of $`Q_j=S^2\times S^1`$, $`Q_jA`$ has $`C_2`$ as a retractor, and hence we can send $`Q_jA`$ to $`C_2`$ by this retract, then extend the map to whole $`๐ฌ๐`$. After $`k_2`$ such operations, we get a quotient space $`X_1=๐ฌ_1๐_1`$ where $`๐ฌ_1=\{Q_1,\mathrm{},Q_{k_1}\}`$, and $`f`$ induces a $`\pi _1`$-surjective map $`X_1M_2`$,
By the maximality of $`๐ฏ`$, each $`Q_i`$ contains no squeeze torus for $`f|_{Q_i}`$, so we have that $`f(Q_i)`$ a fiber of $`M_2`$, and consequently we have the following
Fact (\*)
: each $`Q_i`$ has base orbifold $`S^2`$ and has no more than 3 singular fibers (otherwise there will be a squeeze torus).
So $`f:X_1M_2`$ induces a $`\pi _1`$-surjective map $`X=๐ฎ๐_1M_2`$, $`๐ฎ`$ is a union of $`k_1`$ circle.
Step (3)
We will show that $`r_X<r`$ and reach a contradiction.
If $`g>0`$, there is a non-separating squeeze torus for $`f`$ and clearly $`r_X<r`$.
Below we assume that $`g=0`$. Then every squeeze torus is a separating torus.
Say that $`M_1`$ is of type I, if $`M_1`$ has normal form $`(0;b;2,1;\mathrm{}.;2,1;2\lambda +1,b_k)`$, where $`k4`$ is even, and $`\lambda >1`$, otherwise call $`M_1`$ is of type II. By \[Theorem 1.1 BZ\], the $`r=k2`$ if $`M_1`$ is of type I and $`r=k1`$ if $`M_1`$ is of type II.
Each component $`Q_i`$ of $`๐ฌ`$ must have infinite fundamental group, otherwise $`f(h)`$ is an element of finite order, which must be trivial in $`\pi _1(M_2)`$, and this is forbidden by Step (1). In particular, each $`Q_i`$ contains at least 2 singular fibers, $`i=1,\mathrm{},k_1`$, and $`Q_i`$ contains exactly two singular fibers only if $`Q_i=(0;0;a,b;a,b)`$. Moreover if $`M_1`$ is of type I, then at least one $`Q_i`$ contains 4 singular fibers (since $`\lambda >1`$ and both $`(0;b;2,1;2\lambda +1,b_2)`$ and $`(0;b;2,1;2,1;2\lambda +1,b_3)`$ have finite fundamental groups), which is not possible by Fact\*.
If $`M_1`$ is of type II, then $`r=k1`$ and $`k3`$, but
$$r_Xk_1\frac{k}{2}<k1=r,$$
where the first $``$ is due to the fact that every squeeze torus is separating and the second $``$ is due to every $`Q_i`$ contains at least two singular fibers. โ
Remark 2.4. In Theorem 2.1, the condition โ$`f`$ is $`\pi _1`$-surjectiveโ can be replaced by โ$`f`$ is $`\pi _1`$-finite-indexโ, and the condition โorbifolds are orientableโ can be removed. For details see \[Hu\], where the proof is parallel to the proof above, but involves more complicated case by case argument.
3. On the conjectures of J. Simon on 3-manifold groups.
In this section we study the follwing questions.
###### Question 1.6
Let $`M_i`$ be closed orientable aspherical 3-manifolds. Suppose there is an epimorphism $`\varphi :\pi _1(M_1)\pi _1(M_0)`$.
(A) Is rank $`\pi _1(M_1)>`$ rank $`\pi _1(M_0)`$?
(B) Is Heegaard genus of $`M_1`$ Heegaard genus $`M_0`$?
Moreover given $`M_0`$.
(C1) Is there a number $`N_M`$ such that any sequence of epimorphisms $`\pi _1(M_0)\pi _1(M_1)\mathrm{}.\pi _1(M_n)`$ with $`nN_M`$ contains an isomorphism?
(C2) Does any infinite sequence of epimorphisms $`\pi _1(M_0)\pi _1(M_1)`$ $`\mathrm{}.\pi _1(M_n)\mathrm{}`$ contain an isomorphism?
(D) Are there only finitely many $`M_i`$ with the same first Betti number, or the same $`\pi _1`$-rank, as that of $`M_0`$, for which there is an epimorphism $`\pi _1(M_0)\pi _1(M_i)`$?
We remark that a positive answer for (B) of Question 1.6 implies a positive solution to the Poincare Conjecture. From Example 1.4 of The Introduction the answer to (D) is negative if we remove the condition on first Betti number or $`\pi _1`$-rank on (D) of Question 1.6.
We describe first some examples of non-trivial $`\pi _1`$-surjective maps between two 3-manifolds of the same rank, which give a negative answers of the (A) of Question 1.6. Clearly those examples are all of non-zero degrees.
Example 3.1. Let $`M`$ be a Seifert manifold of normal form $`(0;0;6,b_1;5,b_2;7,b_3)`$. Let $`_2`$ be a cyclic group acting on $`M`$ such that it induces the identity on the base space and standard rotation on each regular fiber. Then one verifies that $`M/_2`$ is a Seifert manifold with normal form $`(0;0;3,b_1;5,2b_2;7,2b_3)`$. Now
$$\pi _1(M)=<s_1,s_2,s_3,h|[s_j,h],s_1^6h^{b_1},s_2^5h^{b_2},s_3^7h^{b_3},s_1s_2s_3>$$
and
$$\pi _1(M/_2)=<t_1,t_2,t_3,h^{}|[t_j,h^{}],t_1^3h_{}^{}{}_{}{}^{b_1},t_2^5h_{}^{}{}_{}{}^{2b_2},t_3^7h_{}^{}{}_{}{}^{2b_3},t_1t_2t_3>$$
The quotient map $`p:MM/_2`$ is a branched covering of degree 2 and $`p_{}`$ sends $`s_jt_j`$ and $`hh_{}^{}{}_{}{}^{2}`$. Since $`(2,b_1)=1`$, $`p_{}`$ is surjective. By \[BZ\] these manifolds have rank 2. $`\mathrm{}`$
Examples 3.2. We now give some examples of $`\pi _1`$-surjective non-zero degree maps between hyperbolic manifolds of the same $`\pi _1`$ ranks.
Let $`M`$ be a closed orientable 3-manifold and $`kM`$ be any hyperbolic fibered knot. Suppose the fiber $`F`$ has genus $`g`$. Let $`M_n`$ be the $`n`$-fold cyclic branched cover of $`M`$ over the knot $`k`$. Then the rank of $`\pi _1(M_n)`$ is bounded by $`2g+1`$ for all $`n`$ and $`M_n`$ is hyperbolic when $`n`$ is large. If $`k|n`$, then $`M_nM_k`$ is a branched cover, which is $`\pi _1`$-surjective. So there are must be infinitely many $`\pi _1`$-surjective branched covering $`M_nM_k`$ between hyperbolic 3-manifolds of the same ranks.
A well studied case is when $`M_n`$ is the n-fold cyclic branched cover of the figure eight knot. Then for $`n3`$ the fundamental groups are all 2-generatorโin fact they are the Fibonacci groups $`F(2,2n)`$, which are all hyperbolic if $`n4`$. By abelianizing $`F(2,2n)`$ we see that all $`M_n`$ have first Betti number zero (see \[MR\] for example).
The next example gives the negative answer of (C1) of Question 1.6.
Example 3.3
(1) Let $`M(n,k)=(0;0;2^k3,b_1;5,2^{nk}b_2;7,2^{nk}b_3)`$. Similar to Example 3.1, we have sequence of degree 2 branched covering $`M(n,n)\mathrm{}M(n,1)M(n,0)`$ of length $`n+1`$, which induces a sequence of epimorphisms of groups $`\pi _1(M(n,n))\mathrm{}\pi _1(M(n,1)))\pi _1(M(n,0))`$ of rank 2. Let $`M`$ be $`\mathrm{\Sigma }_2\times S^1`$. Clearly $`\pi _1(M)`$ surjects onto $`ZZ`$, then we have the sequence of epimorphisms
$$\pi _1(M)\pi _1(M(n,n))\mathrm{}\pi _1(M(n,1))\pi _1(M(n,0))$$
of length $`n+2`$, where $`n`$ can be arbitraily large.
Moreover if we choose $`b_1,b_2,b_3`$ such that the Euler number of $`M(n,n)`$ is non-zero. Since each $`M(N,k)`$ has infinite $`\pi _1`$ and is the image of $`M(n,n)`$ under non-zero degree map, the Euler number of $`M(n,k)`$ is non-zero \[Theorem 2, W\]. It follows $`M(n,k)`$ has neither horizontal or vertical incompressible surface, and therefore all $`M(n,k)`$ are non-Haken \[J\].
(2) Let $`M_n`$ be the $`n`$-fold cyclic branched covering of $`S^3`$ over figure eight knot as in the end of Example 3.2. Then we have sequence of branched coverings of hyperbolic rational homology spheres $`M_{4k}\mathrm{}M_8M_4`$ of length $`l`$ which induces a sequence of epimorphisms of groups $`\pi _1(M_{4k})\mathrm{}\pi _1(M_8)\pi _1(M_4)`$ with rank 2. Let $`M`$ be a hyperbolic 3-manifold with $`\pi _1(M)`$ surjecting $``$ (as in Example 1.4). Then we have the sequence of epimorphisms
$$\pi _1(M)\pi _1(M_{4k})\mathrm{}\pi _1(M_8)\pi _1(M_4)$$
of length $`l+1`$, $`l`$ can be arbitraily large. โ
The next result gives a partial positive answer of (C2) of Question 1.6.
###### Theorem 3.4
Given $`M_0`$, and a sequence $`M_i`$ of closed orientable aspherical Seifert manifolds with epimorphisms $`\pi _1(M_0)\pi _1(M_1)\mathrm{}.\pi _1(M_n)`$ $`\mathrm{}`$, this sequence contains an isomorphism.
###### Demonstration Proof
By passing an infinite subsequence, we may assume that all groups in the sequence have the same rank (each epimorphism in the subsequence is the composition of epimorphisms involved). Then each epimorphism $`\varphi _i:\pi _1(M_i)\pi _1(M_{i+1})`$ in the sequence can be realized by a map $`f_i:M_iM_{i+1}`$ of non-zero degree by Theorem 2.1. Moreover the Seifert fibrations of the $`M_i`$โs can be arranged so that each $`f_i`$ is a fiber preserving. Let $`O_i`$ be the orbifold of $`M_i`$, then $`\chi (O_i)0`$ and we have the induced sequence of epimorphisms
$$\pi _1(O_0)\pi _1(O_1)\mathrm{}.\pi _1(O_n)\mathrm{}$$
of Fuchsian groups. We therefore have a decreasing sequence
$$\chi (M_0)\chi (M_1)\mathrm{}.\chi \pi _1(M_n)\mathrm{}.$$
The $`\{\chi (O)\}`$ form a well-order subset of reals, where $`O`$ runs over compact orbifolds, $`\chi (O_k)=\chi (O_{k+1})`$ for $`k`$ larger than a given $`N`$ (\[Ro2, Lemmas 2.5 and 2.6\] for details). Since there are at most finitely many orbifolds $`O`$ with given $`\chi `$, by passing an infinite sequence, we may assume that all $`O_i`$ are the same.
Let $`O_i=(g;\alpha _{1,},\mathrm{},\alpha _n)`$. Then $`M_i=(g;b_i;\alpha _{1,},\beta _{1,i};\mathrm{};\alpha _{n,},\beta _{n,i})`$.
Since $`0<\beta _{l,i}<\alpha _l`$ for $`l=1,\mathrm{}n`$, by passing a further subsequence, we may assume that $`\beta _{l,i}=\beta _l`$, and finally we get $`M_i=(g;b_i;\alpha _1,\beta _1;\mathrm{};\alpha _n,\beta _n)`$. Moreover we may assume that all $`b_i0`$. Note that by \[p. 680 of LWZ\], all $`M_i`$ have the same first Betti number and the torsion part of $`H_1(M_i,)`$ is unbounded if $`b_i`$ unbounded. Since epimorphisms on $`\pi _1`$ induce epimorphisms on first homology groups, it follows that $`b_i`$โs are bounded. Now we have $`b_i=b_j`$ for some $`i,j`$, then $`M_i=M_j`$ and by the hopfian property of Seifert manifold groups, the epimorphism $`\pi _1(M_i)\pi _1(M_j)`$ is an isomorphism. Then in the sequence above there must be an isomorphism. Theorem 3.4 follows.โ
We have seen that Theorem 2.1 plays important roles for the proof Theorem 3.4. If the answer of Question 1.5 is also YES for hyperbolic 3-manifolds, this will lead to a positive answer for (C2) and (D) for hyperbolic 3-manifolds.
###### Proposition 3.5
Suppose Question 1.5 has a positive answer for hyperbolic 3-manifolds. Then for a given closed orientable hyperbolic 3-manifold $`M_0`$,
(1) any infinite sequence of epimorphisms $`\pi _1(M_0)\pi _1(M_1)\mathrm{}.\pi _1(M_n)\mathrm{}`$ contains an isomorphism, where all $`M_i`$ are closed orientable hyperbolic 3-manifolds.
(2) there are only finitely many closed orientable hyperbolic 3-manifolds $`M_i`$ with the same $`\pi _1`$-rank as that of $`M_0`$, for which there is an epimorphism $`\pi _1(M_0)\pi _1(M_i)`$.
###### Demonstration Proof
(1) By passing an infinite subsequence we may assume all $`\pi _1(M_i)`$ have the same rank. Since we assumed that Question 1.5 has a positive answer for hyperbolic 3-manifolds, this sequence is realized by a sequence of non-zero degree maps
$$M_0M_1\mathrm{}M_n\mathrm{}$$
The rest of the proof is now standard. Since all maps $`f_i:M_iM_{i+1}`$ in the sequence are of non-zero degree, by Gromovโs Theorem \[Chapter 6, Th\], $`v(M_i)v(M_{i+1})`$, where $`v(M_i)`$ is the hyperbolic volume of $`M_i`$. By Thurston-Jรธgensonโs Theorem \[Chapter 6, Th\], $`v(M_k)`$ must be a constant when $`k`$ is larger than a given integer $`N`$. Then by Gromov-Thurstonโs Theorem \[Chapter 6, Th\], $`f_k`$ is homotopic to a homeomorphism, $`k>N`$, so $`f_k`$ is an isomorphism
For (2) since we again assume that Question of 1.5 has a positive answer for hyperbolic 3-manifolds, each $`\varphi _i:\pi _1(M_0)\pi _i(M_i)`$ can be realized by a map of non-zero degree. By Somaโs theorem \[So\], there are only finitely many such $`M_i`$. โ
We also note the following partial positive answer of (D) of Question 1.6 follows easily from the methods of \[RW\].
###### Theorem 3.6
Suppose $`M`$ is a non-Haken hyperbolic 3-manifold. Then there are are only finitely many closed orientable hyperbolic 3-manifolds $`M_i`$ for which there is an epimorphism $`\pi _1(M)\pi _1(M_i)`$. โ
Section 4. A minimal Haken manifold
Let $`E`$ be the complement of trefoil knot with $`m`$ the meridian and $`l`$ the longitude. $`E`$ has a unique Seifert fibration with two singular fiber of indices 2 and 3, over the disc. Via this Seifert structure, we have a presentation
$$\pi _1(E)=<a,b,c,t|a^2t,b^3t,abc>$$
where $`t`$ is the regular Seifert fiber. Let $`E_1`$ and $`E_2`$ be homeomorphic to $`E`$ with meridian and longitudes $`(m_i,l_i)`$, $`i=1,2`$. Now glue $`E_1`$ to $`E_2`$ via a homeomorphism $`h:E_1E_2`$ such that $`h(l_1)=m_2`$ and $`h(m_1)=l_2^1`$. Let $`M`$ denote the resulting manifold, which is a closed graph manifold. The main theorem of this section is:
###### Theorem 4.1
$`M`$ is a minimal closed Haken 3-manifold among all 3-manifolds satisfying Thurstonโs geometric conjecture.
We begin the proof by collecting some elementary facts.
###### Lemma 4.2
(1) For any representation $`\varphi :\pi _1(E)SL(2,)`$, if $`\varphi (t)1`$, then the image $`\varphi (\pi _1(E))`$ is a cyclic group $`<\lambda >`$. Moreover, we must have $`\varphi (a)=\lambda ^2`$, $`\varphi (b)=\lambda ^3`$, $`\varphi (c)=\lambda ^5`$, and $`\varphi (t)=\lambda ^6`$.
(2) In $`\pi _1(T)`$, where $`T=E`$, we have $`m=tc^1`$ and $`l=t^5c^6`$. (Equivalently, $`t=6m+l`$ and $`c=5m+l`$.) Hence $`h(t_1^5c_1^6)=t_2c_2^1`$ and $`h(t_1c_1^1)=t_2^5c_2^6`$.
(3) $`M`$ is an integral homology 3-sphere.
(4) the only 2-sided incompressible surface is the incompressible torus $`T`$, which separates $`M`$ into $`E_1`$ and $`E_2`$.
###### Demonstration Proof
The main part of (1) follows from \[M, Prop. 3\] and the fact that $`H_1(E,)`$ is cyclic. (2) and (3) and the remaining parts of (1) are just direct calculations. Finally to establish (4) we observe the following. Since the trefoil knot is 2-bridge $`E`$ cannot contain a closed embedded essential surface by \[HT\]. If $`M`$ contained an embedded incompressible surface $`T`$, it would follow from the remark above and the gluing homeomorphism that $`E`$ would have a boundary slope $`1/0`$. However \[Theorem 2.0.3, CGLS\] then implies the existence of a closed embedded essential surface in $`E`$. โ
To show that $`M`$ is minimal, we assume not and suppose that there is a degree one map $`f:MN`$, where $`N`$ is irreducible, $`NM`$, and $`NS^3`$. First, since $`M`$ is a graph manifold, its Gromov norm is zero, so $`N`$ cannot be a hyperbolic 3-manifold by \[T, Chapter 6\]. Moreover it is well-known that $`N`$ must be an integer homology sphere (\[Lemma 3.1 RW\]. The proof of Theorem 4.1 will be finished by Lemmas 4.3, 4.4 and 4.5 below.
###### Lemma 4.3
$`N`$ is non-Haken.
###### Demonstration Proof
Suppose $`N`$ is Haken, and let $`FN`$ be an embedded incompressible surface. We may deform $`f`$ so that $`f^1(F)`$ is an incompressible surface in $`M`$. By (4) of Lemma 4.2 $`f^1(F)`$ must consist of parallel copies of $`T`$. By standard 3-manifold topology, we can further deform $`f`$ so that $`f^1(F)=T`$. It follows that $`F`$ is a 2-sphere or torus. Since $`N`$ is irreducible, $`F`$ must be a torus separating $`N`$ into two parts $`N_1`$ and $`N_2`$. Furthermore, the map $`f`$ can be decomposed into two proper degree one map $`f|:E_iN_i`$. However $`E_i`$ is a minimal 3-manifold among knot complements in 3-manifolds via proper degree one maps \[BW\]. Thus, each $`f|`$ is a homeomorphism, and it follows that $`f`$ itself is homotopic to a homeomorphism. โ
###### Lemma 4.4
$`N`$ is not a Seifert manifold with finite fundamental group (other than possibly $`S^3`$).
###### Demonstration Proof
By (3) of Lemma 4.2 if $`N`$ is a Seifert fibered manifold of finite fundamental group and $`NS^3`$, it must be the Poincarรฉ Homology 3-sphere $`P`$. Note $`\pi _1(P)`$ surjects onto $`A_5`$, the alternating group on 5 letters. In particular, (as is well-known) $`A_5`$ is a subgroup of $`PSL(2,)`$โsince $`SO(3)`$ can be identified with $`PSU(2)`$, and the latter is a subgroup of $`PSL(2,)`$. To prove the lemma, it suffices to prove that the image group of any representation of $`\varphi :\pi _1(M)PSL(2,)`$ cannot be $`A_5`$.
Case (1) If $`\varphi (t_1)1`$ and $`\varphi (t_2)1`$, by (1) of Lemma 4.2, the whole image $`\varphi (\pi _1(M))`$ must be a cyclic group (actually trivial).
Case (2) Without loss of generality, we may assume that $`\varphi (t_1)=1`$ and $`\varphi (t_2)1`$. By (1) and (2) of Lemma 4.2, $`\varphi :\pi _1(M)PSL(2,)`$ factors as $`\nu :\pi _1(M)G`$ and $`\mu :GPSL(2,)`$ where $`G`$ is generated by two groups described in (a) and (b) below:
(a) $`\nu (\pi _1(E_1))=<a_1,b_1,c_1|a_1^2,b_1^3,a_1b_1c_1>`$, (b) A cyclic group $`<\lambda _2>`$ such that $`\nu (c_2)=\lambda _2^5`$, $`\nu (t_2)=\lambda _2^6`$.
Since $`h(t_1c_1^1)=t_2^5c_2^6`$, we have $`\nu (h(c_1^1))=\nu (h(t_1c_1^1))=\nu t_2^5c_2^6)=1`$. It follows that
$$G=<a_1,b_1,c_1|a_1^2,b_1^3,a_1b_1c_1,c_1>=<a_1,b_1|a_1^2,b_1^3,a_1b_1>,$$
which is the trivial group.
Case (3) $`\varphi (t_1)=1`$ and $`\varphi (t_2)=1`$. In this case $`\varphi :\pi _1(M)PSL(2,)`$ factors through a group $`G`$ via a map $`\nu :\pi _1(M)G`$, with $`\nu (\pi _1(E_i)`$ is the quotient of $`G_i=<a_i,b_i,c_i|a_i^2,b_i^3,a_ib_ic_i>`$, $`i=1,2`$. Moreover by (2) of Lemma 4.2 we have that in the quotient $`c_1=c_2^6`$ and $`c_2=c_1^6`$. Immediately we have that $`c_1^{37}=1`$ and $`c_2^{37}=1`$ and finally
$$G=<a_i,b_i,c_i,i=1,2|a_i^2,b_i^3,a_ib_ic_i,c_i^{37},c_1=c_2^6,i=1,2>$$
Suppose there is a homomorphism $`\mu :GA_5`$. Since the order of $`c_i`$ is 37, and $`A_5`$ has order $`60`$, under the homomorphism $`\mu `$ the images of $`c_1`$ and $`c_2`$ must be trivial. It follows that $`\mu :GA_5`$ can factor through the group $`G^{}`$,
$$G^{}=<a_1,b_1,|a_1^2,b_1^3,a_1b_1><a_2,b_2,|a_2^2,b_2^3,a_2b_2>,$$
but as above this is trivial. โ
###### Lemma 4.5
$`N`$ is not a Seifert manifold with infinite $`\pi _1`$.
The proof of Lemma 4.5 requires a sequence of additional lemmas. Suppose below $`N`$ is a Seifert manifold of infinite $`\pi _1`$. By Lemma 4.3, we may assume that $`N`$ is non-Haken. Then $`N`$ must be a Seifert manifold with three singular fibers over $`S^2`$.
We begin by establishing:
###### Lemma 4.5.1
(1) Suppose $`\mathrm{\Delta }Iso_+H^2`$ is a triangle group and $`\varphi :\pi _1(2,3,l)\mathrm{\Delta }`$ is of finite index. Then the image of $`\varphi `$ is a hyperbolic triangle group isomorphic to $`\pi _1(2,3,k)`$, where $`k|l`$.
(2) Suppose a Serfert manifolds $`N`$ is an integer homology sphere with infinite $`\pi _1`$ and orbifold $`O=(a_1,a_2,a_3)`$. Then $`gcd(a_i,a_j)=1`$ for $`i,j=1,2,3`$, and $`O`$ is a hyperbolic orbifold.
###### Demonstration Proof
(1) Let $`x^{},y^{}`$ be the order 2 and order 3 elements which generates $`\pi _1(2,3,l)`$ such that $`x^{}y^{}`$ is of order $`l`$. Use $`x`$ and $`y`$ to denote their images in $`Iso_+H^2`$, then $`x`$ and $`y`$ generate the image of $`\varphi `$. Since the image of $`\varphi `$ is of finite index in $`\mathrm{\Delta }`$, it must be co-compact and of rank 2. By well-know fact then the image is a triangle group with $`x^2=y^3=(xy)^k=1`$, where $`k|l`$.
(2) follows from \[p. 680 (d) LWZ\]. โ
###### Lemma 4.5.2
There is a simple closed curve in the kernel of $`f|:TN`$.
###### Demonstration Proof
Since $`\pi _1(N)`$ is torsion free and $`T`$ is a torus, to prove the lemma, we need only that the kernel of $`f|:TN`$.
Suppose $`f(t_1)1`$, otherwise the claim is proved. Note that all elements in $`f(\pi _1(E_1))`$ commute with $`f(t_1)`$. If $`f(t_1)`$ is not the fiber $`t`$ of $`N`$, then either
$$f(\pi _1(E_1))=f(t_1)\text{or}f(\pi _1(E_1))=<f(t_1),f(c_1)>=.$$
The second case is not possible since $`H_1(E_1;)=`$. In the first case we deduce that $`\mathrm{ker}(f|_T)_{}`$ is nontrivial. Similarly if $`f(t_2)`$ is not the fiber $`t`$ of $`N`$, then $`\mathrm{ker}(f|_T)_{}`$ is non-trivial. If $`f(t_1)=t=f(t_2)`$. Since $`t_1`$ and $`t_2`$ do not coincide up to isotopy, still we have $`ker(f|_T)_{}`$ is non-trivial. โ
Let $`C`$ be the simple closed curve provided by Lemma 4.5.2. Suppose $`C=pm_1+ql_1`$ on $`E_1`$, then $`C=qm_2+pl_2`$. By (1) of Lemma 4.2 we have $`pm+gl=(p5q)t+(p+6q)c`$ and $`qm+pl=(q5p)t+(q+6p)c`$. So the degree 1 map $`f`$ factors through $`f:MN_1_{S^1}N_2N`$ where $`N_1`$ and $`E_2^{}`$ are Seifert manifolds whose normal forms are given by $`(2,1;3,1;p+6q,p5q)`$ and $`(2,1;3,1;q+6p,5pq)`$ respectively, and the two cores of the surgery solid tori are identified. If $`f|_{}(\pi _1(N_1))\pi _1(N)`$ and $`f|_{}(\pi _1(N_2))\pi _1(N)`$ then $`\pi _1(N)`$ can be presented as a non-trivial free product with amalgamation by the classical result (see \[CGLS\] for example). It follows that $`N`$ will is Haken contrary to Lemma 4.3. Thus without loss, we assume that $`f|_{}(\pi _1(N_1))=\pi _1(N)`$.
###### Lemma 4.5.3
$`f|_{N_2}`$ is of degree non-zero.
###### Demonstration Proof
Let $`\stackrel{~}{E}`$ be the covering of $`N`$ corresponding to $`f|_{}(\pi _1(N_2))`$. Then $`f:N_2N`$ lifts to $`\stackrel{~}{f}:N_2\stackrel{~}{E}`$, which is $`\pi _1`$-surjective. If $`f|_{}(\pi _1(N_2))\pi _1(N)`$ is of finite index, then $`\stackrel{~}{E}`$ is a closed Seifert manifold. Since both $`\pi _1(N_1)`$ and $`\pi _1(N)`$ are rank 2, $`\pi _1(\stackrel{~}{E}))`$ must be also rank 2. Then $`\stackrel{~}{f}`$ is of non-zero degree by Theorem 2.1. Hence $`f|_{N_2}`$ is non-zero degree.
Below we show $`f|_{}(\pi _1(N_2))\pi _1(N)`$ must be of finite index. Otherwise $`\stackrel{~}{E}`$ is a non-compact, aspherical Seifert manifold, which is known that either the rank of $`H_1(\stackrel{~}{E})`$ is positive or $`\pi _1(\stackrel{~}{E})`$ is trivial. Since $`f|_{}(\pi _1(N_2))`$ is not trivial and $`N_2`$ is a rational homology sphere, all of the above cases are ruled out. So $`f|_{}(\pi _1(N_2)`$ must be of finite index in $`\pi _1(N)`$. โ
Since $`N_1`$ and $`N_2`$ are in symmetry position, we have both $`f|N_1`$ and $`f|N_2`$ are of non-zero degree.
By Lemma 4.5.3, we may assume that $`f|N_i`$ is fiber preserving. Then $`f|N_i`$ induces an homoporphism $`\varphi _i:\pi _1(O_i)\pi _1(O)`$, in particular $`\varphi _1`$ is surjective and $`\varphi _2`$ is finite index, where $`O_1=(2,3,6qp)`$, $`O_2=(2,3,6q+p)`$ and $`O=(a_1,a_2,a_3)`$ are orbifolds of $`N_1`$, $`N_2`$ and $`N`$ respectively. Since $`N`$ is an integer homology sphere of infinite $`\pi _1`$, it follows that $`\pi _1(O)`$ is isomorphic to a hyperbolic triangle group. Since $`\varphi _1:G_1G`$ is surjective, it follows that $`O=(2,3,k)`$, where $`k|6qp`$ by Lemma 4.5.1 (1). Since $`\varphi _2`$ is of finite index, the image of $`\varphi _2`$ is a hyperbolic triangle groups $`\pi _1(2,3,k^{})`$ with $`k^{}|6q+p`$ by Lemma 4.5.1 (1), moreover $`k^{}|k`$. It is easy to see that $`k^{}`$ is a dvisor of both $`12q`$ and $`2p`$. Since $`p`$ and $`q`$ are coprime, the great common divisor of $`12q`$ and $`2p`$ is 12. So $`k^{}`$ is either 2, or 3, or 4, or 6, or 12. Then $`N`$ can not be an integer homology sphere by 4.5.1 (2). โ
REFERENCES
\[BW\] M. Boileau and S. C. Wang Non-zero degree maps and surface bundles over $`S^1`$, J. Diff. Geom. 43 (1996), 789โ908.
\[BZ\] M. Boileau and H. Zieschang, Heegaard genus of closed orientable Seifert 3-manifolds, Invent. Math. 76 (1984) 455โ468.
\[CGLS\] M. Culler, C. McA. Gordon, J. Luecke and P. B. Shalen Dehn surgery on knots, Annals of Math. 125 (1987) 237โ300.
\[CR\] P. J. Callahan and A. W. Reid, Hyperbolic structures on knot complements, Chaos, Solitons and Fractals, 9 (1998), 705โ738.
\[H\] J. Hass, Minimal surfaces in manifolds with $`S^1`$ actions and the simple loop conjecture for Seifert conjecture for Seifert fiber spaces. Proc. AMS, 99, 383-388 (1987)
\[HT\] A. Hatcher and W. Thurston. Incompressible surfaces in 2-bridge knot complements. Invent. Math. 79 (1985) 225โ246.
\[Hu\] C. Huang, Master Thesis, Peking Univ. (1998)
\[J\] W. Jaco, Topology of three manifolds, Regional Conference Series in Mathematics. 43, AMS Math. Soc. Providence, RI 1980.
\[K\] R. Kirby, Problems in low dimensional topology, Geometric Topology, Edited by H. Kazez, AMS/IP 1997.
\[LWZ\] C. Hayat-Legrand, S. C. Wang and H. Zieschang, Minimal Seifert manifolds, Math. Annalen. 308 (1997), 673โ700.
\[MR\] C. Maclachlan and A.W. Reid, Generalised Fibonacci Manifolds, Transformation groups, 2 (1997), 165โ182.
\[M\] K. Motegi, Haken manifolds and representations of their fundamental groups in $`SL(2,)`$, Topology and its Applications 29 (1988), 207โ212.
\[RW\] A.W. Reid and S. C. Wang, Non-Haken 3-manifolds are not large with respect to mappings of non-zero degree, to appear in Comm. Analysis $`\&`$ Geom.
\[Ro1\] Y. Rong, Maps between Seifert fibered spaces of infinite $`\pi _1`$, Pacific J. Math. 160 (1993), 143โ154.
\[Ro2\] Y. Rong, Degree one map between geometric 3-manifolds, Trans. A. M. S. 332 (1992), 411โ436.
\[Sc\] G. P. Scott, The geometries of 3-manifolds, Bull. London Math. Soc. 15 (1983), 401โ487.
\[So\] T. Soma, Maps of non-zero degree to hyperboilc 3-manifolds. To appear J. Diff. Geom.
\[T\] W.P. Thurston, The Geometry and Topology of 3-manifolds, Princeton University mimeographed notes (1979).
\[W\] S.C. Wang, The existance of maps of non-zero degree between aspherical 3-manifolds, Math. Zeit. 208 (1991), 147โ160.
\[ZVC\] H. Zieschang, E. Vogt and H. Coldeway, Surfaces and Planar Discontiuous Group, Lecture Note in Math. 835 Springer-Verlag Berlin, 1980.
Reid
Department of Mathematics,
University of Texas, Austin, TX 78712, USA
Wang
Department of Mathematics,
Peking University, Beijing, 100871, P. R. China.
Zhou
Department of Mathematics,
The East China Normal University,
Shanghai, 200062, P.R. China |
warning/0002/nucl-th0002053.html | ar5iv | text | # Implications of Quark-Lepton Symmetry for Neutrino Masses and Oscillations
(February 22, 2000)
A principal feature of all of the experimentally known fundamental fermions may be summed up in the following manner: The fermions may be grouped into three โfamiliesโ, conventionally assigned by mass, consisting of two color triplets of quarks, one with electric charge $`+2/3`$ and the other $`1/3`$, one lepton with electric charge $`1`$, their antiparticles and a neutrino. Each โfamilyโ of these fermions fills out fifteen $`(\frac{1}{2},0)`$ representations of the Lorentz group. Fourteen of these come in pairs with conjugate color and (electric) charge quantum numbers so that they may be reconstructed into seven Dirac bispinor representations. This is accomplished by using charge conjugation under the Lorentz group<sup>1</sup><sup>1</sup>1This should not be confused with the larger CP operation used earlier in connection with Majorana objects . to convert one member of each pair into the requisite $`(0,\frac{1}{2})`$ representation. Despite this construction, and the fact that these pairs are not conjugate in their electroweak quantum numbers, the Lagrangian of the Standard Model (SM) does not intrinsically violate conservation of the weak interaction quantum numbers, courtesy of the chiral projections included in the interactions.
The exceptional case is that of the neutrino, which has no known partner representation. That such a partner should exist has long been suggested and is especially evident in โvector-likeโ Grand Unified Theories (GUTs). A satisfactory explanation of the stringent bounds on the mass of each of the three different flavors of neutrinos has been developed in this context in terms of the so-called โsee-sawโ mechanism. This mechanism postulates that since the missing partner $`(\frac{1}{2},0)`$ representation carries no SM quantum numbers at all, and so is โsterileโ with respect to all SM interactions, it may naturally acquire a large (GUT scale) Majorana mass. While the usual (active) $`(\frac{1}{2},0)`$ representation could also, in principle, develop a Majorana mass, the associated scalar field must carry weak isospin, $`I_W=1`$, so that mass is usually assumed to be zero. The effect of this is to suppress the induced Majorana mass of the neutrinos active in the SM from the value common for Dirac fermion masses in the SM by a factor of the ratio of such masses to something very roughly on the order of the GUT-scale mass<sup>2</sup><sup>2</sup>2Note that, in the usual discussions of the โsee-sawโ, the mass term that couples the active $`(\frac{1}{2},0)`$ representation to itself as a Lorentz charge conjugate $`(0,\frac{1}{2})`$ representation is referred to as $`m_L`$ while the similar term for sterile neutrinos is referred to as $`m_R`$. Although that notation is natural under the assumption of Dirac neutrinos, here, since the neutrinos under discussion are massive Weyl (Majorana), that identification can lead to confusion. Hence, we use โactiveโ and โsterileโ throughout this paper..
While the prospect of providing a โnaturalโ explanation for the small scale of active neutrino masses is pleasing, there is no principle requiring that the sterile mass be large. Here we discuss an equally valid scenario based on a view of quark-lepton symmetry, which is subject to a general experimental test that will soon be undertaken.
We start from the facts that the known Dirac fermion masses span a range of almost six orders of magnitude and that those of the neutrinos must be at least five to six orders of magnitude smaller still. We allow for the possibility that the true origin of these masses is still not understood and set aside the see-saw. We next recall that it is charge-conservation, for various charges, which eliminates the possibility of Majorana mass terms for each of the fourteen spinor representations that make up the known Dirac bispinors.
Now we recall an old conjecture: That there is indeed a $`(\frac{1}{2},0)`$ representation for a sterile neutrino to form an eighth pair with the known active neutrino for each generation (or family) of fermions. Recalling that all other individual fermion number (baryon number, muon number, etc.) violations seem to be strongly suppressed, we are led to examine the possibility that this is true for neutrinos as well. This leads naturally to the conclusion that, while the sterile neutrino may have a Majorana mass, it should be expected to be small compared to the Dirac mass available to the pair of neutrino representations which can be formed into a Dirac bispinor. We must then simply accept the fact that the Dirac mass for this bispinor is, itself, very small to satisfy experimental constraints, although there are interactions which could exist that would modify the interpretation of these constraints.
Many have conjectured that there should be some parallel (for example, right-chiral interaction) quantum number so that some sort of neutrino number remains. A related point has been made by Cahill, that a Majorana neutrino mass would violate lepton number (L) conservation, and hence also the difference between that and baryon number (B). The experimentally reported suppression of proton decay provides strong support for an assumption that matrix elements for B-L violation are extremely tiny. Thus, we consider it viable to investigate the implications of the point of view that the Majorana mass terms of neutrinos are also quite small compared to Dirac neutrino mass terms.
This very general scenario leads to a quite well-defined class of predictions for neutrino properties and experiments. Denoting the Dirac mass connecting the active and sterile neutrinos by $`M`$ and the Majorana mass of the sterile neutrino by $`m`$, the conjecture that $`m<<M`$ leads to the conclusion that the neutrino states are pseudo-Dirac.<sup>3</sup><sup>3</sup>3Here we use pseudo-Dirac to mean a pair of Majorana neutrinos with masses so nearly degenerate that, for many purposes, the linear combinations appropriate to a particle and an antiparticle are, to a good approximation, eigenstates of the mass matrix. This, we believe, is the general usage. Wolfenstein introduced the term to refer to a particular model in which the two $`(\frac{1}{2},0)`$ representations used to produce the Dirac bispinor were both active under the weak $`SU(2)`$, but were coupled to different charged leptons. That is, the neutrino field eigenstates will be a pair of almost degenerate states and a neutrino will propagate almost as if it were a Dirac fermion. However, with a long oscillation length, it will transform between the active and sterile components, with almost maximal mixing (due to the almost complete degeneracy of the resulting Majorana mass eigenstates). This would be most simply expected to be true for each neutrino type separately. (See also Ref..)
The immediate implication is that when maximal mixing is observed in neutrino oscillations, it will be between active and sterile types. Hence, the long baseline from the Sun leads to the conclusion that if the signal diminution observed in solar neutrino experiments is due to vacuum neutrino oscillations, then the oscillation that is occuring is from active electron neutrinos to sterile (anti)neutrinos. It follows that the SNO experiment is predicted to observe a reduction in the neutral current signal equal to that already found in the charged current signal. This prediction has also been made in Ref..
Similarly, in the observed atmospheric oscillations, the long baseline suggests that the oscillation is from muon to sterile neutrinos. Although this is not favored by the current data set, neither is it inconsistent at present. The scenario discussed here predicts that additional data will find a diminishing signal for active-active oscillation.
We should, however, note a possibility which is difficult to encompass within unified models, but may nonetheless occur: The conjugate partner to the active neutrino representation of one family may be the active neutrino representation of a different family. This would appear to violate quark-lepton symmetry and leaves one uncertain about whether or not there must be sterile partner representations. However, this possibility matches more closely with the preferred interpretation of the observed atmospheric oscillations, if the pair of families involved are those of the muon and tauon. Note that this still implies that the SNO experiment would observe the same reduction of the neutral current signal as of the charged current signal because the solar neutrino oscillation would still necessarily involve a sterile neutrino partner. This conjecture raises the question of whether or not some vestige of quark-lepton symmetry obtains in the form of two additional sterile neutrino representations that mix only with each other, perhaps still forming a pseudo-Dirac bispinor.
We also note that there is the possibility of a modified โsee-sawโ, in which the sterile neutrino mass matrix in family space may have a large scale but a rank less than three. In this case one or two of the families may have neutrinos that are poorly described as pseudo-Dirac without affecting the remaining families. For example, if the rank is one, there are two zero eigenvalues of the $`m`$ matrix, were it diagonalized by itself. The embedding of that matrix in the larger mass matrix for neutrinos can easily change two Dirac neutrinos into pseudo-Dirac neutrinos, or could lead to Majorana neutrinos with masses well-separated on the scale of the Dirac masses.
A priori, no definite predictions are made for flavor oscillations of the type reported to be observed by the LSND collaboration. However, as the Dirac mass terms are larger than the Majorana mass terms in the particular scenario discussed here, and the Dirac mass terms may be presumed (on the basis of quark-lepton symmetry) to be analogous in structure to those found in the quark sector, flavor mixing should be expected to occur with shorter oscillation lengths and modest mixing amplitudes, i.e., much less than maximal. This is certainly consistent with the experimental reports to date. Additional support for this scenario of small flavor mixing between pseudo-Dirac neutrinos may be found in recent discussions of possible interference effects modifying the end point spectrum in Tritium beta decay, when taken in combination with existing limits on Majorana neutrino masses from neutrinoless double beta decay experiments.
Finally, we conclude that the strength of the neutral current signal of the SNO experiment is crucial to determining the viability of any pseudo-Dirac bispinor scenario.
We acknowledge stimulating conversations on this subject with Bill Louis. This research is partially supported by the Department of Energy under contract W-7405-ENG-36, by the National Science Foundation and by the Australian Research Council. |
warning/0002/astro-ph0002434.html | ar5iv | text | # Quantum Perturbative Approach to Discrete Redshift.
## 1 Introduction
### 1.1 evidence of discrete structure
There are at least four types of evidence of discrete structure on the largest scales. The first example of discrete structure is given by discrete redshift, the evidence for this will be discussed in the next paragraph. Redshift comes in discrete values with the characteristic velocity $`v_I=72.2\pm 0.2\mathrm{Km}.\mathrm{s}^1`$, and this leads to a characteristic length $`l_{dr}=v_IH_0^1=3\pm \mathrm{0.8.10}^{22}`$ meters, and a characteristic period $`t_{dr}=v_IH_0^1c^1=3.2\pm \mathrm{0.8.10}^6`$ years. Here this result is used in preference to the other examples of discrete structure because of the advantage of having a qualitative result, namely $`v_I=72.2\pm 0.2\mathrm{Km}.\mathrm{s}.^1`$; the actual techniques used might be applicable to the other cases. A second example is that the Universe appears to consist of superclusters and voids Saunders et al (1991) recent studies of clustering can be found in Cohen (1999) , superclusters and voids occur with apparent regularity (Broadhurst et al (1990) ). Assuming $`q_0=\frac{1}{2}`$ they have a characteristic scale of $`l_{sv}=128\mathrm{h}^1\mathrm{Mpc}=5.3\pm \mathrm{2.6.10}^{24}`$ meters, and hence a characteristic period of $`t_{sv}=\frac{l_{sv}}{c}=5.6\pm 10^8`$ years. Tytler et al (1993) do not confirm Broadhurst et alโs result: instead of Broadhurst et alโs โapparent regularity with a scale of $`128\mathrm{h}^1\mathrm{Mpc}.`$โ they find โThere is no significant periodicity on any scale from $`10`$ to $`210\mathrm{h}^1\mathrm{Mpc}.`$โ The result was also looked at by Willmer et al (1994) . Einasto et al (1997) present evidence for a quiasiregular three-dimensional network of rich superclusters and voids, with regions separated by $``$120Mpc.; and they say that โif this reflects the distribution of all matter, then there must exist some hitherto unknown process that produces regular structure on large scalesโ. A third example of discrete structure is that some normal elliptical galaxies have giant shells surrounding them, Malin and Carter (1980) . These shells probably consist of stars, the most likely method of their formation is from an intergalactic shock wave or an explosive event in the galaxy. The fourth example is of discrete properties from geological time scales. The characteristic periods $`t_{dr}`$ and $`t_{sv}`$ are larger than typical geological frequencies. For example Kortenkamp and Dermott (1998) give a periodicity of $`t_d=10^5`$ years for the accretion rate of interplanetary dust. The rate of accretion of dust by the Earth has varied by a factor of 2 or 3. Extraterrestrial helium-3 concentrations in deep sea cores display a similar periodicity but are $`5.10^4`$ years out of phase. The magnetic polarity time sequence (p.672 of Larson and Birkland (1982) ) gives reversals in the earths magnetic field occurring at intervals $`t_m=10^3\mathrm{to}5.10^3`$ years; and the period of recent glacier advance and retreat is about $`t_g=2.10^3`$ years, (p.488 Larson and Birkland (1982) ). Naidu and Malmgren (1995) give a periodicity of $`t_m=2,200`$ years for the Asian monsoon system, this is found by looking at fluctuations in upwelling intensity in the western Arabian sea. Clearly it would be good if there was an explanation for these characteristic periods and even better if they could be used to predict presently unknown quantities. The obvious factor to try to predict is the size of the Universe as given by the scale factor $`R`$, assuming a Robertson-Walker Universe this is equivalent to deriving a value of the deceleration parameter $`q`$. Now the value of the deceleration parameter is assumed in the derivation of $`t_{sv}`$, thus models using the characteristic period of discrete redshift are studied.
### 1.2 observations of discrete red shift
Perhaps the first paper advocating that redshift can only occur in specific discrete values was Cowan (1969) who found a periodic clustering of redshifts. This was confirmed by Karlsson (1971) who found a number of new peaks in the distribution of quasi-stellar objects; these, together with the peaks at $`z=1.956`$ and $`z=0.061`$ formed a geometric series: $`z_1=1.96,z_2=1.41,z_3=0.96,z_4=0.060,z_5=0.30,z_6=0.06`$; this was supported by Arp et al (1990) . Green and Richstone (1976) did a search for peaks and periodicities in the redshift distribution of a sample of quasars and emission-line galaxies independent of that used in earlier work. In agreement with the results of Burbidge and OโDell (1972) , no statistically significant peak was found at a redshift of $`1.95`$, nor any significant periodicity in redshift in either the sample of quasars alone or the sample of quasars and galaxies together. The strong spectral power peak in their distribution of galaxy redshifts, estimated at a confidence level of $`97.5`$ percent, is completely absent in Green and Richstoneโs analysis; they conclude that the observed redshift distribution is consistent with a random sample of discrete values from a smooth, aperiodic underlying population. Tifft (1976) claimed that well known local galaxies, especially M31, were claimed to consist of two basically opposed streams of outflow material which have an intrinsic difference in red shift of $`7075\mathrm{Km}.\mathrm{s}^1`$. Tifftsโ result was questioned by Monnet and Deharneny (1977) who say that the two opposing streams of material suggest that the best galaxy candidate for a direct test of Tifftsโ result would be a near face-on galaxy: the Doppler effect due to the rotation is minimized and any expansive motion in the spiral arms, as connected with the uneven distribution of neutral gas, which can be of the order of $`1025\mathrm{Km}.\mathrm{s}.^1`$ at most, gives entirely negligible Doppler effect. On the other hand a real intrinsic redshift would still exhibit its full $`7075\mathrm{Km}.\mathrm{s}.^1`$ discontinuity. Monnet and Deharneny chose the nearly face-on galaxy NGC628 and find no intrinsic effects as predicted by Tifft and a very smooth velocity field, with a velocity dispersion $`12\mathrm{Km}.\mathrm{s}.^1`$. In Tifft (1977a) it was claimed that redshift differentials between pairs of galaxies and between galaxies in clusters take preferred values which are various multiples of a basic $`72.5\mathrm{Km}.\mathrm{s}^1`$. In Tifft (1977b) the effect was studied for abnormal galaxies. In Tifft (1978a) it was claimed that the asymmetry in galaxy H1 profiles can be related directly to the properties of discrete redshift. In Tifft (1978b) the concept of discrete redshift was applied to dwarf H1 redshifts and line profiles, and a model of redshift based upon the ultimate discrete levels spaced near $`12\mathrm{Km}.\mathrm{s}^1`$ was developed. Most optical redshift data are not accurate to show discreteness directly, but it is claimed that 21 cm. data on double galaxies show the effect clearly; Tifft (1980) and Tifft (1982a) using the radio data taken by Peterson (1979) at the NRAO 300โ telescope has claimed that the effect is present with a confidence level of $`99.9\mathtt{\%}`$. Cocke and Tifft (1983) have claimed that the effect is present in the 21 cm. data on compact groups of galaxies taken by Haynes (1981) and Heleou et al (1982) at the Arecibo telescope with a confidence level of $`99.5\mathtt{\%}`$. The optical data of Tifft (1982b) also shows the effect strongly. The Fisher and Tully (1981) survey of 21 cm. redshift was found by Tifft and Cocke (1984) to show sharp periodicities at exact multiplies ($`\frac{1}{3}`$ and $`\frac{1}{2}`$) of $`72.45\mathrm{Km}.\mathrm{s}^1`$; the periodicity at $`24.1\mathrm{Km}.\mathrm{s}^1`$ involves galaxies with narrow 21 cm. profiles and the $`\frac{1}{2}`$ periodicity at $`36.2\mathrm{Km}.\mathrm{s}^1`$ involves galaxies with wide profiles, and there appears to be a progression of periods $`243672`$ for galaxies with higher 21 cm. flux levels as the profile width increases. Arp and Sulentic (1985) , using data from the Arecibo telescope of over $`100`$ galaxies in more than $`40`$ groups found that: companion galaxies have a higher redshift than the dominant galaxy of the cluster, and that the difference in redshift between the dominant and companion galaxies occurs in multiplies of $`70\mathrm{Km}.\mathrm{s}^1`$. Sharp (1984) suggested that for double galaxies discrete redshift might be just a statistical effect. Newman et al (1989) suggested that the effect is gaussian random noise and that at least one order of magnitude more data is needed to confirm the effect. Crousdaleโs (1989) study supports Tifftsโ work. Mirzoyan and Vardanyan (1991) claim that the values of the redshifts found preferentially in quasistellar objects essentially coincide with the redshifts for which the strong emission lines of Mg II, C IV, Ly $`\alpha `$, in the spectra of these objects fall close to the maximum sensitivity of the U, B, and V light filters. In this case their effect on the conditions for observing quasars is decisive and causes the quasar redshifts to be discretized. Based on a comparison between the observed quasar redshifts and the expected values assuming that this explanation is correct, they conclude that the observed effect of quasar redshift discretization is caused by observational selection. Guthrie and Napier (1991) confirm the effect in near by galaxies, but whereas Tifft finds 24.2, 36.3, or 72.5 Km.$`\mathrm{s}^1.`$ they find 37.2 Km.$`\mathrm{s}^1`$. Kruogovenko and Orlov (1992) find a periodic cycle of Seyfert and radio galaxies of about $`30\mathrm{h}^1`$Mpc between shells. Holba et al (1994) find a non-negligible region where two quasar samples and the galaxy sample are simultaneously fairly periodic. Guthrie and Napier (1996) confirm redshift periodicity in the local supercluster. Khodyachikhโs (1996) findings contradict the explanation of periodicity by selection effects. Tifft (1996) finds 72 and 36 Km.$`\mathrm{s}^1`$ periodicity for galaxy samples from the Virgo cluster, the Perseus and Cancer supercluster regions, and local space. Tifft (1997) studies the redshift of local galaxies for quantization and finds that ordinary spiral galaxies with 21 cm. profile widths near 200 Km.$`\mathrm{s}^1`$ show periodic redshifts. A review of redshift periodicities has been given by Narlikar (1992) , he finds that different data sets of extragalactic objects including nearby and distant galaxies and quasars that show statistically significant peaks at periodic intervals of redshift; he says at present the data is not complete in any sense but they are substantial enough to make us worry about the fundamental assumption that the Universe is homogeneous on a large scale. Moreover, he claims evidence of this kind has not only persisted in spite of rigorous statistical analysis but has grown with time so that it cannot be altogether ignored. Recent conference proceedings covering some aspects of discrete redshift, Pitucco et al (1996) and Dwari et al (1996) . Newman et al (1994) and Newman and Terzian (1996) analyses the โPower Spectrum Analysisโ (PSA) of Yu and Pebbles (1969) , this is the statistical analysis used in discrete redshift studies. They find that this method generates a sequence of random numbers from observational data which, it was hoped, is exponentially distributed with unit mean and unit variance. The variable derived from this sequence is approximately exponential over much of its range but the tail of the distribution is far removed from an exponential distribution, so that statistical inference and confidence testing based on the tail of the distribution is unreliable. Newman and Terzian go on to say that there are six claims of the PSA method which are wrong or involve some hidden assumptions. For purposes of illustration consider the first of these. Let $`N`$ points $`x_j`$ be distributed in the interval $`0`$ to $`2\pi `$ and let
$$z_n=N^{1/2}\underset{j=1}{\overset{N}{}}\mathrm{exp}(inx_j).$$
(1)
Yu and Peebles claim that the ensemble average of $`z_n(n0)`$ is
$$<z_n>=N^{1/2}<\mathrm{exp}(inx_j)>=N^{1/2}_0^{2\pi }\frac{dx_j}{2\pi }\mathrm{exp}(inx_j)=0$$
(2)
Apparently the correct way to approach this is to suppose that the $`x_j`$ are identically distributed and independent deviates with distribution $`๐ซ(x)`$ and to define a characteristic generating function $`(n)`$ by
$$(n)<\mathrm{exp}(inx)>=_{\mathrm{}}^{\mathrm{}}\mathrm{exp}(inx)๐๐ซ(x),$$
(3)
then
$$<z_n>=N^{1/2}(n).$$
(4)
If the underlying distribution function $`๐ซ(x)`$ is normally distributed (gaussian) with a mean $`\mu `$ and a variance $`\sigma ^2`$ then
$$<z_n>=N^{1/2}\mathrm{exp}(\frac{n^2\sigma ^2}{2}).$$
(5)
Newman and Terzian also say that astronomers choose too small a frequency class interval (bin with) in their frequency histograms, and that the optimal is $`1+\mathrm{log}_2(N)`$, and this is illustrated by their smoother figurers for the larger interval.
### 1.3 theoretical explanations
Theories to explain discrete redshift have been devised by Cocke and Tifft (1983) and Cocke (1985) , Nieto (1986) , and Buitrago (1988) . These theories depend on the introduction of a redshift quantum mechanical operator; this is essentially equivalent to replacing the emitterโs four-momentum $`P^a`$ by an operator. Narlikar and Burbidge (1981) produce an explanation which has a two component model of the Universe with discrete matter. An explanation was devised by Barut et al (1994) where: โIt is shown that the energy distribution in this model is periodic and the periods and density decrease with increasing distance, in striking agreement with experimental data.โ. Discrete redshift of a quantum energy spectrum in an anisotropic universe was found by Lamb et al (1994) . Discrete redshift might be caused be dislocation solutions to equations, such solutions have been described by Edelen (1994) . Arp (1996) argues that there is a โsignal carrierโ for inertial mass, which he calls the machion, and this gives rise to periodicity. Greenberger (1983) and Carvalho (1985) generalizes the quantum commutator to $`[q,p]=i\mathrm{}+i\mathrm{}f(q,p)`$, where $`f(q,p)`$ is a function, Carvalho (1997) uses this to calculate a redshift spectrum with discrete values. Hill et al (1990) constructed three alternative models involving oscillating physics: i) an oscillating gravitational constant, this was also studied by Salgado et al (1996) , ii) oscillating atomic electron mass, this was ruled out by Sudarsky (1992) on the basis of the Braginsky-Panov experiment, and iii) oscillating galactic luminosities. A varying Hubble parameter has been used by Morikawa (1991) as an explanation. A Voroni cellular model was used by van de Weygaert (1991) to explain the Broadhurst et al result. Ikeuchi and Turner (1991) also use a three dimensional Voroni tessellation model to explain voids and walls. Williams et al (1991) find that the Voroni foam model predicts a scale length for the galaxy-galaxy correlation function which is too large. Hill et al (1991) suggest a particular coherent sinusoidal peculiar velocity field of amplitude $`\delta /c3\times 10^3`$ and wavelength $`\lambda 128h^1`$Mpc could explain the result of Broadhurst et al. An oscillating gravitational constant model was constructed by Crittenden and Steinhart (1992) to explain the Broadhurst et al result. Dekel et al argue that the Broadhurst et al result suggests a large scale origin for periodicity. Budinick et al (1995) built another model to explain the Broadhurst et al result. The existence of superclusters and voids is perhaps explained by the cold dark matter theory of White et al (1987) . Redshift periodicity can be used to probe the correctness of general relativity Faroni (1997) . Lui and Hu (1998) suggest an explanation which has the mean free path of heavy elements absorb system varies regularly with cosmic time. Farhi et al (19980 technique might provide an explanation. Some papers on the quantum mechanics of large macroscopic systems are: Greenberger (1983) , Agnese (1984) , DerSarkissan (1984) and (1985) , Carvalho (1985) , Silva (1997) Capozziello et al (1998) , and Carneiro (1998) .
### 1.4 discrete as opposed to continuous energy spectrum
It is a common error in non relativistic quantum mechanics to suppose that quantization implies that there must be a discrete energy spectrum, see for example Schiff (1949) page 34. For an infinite potential $`V(r)`$ there are discrete energy levels; however for a potential with $`E>V`$ there are continuous energy levels when $`E>V`$. Usually however a discrete spectrum is usually an indication of a system with quantum rather than classical properties. By analogy with non relativistic quantum mechanics a closed gravitational interacting system would be expected to display discrete properties at low energies, i.e. when the gravitational field is weakly interacting, rather than at high energies when the gravitational field is strong. This suggests that cosmology and extra-galactic astronomy - as opposed to particle physics, are the subject areas where a single theory combining quantum mechanics and gravity would be necessary. There are many unusual dynamical properties of large scale systems, such as galaxies, see for example Roberts (1991) and references therein, and it might be that these are directly attributable to quantum corrections to classical theory; however here attention is restricted to systems that are characterized by discrete properties rather than unusual dynamics. Bohr-Sommerfeld quantization rule have been applied to gravitational systems: Wereide (1923) applied them to spherically symmetric spacetimes to find the line element of the electron, and Agnese and Testa (1997) applied them to planetary orbits. There are discrete approaches to quantum gravity Loll (1998) , and it might be that discrete structure from near the Planck era is inflated to present day large scale discrete structure. Random walk models suggest a fundamental length of $`l10^{35}`$cm. Sidharth (1999) .
### 1.5 sectional contents
In section 2 the exact solution for the equation of state $`p=(\gamma 1)\mu `$ Robertson-Walker spacetime is derived. This is done because the critical parameter $`2q_0`$ occurs in latter approximations to the Klein-Gordon equation and these exact solutions clarify when the occurrence of this critical parameter is an artifact of the approximations involved; also the standard approach to redshift in Robertson-Walker spacetime is given. Section 3 consists of general remarks on what properties a theory of discrete redshift would be expected to have and discusses various approaches to constructing a theory. Section 4 is devoted to finding approximate solutions to the Klein-Gordon equation in Robertson-Walker spacetimes. Section 5 applies these approximate solutions to theories of discrete redshift where the discreteness of radiation originates from the motion of the emitter. Section 6 discusses how discrete redshift might occur via the massive Klein-Gordon equation. Section 7 uses the solutions for the Klein-Gordon equation in the Einstien static universe to induce weak field metric perturbations of Robertson-Walker spacetime, these weak metric perturbations are used to demonstrate discrete redshift.
### 1.6 conventions
The conventions used are: signature $`+++`$, early latin indices $`a,b,c\mathrm{}=0,1,2,3`$, middle latin indices $`i,j,k\mathrm{}=1,2,3`$, Riemann tensor
$$R_{.bcd}^a=2_{[c}\mathrm{\Gamma }_{d]b}^a+2\mathrm{\Gamma }_{[c|f|}^a\mathrm{\Gamma }_{d]b}^f,$$
(6)
Ricci tensor
$$R_{bd}=R_{.bad}^a,$$
(7)
commutation of covariant derivatives
$$X_{ab;cd}X_{ab:dc}=X_{ae}R_{.bcd}^e+X_{eb}R_{.acd}^e.$$
(8)
Relativistic units are not used; $`c`$, $`G`$, and $`\mathrm{}`$ are put explicitly into equation where appropriate, this is done in order to clarify the relative size of the terms in the approximations used, $`c`$ is included in the definition of the Hubble constant 17
## 2 Robertson-Walker Spacetime.
### 2.1 the line element
The Robertson-Walker line element is
$$ds^2=c^2N(t)^2dt^2+R(t)^2d\mathrm{\Sigma }_3^2,$$
(9)
where $`R`$ is the scale factor and
$`d\mathrm{\Sigma }_3^2`$ $`=`$ $`d\chi ^2+s(\chi )^2(d\theta ^2+\mathrm{sin}^2\theta d\varphi ^2),`$
$`\mathrm{with}`$ (10)
$`s(\chi )`$ $`=`$ $`\{\begin{array}{ccc}\hfill \mathrm{sin}(\chi )& \mathrm{for}& k=+1\hfill \\ \hfill \chi & \mathrm{for}& k=0\hfill \\ \hfill \mathrm{sinh}(\chi )& \mathrm{for}& k=1.\hfill \end{array}`$
The lapse function $`N`$ is arbitrary and depends on the time coordinate used. Three common choices for the lapse function are: i) $`N=1`$ for which the time coordinate is the same as the proper time of a co-moving observer, ii) $`N=R`$ for which the line element is conformal to the Einstein static Universe, and the time coordinate referred to as conformal time, and iii) $`N=R^3`$ for which $`\mathrm{\Gamma }^tg^{ab}\mathtt{\{}_{ab}^t\mathtt{\}}=0`$, and the time coordinate referred to as harmonic time. For $`k=+1`$ using the coordinate transformations, Mรผller (1939)
$`\mathrm{sin}(\alpha )`$ $`=`$ $`\mathrm{sin}(\theta )\mathrm{sin}(\chi ),`$
$`\mathrm{cos}(\beta )`$ $`=`$ $`\sqrt{1+\mathrm{cos}^2(\theta )\mathrm{tan}^2(\chi )},`$
$`\mathrm{cos}(\theta )`$ $`=`$ $`\mathrm{sin}(\beta )\sqrt{\mathrm{sin}(\beta )+\mathrm{tan}^2(\alpha )},`$
$`\mathrm{cos}(\chi )`$ $`=`$ $`\mathrm{cos}(\alpha )\mathrm{cos}(\beta ),`$ (13)
the three-sphere line element becomes
$$d\mathrm{\Sigma }_{3+}^2=d\alpha ^2+\mathrm{cos}^2(\alpha )d\beta ^2+\mathrm{sin}^2(\alpha )d\varphi ^2,$$
(14)
and the three-geodesic distance $`\omega `$ takes the simple form
$$\mathrm{cos}(\omega )=\mathrm{cos}(\alpha )\mathrm{cos}(\alpha ^{})\mathrm{cos}(\beta \beta ^{})+\mathrm{sin}(\alpha )\mathrm{sin}(\alpha ^{})\mathrm{cos}(\varphi \varphi ^{}).$$
(15)
### 2.2 Taylor series expansion
The geodesic distance, or world function is known only for a few special cases, Roberts (1993) . For $`N=1`$, and for times close to an observers time $`t_0`$ the scale factor can be expanded as a Taylor series
$`R`$ $`=`$ $`R_0\left[1+H_0\delta t{\displaystyle \frac{1}{2}}q_0H_0^2\delta t^2+{\displaystyle \frac{1}{6}}j_0H_0^3\delta t^3+O(H_0\delta t)^4\right],`$
$`\mathrm{where}`$
$`\delta t`$ $`=`$ $`tt_0,\dot{R}_tR,`$
$`H`$ $``$ $`{\displaystyle \frac{\dot{R}}{R}},q{\displaystyle \frac{R\ddot{R}}{\dot{R}^2}},j{\displaystyle \frac{\dot{\ddot{R}}R^2}{\dot{R}^3}},`$ (16)
and the the subscripted โ0โ as in $`H_0`$ refers to the value of the object measured by a observer at $`t=t_0`$. For arbitrary lapse and $`c`$ explicit these become
$`H`$ $``$ $`{\displaystyle \frac{\dot{R}}{cNR}},q{\displaystyle \frac{1}{\dot{R}^2}}\left(\ddot{R}{\displaystyle \frac{\dot{N}\dot{R}}{N}}\right),`$
$`j`$ $``$ $`{\displaystyle \frac{R^2}{\dot{R}^3}}\left(\dot{\ddot{R}}3{\displaystyle \frac{\dot{R}\dot{N}}{N}}{\displaystyle \frac{\dot{R}\dot{N}}{N}}+3{\displaystyle \frac{\dot{R}\dot{N}^2}{N^2}}\right).`$ (17)
### 2.3 the field equations
For the line element 9 the Christoffel symbols are
$`\{_{it}^t\}`$ $`=`$ $`\{_{tt}^i\}=0,`$
$`\{_{tt}^t\}`$ $`=`$ $`{\displaystyle \frac{\dot{N}}{N}},\{_{tj}^i\}=\delta _j^i{\displaystyle \frac{\dot{R}}{R}},\{_{ij}^t\}=g_{ij}{\displaystyle \frac{\dot{R}}{c^2N^2R}},`$
$`\{_{\beta \beta }^\alpha \}`$ $`=`$ $`\mathrm{cos}(\alpha )\mathrm{sin}(\alpha ),\{_{\varphi \varphi }^\alpha \}=\mathrm{cos}(\alpha )\mathrm{sin}(\alpha ),`$
$`\{_{\beta \alpha }^\beta \}`$ $`=`$ $`\mathrm{tan}(\alpha ),\{_{\varphi \alpha }^\varphi \}=\mathrm{cot}(\alpha ).`$ (18)
The Riemann tensor is
$`R_{ijkl}^{(3)}`$ $`=`$ $`g_{ik}^{(3)}g_{lj}^{(3)}g_{li}^{(3)}g_{kj}^{(3)},`$
$`R_{ijkl}`$ $`=`$ $`\left({\displaystyle \frac{k}{R^2}}+{\displaystyle \frac{\dot{R}^2}{c^2N^2R^2}}\right)(g_{ik}g_{lj}g_{li}g_{kj}),`$
$`R_{titj}`$ $`=`$ $`\left({\displaystyle \frac{\ddot{R}}{R}}+{\displaystyle \frac{\dot{N}\dot{R}}{NR}}\right)g_{ij}.`$ (19)
using Einsteinโs field equations the density $`\mu `$ and the pressure $`p`$ for a perfect fluid are given by the Friedman equation
$$\frac{8\pi G}{c^4}\mu =3\frac{k}{R^2}+3\frac{\dot{R}^2}{c^2N^2R^2},$$
(20)
and the pressure equation
$$\frac{8\pi G}{c^4}p=\frac{k}{R^2}\frac{1}{c^2N^2}\left(2\frac{\ddot{R}}{R}+\frac{\dot{R}^2}{R^2}2\frac{\dot{N}\dot{R}}{NR}\right).$$
(21)
The conservation equation is
$$\frac{d}{dR}(\mu R^3)=3pR^2.$$
(22)
compare Weinberg equation 15.1.21. For the equation of state
$$p=(\gamma 1)\mu ,$$
(23)
$`(\gamma 1)`$ times the Friedman equation 20 minus the pressure equation 21 is
$$0=\frac{3\gamma 2}{R^2}\left(k+\frac{\dot{R}^2}{c^2N^2}\right)+\frac{2}{c^2N^2}\left(\frac{\ddot{R}}{R}\frac{\dot{N}\dot{R}}{NR}\right),$$
(24)
evaluating this at $`t=t_0`$ and using 17
$$0=k(3\gamma 2)+H_0^2R_0^2(3\gamma 22q_o),$$
(25)
which for given $`\gamma `$ and $`q`$ determines the sign of $`k`$ and hence whether the Robertson-Walker line-element is open or closed, in the particular case $`\gamma =1`$ and $`N=1`$, 25 reduces to Weinberg equation 15.2.5. For the equation of state 23 the conservation equation 22 is
$$\gamma \mu ^{11/\gamma }\frac{d}{dR}\left(\mu ^{\frac{1}{\gamma }}R^3\right)=\mu _{,R}R^3+3\gamma \mu R^2=0.$$
(26)
Integrating
$$\mu ^{\frac{1}{\gamma }}R^3=a,$$
(27)
where $`a`$ is a constant. Thus $`\mu `$ is proportional to $`R^{3\gamma }`$ so that
$$\frac{\mu }{\mu _0}=\left(\frac{R_0}{R}\right)^{3\gamma },$$
(28)
$`\mu _0`$ is also given by the Friedman equation 20 at $`t_0`$, combining with 28 this gives
$$\mu =\frac{3c^4}{8\pi G}\left(H_0^2+\frac{k}{R_0^2}\right)\left(\frac{R_0}{R}\right)^{3\gamma }.$$
(29)
where a factor of $`c`$ is included in our definition of $`H`$ 17. Substituting 29 into 20 the Friedman equation becomes
$$\left(\frac{\dot{R}}{cN}\right)^2+k=\left(\frac{R}{R_0}\right)^{23\gamma }\left(k+H_0^2R_0^2\right).$$
(30)
### 2.4 the general solutions for $`\gamma `$-equation of state
For $`3\gamma =2`$ the solutions of the Friedman equation is the generalized Milne universe
$$N=1,R=R_0(1+H_0t),$$
(31)
for which the scale factor $`R`$ is just given by 16 up to first order. For $`3\gamma 2`$, define
$$R_\rho R_0\left(H_0^2R_0^2+k\right)^{\frac{1}{3\gamma 2}},$$
(32)
then equation 30 fixes the constant in the solution of Vajk (1967)
$`k`$ $`=`$ $`1,cN=R=R_\rho \left[\mathrm{sin}({\displaystyle \frac{3\gamma 2}{2}}\eta )\right]^{\frac{2}{3\gamma 2}},`$ (33)
$`k`$ $`=`$ $`0,N=1,R=R_0\left({\displaystyle \frac{3\gamma cH_0t}{2}}\right)^{\frac{2}{3\gamma }},`$ (34)
$`k`$ $`=`$ $`1,cN=R=R_\rho \left[\mathrm{sinh}({\displaystyle \frac{3\gamma 2}{2}}\eta )\right]^{\frac{2}{3\gamma 2}},`$ (35)
where $`\eta `$ is the time coordinate used when $`N=R`$. These solutions can be expanded around the origin $`t=0`$ of the time coordinate
$$R=R_\rho \left(\frac{3\gamma 2}{2}\right)^{\frac{2}{3\gamma 2}}\left[1\frac{k}{6}\left(\frac{3\gamma 2}{2}\right)\eta \right]^2+O\left(\frac{3\gamma 2}{2}\eta \right)^4.$$
(36)
Using proper time, for $`\gamma =\frac{4}{3}`$
$$R=R_0\left[1+kH_0^2R_0^2\frac{kc^2t^2}{R_0^2}\right]^{\frac{1}{2}},$$
(37)
and for $`\gamma =1`$ there is the expansion round the origin
$$R=R_\rho \left[1\frac{k}{4}\left(\frac{ct}{R_\rho }\right)^2\frac{k}{48}\left(\frac{ct}{R_\rho }\right)^4+O\left(\frac{ct}{R_\rho }\right)^6\right],$$
(38)
and also an expansion in powers of $`t^{\frac{2}{3}}`$. The substitution $`ttt_0`$ merely shifts the origin of the time coordinate, i.e. the singularity of the spacetime moves from $`t=0`$ to $`t=t_0`$. Replacing $`t`$ by
$$t=t_0+\delta t,\delta t=tt_0,$$
(39)
in all of the above gives back the Taylor expansion 16.
### 2.5 the geodesics
The Lagrangian for geodesics is
$$2=g_{ab}\frac{dx^a}{ds}\frac{dx^b}{ds},$$
(40)
for timelike, null, and spacelike geodesics $`2=1,0,+1`$ respectively. A co-moving geodesic is a timelike geodesic with $`\frac{dx^i}{ds}=0`$, the properties and 40 imply that the co-moving velocity vector is always, irrespective of the geometry of the spacetime, given by
$$U^a=\pm \frac{dx^a}{ds}=\pm (\frac{1}{cN},0).$$
(41)
For null geodesics in Robertson-Walker spacetime
$$0=c^2N^2\left(\frac{dt}{d\mathrm{\Omega }}\right)^2+g_{ij}\frac{dx^i}{ds}\frac{dx^j}{ds}.$$
(42)
using the three-sphere distance 15, 42 can be solved to give
$$\mathrm{\Omega }=\frac{cN}{R}๐t+\omega ,$$
(43)
thus the radiation vector is
$$k_a=\mathrm{\Omega },a=(\frac{cN}{R},\omega _i).$$
(44)
The redshift is given by the equations, Ellis (1971)
$$1+z=\frac{ds_0}{ds}=\frac{\lambda _0}{\lambda }=\frac{\nu }{\nu _0}=\frac{U^ak_a}{(U^ak_a)_0},$$
(45)
as in 15 the subscript $`\mathrm{"}0\mathrm{"}`$ refers to the observer. The quantities at the emitter have no subscript. $`\nu `$ and $`\lambda `$ are the frequency and wavelength of the radiation; $`s`$ is the proper time. Equations 41, 44, and 45 give independently of the choice of $`N`$,
$$1+z=\frac{R_0}{R}.$$
(46)
For the Taylor series expansion 15 this becomes
$`z=`$ $``$ $`H_0(tt_0)+\left(1+{\displaystyle \frac{1}{2}}q_0\right)H_0^2(tt_0)^2`$ (47)
$``$ $`\left(1+q+{\displaystyle \frac{1}{6}}j_0\right)H_0^3(tt_0)^3+O\left(H_0(tt_0)\right),`$
the observed values of $`H_0`$ and $`q_0`$ are $`H_0=2.4\pm \mathrm{0.8.10}^{18}\mathrm{sec}.^1`$ and $`q_0=1\pm 1`$, for a recent review see Freedman (1999) . $`z`$ is found to take the discrete values
$$z_n=\frac{nv_I}{c},$$
(48)
where $`n`$ is an integer and $`v_I=72.2\pm 0.2\mathrm{Km}.\mathrm{sec}.^1`$, or perhaps a sixth of this value.
## 3 Remarks on Explanations of Discrete Redshift.
### 3.1 three distinctions categorizing discrete redshift
There are three distinctions that should be made to categorize any theory of discrete redshift. The first is whether the discrete properties enter via the radiation connecting the emitter and observer, or by the motion of either (or both) the emitter and observer. The second is whether the emitter has real discrete motion, or only apparent discrete motion, or neither. The third is whether the effect is due to quantizing the whole system, or part of it, or is not due to quantum mechanics at all.
### 3.2 via connecting radiation
To elaborate on the first distinction, from equation 45 it is apparent that discrete redshift must come from discrete differences in the ratio of the observerโs and emitters proper time $`\frac{ds_{obs}}{ds_{emm}}`$. The underlying structure of spacetime may be so unusual as to forbid the introduction of a time-like four-vector $`U^a`$ or a null tangent vector $`k_a`$; however if this can be done then there are two choices for the origin of discrete properties: either the particles time-like four-vector or the connecting radiation appears to take discrete values. If the metric itself takes discrete values then both of these would presumably occur.
### 3.3 real discrete motion
To elaborate on the second distinction. Real discrete motion means that the emitterโs motion is discrete no matter how it is measured. Apparent discrete motion means that the motion of the emitter merely appears to be discrete to the observer; to put this another way if the observer chooses to observe from a different vantage point the observer might not measure the same discrete motion of the emitter. Real discrete motion has the disadvantage that there would have to be boundaries at the edges of where the emitterโs four-velocity jumps, these boundaries would lead to other effects, such as refraction and reflection on the boundary surface. Such effects have not been observed in association with measurements of discrete redshift. Monet and Deharveny (1977) do not observe what would be expected from real discrete motion, see subsection 1.2. The existence of superclusters and voids, Sanders et al (1991) , however might be an example of an effect of real discrete motion. The quantum perturbations of Grishchuk (1997) , Yamammoto et al (1996) , and Modanese (2000) might give real discrete motion.
### 3.4 whole system quantization
To elaborate on the third distinction. The idea of the whole system being treated quantum mechanically is essentially that of quantum cosmology. If part of the system is quantized it could be either the connecting radiation or the emitting matter. Of course there is the possibility that discrete redshift may not have a quantum origin at all, for example, it might be of fluid dynamical origin.
### 3.5 quantum emitter
In this paper the idea that the discrete properties originate in the quantum treatment of the emitting matter will be pursued. Before proceeding with this some remarks will be made on the possibility of discrete redshift originating in the properties of the connecting radiation from quantum cosmology. There are several possibilities which might give discrete connecting radiation, here two will be mentioned. The first is that the connecting radiation could undergo a scattering process which makes it discrete, the scattering process would have to occur in a wide variety of circumstances in order to explain the many scales and wavelengths over which discrete redshift occurs. The second is that solutions to Maxwellโs equations in Robertson-Walker spacetime involve discrete properties from spherical harmonics. This will be discussed further in Section 5.
### 3.6 quantum cosmology
Quantum cosmology consists of a large amount of theory in which the whole Universe is considered quantum mechanically (see Tipler (1986) and Kiefer (1999) for reviews); and it is the obvious framework in which to start looking for an explanation of discrete redshift. There might be a quantum analog of the Friedman equation 20 which possess solutions in which the metric is discrete. Developing quantum cosmology as it is found in the literature produced nothing along these lines, therefore four more simplistic approaches where attempted: as it would be anticipated that even a simple model of a quantum expanding gas would produce a qualitative result which involved discrete redshift. The dynamics of Robertson-Walker spacetimes are determined by the Friedman equation 20 and the conservation equation 22; the first approach consisted of applying naive operator substitutions to the Friedman equation; by this is meant replacing each symbol in the equation by an operator of the form $`ia_b`$ where $`a`$ is chosen to be the most general combination of $`\mathrm{}`$, $`c`$, $`G`$, $`R`$, $`R_0`$ which provides a dimensionally correct substitution for the symbol, and $`b`$ can be variously considered to be a four-vector index or a time index etcโฆ. This approach failed because it produced inconsistencies when applied to the conservation equation and it was impossible to eliminate $`R`$. A method similar to this which works has been constructed by Rosen (1993) , see also Capozziello et al (1998) . The second approach was to apply nรกive operator substitutions to the Newtonian analog of the Friedman equation (see for example Weinberg (1972) ). It might be anticipated that this would produce a well-defined theory because it involves only a quantum generalization of Newtonian cosmology, however in common with much of Newtonian cosmology $`\dot{R}`$ and $`c`$ occur in places that lead to inconsistencies. The third approach is to note that there are special relativistic gravitational theories, some of which are discussed in North (1965) . It would be hoped that they would lead to equations with $`\dot{R}`$ and $`c`$ in consistent places, however this approach also failed. The fourth approach is that Schrรถdingerโs p.63 (1956) derivation of Robertson-Walker redshift by thermodynamic analogy might be extendable to include quantum mechanics and thence discrete redshift.
### 3.7 constituent quantization
Assuming that the Robertson-Walker spacetime remains correct, it is necessary to identify what constituents of its contents needs to be quantized and by what mechanism. In general relativity the co-moving emitter travels on geodesics for which
$$2=U_aU^a=\frac{1}{m}p_ap^a.$$
(49)
This equation can be nรกively quantized by using the operator substitution
$$p_ai\mathrm{}_a,$$
(50)
giving the Klein-Gordon equation
$$\left(\mathrm{}\frac{m^2c^2}{\mathrm{}^2}\right)\psi =0,$$
(51)
but it is not clear what the interpretation of 51 is in the present context. For example, should the mass in the Klein-Gordon equation be interpreted as the mass of the emitting galaxy or the emitting atom or something in between, and has the single particle theory described by 49 become a many particle theory described by 51? Here what the Klein-Gordon field describes will be discussed later. The Universe will be taken to have closed Robertson-Walker geometry, because it is for closed systems that discrete properties usually occur.
### 3.8 references for the KG equation in RW spacetime
The Klein-Gordon equation in Robertson-Walker spacetimes has been studied for a variety of purposes by: Schrรถdinger (1939) , (1956) , Mรผller (1940) , Lifshitz (1946) , Lifshitz and Khalatnikov (1963) , Ford (1976) , Barrow and Matzner (1980) , and Klainerman and Sarnak (1981) . Maxwellโs equation in Robertson-Walker spacetimes has been studied by Schrรถdinger (1940) , and Mashhoon (1973) . Diracโs equation in Robertson-Walker spacetime has been studied by Schrรถdinger (1938) , (1940) , and Barut and Duru (1987) .
## 4 The Klein-Gordon Equation in Robertson-Walker Spacetime.
### 4.1 spherical harmonics
In Robertson-Walker spacetime the Klein-Gordon equation is
$$0=\frac{1}{NR^3}\left(\frac{R^3\dot{\varphi }}{N}\right)^{}+\frac{c^2}{R^2}K(\varphi )\frac{c^4m^2}{\mathrm{}}\varphi ,$$
(52)
In the coordinates 14, $`K(\varphi )`$ takes the form
$$K(\varphi )=\mathrm{sec}(\alpha )\mathrm{cosec}(\alpha )\left(\mathrm{cos}(\alpha )\mathrm{sin}(\alpha )\varphi _\alpha \right)_\alpha +\mathrm{sec}^2(\alpha )\varphi _{\beta \beta }+\mathrm{cosec}^2(\beta )\varphi _{\alpha \alpha }.$$
(53)
Define
$$Y\mathrm{sin}^{|n|}\alpha )\mathrm{cos}^{|p|}(\alpha )\mathrm{exp}i(|n|\varphi +|p|\beta ),$$
(54)
then
$$\frac{K(\varphi )}{Y}=\left(|n|+|p|\right)\left(|n|+|p|+2\right).$$
(55)
The coordinate ranges $`0<\alpha <\frac{1}{2}\alpha `$ and $`0<\beta ,\varphi <2\pi `$ imply that $`n`$ and $`p`$ are integers, thus defining
$$l|n|+|p|,$$
(56)
gives
$$K(\varphi )=Y_{.|i}^i=l(l+2)Y,$$
(57)
and
$$Y_iY_.^i=l^2\frac{Y^2}{R^2}.$$
(58)
The complex conjugate to equation 57 and 58 also holds, furthermore
$$Y_iY^i.=[n^2(\mathrm{cot}^2(\alpha )+\mathrm{cosec}^2(\alpha ))+p^2(\mathrm{tan}^2(\alpha )+\mathrm{sec}^2(\alpha ))2np]YY^{},$$
(59)
and
$$(YY^{})_{.|i}^i=4\mathrm{sin}^{|2n|}(\alpha )\mathrm{cos}^{|2p|}(\alpha )\left(n^2\mathrm{cot}^2(\alpha )+p^2\mathrm{tan}^2(\alpha )2nppn\right),$$
(60)
also
$`YY_{i|j}`$ $`=`$ $`Y_iY_j,ij,`$ (61)
$`i(YY_\alpha ^{}Y^{}Y_\alpha )`$ $`=`$ $`0,`$
$`i(YY_\beta ^{}Y^{}Y_\beta )`$ $`=`$ $`2pYY^{},`$
$`i(YY_\beta ^{}Y^{}Y_\varphi )`$ $`=`$ $`2nYY^{},`$ (62)
where the dagger $`\mathrm{"}\mathrm{"}`$ denotes complex conjugate. Equations 59, 60, 61, and 62 cause difficulties when considering the stress of a Klein-Gordon field. Defining the dimensionless scalar field
$$\varphi Y=\left(\frac{R}{R_0}\right)^{\frac{3}{2}}\left(\frac{N}{N_0}\right)^{\frac{1}{2}}\psi ,$$
(63)
the Klein-Gordon equation 52 for $`\psi 0`$ becomes
$$0=\frac{\ddot{\psi }}{\psi }+X+c^2l(l+2)\frac{N^2}{R^2}+\frac{c^4m^2N^2}{\mathrm{}},$$
(64)
where
$$X\frac{1}{2}\left(\frac{\ddot{N}}{N}3\frac{\ddot{R}}{R}\right)\frac{3}{4}\left(\frac{\dot{R}}{R}\frac{\dot{N}}{N}\right).$$
(65)
### 4.2 Hillโs equation
For $`R`$ and $`N`$ consisting of trigonometric functions this equation is similar to Hillโs equation, see for example p.406 Whittaker and Watson (1927) . For the Einstein static universe there is the solution
$$\varphi =C_+\mathrm{exp}(i\nu t)+C_{}\mathrm{exp}(i\nu t),\nu ^2=\frac{c^2l(l+2)}{R_0^2}+\frac{m^2c^4}{\mathrm{}^2},$$
(66)
where $`C_+`$ and $`C_{}`$ are constants. For the generalized Milne universe there is the massless solution, Schrรถdinger (1939) ,
$`\varphi `$ $`=`$ $`C_+\tau ^{\frac{1}{2}(1+rt)}+C_{}\tau ^{\frac{1}{2}(1rt)},`$
$`\mathrm{with}\tau `$ $`=`$ $`(2H_0R_0R^2)^1,`$ (67)
$`\mathrm{and}rt^2`$ $`=`$ $`1{\displaystyle \frac{c^2l(l+2)}{H_0^2R_0^2}}.`$
For the closed $`\gamma =\frac{4}{3}`$ spacetime 28 there is the massless solution, Lifshitz (1946) ,
$$\varphi =\mathrm{cosec}(\eta )\left(C_+\mathrm{exp}(+i(l+1)\eta )+C_{}\mathrm{exp}(i(l+1)\eta )\right).$$
(68)
In general 64 is intractable and it is necessary to use approximate WKB solutions to it.
### 4.3 WKB approximation
The WKB approximation is derived as follows, Alvarez (1989) . Assume the differential equation can be put in the form
$$y^{\prime \prime }+\frac{f(x)y}{\mathrm{}^2}=0,$$
(69)
where $`y^{}=\frac{dy}{dx}`$. Let
$$y=\mathrm{exp}(\frac{iz}{\mathrm{}}),$$
(70)
then 69 becomes
$$i\mathrm{}z^{\prime \prime }z^2+f(x)=0,$$
(71)
taking $`\mathrm{}`$ to be small
$$z^{}=\pm \sqrt{f},$$
(72)
which gives the first order approximation
$$y=\mathrm{exp}\left(\frac{i}{\mathrm{}}\sqrt{f}๐x\right).$$
(73)
substituting the derivative of 72, $`z^{\prime \prime }`$, into 71
$$z^2=f\left(1\pm \frac{i\mathrm{}f^{}}{2f^{\frac{3}{2}}}\right),$$
(74)
taking the square root and disregarding terms in $`\mathrm{}^2`$,
$$z^{}=\pm \sqrt{f}+\frac{i\mathrm{}}{4}\frac{f^{}}{f},$$
(75)
where all combinations of sign are possible, choosing the sign in front of the second term to be positive, using equation 61, and integrating
$$y=f^{\frac{1}{4}}\mathrm{exp}\left(\pm \frac{i}{\mathrm{}}\sqrt{f}๐x\right).$$
(76)
The second order approximation is then given by the linear combination
$$y=f^{\frac{1}{4}}\left\{C_+\mathrm{exp}(+\frac{i}{\mathrm{}}\sqrt{f}๐x)+C_{}\mathrm{exp}(\frac{i}{\mathrm{}}\sqrt{f}๐x)\right\}.$$
(77)
Substituting back into 69 this approximation holds if
$$\frac{f}{\mathrm{}^2}>\frac{5}{16}\frac{f^{\prime \prime }}{f^2}.$$
(78)
### 4.4 WKB applied to the KG equation
Applying the WKB approximation to the Klein-Gordon equation 64 implies that $`\varphi `$ is of the form
$$\varphi =A(t)Y(\alpha ,\beta ,\gamma )\left\{C_+\mathrm{exp}(+i\overline{\nu }(t))+C_{}\mathrm{exp}(i\overline{\nu }(t))\right\},$$
(79)
where $`C_+`$ and $`C_{}`$ are dimensionless constants. There are two functions to be determined. The first is the dimensionless frequency $`\nu `$,
$$\nu .t=\overline{\nu }=\left(X+c^2l(l+2)\frac{N^2}{R^2}+\frac{c^4m^2N^2}{\mathrm{}}\right)^{\frac{1}{2}}๐t,$$
(80)
where the constant of integration is a phase factor which is taken to vanish here. When $`t`$ is proper time $`\nu `$ is referred to as the proper frequency. The second is the dimensionless amplitude
$`A`$ $`=`$ $`\left({\displaystyle \frac{R}{R_0}}\right)^{\frac{3}{2}}\left({\displaystyle \frac{N}{N_0}}\right)^{\frac{1}{2}}A_\psi `$ (81)
$`=`$ $`D\left({\displaystyle \frac{R}{R_0}}\right)^{\frac{3}{2}}\left({\displaystyle \frac{N}{N_0}}\right)^{\frac{1}{2}}\left(X+c^2l(l+2){\displaystyle \frac{N^2}{R^2}}+{\displaystyle \frac{c^4m^2N}{\mathrm{}^2}}\right)^{\frac{1}{4}},`$
where $`D`$ is a dimensional constant added in order to keep $`A`$ dimensionless. Both the frequency $`\nu `$ and the amplitude $`A`$ are sensitive to the choice of time coordinate; for example in proper, conformal, and harmonic times respectively
$`N=1`$ $`:`$ $`X={\displaystyle \frac{3}{4}}\left(2{\displaystyle \frac{\ddot{R}}{R}}+{\displaystyle \frac{\dot{R}^2}{R^2}}\right)={\displaystyle \frac{3}{4}}\left({\displaystyle \frac{8\pi Gp}{c^4}}+{\displaystyle \frac{c^2}{R^2}}\right),`$
$`N=R`$ $`:`$ $`X={\displaystyle \frac{R_{,\eta \eta }}{R}},`$
$`N=R^3`$ $`:`$ $`X=0.`$ (82)
The amplitude $`A`$ is computed from straightforward substitution; however the frequency integral 80 usually has to be approximated.
### 4.5 the Taylor series approximation for the frequency
For the Taylor series approximation 16
$$X=\frac{3}{4}\left\{(2q_01)H_0^2+2(j_01)H_0^3(tt_0)+O(tt_0)^2\right\}.$$
(83)
Up to second order in the Taylor series expansion 16 this expression is time dependent, also it depends on the quantity $`(2q_01)`$ from 25 this is just the critical number which determines whether a pressure free $`\gamma =1`$ spacetime is open or closed. Thus any prediction of $`(2q_01)`$ using 83 is just an artifact of the approximations involved rather than a prediction based on $`\gamma =1`$ spacetime, this is the reason that it is necessary to work with exact solutions to Einsteinโs equations rather than use Taylor series approximations. From 83 and 80
$`\nu `$ $`=`$ $`\left\{{\displaystyle \frac{3}{4}}(2q_01)H_0^2+{\displaystyle \frac{c^2l(l+2)}{R_0^2}}+{\displaystyle \frac{c^4m^2}{\mathrm{}^2}}\right\}^{\frac{1}{2}}`$ (84)
$`+`$ $`{\displaystyle \frac{1}{4}}\left\{{\displaystyle \frac{3}{4}}(2q_01)H_0^2+{\displaystyle \frac{c^2l(l+2)}{R_0^2}}+{\displaystyle \frac{c^4m^2}{\mathrm{}^2}}\right\}^{\frac{1}{2}}.\left\{2(j_01)H_0^2{\displaystyle \frac{2c^2l(l+2)}{R_0^2}}\right\}H_0(tt_0)`$
$`+`$ $`O(tt_0)^2.`$
### 4.6 the dimensionless frequency for the $`\gamma `$-solutions
For perfect fluids with $`N=R`$ the solution to Einsteinโs field equations with $`k=1`$ gives
$$X=\frac{3\gamma 2}{2}+\frac{3\gamma 4}{2}\mathrm{cot}^2(\frac{3\gamma 2}{2}\eta ).$$
(85)
For $`m=0`$ the frequency $`\overline{\nu }`$ is
$$\overline{\nu }=\left\{\frac{3\gamma 2}{2}+\frac{3\gamma 4}{2}\mathrm{cot}^2(\frac{3\gamma 2}{2}\eta )+l(l+2)\right\}^{\frac{1}{2}}๐\eta .$$
(86)
Assuming that $`l`$ is large and expanding
$$\overline{\nu }=\sqrt{l(l+2)}\left\{\eta +\frac{1}{2l(l+2)}\left(\eta +\frac{3\gamma 4}{23\gamma }\mathrm{cot}(\frac{3\gamma 2}{2}\eta )\right)\right\}+O(l^4).$$
(87)
For $`m0`$ it is assumed that $`\mathrm{}`$ is small in line with the assumption needed for the WKB approximation. Equivalently, expressions for the wavelength are of the form
$$\frac{1}{\lambda _{total}^2}=\frac{1}{\lambda _{compton}^2}\pm \frac{1}{\lambda _{geometry}^2},$$
(88)
where $`\lambda _{compton}=\frac{\mathrm{}}{mc}`$ is the Compton wavelength. The assumption that $`\mathrm{}`$ is small implies that the total wavelength $`\lambda _{total}`$ is dominated by the Compton term. Usually, as can be seen from 84 the sign in front of the geometric contribution is positive, thus $`\lambda _{total}`$ is shorter than the Compton wavelength. Expanding in $`\mathrm{}`$ the frequency is
$`\overline{\nu }`$ $`=`$ $`{\displaystyle \left\{\frac{3\gamma 2}{2}+\frac{3\gamma 2}{2}\mathrm{cot}^2(\frac{3\gamma 2}{2}\eta )+l(l+2)+\frac{m^2c^2R^2}{\mathrm{}^2}\right\}^{\frac{1}{2}}๐\eta }`$ (89)
$`=`$ $`{\displaystyle \frac{mc}{\mathrm{}}}I_1+{\displaystyle \frac{\mathrm{}}{2mc}}\left\{{\displaystyle \frac{3\gamma 2}{2}}+l(l+2)I_2+{\displaystyle \frac{3\gamma 4}{2}}I_3\right\}+O(\mathrm{}^3),`$
where
$$I_1R๐\eta ,I_2\frac{d\eta }{R},I_3\mathrm{cot}^2(\frac{3\gamma 2}{2}\eta )\frac{d\eta }{R}.$$
(90)
Evaluating these integrals for $`\gamma =\frac{4}{3}`$ and $`\gamma =1`$ gives
$`\overline{\nu }`$ $`=`$ $`{\displaystyle \frac{cmR_0}{\mathrm{}}}\left(H_0^2R_0^2+1\right)^{\frac{1}{2}}\mathrm{cos}(\eta )`$ (91)
$`+`$ $`{\displaystyle \frac{\mathrm{}(l+1)^2}{2cmR_0}}\left(H_0^2R_0^2+1\right)^{\frac{1}{2}}\mathrm{ln}|\mathrm{tan}{\displaystyle \frac{\eta }{2}}|`$
$`+`$ $`O(\mathrm{}^3),`$
and
$`\overline{\nu }`$ $`=`$ $`{\displaystyle \frac{cmR_0}{\mathrm{}}}\left(H_0^2R_0^2+1\right)\left\{{\displaystyle \frac{\eta }{2}}\mathrm{cos}({\displaystyle \frac{\eta }{2}})\mathrm{sin}({\displaystyle \frac{\eta }{2}})\right\}`$ (92)
$`+`$ $`{\displaystyle \frac{\mathrm{}}{2cmR_0}}\left(H_0^2R_0^2+1\right)\mathrm{cot}(\eta )\left\{{\displaystyle \frac{1}{3}}\mathrm{cot}^2({\displaystyle \frac{\eta }{2}}){\displaystyle \frac{1}{2}}(l^2+2l+{\displaystyle \frac{1}{2}})\right\}`$
$`+`$ $`O(\mathrm{}^3),`$
respectively.
### 4.7 further approximation
Using proper time 32 and 82 give for the expansion around the origin
$$X=\frac{3(3\gamma 2)}{4R_\rho ^2}+O(t).$$
(93)
Defining
$$l^2l^2+2l+\frac{3(3\eta 2)}{4},$$
(94)
the proper frequency is
$$\nu =\nu _\rho +O(t),$$
(95)
where
$$\nu _\rho ^2=\frac{c^2l^2}{R_\rho ^2}+\frac{c^4m^2}{\mathrm{}^2}.$$
(96)
For $`\gamma =1`$, using 82 and 38 the integral for the proper frequency is
$$(\overline{\nu }_{,t})^2=\nu _\rho ^2+\frac{c^2l^2}{2R_\rho ^2}\left(\frac{ct}{R_\rho }\right)^2+\frac{11c^2l^2}{48R_\rho ^2}\left(\frac{ct}{R_\rho }\right)^4+O\left(\frac{ct}{R_\rho ^2}\right)^6,$$
(97)
expanding the square root and integrating gives the proper frequency
$$\nu =\nu _\rho \left\{1+\frac{c^2l^2}{12\nu _\rho R_\rho ^2}\left(\frac{ct}{R_\rho }\right)^2+\left(\frac{11}{3}\nu _\rho \frac{c^2l^2}{R_\rho ^2}\right)\frac{c^2l^2}{160\nu _\rho ^2R_\rho ^2}\left(\frac{ct}{R_\rho }\right)^4+O\left(\frac{ct}{R_\rho }\right)^6\right\}.$$
(98)
## 5 Discrete Redshift via the Connecting Radiation.
### 5.1 explanation using the frequency of the connecting radiation
The simplest way to produce a theory of discrete redshift is to note that in 45 $`z`$ has an expansion in terms of the frequency $`\nu `$ and that the solutions to the massive Klein-Gordon equation also involve a frequency. From 82 this frequency depends on the choice of time coordinate, however the $`l`$ dependent term is usually larger than the $`X`$ term so that this choice only makes a small difference; because of this it is sufficient to use the proper frequency. Choosing the massless proper frequency 95, the equation for the redshift 45 becomes
$$\nu _0(1+z)=\frac{cl^{}}{R_\rho }+O(t),$$
(99)
which implies that
$$l^{}(1+z)\frac{\nu _0}{H_0}\left(\frac{3\gamma 2}{2q_0+23\gamma }\right)^{\frac{1}{2}}\left(\frac{2q_0}{2q_0+23\gamma }\right)^{\frac{1}{3\gamma 2}},$$
(100)
as $`\frac{\nu _0}{H_0}10^{26}10^{32}`$, $`l`$ must be a very large number, as noticed by Schrรถdinger (1939) .
### 5.2 preservation of proper frequency
Equation 99 depends on $`z,1,\nu _0,c,H_0,q_0`$ and $`\gamma `$. $`z`$ varies and $`\nu _o,H_0`$ and $`q_0`$ are constants by definition, thus at least one of $`l,c`$ or $`\gamma `$ must also vary. $`c`$ and $`\gamma `$ cannot vary enough to explain $`z`$, therefore suggesting that $`l`$ must vary. By the variables separable assumption 63 $`l`$ is independent of time, however here it is taken that $`l`$ varies slowly so that 63 holds in approximation only. In this section $`\nu `$ is assumed to be the frequency of the electromagnetic connecting radiation and the nature of the variation in $`l`$ is taken to be such that $`\nu `$ maintains the measured value of $`z`$; and this property is a particular example of a property here called the preservation of proper frequency. In general at the most fundamental level a co-moving quantum system would be expected to have basic states dependent upon a proper frequency $`\nu `$ proportional to $`\frac{c^2l(l+2)}{R^2}+\frac{c^4m^2}{\mathrm{}^2}`$. The scale factor $`R`$ is time dependent, in order for the proper frequency to be nearly conserved (or perhaps fixed by some other considerations), it is necessary for $`l`$ to vary. This requirement is here call the preservation of proper frequency. This principle might imply that some quantity of matter on microscopic scales depends on $`R`$ and $`l`$. Precisely what this quantity may be is unclear. The value of the fundamental constants may vary with time, Dirac (1937) , and it could be that these depend on $`R`$ and $`l`$. An alternative way of viewing the preservation of proper frequency is by using the Planck equation; because $`\nu `$ remains nearly a constant the energy $`E=\mathrm{}\nu `$ will remain nearly constant. From 47 and 48 the characteristic time interval corresponding to one unit of discrete redshift is
$$t_{char}=\frac{v_I}{cH_0}=3\pm 1\times 10^6\mathrm{years},$$
(101)
thus making direct measurements depending on $`l`$ unlikely. Suppose that the value of discrete redshift 48 corresponds to $`l`$ varying by a factor of one
$$\delta z=\frac{v_I}{c}=z_{i+1}z_i=\frac{\nu _{i+1}\nu _i}{\nu _0}.$$
(102)
Using the value of $`\nu _0`$ from 99
$$\frac{v_I}{c(1+z)}=\left(\frac{1+4/l+(\frac{9\gamma }{4}+\frac{3}{2})/l^2}{1+2/l+(\frac{9\gamma }{4}\frac{3}{2})/l^2}\right)^{\frac{1}{2}}1$$
(103)
expanding for large $`l`$
$$\frac{v_I}{c(1+z)}=\frac{1}{l}+O(l^2).$$
(104)
From 47 this gives $`l10^7`$. As $`l`$ is large $`ll^{}`$ and substituting for $`l^{}`$ from 100
$$\frac{v_I\nu _0\sqrt{3\gamma 2}}{cH_0}=(2q_0)^{\frac{1}{23\gamma }}\left(2q_0+23\gamma \right)^{\frac{3\gamma }{2(3\gamma 2)}}.$$
(105)
This gives a very large $`q_010^{38}10^{50}`$, well outside the observational limits and comparable in size to $`10^{42}`$ the Dirac (1937) dimensionless constant.
### 5.3 refinements
The above model has scope for refinements: for example by using Maxwellโs equations instead of the Klein-Gordon equation, and more importantly choosing that the field $`\varphi `$ depends on both $`t`$ and $`\chi `$ so that the radiation connects observer and emitter; however this is not pursued as the model has serious problems which it is unlikely that these refinements would overcome. The most important of these is that it predicts a value for $`q_0`$ many orders of magnitude larger than allowed for by observation. It does not give consistent values for $`l`$, if $`\nu `$ is taken to be given by the value given by 41, and 44 then $`l1`$, 100 gives $`l>10^{26}`$, and 104 gives $`l10^7`$. It predicts that $`v_I`$ should depend on frequency and this is not observed, $`v_I`$ has the same value using either optical or radio data. The model is not quantum mechanical because the massless Klein-Gordon equation does not depend explicitly on $`\mathrm{}`$. The positive aspects of the model are that it is simple and predicts apparent discrete motion.
## 6 Discrete Redshift via the Massive Klein-Gordon Equation.
### 6.1 associate Klein-Gordon momenta with the comoving velocity
The scalar field solutions of section 4 can be interpreted as being the wave function for an element of quantized matter in Robertson-Walker spacetime. The stress of the scalar field can be calculated by substitution. Associated with this stress is the momentum density
$$P_a=T_{ab}U^b,$$
(106)
where $`U^b`$ is a time-like vector-field tangential to an observers world-line. In the present case $`U^b`$ can be taken to be the co-moving vector 31. Corresponding to the momentum density 106 there is the velocity vector
$$V_a=\frac{4\pi R^3}{3m}P_a.$$
(107)
In the present case this vector has only the component $`V_t`$. This velocity vector can be taken to be almost co-moving
$$V_aU_a,$$
(108)
where $`U_a`$ is the co-moving vector 31. Two immediate consequences of this assumption are that: the redshift 35 constructed using $`V_a`$ is almost the same as that constructed from $`U_a`$, and that the redshift is discrete. The almost co-moving stress assumption can be used to calculate the size of the scale factor by methods similar to those of the previous section. It predicts a scale factor of microscopic size. A variant of this approach is to take the vector
$$W_a=\frac{\varphi _a}{\sqrt{\varphi _c\varphi ^c}},$$
(109)
to be almost co-moving. This variant is similar to the co-moving stress approach, but it also predicts angular terms because of non-vanishing of $`W_i`$. The reason that the moving stress approach gives the wrong result is that it is unrealistic to consider the stress given by a single scalar field. It might be hoped that a statistical ensemble of such wave functions would produce the correct co-moving stress. It is not immediately apparent how how to construct such ensembles. For a review of statistical mechanics in curved spacetime see Ehlers (1971) . It can be assumed that such statistical ensembles give back the solutions of section 1 and that the single Klein-Gordon equation produces a quantum perturbation of this. This approach is pursued in section 7.
### 6.2 other massive Klein-Gordon approaches
There are several other approaches based on solutions to the massive Klein-Gordon equation. In non-relativistic quantum mechanics there is the equation
$$<p_a>=i\mathrm{}๐x^a\varphi ^{}_a\varphi ,$$
(110)
see for example Schiff (1949) equation 7.8. This equation would give non-vanishing $`p_i`$ components, so that the requirement that $`\frac{p_a}{m}`$ is approximately co-moving again would not hold; also it is not clear what the interpretation of 110 is in relativistic quantum mechanics. Another approach is to use the group velocity $`c_{group}`$ of $`\varphi `$, as defined by Schrรถdinger (1939) , to define an effective energy and hence $`p_t`$ component; however the existence of a non-negligible group velocity $`c_{group}`$ would again imply non-negligible $`p_i`$ components. There is an entirely different approach which consists of investigating how the hydrogen atom behaves in a Robertson-Walker universe, see for example Trees (1956) ; the problem with this approach is that it requires the value of $`l`$ to be small and is thus unsuitable for the problem in hand.
## 7 Discrete Redshift via Quantum Perturbation of Classical Solutions.
### 7.1 general weak metric perturbation
The metric is taken to be of the form
$$g_{ab}=\overline{g}_{ab}+h_{ab},g^{ab}=\overline{g}^{ab}h^{ab},$$
(111)
where $`\overline{g}_{ab}`$ is a given background field, in the present case this is the Robertson-Walker metric 9, and $`h_{ab}`$ is a small perturbative term. The connection is
$$\mathrm{\Gamma }_{bc}^a\{_{bc}^a\}+\frac{1}{2}K_{.bc}^a,$$
(112)
where the contorsion is
$$K_{abc}=h_{ba;c}+h_{ca;b}h_{bc;a},$$
(113)
and $`\{_{bc}^a\}`$ is the Christoffel symbol of the background field $`\overline{g}_{ab}`$, $`\mathrm{"};\mathrm{"}`$ is the covariant derivative with respect to the background field $`\overline{g}_{ab}`$. For any connection which is a sum of the Christoffel connection and a contorsion tensor, the Riemann tensor is
$$R_{.bcd}^a=\overline{R}_{.bcd}^a+K_{.[d|b|;c]}^a+K_{.eb}^aS_{.cd}^e+\frac{1}{2}K_{.[c|e|}^aK_{.d]b}^e.$$
(114)
In the present case the torsion $`S_{.bc}^a`$ and the cross terms in $`h_{ab}`$ are taken to vanish, then after using 8 for the commutation of covariant derivatives, the Riemann tensor becomes
$$R_{.bcd}^a=\overline{R}_{.bcd}^a\frac{1}{2}h_{.e}^a\overline{R}_{.bcd}^e\frac{1}{2}h_{be}\overline{R}_{..cd}^{ea}+\frac{1}{2}(h_{d.;bc}^ah_{db;.c}^ah_{c.;bd}^a+h_{cb;.d}^a),$$
(115)
Contracting and again using 8
$$R_{bd}=\overline{R}_{bd}h^{fe}\overline{R}_{ebfd}+\frac{1}{2}h_{be}\overline{R}_{.d}^e+\frac{1}{2}h_{de}\overline{R}_{.b}^e+\frac{1}{2}(h_{d.,cb}^c+h_{b.;cd}^c\mathrm{}h_{db}h_{;bd}),$$
(116)
where $`h=h_{.a}^a`$. Equations 114, 115, and 116 differ from (I4) and (I5) of Lifshitz and Khalatnikov (1963) as they leave out all the terms involving products of $`h_{ab}`$ and $`\overline{R}_{cdef}`$.
### 7.2 perturbing the stress of a perfect fluid
Perturbations of the stress usually (see for example Lifshitz and Khalatnikov (1963) and Sacks and Wolfe (1967) ) are of a perfect fluid which obeys
$$R_{ab}=(\mu +p)U_aU_b+\frac{1}{2}(\mu p)g_{ab},$$
(117)
the first conservation equation
$$\mu _aU^a+(\mu +p)U_{.;a}^a=0,$$
(118)
and the second conservation equation
$$(\mu +p)\dot{U}_a+(g_a^b+U_aU^b)p_{,b}=0,$$
(119)
where
$$\dot{U}_a=U_{a;b}U^a.$$
(120)
The perfect fluid is linearly perturbed thus
$$\mu =\overline{\mu }+\delta \mu ,p=\overline{p}+\delta p,U_a=\overline{U}_a+\delta U_a,g_{ab}=\overline{g}_{ab}+\delta g_{ab}.$$
(121)
Identifying
$$\delta \mu =\varphi _c\varphi ^c+V(\varphi ),\delta p=\varphi _c\varphi ^cV(\varphi ),\delta U_a=\frac{\varphi _a}{\sqrt{\varphi _c\varphi ^c}},\delta g_{ab}=h_{ab},$$
(122)
the Klein-Gordon equation and the scalar field stress are recovered at the second and third orders. The solution to the Klein-Gordon equation given in Section 4 are not compatible with this linearization because the first order perturbation produces cross equations in the perfect fluid and scalar field which are not obeyed. This can be readily verified by investigating the time component of the first order perturbation of the second conservation equation 119 for the Einstein static universe this gives
$$\left(\frac{\varphi _t}{\sqrt{\varphi _c\varphi ^c}}\right)_{,t}=0,$$
(123)
which is incompatible with the solution 66. Usually perturbation theory fixes the values of $`\delta \mu `$ and $`\delta p`$, thus not giving the freedom necessary to replace them with scalar fields.
### 7.3 perturbing the stress of a scalar field
Here perturbation are taken to given by the scalar field stress 79, thus
$$R_{ab}=\overline{R}_{ab}+\varphi _a\varphi _b^{}+\varphi _a^{}\varphi _b+g_{ab}\frac{m^2c^2}{\mathrm{}^2}\varphi \varphi ^{},$$
(124)
The components of the Ricci tensor follow immediately after noting
$`\varphi \varphi ^{}`$ $`=`$ $`A^2YY^{}\left(C_+^2+C_{}^2+2C_+C_{}\mathrm{cos}(2\overline{\nu })\right)`$
$`\varphi _t\varphi _t^{}`$ $`=`$ $`A_t^2YY^{}\left(C_+^2+C_{}^2+2C_+C_{}\mathrm{cos}(2\overline{\nu })\right)`$
$`4AA_tYY^{}\mathrm{sin}(2\overline{\nu })`$
$`+A^2YY^{}\left(C_+^2+C_{}^22C_+C_{}\mathrm{cos}(2\overline{\nu })\right)`$
$`\varphi _i\varphi _t^{}+\varphi _i^{}\varphi _t`$ $`=`$ $`AA_t(YY_i^{}+Y_iY^{})\left(C_+^2+C_{}^2+2C_+C_{}\mathrm{cos}(2\overline{\nu })\right)`$
$`+iA^2(YY_i^{}Y_iY^{})(C_+^2C_{}^2)\overline{\nu }_t`$
$`2(YY_i^{}+Y_iY^{})C_+C_{}\overline{\nu }_t\mathrm{sin}(2\overline{\nu }),`$
$`\varphi _i\varphi _j^{}+\varphi _i^{}\varphi _i`$ $`=`$ $`A^2(Y_iY_j^{}+Y_i^{}Y_j)\left(C_+^2+C_{}^2+2C_+C_{}\mathrm{cos}(2\overline{\nu })\right).`$ (125)
### 7.4 the $`R_{tt}`$ component
In general the equations resulting from 124 are intractable and therefore attention is restricted to the Einstein static universe. The perturbations are taken to be in the harmonic gauge
$$h_{a.;b}^b=\frac{1}{2}h_{,a},$$
(126)
the $`R_{tt}`$ component of 124 is
$`R_{tt}\overline{R}_{tt}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{}h_{tt}`$ (127)
$`=`$ $`YY^{}\{(C_+^2+C_{}^2)(2\nu ^2{\displaystyle \frac{c^2m^2}{\mathrm{}^2}}(c^2h_{tt})`$
$`+2C_+C_{}\mathrm{cos}(2\nu t)(2\nu ^2{\displaystyle \frac{c^2m^2}{\mathrm{}^2}}(c^2h_{tt}))\}.`$
due to the presence of $`YY^{}`$, which obeys 60, this equation appears to be intractable. Now $`YY^{}=\mathrm{sin}^{|2n|}(\alpha )\mathrm{cos}^{|2n|}(\alpha )`$, expanding the trigonometrical functions and taking $`n=0`$ (implying $`p=l`$) gives $`YY^{}`$ is approximately one for small $`\alpha `$. Taking the cross term $`m^2h_{tt}`$ to be negligible and assuming that $`h_{tt}`$ is only a function of $`t`$, 127 reduces to
$`h_{tt,tt}`$ $`=`$ $`2C^2(C_+^2+C_{}^2)\left(2\nu ^2{\displaystyle \frac{c^4m^2}{\mathrm{}^2}}\right)`$ (128)
$`4C^2C_+C_{}\left(2\nu ^2+{\displaystyle \frac{c^4m^2}{\mathrm{}^2}}\right)\mathrm{cos}(2\nu t),`$
which has solution
$`h_{tt}`$ $`=`$ $`C_2(C_+^2+C_{}^2)\left(2\nu ^2{\displaystyle \frac{c^4m^2}{\mathrm{}^2}}\right)t^2`$ (129)
$`+C_+C_{}{\displaystyle \frac{2\nu ^2+c^4m^2\mathrm{}^2}{c^2\nu ^2}}\mathrm{cos}(2\nu t)+B_1+B_2,`$
where $`B_1`$ and $`B_2`$ are constants.
### 7.5 the $`R_{ti}`$ component
The $`R_{ti}`$ component of 124 is
$`R_{ti}\overline{R}_{ti}`$ $`=`$ $`h_{ti}{\displaystyle \frac{1}{2}}\mathrm{}h_{ti}`$ (130)
$`=`$ $`2(YY_i^{}+Y_iY^{})C_+C_{}\nu \mathrm{sin}(2\nu t)`$
$`+i(YY_i^{}Y_iY^{})(C_+^2C_{}^2)\nu `$
$`+{\displaystyle \frac{c^2m^2}{\mathrm{}^2}}h_{ti}YY^{}(C_+^2+C_{}^2+2C_+C_{}\mathrm{cos}(2\nu t)).`$
The $`i=\beta `$ and $`i=\varphi `$ components of the first term vanish, however the $`i=\alpha `$ term does not as $`YY_\alpha ^{}+Y_\alpha Y^{}=2(n\mathrm{cot}(\alpha )p\mathrm{tan}(\alpha ))YY^{}`$; again taking $`n=0`$ and expanding the trigonometrical functions gives that this term vanishes for small $`\alpha `$. By equations 62 the angular part of the second term is non-vanishing and remains so after expanding for small $`\alpha `$, by assuming that $`C_+^2`$ is of the same magnitude as $`C_{}^2`$ this term can be taken to vanish. Similarly to the $`R_{tt}`$ component the $`m^2h_{ti}`$ term can be taken to be negligible. Hence all of the left hand side of 130 can be taken to vanish, thus $`h_{ti}=0`$ is an approximate solution to 130.
### 7.6 the $`R_{ij}`$ component
The $`R_{ij}`$ component of 124 is
$`R_{ij}\overline{R}_{ij}`$ $`=`$ $`g_{ij}h_{.k}^kR_0^2+3h_{ij}{\displaystyle \frac{1}{2}}h_{ij}`$ (131)
$`=`$ $`\left(Y_iY_j^{}+Y_i^{}Y_i^{}+(\overline{g}_{ij}+h_{ij}){\displaystyle \frac{m^2c^2}{\mathrm{}^2}}YY^{}\right)`$
$`\times \left(C_+^2+C_{}^2+2C_+C_{}\mathrm{cos}(2\nu t)\right).`$
Now $`Y_iY_j^{}+Y_i^{}Y_j`$ is non-vanishing for $`i=j=\alpha ,\beta ,\varphi `$ and $`i=\beta `$, $`j=\varphi `$ or $`i=\varphi `$, $`j=\beta `$. Taking $`n=0`$ and expanding there remains just the $`i=j=\varphi `$ component and it is of size $`l^2`$; the spatial axes can be rotated so that
$$Y_iY_j^{}+Y_iY_j^{}=l^2\overline{g}_{ij}^{(3)}.$$
(132)
Similarly to the $`R_{tt}`$ component $`YY^{}`$ is taken to be approximately one and $`m^2h_{ij}`$ is taken to be negligible. Subject to the ansatz
$$h_{ij}\sigma (t)\overline{g}_{ij}^{(3)},$$
(133)
131 reduces to
$$\sigma _{,tt}=2c^2\left(l^2+\frac{m^2c^2R_0^2}{\mathrm{}^2}\right)\left(C_+^2+C_{}^2+2C_+C_{}\mathrm{cos}(2\nu t)\right),$$
(134)
which has solution
$`\sigma `$ $`=`$ $`+c^2(C_+^2+C_{}^2)\left(l^2+{\displaystyle \frac{c^2m^2R_0^2}{\mathrm{}^2}}\right)l^2`$ (135)
$`c^2C_+C_{}{\displaystyle \frac{l^2+c^2m^2R_0^2\mathrm{}^2}{\nu ^2}}\mathrm{cos}(2\nu t)+B_3t+B_4,`$
where $`B_3`$ and $`B_4`$ are constants.
### 7.7 the harmonic gauge condition
The solutions for $`h_{ab}`$ given by 129, 133, and 135 do not obey the harmonic gauge condition 126. This is a result of the approximations made for $`Y`$. From 132 $`h_{,t}`$ will depend on $`l`$ but from 129 $`h_{t.;t}^t`$ will not, thus violating the time component of the harmonic gauge condition. No approximation for $`Y`$ which allows the harmonic gauge condition to be preserved are known. 129, 133, and 135 give the weak field metric perturbations
$$N^2=\overline{N}^2c^2h_{tt},R^2=\overline{R}^2+\sigma .$$
(136)
Note that the equations 127, 130 and 131 are not equivalent to the differential equations that arise if the substitutions
$$N=\overline{N}+ฯต_1,R=\overline{R}+ฯต_2,$$
(137)
are used in 20 and 21, because for example, there are no second derivatives of $`N`$ in 20 and 21.
### 7.8 the change in redshift
In general weak metric perturbations induce a complicated change in the redshift. This has been calculated for the conformally flat ($`k=0`$) case by Sacks and Wolfe (1967) . The present case is much simplified because the metric perturbations are of the form 136, the new values of $`N`$ and $`R`$ can be used for the vectors 41 and 44 to give the redshift of the form 46, thus
$$1+z=\frac{R_0}{R}=\frac{R_0}{\overline{R}}\left(1+\frac{\sigma }{\overline{R}^2}\right)\frac{R_0}{\overline{R}}\left(1\frac{\sigma }{2\overline{R}^2}\right).$$
(138)
Appealing to the principle of the preservation of proper frequency, introduced in section 5, a change in the value of $`l`$ by a factor of one is taken to correspond to one unit of discrete redshift
$$\frac{v_I}{c}=|\delta z|,\delta zz_{l+1}z_l\frac{1+z}{2\overline{R}^2}(\sigma _l\sigma _{l+1}).$$
(139)
This equation allows rough estimates to be made of the size of $`R_0`$; as such an estimate can only be made of the order of magnitude of $`R`$ it can be assumed that $`\frac{1+z}{2\overline{R}^2}R_0^2`$, this still leaves $`C_+,C_{},l,t`$ and $`R`$ of unknown size. The frequency $`\nu `$ is large compared to the time scales involved in 135, as for example the Compton frequency of the electron $`\nu _e=\frac{m_ec^2}{\mathrm{}}10^{21}\mathrm{sec}^1`$; thus the $`t^2`$ term in 135 is larger than the $`\mathrm{cos}`$ term and 139 becomes
$$R_0^2=(C_+^2+C_{}^2)(2l+1)t^2c^3v_I^1.$$
(140)
Using the $`T_{.t}^t`$ component of the scalar fieldโs stress
$$C_+^2+C_{}^2=\frac{8\pi G}{c^2}\left(\frac{2l}{R_0^2}+\frac{c^2m^2}{\mathrm{}^2}\right)^1,$$
(141)
where $`\mu _S`$ is the density of the scalar field, and the $`\mathrm{cos}`$ term is taken to be negligible and $`YY^{}1`$, thus 140 becomes
$$R_0^2=8\pi G\mu _S\left(\frac{2l}{R_0^2}+\frac{c^2m^2}{\mathrm{}^2}\right)(2l+1)t^2\frac{c}{v_I}.$$
(142)
It has been assumed that $`C_+^2+C_{}^2`$ is a constant and that $`\mu _S`$ is $`l`$ dependent, this implies that the substitution for $`C_+^2+C_{}^2`$ takes place after equation 139 has been applied.
### 7.9 some incompatible conditions
Rather than deriving a value of $`R_0`$ from 142, it is shown what values of $`R_0,l,t`$ and $`s`$ are compatible with this equation. First it is proved that the following conditions are incompatible:
i)142 holds,
ii) the $`\mathrm{"}t\mathrm{"}`$ in 142 is less than $`H_0^1`$,
iii)$`R>10^{22}`$ meters (this is a typical distance between galaxies),
iv)$`\mu _s<10^5\mathrm{Kg}\mathrm{m}^3`$ (this is an extremely high density compared with
a typical stellar interior density)
$`\mu _c=10^{26}\mathrm{Kg}\mathrm{m}^3`$ the critical density for a Robertson-Walker universe,
$`\mu =10^5\mathrm{Kg}\mathrm{m}^3`$ a typical photosphere density (see p.163 Allen (1973)),
$`\mu =10^{+3}\mathrm{Kg}\mathrm{m}^3`$ an average stellar density; it would be expected that $`\mu _s<\mu _c`$ for the weak metric approximations used in deriving 142 to work),
v)$`\nu _{geomerty}<<\nu _e`$ (this is necessary if the scalar field $`\varphi `$ is chosen to represent a known field).
Proof: v) implies that the $`\frac{l}{R}`$ term in 142 can be neglected thus
$$R_0^2=\frac{8\pi G}{c^2}\frac{\mathrm{}^2}{m^2c^2}(2l+1)\frac{t^2c^2}{v_I},$$
(143)
using ii)
$$R_0^2<\frac{8\pi G\mu _S\mathrm{}^2(2l+1)}{m^2cv_IH_0^2},$$
(144)
again using v)
$$R_0<\frac{16\pi G\mu _s\mathrm{}}{mv_IH_0^2},$$
(145)
in $`SI`$ units this is $`R_0<10^{18}\mu _S`$, from which conditions ii) and iii) can be seen to be incompatible.
### 7.10 some compatible conditions
The most realistic compatible conditions are:
i) 142 holds,
ii) the $`\mathrm{"}t\mathrm{"}`$ in 142 equals $`H_0^1`$,
iii) $`R_0=10^{28}`$ meters (this implies that $`2q_01=10^4`$),
iv) $`\mu _S=10^{13}\mathrm{Kg}\mathrm{m}^3`$ (this is well above the critical density but below
a typical photosphere density),
v) $`\nu =t_p^1`$, where $`t_p^1=\mathrm{}^{\frac{1}{2}}G^{\frac{1}{2}}c^{\frac{5}{2}}10^{44}\mathrm{sec}.`$ is the Planck time,
(this forces the scalar field $`\varphi `$ to be a hypothetical field rather than a known field,
together with the above value of $`R_0`$ it implies that $`l=10^{63}`$ which is the Dirac
dimensionless constant to the power of $`\frac{3}{2}`$).
Proof: Re-arranging 142
$$\frac{c^2m^2R_0^2}{\mathrm{}^2}=2l(8\pi G\mu _St^2cv_I^11)+8\pi G\mu _St^2cv_I^1,$$
(146)
using ii) for the value of $`t`$ and using $`SI`$ units
$$10^{25}R_0^22l(10^{31}1)+10^{31}.$$
(147)
Now v) implies that $`\nu _{geometry}>>\nu _{compton}`$ giving
$$l10^{35}R_010^6\mu _S^1R_0,$$
(148)
and the values of $`R_0,\mu _S`$ and $`l`$ given above can be shown to obey 148 by substitution.
### 7.11 summary
To summarize some of the deficiencies of the above model. There are at least three technical deficiencies: there is no proof that the dynamical equations 124 are consistent; the Einstein static universe has been assumed in order to solve the perturbation equations, but the background metric is time dependent; and various approximations have been made for the angular terms, in particular the approximations 132 result in the loss of the initially assumed harmonic gauge condition. There are at least three physical deficiencies: the equation for discrete redshift in the form 139 depends on $`z`$ contrary to observation; the most realistic compatible conditions for equation 142 require a hypothetical field with the unusual property of a frequency of the order of the inverse Planck time; and the result requires a scalar field density above the critical density of a Robertson-Walker-Friedman universe.
## 8 Conclusion.
Properties of Robertson-Walker spacetimes can be discrete if they depend on the spherical harmonic integer $`l`$, in particular redshift is discrete if there is a mechanism to connect it to this integer. Density perturbations have been known for a long time to depend on this integer and thus integer dependent redshift could have been predicted before it was observed. The problems with introducing discrete redshift via density perturbations include: density perturbations are irregular whereas the value of discrete redshift is constant irrespective of other conditions, and more importantly density perturbations provide no mechanism which will alter $`l`$. Solutions to the Klein-Gordon equation and other field equations also depend on the spherical harmonic integer $`l`$. The requirement that the energy of these fields is almost conserved implies that, as discussed in section 5, the proper frequency is preserved; this in turn implies that the value of $`l`$ changes in a regular manner proportional to the increase in the scale factor $`R`$. In principle all quantum fields have the Universe as an ultimate boundary condition and are thus presumably $`l`$ dependent. Here it was found that theories using only solutions to the Klein-Gordon equation predict a microscopic value for the scale factor $`R`$. It was suggested that the large scale behaviour of Robertson-Walker spacetime is governed, as it is classically, by the Friedmann equation and that the Klein-Gordon solutions in this background induce weak metric perturbations of the spacetime. It might be that solutions for fields involve other integers, apart from the spherical harmonic integer, and this could also lead to discrete redshift via induced metric perturbations. From 122 it is not clear that this coupling is well-defined as the interactions between the Klein-Gordon field and the background fluid may be non-negligible. To produce a realistic prediction of the size of the scale factor using induced metric perturbations, it was necessary to make some very coarse technical and physical assumptions, including the requirement that the scalar field has a frequency approximately equal to the inverse Planck time, this precludes the scalar field representing a known particle. Any theory based upon metric perturbations will predict what in section 3 is called real discrete motion; this implies that there should be boundary effects where the value of the discrete redshift jumps. It is hoped that using the theory of quantum fields on curved spacetimes will remove, or reduce the bounds on, the free parameters such as $`C_+`$ and $`C_{}`$ in the induced weak metric perturbations, and will give a more rigorously defined theory. <sup>1</sup><sup>1</sup>1The referee suggests:โ- take away โ$`2q_01=10^4`$โ from the Abstract and add in the Conclusions that the value $`q_0=1/2`$ is obtained in the absence of the cosmological constant. Taking it into account, the value of $`q_0`$ would be in agreement with the recent observations on the Ia type supernovae.โ
## 9 Acknowledgement
I would like to thank Tony Fairall and Tim Gebbie for reading and commenting on this paper.. This work was supported in part by the South African Foundation for Research and Development (FRD). |
warning/0002/math0002089.html | ar5iv | text | # 0 Introduction.
## 0 Introduction.
Core models were constructed in the papers , , , and , (see also ), , and . We refer the reader to , , and for less painful introductions into core model theory.
A core model is intended to be an inner model of set theory (that is, a transitive class-sized model of $`\mathrm{๐น๐ฅ๐ข}`$) which meets two requirements:
F$`_1`$ It is close to $`V`$ ($`=`$ the universe of all sets), and
F$`_2`$ it can be analyzed in great detail.
Both requirements should be formulated more precisely, of course. However, as โcore modelโ is no formal concept, we canโt expect a thorough general definition. Letโs try to give some hints.
As for F<sub>1</sub>, a core model, call it $`K`$, should reflect the large cardinal situation of $`V`$ (for example, the large cardinals that exist in $`V`$ should be found in $`K`$ as well up to a certain size), it should satisfy certain forms of covering (for example, $`K`$ should compute successors of singular cardinals correctly), and it should be absolute for set-forcings (i.e., the definition of $`K`$ should determine the same object in any set-forcing extension of $`V`$). As for F<sub>2</sub>, a core model should be a fine structural inner model (a class-sized premouse, technically speaking), which satisfies certain forms of condensation (which in turn typically follow from a certain amount of iterability of $`K`$, and make $`K`$ amenable for combinatorial studies).
For the purposes of this introduction weโll say that $`K`$ exists if there is an inner model satisfying appropriate versions of F<sub>1</sub> and F<sub>2</sub> above. The key strategy for constructing core models has not changed through time. As a matter of fact, in order to have any chance to build $`K`$ satisfying F<sub>1</sub> and F<sub>2</sub> one has to work in a theory
$$\mathrm{๐น๐ฅ๐ข}+\neg \mathrm{\Psi },$$
where $`\neg \mathrm{\Psi }`$ denotes an anti large cardinal assumption (for example, โthere is no inner model with a certain large cardinalโ). The following table lists the achievements of , , , and , , and , and of their forerunner, Gรถdel.
| Author(s) | $`\neg \mathrm{\Psi }`$ |
| --- | --- |
| Gรถdel | $`\neg 0^\mathrm{\#}`$ |
| Dodd + Jensen | $`\neg 0^{}`$ |
| Koepke | $`\neg 0^{long}`$ |
| Jensen | $`\neg 0^{sword}`$ |
| Mitchell | $`\neg o(\kappa )=\kappa ^{++}`$ |
| Jensen | $`\neg 0^{\mathrm{}}`$ |
| Steel | $`\neg M_1^\mathrm{\#}`$ |
Except for the last entry this means that the person(s) listed on the left hand side has (have) shown $`K`$ to exist in the theory $`\mathrm{๐น๐ฅ๐ข}+\neg \mathrm{\Psi }`$. (The definitions of the respective anti large cardinal assertions can be found in the above cited papers.) It must however be mentioned that it was Jensen who proved the Covering Lemma for $`L`$ and who also developed the fine structure theory for $`L`$.
What Steel does in is to prove the existence of $`K`$ in the theory
$$\mathrm{๐น๐ฅ๐ข}+\neg M_1^\mathrm{\#}+\mathrm{there}\mathrm{is}\mathrm{a}\mathrm{measurable}\mathrm{cardinal},$$
where by $`\neg M_1^\mathrm{\#}`$ we mean that there is no sharp for an inner model with a Woodin cardinal. For certain applications of $`K`$ this is a deficit (see the discussion in ยง0 of ): Core model theory is typically applied as follows. If $`\mathrm{\Phi }`$ is any statement whatsoever, then one can try to lead
$$\mathrm{๐น๐ฅ๐ข}+\neg \mathrm{\Psi }+\mathrm{\Phi }$$
to a contradiction by using
$$\mathrm{๐น๐ฅ๐ข}+\neg \mathrm{\Psi }K\mathrm{exists}$$
and showing that $`\mathrm{\Phi }`$ contradicts the existence of $`K`$ โbelow $`\mathrm{\Psi }`$.โ If this is the case, one has arrived at finding a lower bound in terms of the consistency strength of $`\mathrm{\Phi }`$, that is, at proving
$$\mathrm{๐น๐ฅ๐ข}+\mathrm{\Phi }\mathrm{\Psi }.$$
Consequently, from we can (often, not always) only hope to get a theorem of the form
$$\mathrm{๐น๐ฅ๐ข}+\mathrm{\Phi }+\mathrm{there}\mathrm{is}\mathrm{a}\mathrm{measurable}\mathrm{cardinal}$$
$$\mathrm{there}\mathrm{is}\mathrm{a}\mathrm{sharp}\mathrm{for}\mathrm{an}\mathrm{inner}\mathrm{model}\mathrm{with}\mathrm{a}\mathrm{Woodin}\mathrm{cardinal}.$$
(It should be noted that there is an important variation of this use of core model theory. Instead of trying to lead
$$\mathrm{๐น๐ฅ๐ข}+\neg \mathrm{\Psi }+\mathrm{\Phi }$$
to an outright contradiction one can try to show that $`\mathrm{๐น๐ฅ๐ข}+\mathrm{\Phi }`$ implies that such-and-such large cardinals exist in the $`K`$ โbelow $`\mathrm{\Psi }`$.โ)
At the time of writing it is not known how to develop the theory of $`K`$ just assuming โ$`\mathrm{๐น๐ฅ๐ข}+`$ there is no inner model with a Woodin cardinal.โ The present paper solves the puzzle of constructing the core model for many strong cardinals. Specifically, we improve Jensenโs notes by establishing that $`K`$ can be shown to exist in the theory
$`\mathrm{๐น๐ฅ๐ข}+`$ there is no sharp for an inner model with
a proper class of strong cardinals,
or
$$\mathrm{๐น๐ฅ๐ข}+0^{^{}}(\mathrm{zero}\mathrm{hand}\mathrm{grenade})\mathrm{does}\mathrm{not}\mathrm{exist},$$
as we shall say. It has turned out that the means by which can be improved in this direction are far from being straightforward generalizations of the means provided by . We also want to emphasize that the work done here yields applications which could not be obtained before (cf. section 9).
The main new idea here is that if $`0^{^{}}`$ does not exist then normal iteration trees are simple enough (โalmost linearโ) so that $`K^c`$ can (still) be built by just requiring new extenders in its recursive definition to be countably complete (rather than having certificates). This will allow a proof of weak covering for $`K^c`$, and finally a proof of the existence of $`K`$ in the theory $`\mathrm{๐น๐ฅ๐ข}+\neg 0^{^{}}`$. There are serious obstacles to proving the existence of $`K`$ in a theory weaker than $`\mathrm{๐น๐ฅ๐ข}+\neg 0^{^{}}`$.
We believe that the present paper raises (and answers) questions and develops techniques which are interesting from the point of view of inner model theory, and which should be useful also outside the realm of this paper.
The paper is organized as follows. Section 1 discusses fine structure. In particular, we shall present our notation and state what weโll have to assume familiarity with. Section 2 introduces $`0^{^{}}`$, and proves that any normal iteration tree of a premouse which is below $`0^{^{}}`$ is โalmost linear.โ Section 3 builds $`K^c`$, the countably complete core model below $`0^{^{}}`$. $`K^c`$ is a preliminary version of $`K`$; it is not known how to develop the theory of $`K`$ without going through $`K^c`$ and proving weak covering for it. Section 4 proves that $`K^c`$ is maximal in the mouse order, and it shows a โgoodnessโ property of $`K^c`$. Sections 5 and 6 establish technical tools for proving the โweak covering lemma for $`K^c`$โ in section 7: we have to study bicephali which are somewhat more liberal than any bicephalus studied so far, and we have to show a โmaximalityโ property of $`K^c`$. Section 8 finally isolates $`K`$.
Section 9 contains an application, and motivation, for this work. Weโll show that the results of the present paper provide the last step for determining the exact consistency strength of the assumption in the statement of the $`12^{\mathrm{th}}`$ Delfino problem (cf. ). Section 9 also contains a brief description of that problem.
Historical note. The work described in sections 2, 3, 4, 8, and 9 was done in Berkeley in the fall of 1997 (parts of section 8 are due to John Steel). The results in sections 5, 6, and 7 were obtained in Vienna in the fall and winter of 1999/2000.
## 1 Preliminaries.
We assume familiarity to a certain extent with Jensenโs classical fine structure theory (cf. , or \[32, Chapter 1\]), with inner model theory as presented in or , and with core model theory as developed in , , , or . The policy of the present paper is that we have elaborated only on new ideas. In most cases we wonโt give proofs for things which can be found elsewhere.
This paper builds upon the concept of a โ(Friedman-Jensen) premouseโ proposed in . Let us first briefly recall essentials from \[9, ยงยง1 and 4\]. We explicitly warn the reader that this part of the current section of our present paper does not give any complete definitions. The reader may find any missing details in the monograph .
An extender is a partial map $`F:๐ซ(\kappa )๐ซ(\lambda )`$ for some ordinals $`\kappa <\lambda `$ (cf. \[9, ยง1 p.2\]) rather than a system of hyper-measures. (We shall sometimes abuse the notation by writing $`F`$ for the partial map $`๐ซ([\kappa ]^{<\omega })๐ซ([\lambda ]^{<\omega })`$ induced by $`F`$ via soft coding \[cf. \[9, ยง1 p.2\]\]; and weโll often write $`i_F`$ for an ultrapower map $`๐ฉ`$ induced by $`F`$.) Premice will be $`J`$-structures constructed from certain well-behaved extender sequences.
A pre-premouse (cf. \[9, ยง4 p.1\]) is an acceptable $`J`$-structure of the form $`=(J_\alpha [\stackrel{}{E}];,\stackrel{}{E},E_\alpha )`$, where
$$\stackrel{}{E}=\{(\nu ,\xi ,X):\xi <\nu <\alpha \mathrm{and}\xi E_\nu (X)\}$$
codes a sequence of extenders, with the following two properties:
(a) If $`E_\nu \mathrm{}`$ for $`\nu \alpha `$, then $`E_\nu `$ is an extender whose domain is $`๐ซ(\kappa )J_\nu [\stackrel{}{E}]`$ for some $`\kappa <\nu `$, $`(J_\nu [\stackrel{}{E}\nu ];,\stackrel{}{E}\nu )`$ is the ($`\mathrm{\Sigma }_0`$-) ultrapower of $`(J_{\kappa ^+}[\stackrel{}{E}\kappa ^+];,\stackrel{}{E}\kappa ^+)`$ where $`\kappa ^+`$ is calculated in $`J_\nu [\stackrel{}{E}\nu ]`$, and if
$$i_{E_\nu }:(J_\nu [\stackrel{}{E}\nu ];,\stackrel{}{E}\nu )_{E_\nu }N$$
is the ($`\mathrm{\Sigma }_0`$-) ultrapower map then $`\stackrel{}{E}^N\nu =\stackrel{}{E}\nu `$ and $`E_\nu ^N=\mathrm{}`$, and
(b) Proper initial segments of $``$ are sound.
Condition (a) is often referred to as โcoherence.โ The ultrapower formed in (a) is according to the upward extension of embeddings technique using $`E_\nu `$ as a fragment of the $`i_{E_\nu }`$ to be formed. We always suppose the well-founded part of a model to be transitive. (In our formulation of (a), we suppose that $`\nu +1`$ is a subset of the well-founded part.) The concept of โsoundnessโ in (b) is according to Jensenโs fine structure (see below).
For pre-premice $`=(J_\alpha [\stackrel{}{E}];,\stackrel{}{E},E_\alpha )`$ as above and for $`\nu \alpha `$ we adopt the notation of \[20, Definition 5.0.4\] and write $`๐ฅ_\nu ^{}`$ to mean $`(J_\nu [\stackrel{}{E}\nu ];,\stackrel{}{E}\nu ,E_\nu )`$. Sometimes we shall confuse $`๐ฅ_\nu ^{}`$ with its underlying universe, $`J_\nu [\stackrel{}{E}\nu ]`$.
Let $``$ be a pre-premouse as above. Let $`F=E_\nu \mathrm{}`$ be an extender with critical point $`\kappa `$ for some $`\kappa <\nu `$ and $`\nu \mathrm{OR}`$. For any $`\lambda F(\kappa )`$ we define $`F|\lambda `$ by setting $`(X,y)F|\lambda `$ if and only if $`Y((X,Y)F\sigma (y)=Y)`$ where $`\sigma `$ the inverse of the collapse of the $`\mathrm{\Sigma }_1`$ hull of $`\lambda `$ inside $`\mathrm{Ult}(๐ฅ_\nu ^{};F)`$. Let $`C_\nu ^M`$ be the set of all $`\lambda (\kappa ,F(\kappa ))`$ such that $`F|\lambda `$ is its own trivial completion, i.e., such that
$$\alpha <\lambda f{}_{}{}^{\kappa }\kappa ๐ฅ_\nu ^{}i_F(f)(\alpha )<\lambda .$$
(For such $`\lambda `$ weโll have that $`(F|\lambda )(X)=F(X)\lambda `$ for $`X\mathrm{dom}(F)`$.) The pre-premouse $``$ is now called a premouse (cf. \[9, ยง4 p.2\]) provided we always have
$$\lambda C_\nu ^MF|\lambda ๐ฅ_\nu ^{}.$$
This clause is called the initial segment condition (cf. \[10, I\], which gives the correction to the initial segment condition proposed in ).
If $``$ is a premouse with $`C_\nu ^{}\mathrm{}`$ for some $`\nu \mathrm{OR}`$ then $``$ is not โbelow superstrong,โ in that $``$ has then an extender on its sequence witnessing that $`๐ฅ_\nu ^{}`$ โthere is a superstrong cardinal.โ That is, below superstrongs does the initial segment condition collapse to the requirement that we have $`C_\nu ^{}=\mathrm{}`$ for all $`\nu \mathrm{OR}`$.
Inner model theory has to iterate premice. Fine iterations in are based on Jensenโs smooth $`\mathrm{\Sigma }^{}`$-theory. Jensen has developed a machinery for taking fine ultrapowers $`\pi :_F^{}๐ฉ`$ in such a way that $`\pi `$ will be what he calls $`\mathrm{\Sigma }_0^{(n)}`$-elementary for all $`n<\omega `$ with $`\rho _n()>\mathrm{c}.\mathrm{p}.(\pi )`$. As presented in , in order to develop the theory of $`K`$ one would only need $`\mathrm{\Sigma }_0^{(n)}`$-elementarity here for all those $`n<\omega `$ such that $``$ is $`n`$-sound. When following this latter route and defining the concept of a fine ultrapower ร la , one can moreover in fact stick to the more traditional master codes, that is to โcoding $``$ onto $`\rho _n()`$, taking a $`\mathrm{\Sigma }_0`$-ultrapower of the coded structure, and then decodingโ ( p. 40; see also the discussion in the introduction to ). Being pragmatic, it is this latter approach which we shall follow here. It has a couple of advantages: it suffices for our purposes, and it will simplify our iterability proof (cf. 3.3 below) and will make it accessible for people ignorant of Jensenโs $`\mathrm{\Sigma }^{}`$-theory.
Let $``$ be an acceptable $`J`$-structure. We shall write (cf. \[32, Chapter 1\]):
$``$ $`\rho _n()`$ for the $`n^{\mathrm{th}}`$ projectum of $``$,
$``$ $`P_{}^n`$ for the set of good parameters (i.e., for the set of parameters witnessing $`\rho _n()`$ is the $`n^{\mathrm{th}}`$ projectum),
$``$ $`p_{,n}`$ for the $`n^{\mathrm{th}}`$ standard parameter of $``$ (i.e., the least element of $`P_{}^n[\mathrm{OR}]^{<\omega }`$ where โleastโ refers to the canonical well-order of $`[\mathrm{OR}]^{<\omega }`$),
$``$ $`W_{}^\nu `$ for the witness for $`\nu p_{,n}`$,
$``$ $`A_{}^{n,p}`$ for the $`n^{\mathrm{th}}`$ standard code determined by $`p`$,
$``$ $`^{n,p}`$ for the $`n^{\mathrm{th}}`$ reduct determined by $`p`$,
$``$ $`^n`$ for $`^{n,p_{,n}}`$,
$``$ $`h_{}^{n,p}`$ ($`h_{}^n`$) for the canonical $`\mathrm{\Sigma }_1`$ Skolem function of $`^{n,p}`$ (of $`^n`$),
$``$ $`R_{}^n`$ for the set of very good parameters (i.e., for the set of $`p`$ such that $`^{n1,p}`$ is generated by $`h_{}^{n1,p}`$ from $`\rho _n()\{p\}`$), and
$``$ $`_n()`$ for the $`n`$-core of $``$ (i.e., for that $`\overline{}`$ such that $`\overline{}^k`$ is the transitive collapse of what is generated by $`h_{}^k`$ from $`\rho _{k+1}()\{p_{^k,1}\}`$ for all $`k<n`$).
By definition, $``$ is $`n`$-sound if and only if $`P_{}^n=R_{}^n`$ (we thereby follow and emphasize that this is in contrast to \[20, Defition 2.8.3\]). As a matter of fact, $``$ is $`n`$-sound if and only if $`p_{,k}R_{}^k`$ for all $`kn`$ if and only if $`=_n()`$. Moreover, if $``$ is $`n`$-sound then $`p_{,k}=p_{,n}k`$ for all $`kn`$.
Let $``$ and $`๐ฉ`$ be acceptable $`J`$-structures, and let $`n\omega `$. We call $`\pi :๐ฉ`$ an $`n`$-embedding (cf. \[20, Definition 2.8.4\]) if and only if
$``$ $``$ and $`๐ฉ`$ are both $`n`$-sound,
$``$ $`\pi (p_{,n})=p_{๐ฉ,n}`$,
$``$ $`\pi (\rho _k())=\rho _k(๐ฉ)`$ for all $`k<n`$ (by convention, $`\pi (\mathrm{OR})=๐ฉ\mathrm{OR}`$),
$``$ $`\pi `$$`\rho _n()`$ is cofinal in $`\rho _n(๐ฉ)`$, and
$``$ $`\pi ^k:^k_{\mathrm{\Sigma }_1}๐ฉ^k`$ for all $`kn`$.
The significance of $`n`$-embeddings is that they are generated by taking $`n`$-ultra-powers. Let $``$ be an $`n`$-sound premouse where $`n\omega `$. Suppose that $`F`$ is an extender which measures all the subsets of its critical point which are in $`^n`$, and let
$$i_F:^n_F\overline{๐ฉ}$$
be the $`\mathrm{\Sigma }_0`$-ultrapower map. There will be at most one transitive structure $`๐ฉ`$ together with a $`\mathrm{\Sigma }_1`$-elementary embedding $`\pi :๐ฉ`$ such that $`\pi (p_{,n})R_{}^n`$, $`\overline{๐ฉ}=๐ฉ^{n,\pi (p_{,n})}`$, and $`\pi ^n=i_F`$. This follows from the upward extension of embeddings lemma (cf. \[32, Chapter 1\]). By forming a term model, one can always produce a (not necessarily well-founded) unique (up to isomorphism) candidate for such an $`๐ฉ`$. We shall denote this situation by
$$\pi :_F^n๐ฉ=\mathrm{Ult}_n(;F),$$
and call it the $`n`$-ultrapower of $``$ by $`F`$. As a matter of fact, if $`F`$ is โclose toโ $``$ (cf. \[20, Definition 4.4.1\]) then weโll have that $`\rho _{n+1}()=\rho _{n+1}(๐ฉ)`$.
Now if $`\pi :_F^n๐ฉ`$ and $`๐ฉ`$ is transitive then $`\pi `$ is an $`n`$-embedding if and only if $`\pi (p_{,n})=p_{๐ฉ,n}`$ and $`\pi (\rho _k())=\rho _k(๐ฉ)`$ for all $`k<n`$. This in turn will follow from a certain solidity and universality of $`p_{,n}`$. We call $``$ $`n`$-solid just in case that for all $`\nu p_{,n}`$ we have that $`W_{}^\nu `$. We call $``$ $`n`$-universal if and only if
$$๐ซ(\rho _n())=\{h_{}^n(i,(\stackrel{}{\alpha },p_{,n}))\rho _n():i<\omega \stackrel{}{\alpha }\rho _n()\}.$$
We have that $`\pi (p_{,n})=p_{๐ฉ,n}`$ and $`\pi (\rho _k())=\rho _k(๐ฉ)`$ for all $`k<n`$ hold true in the above setting if $``$ is $`n`$-solid and $`n`$-universal. One of the key theorems of fine structure theory (cf. \[27, ยง8\], or \[9, ยง7\]) will tell us that if $``$ is $`n`$-sound and $`n`$-iterable then $``$ is $`(n+1)`$-solid and $`(n+1)`$-universal. Here, โ$`n`$-iterabilityโ has to be explained.
Before turning to that it remains to introduce a weakened version of $`n`$-embeddings which will come up in the proofs of 3.2 and 5.10. Let $``$ and $`๐ฉ`$ be acceptable $`J`$-structures, and let $`n\omega `$. We call $`\pi :๐ฉ`$ a weak $`n`$-embedding (cf. \[20, p. 52 ff.\]) if and only if
$``$ $``$ and $`๐ฉ`$ are both $`n`$-sound,
$``$ $`\pi (p_{,n})=p_{๐ฉ,n}`$,
$``$ $`\pi (\rho _k())=\rho _k(๐ฉ)`$ for all $`k<n`$ (by convention, $`\pi (\mathrm{OR})=๐ฉ\mathrm{OR}`$),
$``$ $`\mathrm{sup}(\pi `$$`\rho _n())\rho _n(๐ฉ)`$,
$``$ $`\pi ^k:^k_{\mathrm{\Sigma }_1}๐ฉ^k`$ for all $`k<n`$, and
$``$ there is a set $`X`$ with $`\{p_{n,},\rho _n()\}X`$, $`X\rho _n()`$ is cofinal in $`\rho _n()`$, and $`\pi ^n:^n๐ฉ^n`$ is $`\mathrm{\Sigma }_1`$-elementary on parameters from $`X`$.
There is a generalization of taking an ultrapower by an extender, namely, taking an ultrapower by a โlong extenderโ (cf. for example \[19, ยง2.5\]). If $``$ is a premouse, if $`\nu `$ is a regular cardinal of $``$, and if $`\pi :๐ฅ_\nu ^{}๐ฉ`$ is an embedding then we shall write
$$\mathrm{Ult}_n(;\pi )$$
for the $`n`$-ultrapower of $``$ by the long extender derived from $`\pi `$ (see \[19, ยง2.5\]), also called the lift up of $`๐ฉ`$ by $`\pi `$. Weโll use this technique in sections 6 and 7.
Let us now discuss iterations of premice. Premice are iterated by forming iteration trees. When $``$ is a premouse then a typical iteration tree on $``$ will have the form
$$๐ฏ=((_\alpha ^๐ฏ,\pi _{\alpha \beta }^๐ฏ:\alpha T\beta <\theta ),(E_\alpha ^๐ฏ:\alpha +1<\theta ),T)$$
where $`T`$ is the tree structure, $`_\alpha ^๐ฏ`$ are the models, $`\pi _{\alpha \beta }^๐ฏ`$ are the embeddings, and $`E_\alpha ^๐ฏ`$ are the extenders used. When dealing with iteration trees we shall use the notation from . We call an iteration tree $`๐ฏ`$ as above normal (cf. \[9, ยง4 p. 4\]) if and only if the following three requirements are met.
$``$ The index of $`E_\beta ^๐ฏ`$ is (strictly) greater than the index of $`E_\alpha ^๐ฏ`$ for all $`\alpha <\beta <\theta `$,
$``$ $`T`$-pred$`(\alpha +1)=`$ the least $`\beta \alpha `$ such that
$$\mathrm{c}.\mathrm{p}.(E_\alpha ^๐ฏ)<E_\beta ^๐ฏ(\mathrm{c}.\mathrm{p}.(E_\beta ^๐ฏ)),\mathrm{and}$$
$``$ setting $`\alpha ^{}=T`$-pred$`(\alpha +1)`$, $`\mathrm{dom}(\pi _{\alpha ^{}\alpha +1}^๐ฏ)=๐ฅ_\eta ^{_\alpha ^{}^๐ฏ}`$ where $`\eta `$ is largest such that $`E_\alpha ^๐ฏ`$ measures all the subsets of its critical point which are in $`๐ฅ_\eta ^{_\alpha ^{}^๐ฏ}`$, for all $`\alpha +1<\theta `$.
We shall build our fine iteration trees by forming $`n`$-ultrapowers. That is, if
$$๐ฏ=((_\alpha ^๐ฏ,\pi _{\alpha \beta }^๐ฏ:\alpha T\beta <\theta ),(E_\alpha ^๐ฏ:\alpha +1<\theta ),T)$$
is an iteration tree weโll have that for all $`\alpha +1<\theta `$, setting $`\alpha ^{}=T`$-pred$`(\alpha +1)`$,
$$\pi _{\alpha ^{}\alpha }^๐ฏ:\mathrm{dom}(\pi _{\alpha ^{}\alpha }^๐ฏ)_{E_\alpha ^๐ฏ}^n_{\alpha +1}^๐ฏ$$
for some $`n\omega `$. In this case, we shall denote $`n`$ by $`\mathrm{deg}^๐ฏ(\alpha +1)`$. Notice that if $`\mathrm{deg}^๐ฏ(\alpha +1)=n`$ then $`\mathrm{dom}(\pi _{\alpha ^{}\alpha }^๐ฏ)`$ has to be $`n`$-sound, and $`_{\alpha +1}^๐ฏ`$ will be $`n`$-sound by design. We call an iteration tree $`๐ฏ`$ on a premouse $``$ $`n`$-bounded if $`\mathrm{deg}^๐ฏ(\alpha +1)n`$ for all $`\alpha +1<\theta `$ such that $`๐^๐ฏ(0,\alpha +1]_T=\mathrm{}`$, i.e., there is no drop along $`[0,\alpha +1]_T`$.
We call $`๐ฏ`$ a putative iteration tree if $`๐ฏ`$ is an iteration tree except for the fact that if $`๐ฏ`$ has a last model, $`_{\mathrm{}}^๐ฏ`$, then $`_{\mathrm{}}^๐ฏ`$ does not have to be well-founded.
Let $``$ be a premouse. For our purposes, we may think of $``$ as being normally $`n`$-iterable if all putative normal $`n`$-bounded iteration trees on $``$ of successor length have the property that $`_{\mathrm{}}^๐ฏ`$, the last model of $`๐ฏ`$, is transitive, and $`๐^๐ฏ(0,\mathrm{}]_T`$ is finite.<sup>1</sup><sup>1</sup>1The informed reader might miss an assertion about the existence of well-founded branches for trees of limit length; the reason for not including such a thing is that by virtue of 2.2 weโll only have to deal with trees with exactly one cofinal branch. In general weโd have to say that there exists an iteration strategy for $``$, i.e., a certain partial function picking cofinal well-founded branches through certain trees on $``$ of limit length. In a comparison process weโll typically use maximal trees. For $`n\omega `$ we call $`๐ฏ`$ $`n`$-maximal (cf. \[27, Definition 6.1.2\]) if the following two requirements are met.
$``$ if $`๐^๐ฏ(0,\alpha +1]_T=\mathrm{}`$ then $`\mathrm{deg}^๐ฏ(\alpha +1)`$ is the largest $`kn`$ such that $`\mathrm{c}.\mathrm{p}.(E_\alpha ^๐ฏ)<\rho _k(\mathrm{dom}(\pi _{\alpha ^{}\alpha +1}^๐ฏ))`$ where $`\alpha ^{}=T`$-pred$`(\alpha +1)`$, and
$``$ if $`๐^๐ฏ(0,\alpha +1]_T\mathrm{}`$ then $`\mathrm{deg}^๐ฏ(\alpha +1)`$ is the largest $`k\omega `$ such that $`\mathrm{c}.\mathrm{p}.(E_\alpha ^๐ฏ)<\rho _k(\mathrm{dom}(\pi _{\alpha ^{}\alpha +1}^๐ฏ))`$ where $`\alpha ^{}=T`$-pred$`(\alpha +1)`$.
That $``$ is $`n`$-iterable (see \[20, Definition 5.1.4\]) will say that every stack
$$๐ฏ_0{}_{}{}^{}๐ฏ_{1}^{}{}_{}{}^{}๐ฏ_{2}^{}{}_{}{}^{}\mathrm{}$$
of normal iteration trees is well-behaved, where $`๐ฏ_{i+1}`$ is a tree on an initial segment of the last model of $`๐ฏ_i`$. More precisely, $``$ is $`n`$-iterable provided the following holds true: if $`0<k\omega `$ and $`(๐ฏ_i:i<k)`$ is a sequence of normal iteration trees such that
$``$ $`_0^{๐ฏ_0}=`$,
$``$ $`๐ฏ_0`$ is $`n`$-bounded,
$``$ $`๐ฏ_i`$ has successor length $`\theta _i`$ for all $`i<k`$,
$``$ $`\gamma _i_0^{๐ฏ_{i+1}}=๐ฅ_{\gamma _i}^{_{\theta _i}^{๐ฏ_i}}`$ for all $`i+1<k`$, and
$``$ $`๐ฏ_{i+1}`$ is $`n(i)`$-bounded, where $`n(i)`$ is maximal such that $`_0^{๐ฏ_{i+1}}`$ is $`n(i)`$-sound and $`ji[\gamma _j=_{\theta _j}^{๐ฏ_j}\mathrm{OR}๐^{๐ฏ_j}(0,\theta _j]_{T_j}=\mathrm{}]n(i)n`$, for all $`i+1<k`$,
then either $`k<\omega `$ and $`_{\theta _{k1}}^{๐ฏ_{k1}}`$ is normally $`\overline{n}`$-iterable, where $`\overline{n}`$ is maximal such that $`_0^{๐ฏ_{k1}}`$ is $`\overline{n}`$-sound and $`j<k1[\gamma _j=_{\theta _j}^{๐ฏ_j}\mathrm{OR}๐^{๐ฏ_j}(0,\theta _j]=\mathrm{}]\overline{n}n`$, or else $`k=\omega `$ and $`๐^{๐ฏ_i}(0,\theta _i]_{T_i}=\mathrm{}`$ for all but finitely many $`i<\omega `$ and the direct limit of the $`_0^{๐ฏ_i}`$โs together with the obvious maps (for sufficiently large $`i`$) is well-founded.
The following lemma, which we might call the โstrong initial segment conditionโ for Friedman-Jensen mice, is an immediate consequence of the condensation lemma \[9, ยง8 Lemma 4\]. It is just a slight strengthening of \[9, ยง8 Corollary 4.2\].
###### Lemma 1.1
Let $`=(J_\alpha [\stackrel{}{E}];,\stackrel{}{E},E_\alpha )`$ be a $`0`$-iterable premouse, where $`E_\alpha \mathrm{}`$. Suppose that for no $`\mu \mathrm{OR}`$ do we have $`๐ฅ_\mu ^{}\mathrm{๐น๐ฅ๐ข}+`$ โthere is a Woodin cardinal.โ Set $`\kappa =\mathrm{c}.\mathrm{p}.(E_\alpha )`$, and let $`\xi (\kappa ,\rho _1())`$. Then one of (a) or (b) below holds:
(a) There is some $`\eta <\alpha `$ such that $`E_\alpha |\xi =E_\eta ^{}`$.
(b) $`\xi `$ is not a cardinal in $``$, and there are $`\mu <\eta <\mathrm{OR}`$ and $`k<\omega `$ such that $`E_\alpha |\xi `$ is the top extender of $`\mathrm{Ult}_k(๐ฅ_\eta ^{};E_\mu ^{})`$ and $`(\eta ,k)`$ is $`<_{lex}`$-least with $`\eta \mu `$ and $`\rho _{k+1}(๐ฅ_\eta ^{})\mathrm{c}.\mathrm{p}.(E_\mu ^{})`$.
Moreover, in any case there is some $`E_{\stackrel{~}{\nu }}^{}`$ with $`\xi <\stackrel{~}{\nu }<\xi ^+`$ and $`\mathrm{c}.\mathrm{p}.(E_{\stackrel{~}{\nu }}^{})=\kappa `$.
Proof. As $`E_\alpha |\xi =E_\alpha |(\xi \kappa ^+)`$ for $`\xi >\kappa `$, let us assume without loss of generality that $`\xi \kappa ^+`$. Consider
$$\sigma :\overline{}H^{}(\xi )_{\mathrm{\Sigma }_1},$$
where $`H^{}(\xi )`$ denotes the hull generated from $`\xi `$ by $`h_{}^0`$, and $`\overline{}`$ is transitive. Let $`\nu \overline{}\mathrm{OR}`$ be maximal with $`\sigma \nu =\mathrm{id}`$. We have that $`\rho _1(\overline{})\xi `$, and $`\overline{}`$ is trivially $`1`$-sound above $`\nu \xi `$ (i.e., $`\overline{}`$ is the hull generated from $`\nu \{p_{\overline{},1}\}`$ by $`h_\overline{}^0`$). Because $`\rho _1(\overline{})\xi <\rho _1()`$, by the condensation lemma \[9, ยง8 Lemma 4\] there remain just two possibilities:
(i) $`\overline{}=๐ฅ_\eta ^{}`$ for some $`\eta <\mathrm{OR}`$, or
(ii) $`\overline{}=\mathrm{Ult}_k(๐ฅ_\eta ^{};E_\mu ^{})`$ where
$$\rho _\omega (๐ฅ_\eta ^{})<\nu =(\mathrm{c}.\mathrm{p}.(E_\mu ^{}))^{+๐ฅ_\mu ^{}}=(\mathrm{c}.\mathrm{p}.(E_\mu ^{}))^{+๐ฅ_\eta ^{}}<\mu \eta <\mathrm{OR},$$
$`E_\mu ^{}`$ is generated by its critical point, and $`k<\omega `$ is such that $`\rho _{k+1}(๐ฅ_\eta ^{})<\nu \rho _k(๐ฅ_\eta ^{})`$.
Now let $`F`$ denote the top extender of $`\overline{}`$. It is easy to see that $`F|\xi =E_\alpha |\xi `$. Moreover, all the generators of $`F`$ are below $`\xi `$, so that in fact $`F=F|\xi =E_\alpha |\xi `$.
Let us first assume that (i) holds. Then $`F=E_\eta ^{}`$, so that (a) in the statement of 1.1 holds.
Let us then assume that (ii) holds. Then $`F`$ has to be the top extender of $`\mathrm{Ult}_k(๐ฅ_\eta ^{};E_\mu ^{})`$. It is also easy to see that $`\eta >\mu `$. Set $`\overline{\mu }=\mathrm{c}.\mathrm{p}.(E_\mu ^{})`$. As $`\rho _\omega (๐ฅ_\eta ^{})\overline{\mu }`$ and $`\kappa ^{+\overline{}}=\kappa ^+`$, we must have $`\overline{\mu }>\kappa `$. Hence by (ii), $`๐ฅ_\eta ^{}`$ has to have a top extender with critical point $`\kappa `$, namely $`E_\eta ^{}`$. Notice that $`\eta >\nu \xi `$.
Now suppose that $`\xi `$ is a cardinal in $``$. Weโll then have that $`\rho _\omega (๐ฅ_\eta ^{})\overline{\mu }`$ and $`\eta <\mathrm{OR}`$ imply that $`\xi \overline{\mu }`$. But taking the ultrapower $`\mathrm{Ult}_k(๐ฅ_\eta ^{};E_\mu ^{})`$ is supposed to give $`\overline{}`$ and it adds $`\overline{\mu }`$ as a generator of its top extender, $`F`$. However, all the generators of $`F`$ are below $`\xi `$. Contradiction!
We have shown that if (ii) holds then (b) in the statement of 1.1 holds.
Finally, the careful reader will have observed that we also have established the โmoreoverโ clause in the statement of 1.1: notice that weโll always have that $`\overline{}`$, and that $`\overline{}`$ has the same cardinality as $`\xi `$ inside $``$.
$`\mathrm{}`$ (1.1)
1.1 should be compared with the initial segment condition of \[20, Definition 1.0.4 (5)\] (cf. also \[26, Definition 2.4\]). However, as Friedman-Jensen extenders allow more room between the sup of their generators and their index, one has to add the hypothesis that $`\xi <\rho _1()`$ in 1.1.
###### Definition 1.2
Let $``$ be a premouse, and let $`\kappa <\tau \mathrm{OR}`$. Then $`\kappa `$ is said to be $`<\tau `$-strong in $``$ if for all $`\alpha <\tau `$ there is some extender $`F`$ with $`\mathrm{dom}(F)=๐ซ(\kappa )`$, $`\mathrm{c}.\mathrm{p}.(F)=\kappa `$, and $`๐ฅ_\alpha ^{}\mathrm{Ult}_0(;F)`$. Furthermore, $`\kappa `$ is said to be $`<\tau `$-strong in $``$ as witnessed by $`\stackrel{}{E}^{}`$ if for all $`\alpha <\tau `$ there is some $`E_\beta ^{}\mathrm{}`$ with $`\mathrm{c}.\mathrm{p}.(E_\beta ^{})=\kappa `$ and $`\alpha <\beta <\tau `$.
We may now phrase an immediate corollary to 1.1 as follows.
###### Corollary 1.3
Let $`=(J_\alpha [\stackrel{}{E}];,\stackrel{}{E},E_\alpha )`$ be a $`0`$-iterable premouse, where $`E_\alpha \mathrm{}`$. Suppose that for no $`\mu \mathrm{OR}`$ do we have $`๐ฅ_\mu ^{}\mathrm{๐น๐ฅ๐ข}+`$ โthere is a Woodin cardinal.โ Set $`\kappa =\mathrm{c}.\mathrm{p}.(E_\alpha )`$, and let $`\tau (\kappa ^+,\rho _1()]`$ be a cardinal in $``$. Then $`\kappa `$ is $`<\tau `$-strong in $``$ as witnessed by $`\stackrel{}{E}^{}`$.
## 2 Almost linear iterations and $`0^{^{}}`$.
In this section we show that for premice which do not encompass a measurable limit of strong cardinals (i.e., which are โbelow $`0^{^{}}`$,โ as we shall say) almost linear iterability โ in a sense to be made precise โ suffices for comparison. It appears that T. Dodd already knew that such small mice can be linearily compared, at least the mice of his time (see ).
###### Definition 2.1
An iteration tree $`๐ฏ`$ is called almost linear<sup>2</sup><sup>2</sup>2\[30, Definition 1.10\] introduces a different concept of almost linearity. If a Mitchell-Steel premouse $``$ is below $`0^{\mathrm{}}`$ then every iteration tree on $``$ is almost linear in the sense of \[30, Definition 1.10\]. if the following hold true.
(a) For all $`i+1<\mathrm{lh}(๐ฏ)`$ we have that $`๐ฏ`$-$`\mathrm{pred}(i+1)[0,i]_T`$, and
(b) any $`i<\mathrm{lh}(๐ฏ)`$ only has finitely many immediate $`๐ฏ`$-successors.
We shall not have to worry about finding cofinal branches for almost linear iteration trees of limit length.
###### Lemma 2.2
Let $`๐ฏ`$ be an almost linear iteration tree. Let $`\lambda \mathrm{lh}(๐ฏ)`$ be a limit ordinal (or $`\lambda =\mathrm{OR}`$). Then $`๐ฏ\lambda `$ has a unique cofinal branch $`b`$, and if $`\lambda <\mathrm{lh}(๐ฏ)`$ then $`b=[0,\lambda )_T`$.
Proof. Obviously, the second part follows from the first. We get $`b`$ by the following simple recursion. If $`ib`$, $`i+1<\lambda `$, then the immediate $`๐ฏ`$-successor of $`i`$ in $`b`$ is the maximal (as an ordinal) $`i^{}>i`$ being immediate $`๐ฏ`$-successor of $`i`$. And if $`\overline{\lambda }<\lambda `$ is a limit ordinal and cofinally many $`i<\overline{\lambda }`$ are in $`b`$ then $`\overline{\lambda }`$ is in $`b`$, too. 2.1 is just what is needed to see that this works.
$`\mathrm{}`$ (2.2)
###### Definition 2.3
Let $``$ be a premouse. $``$ is said to be below $`0^{^{}}`$ (pronounced โzero hand-grenadeโ) if there is no $`\kappa `$ which is the critical point of an extender $`E_\nu ^{}\mathrm{}`$ with the property that
$$\{\mu <\kappa :\mu \mathrm{is}<\kappa \mathrm{strong}\mathrm{in}\}\mathrm{is}\mathrm{unbounded}\mathrm{in}\kappa .$$
We admit that $`{}_{}{}^{}`$ doesnโt seem to look like a hand-grenade. The mice of inner model theory donโt resemble the mice in our backyard as well. โHand-grenadeโ is just another math term in the tradition of daggers, swords, and pistols.
We want to emphasize that if the premouse $``$ is below $`0^{^{}}`$ then $``$ is of course below superstrong as well, so that weโll have that $`C_\nu ^{}=\mathrm{}`$ for all $`\nu \mathrm{OR}`$.
The following lemma shows the benefit of life without hand-grenades.
###### Lemma 2.4
Let $``$ be a premouse which is below $`0^{^{}}`$. Then any normal iteration of $``$ is almost linear.
Proof. Let us fix $``$. We first aim to prove:
Claim 1. For all normal $`๐ฏ`$ on $``$ of successor length $`\vartheta +1`$ do we have the following. Let $`i+1(0,\vartheta ]_T`$, and set $`\kappa =\mathrm{c}.\mathrm{p}.(E_i^๐ฏ)`$, and $`\lambda =E_i^๐ฏ(\kappa )`$. Then there is no $`G=E_\nu ^{_\vartheta ^๐ฏ}\mathrm{}`$ with $`\mathrm{c}.\mathrm{p}.(G)[\kappa ,\lambda )`$ and $`\nu \lambda `$.
Proof. Let $`๐ฏ`$ be a normal iteration tree on $``$ of length $`\vartheta +1`$, and let us consider
$$\pi _{i^{}i+1}^๐ฏ:๐ฅ_\eta ^{_i^{}^๐ฏ}_F_{i+1}^๐ฏ,$$
where $`F=E_i^๐ฏ`$, $`i^{}=T`$-$`\mathrm{pred}(i+1)`$, $`๐ฅ_\eta ^{_i^{}^๐ฏ}=\mathrm{dom}(\pi _{i^{}i+1}^๐ฏ)`$, and $`i+1T\vartheta `$. We put $`\kappa =\mathrm{c}.\mathrm{p}.(F)`$, and $`\lambda =F(\kappa )`$. Suppose that there is some $`G=E_\nu ^{_\vartheta ^๐ฏ}\mathrm{}`$ with $`\mathrm{c}.\mathrm{p}.(G)[\kappa ,\lambda )`$ and $`\nu \lambda `$. Let $`\mu =\mathrm{c}.\mathrm{p}.(G)`$.
Subclaim. $`\rho _\omega (๐ฅ_\nu ^{_\vartheta ^๐ฏ})\lambda `$.
Proof. Of course, $`\lambda `$ is a cardinal of all $`_k^๐ฏ`$ with $`ki+1`$. This trivially implies $`\rho _\omega (๐ฅ_\nu ^{_\vartheta ^๐ฏ})\lambda `$ if $`๐^๐ฏ(i+1,\vartheta ]_T\mathrm{}`$ or if $`G`$ is not the top extender of $`_\vartheta ^๐ฏ`$.
Let us suppose that $`๐^๐ฏ(i+1,\vartheta ]_T=\mathrm{}`$ and that $`G`$ is the top extender of $`_\vartheta ^๐ฏ`$. By the normality of $`๐ฏ`$ we will have that $`\pi _{i+1\vartheta }^๐ฏ\lambda =\mathrm{id}`$. But $`\pi _{i+1\vartheta }^๐ฏ:_{i+1}^๐ฏ_\vartheta ^๐ฏ`$ is sufficiently elementary to yield that then $`_{i+1}^๐ฏ`$ must have a top extender with critical point $`\mu =\mathrm{c}.\mathrm{p}.(G)`$. In particular, $`\mu \mathrm{ran}(\pi _{i^{}i+1}^๐ฏ)`$. However, $`\mu [\kappa ,\lambda )`$ clearly implies $`\mu \mathrm{ran}(\pi _{i^{}i+1}^๐ฏ)`$. Contradiction!
$`\mathrm{}`$ (Subclaim)
Now by the Subclaim we get that $`G|\xi ๐ฅ_\nu ^{_\vartheta ^๐ฏ}`$ for all $`\xi (\mu ,\lambda )`$. Thus
$$(G|\xi :\xi (\mu ,\lambda ))$$
witnesses that $`\mu `$ is $`<\lambda `$-strong in $`_\vartheta ^๐ฏ`$. But as $`๐ฅ_\lambda ^{_\vartheta ^๐ฏ}=๐ฅ_\lambda ^{_{i+1}^๐ฏ}`$ and $`\lambda `$ is a cardinal in both of these models, this says that $`\mu `$ is $`<\lambda `$-strong in $`_{i+1}^๐ฏ`$.
Let $`\alpha <\kappa `$ be arbitrary. We have seen that
$$_{i+1}^๐ฏ\mathrm{`}\mathrm{`}\stackrel{~}{\mu }(\alpha ,\lambda )(\stackrel{~}{\mu }\mathrm{is}<\lambda \mathrm{strong}).\mathrm{"}$$
As $`\pi _{i^{}i+1}^๐ฏ`$ is sufficiently elementary, we can deduce that
$$_i^{}^๐ฏ\mathrm{`}\mathrm{`}\stackrel{~}{\mu }(\alpha ,\kappa )(\stackrel{~}{\mu }\mathrm{is}<\kappa \mathrm{strong}).\mathrm{"}$$
But $`๐ฅ_\kappa ^{_i^{}^๐ฏ}=๐ฅ_\kappa ^{_i^๐ฏ}`$ and $`\alpha <\kappa `$ was arbitrary, so that
$$\{\stackrel{~}{\mu }<\kappa :\stackrel{~}{\mu }\mathrm{is}<\kappa \mathrm{strong}\mathrm{in}_i^๐ฏ\}\mathrm{is}\mathrm{unbounded}\mathrm{in}\kappa .$$
Because $`F=E_i^๐ฏ`$ has critical point $`\kappa `$, this shows that $`_i^๐ฏ`$ is not below $`0^{^{}}`$. But this implies that $``$ was not below $`0^{^{}}`$ to begin with. Contradiction!
$`\mathrm{}`$ (Claim 1)
Claim 2. For all normal $`๐ฏ`$ on $``$ of double successor length $`\vartheta +2`$ do we have the following. Let $`i<j+1<\vartheta +1`$ be such that $`j+1(i,\vartheta ]_T`$, and $`i=T`$-pred$`(j+1)=T`$-pred$`(\vartheta +1)`$. Then $`\mathrm{c}.\mathrm{p}.(E_j^๐ฏ)>\mathrm{c}.\mathrm{p}.(E_\vartheta ^๐ฏ)`$.
Proof. This easily follows from Claim 1.
$`\mathrm{}`$ (Claim 2)
Let us finally prove that every normal $`๐ฏ`$ on $``$ is almost linear. Suppose not, and let $`๐ฏ`$ be a normal iteration tree on $``$ of length $`\vartheta `$ such that $`๐ฏ\theta `$ is almost linear for all $`\theta <\vartheta `$, whereas $`๐ฏ`$ is not almost linear.
Case 1. $`\vartheta `$ is a limit ordinal.
Then there is some $`i<\vartheta `$ such that $`i`$ has infinitely many immediate $`T`$-successors, say $`i_k+1`$ for all $`k<\omega `$. We may then apply Claim 2 successively to the almost linear trees $`๐ฏi_k+2`$ to get an infinite descending sequence of ordinals, namely $`(\mathrm{c}.\mathrm{p}.(E_{i_k}^๐ฏ):k\omega )`$. Contradiction!
Case 2. $`\vartheta `$ is a successor ordinal.
By 2.2 we then know that $`\mathrm{lh}(๐ฏ)`$ must be a double successor, say $`\mathrm{lh}(๐ฏ)=\vartheta +2`$, and $`๐ฏ\vartheta +1`$ is almost linear. For $`i<\vartheta +1`$ let $`\kappa _i=\mathrm{c}.\mathrm{p}.(E_i^๐ฏ)`$, and $`\lambda _i=E_i^๐ฏ(\kappa _i)`$. Because $`๐ฏ`$ is normal, we must have $`i<j<\vartheta +1\lambda _i<\lambda _j`$.
Set $`j=๐ฏ`$-$`\mathrm{pred}(\vartheta +1)`$. Note that $`j[0,\vartheta ]_T`$ (and so in particular $`j<\vartheta `$) by our assumption on $`๐ฏ`$. Let then $`k<j`$ be maximal such that $`k[0,j)_T`$ as well as $`k[0,\vartheta )_T`$.
By the normality of $`๐ฏ`$, $`j`$ is least such that $`\kappa _\vartheta <\lambda _j`$. In particular,
(1) $`i(ji<\vartheta +1\kappa _\vartheta <\lambda _j\lambda _i).`$
Let $`i+1`$ be minimal in $`(k,\vartheta ]_T`$, so that $`E_i^๐ฏ`$ is the first extender used on $`[k,\vartheta )_๐ฏ`$. As $`๐ฏ\vartheta +1`$ is almost linear, we must have that $`i+1>j`$, i.e., $`ij`$. Hence by (1),
(2) $`\kappa _\vartheta <\lambda _j\lambda _i\lambda _\vartheta .`$
Recall that $`k<j`$. As $`k=๐ฏ`$-$`\mathrm{pred}(i+1)`$ we must have $`\kappa _i<\lambda _k`$. We thus have that
(3) $`\kappa _i<\kappa _\vartheta ,`$
because otherwise $`k<j`$ would be such that $`\kappa _\vartheta \kappa _i<\lambda _k`$; but $`j=๐ฏ`$-$`\mathrm{pred}(\vartheta +1)`$ is least with $`\kappa _\vartheta <\lambda _j`$.
But (2) and (3) contradict Claim 1!
$`\mathrm{}`$ (2.4)
For future reference, let us state an immediate consequence of the proof of Claim 1 in the proof of 2.4.
###### Lemma 2.5
Let $``$ be a premouse which is below $`0^{^{}}`$, and let $`๐ฏ`$ be a normal tree on $``$ with last model $`๐ฉ=_{\mathrm{}}^๐ฏ`$. Let $`i+1<\mathrm{lh}(๐ฏ)`$, $`F=E_i^๐ฏ`$, $`\kappa =\mathrm{c}.\mathrm{p}.(F)`$, and $`\lambda =F(\kappa )`$. Then for no $`\stackrel{~}{\lambda }\lambda `$ and $`\stackrel{~}{\kappa }[\kappa ,\lambda )`$ do we have that $`\stackrel{~}{\kappa }`$ is $`<\stackrel{~}{\lambda }`$-strong in $`๐ฉ`$.
Using 1.3 we could have shown the following. Let $``$ be a $`0`$-iterable premouse which is such that for no $`E_\nu ^{}\mathrm{}`$ with critical point $`\kappa `$ do we have
$$\{\mu <\kappa :\mu \mathrm{is}<\kappa \mathrm{strong}\mathrm{in}\mathrm{as}\mathrm{witnessed}\mathrm{by}\stackrel{}{E}^{}\}\mathrm{is}\mathrm{unbounded}\mathrm{in}\kappa .$$
Then any normal iteration tree on $``$ is almost linear. The advantage of 2.4 is that it can be applied to premice of which we donโt (yet) know that they are iterable.
The insight which gives 2.4 also shows that normal iterations of phalanxes below $`0^{^{}}`$ are of a simple form.
###### Definition 2.6
Let $`\stackrel{}{๐ซ}=((๐ซ_i:i<\alpha +1),(\mu _i:i<\alpha ))`$ be a phalanx. $`\stackrel{}{๐ซ}`$ is said to be below $`0^{^{}}`$ if every $`๐ซ_i`$ for $`i<\alpha +1`$ is below $`0^{^{}}`$.
###### Lemma 2.7
Let the phalanx $`\stackrel{}{๐ซ}=((๐ซ_i:i<\alpha +1),(\mu _i:i<\alpha ))`$ be below $`0^{^{}}`$, and let $`๐ฏ`$ be a normal iteration of $`\stackrel{}{๐ซ}`$. Then there are $`i_n<i_{n1}<\mathrm{}<i_1<i_0=\alpha `$ and $`\alpha =\beta _0<\beta _1<\mathrm{}<\beta _{n1}<\beta _n<\beta _{n+1}=\mathrm{lh}(๐ฏ)`$ suchthat for all $`\gamma <\mathrm{lh}(๐ฏ)`$ and for all $`kn`$ do we have that
$$\gamma [\beta _k,\beta _{k+1})i_kT\gamma .$$
The lemma says, among other things, that we can write $`๐ฏ`$ as
$$๐ฏ=๐ฏ_0{}_{}{}^{}๐ฏ_{1}^{}{}_{}{}^{}\mathrm{}{}_{}{}^{}๐ฏ_{n}^{}$$
where $`๐ฏ_0`$ is an iteration of $`๐ซ_{i_0}=๐ซ_\alpha `$ (and may be trivial) and $`๐ฏ_{k+1}`$ is an iteration of $`๐ซ_{i_{k+1}}`$ (except for the fact that its first extender is taken from the last model of $`๐ฏ_k`$). However, the proof of 2.7 is straightforward in the light of the proof of 2.4 and may be left to the reader.
Lemma 2.4 is optimal in the following sense.
###### Lemma 2.8
Let $``$ be a premouse with top extender $`F`$ such that for $`\kappa =\mathrm{c}.\mathrm{p}.(F)`$, there are arbitrary large $`\mu <\kappa `$ such that $`\mu `$ is $`<\kappa `$ strong in $``$ as witnessed by $`\stackrel{}{E}^{}`$. (In particular, $``$ is not below $`0^{^{}}`$.) Then assuming that $``$ is sufficiently iterable there is a normal iteration tree $`๐ฏ`$ on $``$ of length $`4`$ which is not almost linear.
Proof. Let $`\lambda `$ be the largest cardinal of $``$. Let $`\nu _0<\lambda `$ be such that $`E_{\nu _0}^{}\mathrm{}`$ is total on $``$ and has critical point $`\kappa _0>\kappa `$. Let $`\pi _{01}^๐ฏ:_{E_{\nu _0}^{}}_1`$. Let $`F^{}`$ be the top extender of $`_1`$, and let $`\pi _{02}^๐ฏ:_F^{}_2`$. Let $`F^{\prime \prime }`$ be the top extender of $`_2`$.
Note that $`F^{}(\kappa )=\pi _{01}^๐ฏ(\lambda )`$ is the critical point of $`F^{\prime \prime }`$, and that $`F^{}(\kappa )\lambda >\nu _0`$, so that we may pick some $`\mu (\nu _0,F^{}(\kappa ))`$ which is $`<F^{}(\kappa )`$-strong in $`_1`$ as witnessed by $`\stackrel{}{E}^_1`$. Using $`F^{\prime \prime }`$, $`\mu `$ is also $`<F^{\prime \prime }(F^{}(\kappa ))`$-strong in $`_2`$ as witnessed by $`\stackrel{}{E}^_2`$. We may thus pick $`\nu _2>_1\mathrm{OR}`$ such that $`E_{\nu _2}^_2\mathrm{}`$ is total on $`_2`$ and has critical point $`\mu `$. Let $`\pi _{13}^๐ฏ:_1_{E_{\nu _2}^_2}_3`$.
It is easy to see that we have constructed a normal iteration tree on $``$ of length $`4`$ which is not almost linear.
$`\mathrm{}`$ (2.8)
One can easily generalize 2.4 and 2.8.
###### Definition 2.9
Let $``$ be a premouse, and let $`\mu <\kappa \mathrm{OR}`$. We call $`\mu `$ $`<\kappa `$-$`0`$-strong in $``$ if $`\mu `$ is a measurable cardinal in $`๐ฅ_\kappa ^{}`$ as witnessed by $`\stackrel{}{E}^{}`$, i.e, there is an extender $`E_\nu ^{}\mathrm{}`$ with critical point $`\mu `$ and such that $`\nu >\mu ^{+๐ฅ_\kappa ^{}}`$. For $`n<\omega `$ we call $`\mu `$ $`<\kappa `$-$`(n+1)`$-strong in $``$ if $`\mu `$ is $`<\kappa `$-strong in $``$ as witnessed by $`\stackrel{}{E}^{}`$ and there are arbitrary large $`\overline{\mu }<\mu `$ such that $`\overline{\mu }`$ is $`<\mu `$-$`n`$-strong in $``$.
###### Definition 2.10
Let $``$ be a premouse, and let $`n<\omega `$. $``$ is said to be below $`0^n^{}`$ if there is no $`\kappa `$ which is the critical point of an extender $`E_\nu ^{}\mathrm{}`$ with the property that
$$\{\mu <\kappa :\mu \mathrm{is}<\kappa n\mathrm{strong}\mathrm{in}\}\mathrm{is}\mathrm{unbounded}\mathrm{in}\kappa .$$
Notice that $``$ is below $`0^1^{}`$ if and only if $``$ is below $`0^{^{}}`$. The following lemma, whose easy proof we omit, generalizes 2.8.
###### Lemma 2.11
Let $``$ be a premouse which is not below $`0^n^{}`$. Then โ provided all the involved ultrapowers are well-founded โ one can build a normal alternating chain on $``$ of length $`3+n`$, i.e., a normal iteration tree with tree structure $`T`$ given by $`iTj`$ if and only if $`i=0`$ or $`(ijij(mod2))`$ for $`ij<3+n`$.
This is tight in the sense that there are no iteration trees with so much โjumping from branch to branchโ if every premouse is below $`0^n^{}`$. We leave it to the reader to formulate a precise generalization of 2.4.
Now suppose that $``$ is a premouse with top extender $`F`$, which is such that, setting $`\kappa =\mathrm{c}.\mathrm{p}.(F)`$,
$$\{\mu <\kappa :\mu \mathrm{is}<\kappa \mathrm{strong}\mathrm{in}\mathrm{as}\mathrm{witnessed}\mathrm{by}\stackrel{}{E}^{}\}\mathrm{is}\mathrm{stationary}.$$
Then for no $`n<\omega `$ is $``$ below $`0^n^{}`$, and โ again provided all the involved ultrapowers are well-founded โ for every $`n<\omega `$ can one build a normal alternating chain on $``$ of length $`3+n`$. Recall that (by \[20, Theorem 6.1\]; see also \[9, ยง6\]) to build an alternating chain of length $`\omega `$ requires an assumption at the level of a (definably) Woodin cardinal.
We also want to mention an extension of 2.4 into another direction, namely by revising the definition of โnormal trees.โ
###### Definition 2.12
Let $``$ be a premouse. $``$ is said to be below $`0\mathrm{}`$ if for no $`\kappa `$ which is the critical point of some extender $`F=E_\nu ^{}\mathrm{}`$ do we have that, setting $`\lambda =F(\kappa )`$, there is some $`\mu (\kappa ,\lambda )`$ with
$$๐ฅ_\lambda ^{}\mathrm{`}\mathrm{`}\mu \mathrm{is}\mathrm{a}\mathrm{strong}\mathrm{cardinal},\mathrm{"}$$
and $`F`$ has a generator (strictly) above $`\mu `$.
If we were to let $`T`$-$`\mathrm{pred}(i+1)`$ be the least $`ji`$ such that $`\mathrm{c}.\mathrm{p}.(E_i^๐ฏ)<`$ the $`\mathrm{sup}`$ of the generators of $`E_j^๐ฏ`$ in a D-normal iteration tree $`๐ฏ`$ then we would get that any D-normal iteration tree of a premouse which is below $`0\mathrm{}`$ is almost linear.<sup>3</sup><sup>3</sup>3$`0\mathrm{}`$ is mentioned in the introduction to . Cf. also \[11, ยง3\]. This observation can be used to develop the theory of $`K`$ up to the level of $`0\mathrm{}`$ (cf. ).
## 3 Iterability, and the existence of $`K^c`$.
This section introduces $`K^c`$, a preliminary version of $`K`$. $`K^c`$ is constructed recursively, exactly as in \[9, ยง11\] (see also \[27, p. 6 f.\]), except for the fact that we only require new extenders to be countably complete (rather than โcertifiableโ as in \[9, ยง11\] when they are put onto the sequence. Recall that an extender $`F`$ with $`\kappa =\mathrm{c}.\mathrm{p}.(F)`$ is called countably complete if for all $`(a_n,X_n:n<\omega )`$ such that $`a_n[F(\kappa )]^{<\omega }`$ and $`a_nF(X_n)`$ for all $`n<\omega `$ there is some order preserving $`\tau :_{n<\omega }a_n\kappa `$ with $`\tau \mathrm{"}a_nX_n`$ for every $`n<\omega `$.
Let us define premice $`๐ฉ_\xi `$ and $`_\xi `$ by induction on $`\xi \mathrm{OR}\{\mathrm{OR}\}`$ (cf. \[9, ยง10 p. 9 f.\]). We let $`๐ฉ_0=(J_\omega ;,\mathrm{})`$. Having defined $`๐ฉ_\xi `$, we let the construction break down unless $`๐ฉ_\xi `$ is reliable, i.e., unless for all $`n\omega `$, $`_n(๐ฉ_\xi )`$ is $`n`$-iterable. If $`๐ฉ_\xi `$ is โreliableโ then we continue by setting
$$_\xi =_\omega (๐ฉ_\xi ).$$
Now suppose that $`_\xi `$ has been defined.
Case 1. $`_\xi =(J_\alpha [\stackrel{}{E}];,\stackrel{}{E})`$ is passive, and there is a unique countably complete extender $`F`$ such that
$$(J_\alpha [\stackrel{}{E}];,\stackrel{}{E},F)$$
is a premouse below $`0^{^{}}`$. In this case we set $`๐ฉ_{\xi +1}=(J_\alpha [\stackrel{}{E}];,\stackrel{}{E},F)`$.
Case 2. Otherwise. Then we just construct one more step. I.e., if $`_\xi =(J_\alpha [\stackrel{}{E}];,\stackrel{}{E})`$ then we let $`๐ฉ_{\xi +1}=(J_{\alpha +1}[\stackrel{}{E}];,\stackrel{}{E})`$, and if $`_\xi =(J_\alpha [\stackrel{}{E}];,\stackrel{}{E},F)`$ (with $`F\mathrm{}`$) then we let $`๐ฉ_{\xi +1}=(J_{\alpha +1}[\stackrel{}{E}^{}F];,\stackrel{}{E}^{}F)`$.
Now suppose that $`_\xi `$ has been defined for all $`\xi <\lambda `$ where $`\lambda `$ is a limit ordinal. We may let
$$\mu =\mathrm{sup}_{\xi <\lambda }\mathrm{min}_{\xi <\zeta <\lambda }\rho _\omega (_\zeta )^{+_\xi },$$
and let $`๐ฉ_\lambda `$ be that passive premouse of height $`\mu `$ such that for all $`\overline{\mu }<\mu `$, $`๐ฅ_{\overline{\mu }}^{๐ฉ_\lambda }`$ is the eventual value of $`๐ฅ_{\overline{\mu }}^{๐ฉ_\xi }`$ as $`\xi \lambda `$ (cf. \[9, ยง10 p. 10\], where it is shown that $`๐ฉ_\lambda `$ is well-defined if all $`_\xi `$, $`\xi <\lambda `$, are well-defined).
Notice that, whereas there is no restriction on $`\mathrm{cf}^V(\mathrm{c}.\mathrm{p}.(F))`$ in Case 1 above, it is automatic that $`\mathrm{cf}^V(\mathrm{c}.\mathrm{p}.(F))>\omega `$, because no countably complete extender can have a critical point with cofinality $`\omega `$.
We are now going to prove that the construction never breaks down. For this purpose, we need the resurrection maps provided by the following lemma. (We handle the issue of resurrecting slightly different than both \[20, ยง12\] and \[27, ยง9\].)
###### Lemma 3.1
and Definition Suppose that $`๐ฉ_{\xi _0}`$ exists for some $`\xi _0\mathrm{OR}`$. Then there are
$$(\tau _{(\xi ,\eta )}:\xi \xi _0\eta ๐ฉ_\xi \mathrm{OR})$$
and
$$(\phi (\xi ,\eta ):\xi \xi _0\eta ๐ฉ_\xi \mathrm{OR})$$
such that for every pair $`(\xi ,\eta )`$ with $`\xi \xi _0`$ and $`\eta ๐ฉ_\xi \mathrm{OR}`$ do we have that
$$\tau _{(\xi ,\eta )}:๐ฅ_\eta ^{๐ฉ_\xi }๐ฉ_{(\phi _{(\xi ,\eta )})},$$
where, for every $`n<\omega `$,
$$\tau _{(\xi ,\eta )}(๐ฅ_\eta ^{๐ฉ_\xi })^n:(๐ฅ_\eta ^{๐ฉ_\xi })^n_{\mathrm{\Sigma }_0}(๐ฉ_{(\phi (\xi ,\eta ))})^n;$$
the map $`\tau _{(\xi ,\eta )}`$ is called a resurrection map. Moreover, whenever $`\eta <\eta ^{}๐ฉ_\xi \mathrm{OR}`$ then, setting $`\rho =\mathrm{min}\{\rho _\omega (๐ฅ_{\overline{\eta }}^{๐ฉ_\xi }):\eta \overline{\eta }<\eta ^{}\}`$, we have that $`\tau _{(\xi ,\eta )}`$ agrees with $`\tau _{(\xi ,\eta ^{})}`$ up to $`\rho `$, i.e.,
$$๐ฅ_{\tau _{(\xi ,\eta ^{})}(\rho )}^{๐ฉ_{(\phi (\xi ,\eta ))}}=๐ฅ_{\tau _{(\xi ,\eta ^{})}(\rho )}^{๐ฉ_{(\phi (\xi ,\eta ^{}))}},\mathrm{and}$$
$$\tau _{(\xi ,\eta )}๐ฅ_\rho ^{๐ฉ_\xi }=\tau _{(\xi ,\eta ^{})}๐ฅ_\rho ^{๐ฉ_\xi }.$$
Proof sketch. The maps $`\tau _{(\xi ,\eta )}`$ and their target models are chosen by induction on $`\xi `$. Suppose $`\tau _{(\xi ,\eta )}`$ has been defined for all $`\eta ๐ฉ_\xi \mathrm{OR}`$. We let $`\tau _{(\xi +1,๐ฉ_{\xi +1}\mathrm{OR})}=\mathrm{id}`$. If $`\eta <๐ฉ_{\xi +1}\mathrm{OR}`$ then we have that $`๐ฅ_\eta ^{๐ฉ_{\xi +1}}=๐ฅ_\eta ^_\xi `$. Let $`\pi :_\xi ๐ฉ_\xi `$ be the core map. We then let
$$\phi (\xi +1,\eta )=\phi (\xi ,\pi (\eta )),\mathrm{and}$$
$$\tau _{(\xi +1,\eta )}=\tau _{(\xi ,\pi (\eta ))}\pi ๐ฅ_\eta ^_\xi .$$
It is straightforward to verify the required agreement between appropriate $`\tau _{(\xi +1,\eta )}`$ and $`\tau _{(\xi +1,\eta ^{})}`$.
At limit stages $`\lambda `$ we exploit the fact that any proper initial segment of $`๐ฉ_\lambda `$ is an initial segment of some $`๐ฉ_\xi `$ for $`\xi <\lambda `$. We leave the details to the reader.
$`\mathrm{}`$ (3.1)
The following two lemmas, 3.2 and 3.3, proving normal and full $`n`$-iterability of $`_n(๐ฉ_\xi )`$, are shown simultaneously by induction on $`n<\omega `$.
###### Lemma 3.2
Suppose that $`๐ฉ_\xi `$ exists for some $`\xi \mathrm{OR}`$ and $`n<\omega `$. Then $`_n(๐ฉ_\xi )`$ is normally $`n`$-iterable.
Proof. Let us fix some $`n<\omega `$ and assume that $`_m(๐ฉ_\xi )`$ is $`m`$-iterable for all $`m<n`$. By \[20, ยง4, and Lemmas 6.1.5 and 8.1\], this will buy us that we already know that for all $`kn`$, any $`k`$-bounded iteration of $`_n(๐ฉ_\xi )`$ will move $`p_{_n(๐ฉ_\xi ),k}`$ as well as all standard parameters of all non-simple iterates of $`_n(๐ฉ_\xi )`$ correctly.
By 2.2 and 2.4, if the lemma fails then there is a putative normal $`n`$-bounded iteration tree $`๐ฐ`$ of minimal length $`\beta `$ on $`_n(๐ฉ_\xi )`$ such that either $`๐ฐ`$ has a last ill-founded model, or else $`๐ฐ`$ has limit length and $`๐^๐ฐb^{}`$ is infinite where $`b^{}`$ is the unique cofinal branch through $`U`$. A standard argument yields a (fully elementary) embedding
$$\sigma :_n(๐ฉ_\xi )$$
such that $``$ is countable and there is a putative normal $`n`$-bounded iteration tree $`๐ฏ`$ of minimal length $`\alpha `$ on $``$ such that $`\alpha <\omega _1`$ and either $`๐ฏ`$ has a last ill-founded model, or else $`๐ฏ`$ has limit length and $`๐^๐ฏb`$ is infinite where $`b`$ is the unique cofinal branch through $`T`$. Let
$$๐ฏ=((_i^๐ฏ,\pi _{ij}^๐ฏ:iTj<\alpha +1),(E_i^๐ฏ:i<\alpha ),T).$$
For $`i<\alpha `$ we set $`\kappa _i=\mathrm{c}.\mathrm{p}.(E_i^๐ฏ)`$, $`\lambda _i=E_i^๐ฏ(\kappa _i)`$, and $`\eta _i=\mathrm{dom}(\pi _{i^{}i+1}^๐ฏ)\mathrm{OR}`$ where $`i^{}=T`$-pred$`(i+1)`$. We set $`n(0)=n`$, and for $`i>0`$ we let $`n(i)`$ be such that
$$j(j<_Tik+1(j,i]_T_{k+1}^๐ฏ=\mathrm{Ult}_{n(i)}(๐ฅ_{\eta _k}^{_{Tpred(k+1)}^๐ฏ};E_k^๐ฏ)).$$
(That is, $`n(i)=\mathrm{deg}^๐ฏ(i)`$ for successor ordinals $`i`$.) Notice that for all $`i<\alpha +1`$ do we have that $`_i^๐ฏ`$ is $`n(i)`$-sound.
We shall now pick for all $`i\alpha `$ some $`\xi (i)\xi `$ together with a weak $`n(i)`$-embedding
$$\stackrel{~}{\sigma }_i:_i^๐ฏ_{n(i)}(๐ฉ_{\xi (i)}).$$
The maps $`\stackrel{~}{\sigma }_i`$ may be obtained by composing some $`\sigma _i`$ with (the inverse of) a core map. We shall therefore just recursively pick $`\xi (i)`$ and $`\sigma _i`$. Weโll inductively maintain that the following three requirements are met.
R 1<sub>i</sub> $`k[0,i]_T`$ do we have that $`\sigma _k:_k^๐ฏ๐ฉ_{\xi (k)}`$ is the extension of
$$\sigma _k(_k^๐ฏ)^{n(k)}:(_k^๐ฏ)^{n(k)}_{\mathrm{\Sigma }_0}(๐ฉ_{\xi (k)})^{n(k)}$$
given by the downward extension of embeddings lemma.
R 2<sub>i</sub> $`k+1(0,i]_T`$, setting $`k^{}=๐ฏ`$-$`\mathrm{pred}(k+1)`$, we have that $`\sigma _i`$ agrees with $`\tau _{(\xi (k^{}),\sigma _k^{}(\eta _k))}\sigma _k^{}`$ up to $`\kappa _k`$, i.e.,
$$๐ฅ_{\sigma _i(\kappa _k)}^{๐ฉ_{\xi (i)}}=๐ฅ_{\sigma _i(\kappa _k)}^{๐ฉ_{\phi (\xi (k^{}),\sigma _k^{}(\eta _k))}},\mathrm{and}$$
$$\sigma _i๐ฅ_{\kappa _k}^{_i^๐ฏ}=\tau _{(\xi (k^{}),\sigma _k^{}(\eta _k))}\sigma _k^{}๐ฅ_{\kappa _k}^{_k^{}^๐ฏ}.$$
R 3<sub>i</sub> if $`kTj(0,i]_T`$ and $`๐^๐ฏ[k,j]_T=\mathrm{}`$ then $`\sigma _j\pi _{kj}^๐ฏ=\sigma _k`$.
To commence, we set $`\xi (0)=\xi `$ and $`\sigma _0=\pi \sigma `$ where $`\pi :_n(๐ฉ_\xi )๐ฉ_\xi `$ is the core map. It is easy to see that R 1<sub>0</sub> holds. Moreover, R 2<sub>0</sub> and R 3<sub>0</sub> are vacuously true.
Now suppose we have defined $`\sigma _j:_j^๐ฏ๐ฉ_{\xi (j)}`$ for all $`ji`$ in such a way that A 1<sub>i</sub>, A 2<sub>i</sub>, and A 3<sub>i</sub> hold. Suppose that $`i+1<\alpha `$, so that the $`n(i+1)`$-ultrapower
$$\pi _{i^{}i+1}^๐ฏ:๐ฅ_{\eta _i}^{_i^{}^๐ฏ}_F_{i+1}^๐ฏ$$
exists with $`F=E_i^๐ฏ`$ and $`i^{}=๐ฏ`$-$`\mathrm{pred}(i+1)`$. It is convenient to split the construction into two cases.
Case 1. $`i^{}<i`$.
In this case, $`\lambda _i^{}`$ is a cardinal in $`_i^๐ฏ`$, so that by $`\kappa _i<\lambda _i^{}`$ we have that $`F`$ is a total extender on $`_i^๐ฏ`$. Hence $`G=\sigma _i(F)`$ is countably complete, being total on $`๐ฉ_{\xi (i)}`$. So we may pick $`\rho :\sigma _i(\lambda _i)\mathrm{ran}(\sigma _i)\sigma _i(\kappa _i)`$ order preserving such that
$$aG(X)\rho \mathrm{"}aX$$
for all appropriate $`a`$, $`X\mathrm{ran}(\sigma _i)`$. Let $`\xi (i+1)=\phi (\xi (i^{}),\sigma _i^{}(\eta _i))`$, i.e., $`_{\xi (i+1)}`$ is the target model of the resurrection map
$$\tau _{(\xi (i^{}),\sigma _i^{}(\eta _i))}:๐ฅ_{\sigma _i^{}(\eta _i)}^{_{\xi (i^{})}}๐ฉ_{\xi (i+1)}.$$
By 2.4, $`i^{}[0,i)_๐ฏ`$. Let $`k+1`$ be minimal in $`(i^{},i]_๐ฏ`$. In particular, $`๐ฏ`$-$`\mathrm{pred}(k+1)=i^{}`$. Now R 2<sub>i</sub> gives us that $`\sigma _i`$ agrees with $`\tau _{\xi (i^{}),\sigma _i^{}(\eta _k)}\sigma _i^{}`$ up to $`\kappa _k`$. But by the proof of 2.4 we have $`\kappa _i<\kappa _k`$, so that $`\eta _i\eta _k`$ and
$$\xi [\eta _k,\eta _i)(\rho _\omega (๐ฅ_\xi ^{_i^{}^๐ฏ})>\kappa _i),\mathrm{and}\mathrm{thus}$$
$$\xi [\sigma _i^{}(\eta _k),\sigma _i^{}(\eta _i))(\rho _\omega (๐ฅ_\xi ^{๐ฉ_{\xi (i^{})}})>\sigma _i^{}(\kappa _i)),$$
which by 3.1 implies that $`\tau _{(\xi (i^{}),\sigma _i^{}(\eta _i))}`$ agrees with $`\tau _{(\xi (i^{}),\sigma _i^{}(\eta _k))}`$ up to $`\sigma _i^{}(\kappa _i^+)`$ (where $`\kappa _i^+`$ is calculated in $`๐ฅ_{\eta _i}^{_i^{}^๐ฏ}`$). These agreements combined easily give that $`\sigma _i`$ agrees with $`\tau _{(\xi (i^{}),\sigma _i^{}(\eta _i))}\sigma _i^{}`$ up to $`\kappa _i^+`$. Let us write $`\tau =\tau _{(\xi (i^{}),\sigma _i^{}(\eta _i))}`$.
Let
$$p=p_{๐ฅ_{\eta _i}^{_i^{}^๐ฏ},n(i+1)}.$$
By R 1<sub>i</sub> we will have that
$$\sigma _i^{}(p)=p_{๐ฅ_{\sigma _i^{}(\eta _i)}^{๐ฉ_{\xi (i^{})}},n(i+1)}.$$
Moreover, by 3.1 we then get that
$$\tau \sigma _i^{}(p)=p_{๐ฉ_{\xi (i+1)},n(i+1)},\mathrm{and}$$
$$\tau (๐ฅ_{\sigma _i^{}(\eta _i)}^{๐ฉ_{\xi (i^{})}})^{n(i+1),\sigma _i^{}(p)}:(๐ฅ_{\sigma _i^{}(\eta _i)}^{๐ฉ_{\xi (i^{})}})^{n(i+1),\sigma _i^{}(p)}_{\mathrm{\Sigma }_0}(๐ฉ_{\xi (i+1)})^{n(i+1),\tau \sigma _i^{}(p)}.$$
We may now define
$$\overline{\sigma }:(_{i+1}^๐ฏ)^{n(i+1),\pi _{i^{}i+1}^๐ฏ(p)}_{\mathrm{\Sigma }_0}(๐ฉ_{\xi (i+1)})^{n(i+1),\tau \sigma _i^{}(p)}=(๐ฉ_{\xi (i+1)})^{n(i+1)}$$
by setting
$$\overline{\sigma }([a,f])=\tau \sigma _i^{}(f)(\rho \mathrm{"}\sigma _i(a)).$$
To show that this is well-defined and $`\mathrm{\Sigma }_0`$-elementary we may reason as follows. Let $`\mathrm{\Phi }`$ be a $`\mathrm{\Sigma }_0`$ formula. Then
$$(_{i+1}^๐ฏ)^{n(i+1),\pi _{i^{}i+1}^๐ฏ(p)}\mathrm{\Phi }([a_1,f_1],\mathrm{},[a_k,f_k])$$
$$(a_1,\mathrm{},a_k)F(\{(u_1,\mathrm{},u_k):(๐ฅ_{\eta _i}^{_i^{}^๐ฏ})^{n(i+1),p}\mathrm{\Phi }(f_1(u_1),\mathrm{},f_k(u_k))\})$$
$$(\sigma _i(a_1),\mathrm{},\sigma _i(a_k))G(\sigma _i(\{(u_1,\mathrm{},u_k):(๐ฅ_{\eta _i}^{_i^{}^๐ฏ})^{n(i+1),p}\mathrm{\Phi }(f_1(u_1),\mathrm{},f_k(u_k))\})),$$
which, by the amount of agreement of $`\sigma _i`$ with $`\tau \sigma _i^{}`$, holds if and only if
$`(\sigma _i(a_1),\mathrm{},\sigma _i(a_k))`$
$$G(\tau \sigma _i^{}(\{(u_1,\mathrm{},u_k):(๐ฅ_{\eta _i}^{_i^{}^๐ฏ})^{n(i+1),p}\mathrm{\Phi }(f_1(u_1),\mathrm{},f_k(u_k))\}))$$
$`(\sigma _i(a_1),\mathrm{},\sigma _i(a_k))`$
$$G(\{(u_1,\mathrm{},u_k):(๐ฉ_{\xi (i+1)})^{n(i+1),\tau \sigma _i^{}(p)}\mathrm{\Phi }(\tau \sigma _i^{}(f_1)(u_1),\mathrm{},\tau \sigma _i^{}(f_k)(u_k))\})$$
$$(๐ฉ_{\xi (i+1)})^{n(i+1),\tau \sigma _i^{}(p)}\mathrm{\Phi }(\tau \sigma _i^{}(f_1)(\rho \mathrm{"}\sigma _i(a_1),\mathrm{},\tau \sigma _i^{}(f_k)(\rho \mathrm{"}\sigma _i(a_k))).$$
We let $`\sigma _{i+1}`$ be the extension of $`\overline{\sigma }`$ given by the downward extension of embeddings lemma. By the remark in the first paragraph of this proof, weโll have that
$$\pi _{i^{}i+1}^๐ฏ(p)=p_{_{i+1}^๐ฏ,n(i+1)}.$$
Also, by the definition of $`\sigma _{i+1}`$,
$$\sigma _{i+1}(\pi _{i^{}i+1}^๐ฏ(p))=\tau \sigma _i^{}(p)=p_{๐ฉ_{\xi (i+1)},n(i+1)},$$
and hence
$$\sigma _{i+1}(p_{_{i+1}^๐ฏ,n(i+1)})=p_{๐ฉ_{\xi (i+1)},n(i+1)}.$$
Hence we have established R 1<sub>i+1</sub>.
Let us verify R 2<sub>i+1</sub>. It is clear by construction that $`\sigma _{i+1}`$ agrees with $`\tau \sigma _i^{}`$ up to $`\kappa _i`$. So let $`k+1(0,i+1]_T`$ be such that $`k<i`$. Then $`\lambda _k\kappa _i`$ is a cardinal in $`_i^{}^๐ฏ`$, so that
$$\xi [\lambda _k,\mathrm{OR}_i^{}^๐ฏ)\rho _\omega (๐ฅ_\xi ^{_i^{}^๐ฏ})\lambda _k,\mathrm{thus}$$
$$\xi [\sigma _i^{}(\lambda _k),\mathrm{OR}๐ฉ_{\xi (i^{})})\rho _\omega (๐ฅ_\xi ^{๐ฉ_{\xi (i^{})}})\sigma _i^{}(\lambda _k),$$
which implies that
$$\tau _{(\xi (i^{}),\sigma _i^{}(\eta _i))}\sigma _i^{}(\lambda _k)=\mathrm{id}.$$
Combining this with R 2$`_i^{}`$ we get that $`\tau _{(\xi (k^{}),\sigma _k^{}(\eta _k))}\sigma _k^{}`$ agrees with $`\tau _{(\xi (i^{}),\sigma _i^{}(\eta _i))}\sigma _i^{}`$ up to $`\kappa _k`$, which in turn agrees with $`\sigma _{i+1}`$ up to $`\kappa _k`$ by the construction of $`\sigma _{i+1}`$. This proves R 2<sub>i+1</sub>.
It is now straightforward to verify R 3<sub>i+1</sub>.
Case 2. $`i^{}=i`$.
In this case, $`F`$ may be partial on $`_i^๐ฏ`$. Let $`\xi (i+1)=\phi (\xi (i),\sigma _i(\eta _i))`$, i.e., $`_{\xi (i+1)}`$ is the target model of the resurrection map
$$\tau _{(\xi (i),\sigma _i(\eta _i))}:๐ฅ_{\sigma _i(\eta _i)}^{_{\xi (i)}}๐ฉ_{\xi (i+1)},$$
and let us write $`\tau =\tau _{(\xi (i),\sigma _i(\eta _i))}`$. We then have that $`G=\tau \sigma _i(F)`$ is countably complete, being total on $`๐ฉ_{\xi (i+1)}`$. So we may pick an orderpreserving $`\rho :\tau \sigma _i(\kappa _i)\mathrm{ran}(\tau \sigma _i)\tau \sigma _i(\kappa _i)`$ such that
$$aG(X)\rho \mathrm{"}aX$$
for all appropriate $`a`$, $`X\mathrm{ran}(\tau \sigma _i)`$.
Let
$$p=p_{๐ฅ_{\eta _i}^{_i^๐ฏ},n(i+1)}.$$
We may define
$$\overline{\sigma }:(_{i+1}^๐ฏ)^{n(i+1),\pi _{ii+1}^๐ฏ(p)}_{\mathrm{\Sigma }_0}(๐ฉ_{\xi (i+1)})^{n(i+1),\sigma _i(p)}$$
by setting
$$\overline{\sigma }([a,f])=\tau \sigma _i(f)(\rho \mathrm{"}\tau \sigma _i(a)).$$
To show that this is well-defined and $`\mathrm{\Sigma }_0`$-elementary we may reason as follows. Let $`\mathrm{\Phi }`$ be a $`\mathrm{\Sigma }_0`$ formula. Then
$$(_{i+1}^๐ฏ)^{n(i+1),\pi _{ii+1}^๐ฏ(p)}\mathrm{\Phi }([a_1,f_1],\mathrm{},[a_k,f_k])$$
$$(a_1,\mathrm{},a_k)F(\{(u_1,\mathrm{},u_k):(๐ฅ_{\eta _i}^{_i^๐ฏ})^{n(i+1),p}\mathrm{\Phi }(f_1(u_1),\mathrm{}f_k(u_k))\})$$
$`(\tau \sigma _i(a_1),\mathrm{},\tau \sigma _i(a_k))`$
$$G(\tau \sigma _i(\{(u_1,\mathrm{},u_k):(๐ฅ_{\eta _i}^{_i^๐ฏ})^{n(i+1),p}\mathrm{\Phi }(f_1(u_1),\mathrm{},f_k(u_k))\}))$$
$`(\tau \sigma _i(a_1),\mathrm{},\tau \sigma _i(a_k))`$
$$G(\{(u_1,\mathrm{},u_k):(๐ฉ_{\xi (i+1)})^{n(i+1),\tau \sigma _i(p)}\mathrm{\Phi }(\tau \sigma _i(f_1)(u_1),\mathrm{},\tau \sigma _i(f_k)(u_k))\})$$
$$(๐ฉ_{\xi (i+1)})^{n(i+1),\tau \sigma _i(p)}\mathrm{\Phi }(\tau \sigma _i(f_1)(\rho \mathrm{"}\tau \sigma _i(a_1)),\mathrm{},\tau \sigma _i(f_k)(\rho \mathrm{"}\tau \sigma _i(a_k))).$$
We may now let $`\sigma _{i+1}`$ be the extension of $`\overline{\sigma }`$ given by the downward extension of embeddings lemma. We can then argue exactly as in Case 1 to see that we have established R 1<sub>i+1</sub>.
As for R 2<sub>i+1</sub> and R 3<sub>i+1</sub>, the proofs are similar as in the previous case.
Now suppose that we have defined $`\sigma _i:_i^๐ฏ๐ฉ_{\xi (i)}`$ for all $`i<\lambda `$ where $`\lambda <\alpha `$ is a limit ordinal. By our minimality assumption on $`\alpha `$ we have that $`๐^๐ฏ(0,\lambda ]_T`$ is finite. Hence we may use that for all $`i<\lambda `$ the statement R 3<sub>i</sub> holds to define
$$\sigma _\lambda (x)=\sigma _i(\pi _{i\lambda }^๐ฏ)^1(x)$$
where $`i<_T\lambda `$ is large enough. It is straightforward to verify that with this definition we have R 1<sub>ฮป</sub>, R 2<sub>ฮป</sub>, and R 3<sub>ฮป</sub>.
Now if $`๐ฏ`$ has a last model $`_{\alpha 1}^๐ฏ`$ then it follows from R 1<sub>ฮฑ-1</sub> that $`_{\alpha 1}^๐ฏ`$ canโt be ill-founded. Also, if $`๐ฏ`$ has limit length and $`๐^๐ฏb`$ is infinite where $`b`$ is the unique cofinal branch through $`T`$, then the indices of the target models of $`\sigma _i`$ for $`ib`$ yield an infinite descending sequence of ordinals. We have reached a contradiction!
$`\mathrm{}`$ (3.2)
We can now exploit the previous argument a bit further and arrive at the following.
###### Lemma 3.3
Suppose that $`_n(๐ฉ_\xi )`$ exists for some $`\xi <\mathrm{}`$ and $`n<\omega `$. Then $`_n(๐ฉ_\xi )`$ is $`n`$-iterable.
Proof sketch. This time, we have to deal with a sequence of iteration trees
$$๐ฏ_0{}_{}{}^{}๐ฏ_{1}^{}{}_{}{}^{}๐ฏ_{2}^{}{}_{}{}^{}\mathrm{}$$
of length $`\omega `$, where $`๐ฏ_0`$ is on $``$, and $`๐ฏ_{i+1}`$ is on the last model of $`๐ฏ_i`$. We may now repeatedly apply the proof of 3.2 to see that the last model of any $`๐ฏ_i`$ can be embedded into some $`๐ฉ_\xi `$. This will give the desired conclusion.
$`\mathrm{}`$ (3.3)
It is now easy to see that $`_{\mathrm{OR}}=๐ฉ_{\mathrm{OR}}`$ is a model of height $`\mathrm{OR}`$.
###### Definition 3.4
We write $`K^c`$ for $`_{\mathrm{OR}}=๐ฉ_{\mathrm{OR}}`$. $`K^c`$ is called the countably complete core model below $`0^{^{}}`$.
The method of the proof of 3.3 also yields the following.
###### Lemma 3.5
Let $`_\xi =(J_\alpha [\stackrel{}{E}];,\stackrel{}{E})`$ be passive. Then there is at most one countably complete extender $`F`$ such that $`(J_\alpha [\stackrel{}{E}];,\stackrel{}{E},F)`$ is a premouse.
Proof sketch. Suppose $`F`$, $`F^{}`$ are both countably complete and such that $`(J_\alpha [\stackrel{}{E}];,\stackrel{}{E},F)`$ as well as $`(J_\alpha [\stackrel{}{E}];,\stackrel{}{E},F^{})`$ is a premouse. We may then form the prebicephalus
$$=(J_\alpha [\stackrel{}{E}];,\stackrel{}{E},F,F^{}).$$
(Cf. \[20, ยง9\] for a paradigmatic theory of bicephali.) The proof of 3.2 shows that $``$ is $`0`$-iterable in the obvious sense. By coiterating $``$ against itself we may then as usual conclude that in fact $`F=F^{}`$.
$`\mathrm{}`$ (3.5)
Lemma 3.5 may also be derived from our more general theory of bicephali which we shall develop in section 5 below.
Notice that in Case 1 of the recursive construction of $`๐ฉ_{\xi +1}`$ we had $`F`$ required to be unique. 3.5 now says that this requirement poses no restriction at all.
For the purposes of isolating $`K`$ we shall in fact need a variant of $`K^c`$. Let $`\mathrm{\Gamma }\mathrm{}`$ be a class of regular cardinals, $`\omega \mathrm{\Gamma }`$. We then define $`๐ฉ_\xi `$ and $`_\xi `$ exactly as before except that in Case 1 of the definition of $`๐ฉ_{\xi +1}`$ we also require that $`\mathrm{cf}^V(\mathrm{c}.\mathrm{p}.(F))\mathrm{\Gamma }`$ (here, $`F`$ is the top extender of $`๐ฉ_{\xi +1}`$). Let us write $`๐ฉ_\xi ^\mathrm{\Gamma }`$ and $`_\xi ^\mathrm{\Gamma }`$ instead of $`๐ฉ_\xi `$ and $`_\xi `$. The proofs of the current section show that they all exist, and that $`_n(๐ฉ_\xi ^\mathrm{\Gamma })`$ is always $`n`$-iterable.
###### Definition 3.6
Let $`\mathrm{\Gamma }\mathrm{}`$ be a class of regular cardinals, $`\omega \mathrm{\Gamma }`$. We write $`K_\mathrm{\Gamma }^c`$ for $`_{\mathrm{OR}}^\mathrm{\Gamma }=๐ฉ_{\mathrm{OR}}^\mathrm{\Gamma }`$. $`K^c`$ is called the (countably complete) $`\mathrm{\Gamma }`$-full core model below $`0^{^{}}`$.
## 4 Universality of $`K^c`$.
###### Definition 4.1
Let $`W`$ be a weasel. $`W`$ is called full provided the following holds. Let $`\alpha <\beta `$ be any cardinals of $`W`$. There are no $`\stackrel{~}{\tau }>\beta `$ and a countably complete $`\stackrel{~}{F}:๐ซ(\alpha )W๐ซ(\beta )๐ฅ_{\stackrel{~}{\tau }}^W`$ such that $`(๐ฅ_{\stackrel{~}{\tau }}^W,\stackrel{~}{F})`$ is a pre-premouse, $`\mathrm{c}.\mathrm{p}.(\stackrel{~}{F})=\alpha `$, and $`\stackrel{~}{F}(\alpha )=\beta `$.
We shall from now on denote by โ$`\neg 0^{^{}}`$โ the assumption that every premouse is below $`0^{^{}}`$.
###### Lemma 4.2
$`(\neg 0^{^{}})`$ $`K^c`$ is full.
Proof. Suppose otherwise. Let $`\alpha <\beta `$ be cardinals of $`K^c`$, let $`\stackrel{~}{\tau }>\beta `$ and let $`\stackrel{~}{F}:๐ซ(\alpha )K^c๐ซ(\beta )๐ฅ_{\stackrel{~}{\tau }}^{K^c}`$ be countably complete and such that $`(๐ฅ_{\stackrel{~}{\tau }}^{K^c},\stackrel{~}{F})`$ is a pre-premouse, $`\mathrm{c}.\mathrm{p}.(\stackrel{~}{F})=\alpha `$ and $`\stackrel{~}{F}(\alpha )=\beta `$.
Let $`\lambda \beta `$ be least such that $`\stackrel{~}{F}(f)(\xi )<\lambda `$ whenever $`f{}_{}{}^{\alpha }\alpha K^c`$ and $`\xi <\lambda `$. Set $`F=\stackrel{~}{F}|\lambda `$, and $`\tau =\alpha ^{+K^c}`$.
There is $`\pi :\mathrm{Ult}(๐ฅ_\tau ^{K^c},F)๐ฅ_{\stackrel{~}{\tau }}^{K^c}`$ defined by $`\pi (F(f)(a))=\stackrel{~}{F}(f)(a)`$, and of course $`\pi \lambda =\mathrm{id}`$ and $`\stackrel{~}{F}=\pi F`$. By the choice of $`\lambda `$, we also easily get $`F(\alpha )=\lambda `$.
As $`\lambda `$ is a limit cardinal in $`\mathrm{Ult}(๐ฅ_\tau ^{K^c},F)`$, by $`\pi \lambda =\mathrm{id}`$ we then get that in fact $`\lambda `$ is a (limit) cardinal in $`K^c`$. Also, $`\lambda `$ is the largest cardinal in $`\mathrm{Ult}(๐ฅ_\tau ^{K^c},F)`$, so that applying the condensation lemma \[9, Lemma 4\] to cofinally many restrictions of $`\pi `$ gives that actually $`\mathrm{Ult}(๐ฅ_\tau ^{K^c},F)=๐ฅ_\gamma ^{K^c}`$ for some $`\gamma `$.
Notice that $`(๐ฅ_\gamma ^{K^c},F)`$ is a premouse. (The initial segment condition for $`F`$ is vacuously true by the choice of $`\lambda `$.) Moreover, we have that $`๐ฅ_\lambda ^{K^c}=_\eta `$ for some $`\eta `$, as $`\lambda `$ is a cardinal in $`K^c`$. It is also straightforward to check that because $`\lambda `$ is the largest cardinal of $`๐ฅ_\gamma ^{K^c}`$ there is $`\eta ^{}>\eta `$ with $`๐ฅ_\gamma ^{K^c}=_\eta ^{}`$.
We thus have $`_{\eta ^{}+1}=(๐ฅ_\gamma ^{K^c},F)`$ by 3.5, so that in particular $`\rho _\omega (_{\eta ^{}+1})<\lambda `$. But the fact that $`\lambda `$ is a cardinal in $`K^c`$ implies that $`\rho _\omega (_{\stackrel{~}{\eta }})\lambda `$ for all $`\stackrel{~}{\eta }\eta `$. Contradiction!
$`\mathrm{}`$ (4.2)
###### Definition 4.3
Let $`W`$ be a weasel. $`W`$ is called universal if $`W`$ is iterable, and whenever $`๐ฏ`$, $`๐ฐ`$ are iteration trees arising from the coiteration of $`W`$ with some (set- or class-sized) premouse $``$ such that $`\mathrm{lh}(๐ฏ)=\mathrm{lh}(๐ฐ)=\mathrm{OR}+1`$, then $``$ is a weasel, $`๐^๐ฐ(0,\mathrm{OR}]_U=\mathrm{}`$, $`\pi _0\mathrm{}^๐ฐ\mathrm{"}\mathrm{OR}\mathrm{OR}`$, and $`_{\mathrm{}}^๐ฐ\mathrm{}_{\mathrm{}}^๐ฏ`$.
###### Lemma 4.4
$`(\neg 0^{^{}})`$ Every full weasel is universal.
Proof. This is shown by varying an argument which is due to Jensen, cf. the Addendum to \[9, ยง3 Theorem 5\].
Fix a full weasel $`W`$, and suppose that $`W`$ is not universal. Let $`๐ฉ`$ witness this, i.e., $`๐ฉ`$ is a (set or proper class sized) premouse, and if $`(๐ฏ,๐ฐ)`$ denotes the coiteration of $`W`$ with $`๐ฉ`$ then both $`๐ฏ`$ and $`๐ฐ`$ have length $`\mathrm{OR}+1`$ and there is a club $`C\mathrm{OR}`$ of strong limit cardinals such that $`D^๐ฐ[\mathrm{min}(C),\mathrm{OR}]_U=\mathrm{}`$, $`\pi _\alpha \mathrm{}^๐ฐ\alpha =\mathrm{id}`$, and $`\pi _{\alpha \beta }^๐ฐ(\alpha )=\beta `$ for all $`\alpha \beta C`$, and $`D^๐ฏ[0,\mathrm{OR}]_T=\mathrm{}`$ and $`\pi _0\mathrm{}^๐ฏ\mathrm{"}\beta \beta `$ for all $`\beta C`$.
Case 1. There is some $`\beta C`$ such that $`\pi _{0\beta }^๐ฏ(\beta )>\beta `$.
Fix such $`\beta C`$. As $`\pi _{0\beta }^๐ฏ\mathrm{"}\beta \beta `$ but $`\pi _{0\beta }^๐ฏ(\beta )>\beta `$, we must have that $`\mu =\mathrm{cf}^W(\beta )<\beta `$ is measurable in $`W`$ and that $`\mu `$ is used in the iteration giving $`\pi _{0\beta }^๐ฏ`$. Let $`\alpha `$ be least in $`[0,\beta )_T`$ such that $`\stackrel{~}{\mu }=\pi _{0\alpha }^๐ฏ(\mu )`$ is the critical point of $`\pi _{\alpha \beta }^๐ฏ`$. Then
$$\mathrm{cf}^{_\alpha ^๐ฏ}(\beta )=\stackrel{~}{\mu }.$$
Pick $`f_\alpha ^๐ฏ`$, $`f:\stackrel{~}{\mu }\beta `$ cofinal. Then $`\pi _{\alpha \mathrm{OR}}^๐ฏ(f)\stackrel{~}{\mu }:\stackrel{~}{\mu }\beta `$ is cofinal, too. (For $`\xi <\stackrel{~}{\mu }`$ do we have $`\pi _\alpha \mathrm{}^๐ฏ(f)(\xi )=\pi _\alpha \mathrm{}^๐ฏ(f(\xi ))<\beta `$, as $`\pi _0\mathrm{}^๐ฏ\mathrm{"}\beta \beta `$; and of course $`\pi _\alpha \mathrm{}^๐ฏ(f(\xi ))f(\xi )`$.) But $`\pi _\alpha \mathrm{}^๐ฏ(f)\stackrel{~}{\mu }_{\mathrm{}}^๐ฏ`$, and thus $`\beta `$ is singular in $`_{\mathrm{}}^๐ฏ`$.
On the other hand, $`\beta `$ is of course inaccessible in $`_{\mathrm{}}^๐ฐ`$, and $`_{\mathrm{}}^๐ฐ\mathrm{}_{\mathrm{}}^๐ฏ`$. This is a contradiction!
We may hence assume that:
Case 2. All $`\beta C`$ are such that $`\pi _{0\beta }^๐ฏ(\beta )=\beta `$.
Claim 1. There are a club class $`D^{}\mathrm{OR}`$ and a commutative system $`(\sigma _{\alpha \beta }:\alpha \beta D^{})`$ such that $`\sigma _{\alpha \beta }:๐ฅ_{\alpha ^{+W}}^W_{\mathrm{\Sigma }_0}๐ฅ_{\beta ^{+W}}^W`$ cofinally with $`\sigma _{\alpha \beta }\alpha =\mathrm{id}`$ and $`\sigma _{\alpha \beta }(\alpha )=\beta `$ for all $`\alpha \beta D^{}`$.
Proof. Let us write $`W_\alpha =๐ฅ_{\alpha ^{+W}}^W`$ and $`W_\alpha ^{}=\pi _{0\alpha }^๐ฏ(W_\alpha )`$ for $`\alpha C`$. Notice that
$$W_\alpha ^{}=๐ฅ_{\alpha ^+}^{_\alpha ^๐ฐ}\mathrm{where}\alpha ^+=\alpha ^{+_\alpha ^๐ฐ},$$
so that we have that
$$\pi _{\alpha \beta }^๐ฐW_\alpha ^{}:W_\alpha ^{}_{\mathrm{\Sigma }_0}W_\beta ^{}\mathrm{cofinally},$$
for $`\alpha \beta C`$. But we clearly also have that
$$\pi _{0\beta }^๐ฏW_\beta :W_\beta _{\mathrm{\Sigma }_0}W_\beta ^{}\mathrm{cofinally}$$
for $`\beta C`$. We aim to show that typically $`\mathrm{ran}(\pi _{\alpha \beta }^๐ฐW_\alpha ^{}\pi _{0\alpha }^๐ฏW_\alpha )\mathrm{ran}(\pi _{0\beta }^๐ฏW_\beta )`$, so that $`(\pi _{0\beta }^๐ฏW_\beta )^1\pi _{\alpha \beta }^๐ฐW_\alpha ^{}\pi _{0\alpha }^๐ฏW_\alpha `$ makes sense.
Let $`\tau _\alpha =\mathrm{OR}W_\alpha `$, $`\alpha C`$. The point is that $`\mathrm{cf}(\tau _\alpha )=\mathrm{cf}(\tau _\beta )`$ for all $`\alpha `$, $`\beta C`$. Let $`\gamma =\mathrm{cf}(\tau _\alpha )`$, $`\alpha C`$, and pick $`X_\alpha =\{\xi _0^\alpha <\xi _1^\alpha <\mathrm{}<\xi _i^\alpha <\mathrm{}:i<\gamma \}\tau _\alpha `$ unbounded in $`\tau _\alpha `$ and of order type $`\gamma `$ for all $`\alpha C`$. We may define a regressive function $`\delta :C\{\alpha :\mathrm{cf}(\alpha )>\gamma \}\mathrm{OR}`$ by letting $`\delta (\alpha )`$ be the least $`\delta <\alpha `$ such that
$$\pi _{0\alpha }^๐ฏ\mathrm{"}X_\alpha \mathrm{ran}(\pi _{\delta \alpha }^๐ฐ).$$
We may further define a regressive function $`h:C\{\alpha :\mathrm{cf}(\alpha )>\gamma \}\mathrm{OR}`$ by letting
$$h(\alpha )=(\pi _{\delta (\alpha )\alpha }^๐ฐ)^1\mathrm{"}\pi _{0\alpha }^๐ฏ\mathrm{"}X_\alpha .$$
By Fodor, there are an ordinal $`\delta _0`$, some $`Y\mathrm{OR}`$, and some unbounded $`DC`$ with $`\delta (\alpha )=\delta _0`$ and $`h(\alpha )=Y`$, all $`\alpha D`$. We then have that
$$\pi _{\alpha \beta }^๐ฐ\pi _{0\alpha }^๐ฏ(\xi _i^\alpha )=\pi _{0\beta }^๐ฏ(\xi _i^\alpha )$$
for all $`\alpha \beta D`$ and $`i<\gamma `$.
We now verify that
$$\pi _{\alpha \beta }^๐ฐW_\alpha ^{}\pi _{0\alpha }^๐ฏW_\alpha \pi _{0\beta }^๐ฏW_\beta $$
for all $`\alpha \beta D`$.
Well, let $`\alpha \beta D`$ and $`xW_\alpha `$. Then $`x๐ฅ_{\xi _i^\alpha }^W`$ for some $`i<\gamma `$. Let $`fW_\alpha `$ be the least (in the order of constructibility) surjection $`f:\alpha ๐ฅ_{\xi _i^\alpha }^W`$. So $`x=f(\nu )`$ , some $`\nu <\alpha `$. We have that $`\pi _{\alpha \beta }^๐ฐ\pi _{0\alpha }^๐ฏ(f)=`$ the least surjection
$$g:\beta ๐ฅ_{\pi _{\alpha \beta }^๐ฐ\pi _{0\alpha }^๐ฏ(\xi _i^\alpha )}^{W_\beta ^{}}=๐ฅ_{\pi _{0\beta }^๐ฏ(\xi _i^\beta )}^{W_\beta ^{}},$$
so that in particular $`\pi _{\alpha \beta }^๐ฐ\pi _{0\alpha }^๐ฏ(f)\mathrm{ran}(\pi _{0\beta }^๐ฏ)`$, say $`\pi _{\alpha \beta }^๐ฐ\pi _{0\alpha }^๐ฏ(f)=\pi _{0\beta }^๐ฏ(g^{})`$.
We then get $`\pi _{\alpha \beta }^๐ฐ\pi _{0\alpha }^๐ฏ(x)=\pi _{\alpha \beta }^๐ฐ\pi _{0\alpha }^๐ฏ(f(\nu ))=\pi _{0\beta }^๐ฏ(g^{})(\pi _{0\alpha }^๐ฏ(\nu ))`$, as $`\pi _{0\alpha }^๐ฏ\mathrm{"}\alpha \alpha `$ and $`\pi _{\alpha \beta }^๐ฐ\alpha =\mathrm{id}`$, $`=\pi _{0\beta }^๐ฏ(g^{})(\pi _{0\beta }^๐ฏ(\nu ))`$, as $`\pi _{\alpha \beta }^๐ฏ\alpha =\mathrm{id}`$, $`=\pi _{0\beta }^๐ฏ(g^{}(\nu ))`$.
Now let
$$D^{}=D\{\alpha :\lambda \mathrm{Lim}\alpha =\mathrm{sup}(D\lambda )\}.$$
$`D^{}`$ is club. In order to finish the proof of Claim 1 it suffices to verify the following.
Subclaim. If $`\{\beta <\alpha :\beta D\}`$ is unbounded in $`\alpha `$, and if $`(\stackrel{~}{W},\sigma _\beta :\beta D\alpha )`$ is the direct limit of the system
$$(W_\beta ,((\pi _{0\beta }^๐ฏ)^1\pi _{\gamma \beta }^๐ฐ\pi _{0\gamma }^๐ฏ:\gamma \beta D\alpha )$$
then we have that $`\sigma _\beta :W_\beta W_\alpha `$ is cofinal with $`\sigma _\beta \beta =\mathrm{id}`$ and $`\sigma _\beta (\beta )=\alpha `$.
Proof. Let $`\delta =\mathrm{min}(D\alpha +1)`$. There is $`\sigma :\stackrel{~}{W}W_\delta `$ defined by
$$x(\pi _{0\delta }^๐ฏ)^1\pi _{\beta \delta }^๐ฐ\pi _{0\beta }^๐ฏ(\overline{x})\mathrm{where}\overline{x}=\sigma _\beta ^1(x).$$
It is clear that $`\sigma \alpha =\mathrm{id}`$, and $`\sigma (\alpha )=\delta `$. By \[9, ยง8 Lemma 4\] we hence have that ($`\stackrel{~}{W}`$ is transitive and) $`\stackrel{~}{W}\mathrm{}W_\alpha `$.
Now consider the two maps $`\pi _{0\alpha }^๐ฏ:W_\alpha W_\alpha ^{}`$ and $`\stackrel{~}{\sigma }:\stackrel{~}{W}W_\alpha ^{}`$, where
$$\stackrel{~}{\sigma }(x)=\pi _{\beta \alpha }^๐ฐ\pi _{0\beta }^๐ฏ(\overline{x})\mathrm{for}\overline{x}=\sigma _\beta ^1(x).$$
We may define $`\tau :W_\alpha ^{}\pi _{0\alpha }^๐ฏ(\stackrel{~}{W})`$ by setting
$$\tau (\stackrel{~}{\sigma }(f)(\xi ))=\pi _{0\alpha }^๐ฏ(f)(\xi )$$
for $`\xi <\alpha `$ and appropriate $`f`$ with $`\mathrm{dom}(f)`$ bounded below $`\alpha `$. $`\tau `$ is easily seen to be well-defined, and in fact surjective.
But then $`\pi _{0\alpha }^๐ฏ(\stackrel{~}{W})=W_\alpha ^{}`$, and so $`W_\alpha =\stackrel{~}{W}`$.
$`\mathrm{}`$ (Subclaim)
$`\mathrm{}`$ (Claim 1)
Now fix $`D^{}`$ and maps $`\sigma _{\alpha \beta }`$ as in Claim 1. For each $`\alpha D^{}`$, let $`\alpha ^{}=\mathrm{min}(D^{}\alpha +1)`$, and let $`F_\alpha `$ the extender derived from $`\sigma _{\alpha \alpha ^{}}`$.
Claim 2. There is some $`\beta D^{}`$ such that $`F_\beta `$ is countably complete.
Proof. Let $`D^{\prime \prime }=\{\alpha D^{}:D^{}\alpha `$ is unbounded in $`\alpha \}`$. Notice that $`D^{\prime \prime }`$ is club. Suppose that $`F_\alpha `$ is not countably complete for any $`\alpha D^{\prime \prime }`$. Pick for any $`\alpha D^{\prime \prime }`$ sequences $`(a_n^\alpha :n<\omega )`$ and $`(X_n^\alpha :n<\omega )`$ witnessing this. This means that $`a_n^\alpha F_\alpha (X_n^\alpha )`$ for all natural numbers $`n`$, but there is no orderpreserving function $`\tau :_{n<\omega }a_n^\alpha \alpha =\mathrm{c}.\mathrm{p}.(F_\alpha )`$ with the property that for all natural numbers $`n`$, we have $`\tau \mathrm{"}a_n^\alpha X_n^\alpha `$.
Let $`\alpha D^{\prime \prime }`$. Let $`g^\alpha :_{n<\omega }a_n^\alpha \mathrm{otp}(_{n<\omega }a_n^\alpha )<\omega _1`$ denote the transitive collapse, and let
$$g(\alpha )=(\mathrm{otp}(\underset{n<\omega }{}a_n^\alpha ),(g^\alpha \mathrm{"}a_n^\alpha :n<\omega )),$$
i.e., $`g(\alpha )`$ tells us how the $`a_n^\alpha `$ sit inside $`_{n<\omega }a_n^\alpha `$. Let $`\beta (\alpha )`$ be the least $`\beta D^{}`$ such that for all natural numbers $`n`$ do we have that $`X_n^\alpha \mathrm{ran}(\pi _{\beta \alpha })`$. Notice that $`\beta (\alpha )<\alpha `$, as $`\mathrm{cf}(\alpha )>\omega `$.
We may now apply Fodorโs lemma to the function $`F`$ with $`\mathrm{dom}(F)=D^{\prime \prime }`$ defined by
$$F(\alpha )=(\beta (\alpha ),(\pi _{\beta (\alpha )\alpha }^1(X_n^\alpha ):n<\omega ),g(\alpha ))$$
to get some unbounded $`ED^{\prime \prime }`$ on which $`F`$ is constant. Let $`\alpha <\gamma E`$. As $`g(\gamma )=g(\alpha )`$, we have that $`\tau =(g^\alpha )^1g^\gamma :_{n<\omega }a_n^\gamma _{n<\omega }a_n^\alpha `$ is order preserving with $`\tau \mathrm{"}a_n^\gamma =a_n^\alpha `$ for all natural numbers $`n`$. But we now get, by $`F(\gamma )=F(\alpha )`$ together with the commutativity of the maps $`\pi _{\beta \beta ^{}}`$, that for all $`n<\omega `$,
$$\tau \mathrm{"}a_n^\gamma =a_n^\alpha F_\alpha (X_n^\alpha )=F_\alpha (\pi _{\alpha \gamma }^1(X_n^\gamma ))=X_n^\gamma [F_\alpha (\alpha )]^{<\omega },$$
which contradicts the choice of $`(a_n^\gamma ,X_n^\gamma ,n<\omega )`$.
$`\mathrm{}`$ (Claim 2)
But Claim 2 now plainly contradicts the fact that $`W`$ is supposed to be full.
$`\mathrm{}`$ (4.4)
The following is an immediate corollary to 4.2 and 4.4.
###### Corollary 4.5
$`(\neg 0^{^{}})`$ $`K^c`$ is universal.
Let $`\mathrm{\Gamma }\mathrm{}`$ be a class of regular cardinals, $`\omega \mathrm{\Gamma }`$. In the verification of Claim 2 during the proof of 4.4 we might have considered the stationary class $`D^{\prime \prime }\mathrm{\Gamma }`$ rather than $`D^{\prime \prime }`$, thus deducing that there is some $`\beta D^{}`$ such that $`F_\beta `$ is countably complete and $`\mathrm{cf}^V(\mathrm{c}.\mathrm{p}.(F_\beta ))\mathrm{\Gamma }`$. This leads to the following.
###### Corollary 4.6
$`(\neg 0^{^{}})`$ Let $`\mathrm{\Gamma }\mathrm{}`$ be a class of regular cardinals, $`\omega \mathrm{\Gamma }`$. Then $`K_\mathrm{\Gamma }^c`$ is universal.
It is however not true that an initial segment of $`K^c`$ whose height is a cardinal in $`V`$ greater than $`\mathrm{}_1`$ is universal for coiterable premice of at most the same height.<sup>4</sup><sup>4</sup>4Compare with \[22, Theorem 3.4\]. Anticipating the theory of $`K`$, here is an example. Suppose that $`K`$ has a measurable cardinal, and let $`\mu `$ be the least one. Let $`\lambda >\mu `$ be, in $`K`$, a singular cardinal of cofinality $`\mu `$. Suppose that $`V=K^{Col(\omega ,\mu )}`$. Then $`๐ฅ_{\lambda ^+}^{K^c}`$ does not win the coiteration against $`๐ฅ_{\lambda ^+}^K`$, because if $`F`$ is an extender on $`K`$ with critical point $`\mu `$ then $`\mathrm{Ult}_0(๐ฅ_{\lambda ^+}^K;F)`$ will have height $`>\lambda ^+`$. We leave the further (easy) details to the reader.
###### Lemma 4.7
(Goodness of $`K^c`$.) Let $`\mathrm{\Gamma }\mathrm{}`$ be a class of regular cardinals, $`\omega \mathrm{\Gamma }`$. Let $`\kappa `$ be a cardinal of $`K_\mathrm{\Gamma }^c`$, and let $`๐ซ`$ be an iterate of $`K_\mathrm{\Gamma }^c`$ above $`\kappa `$, i.e., there is an iteration tree $`๐ฏ=((_\alpha ^๐ฏ,\pi _{\alpha \beta }^๐ฏ:\alpha T\beta \theta +1),(E_\alpha ^๐ฏ:\alpha <\theta ),T)`$ on $`K_\mathrm{\Gamma }^c`$ such that $`\mathrm{c}.\mathrm{p}.(E_\alpha ^๐ฏ)\kappa `$ for every $`\alpha <\theta `$ and $`๐ซ=_\theta ^๐ฏ`$. Let $`F=E_\alpha ^๐ซ\mathrm{}`$ be such that $`\alpha >\kappa `$, and $`\mu =\mathrm{c}.\mathrm{p}.(F)<\kappa `$ (notice that $`\alpha >\mu ^{+K_\mathrm{\Gamma }^c}`$ and $`F`$ is total on $`๐ซ`$).
Then $`F`$ is countably complete.
Proof. This is shown by revisiting the arguments for 3.2 and 3.3. Suppose without loss of generality that $`K_\mathrm{\Gamma }^c=K^c`$. Let $`((a_n,X_n):n<\omega )`$ be given such that
$$a_nF(X_n)\mathrm{for}\mathrm{all}n<\omega .$$
Let $`\lambda ^+`$, a successor cardinal in $`V`$, be large enough such that $`๐ฏ`$ and $`F`$ are both hereditarily smaller than $`\lambda ^+`$. We may then construe $`๐ฏ`$ as an iteration tree on $`๐ฅ_{\lambda ^+}^{K^c}`$, and we shall have that $`F`$ is on the final model of $`๐ฏ`$ so construed. Let $`\xi `$ be such that $`๐ฅ_{\lambda ^+}^{K^c}=_\xi `$.
Running the proofs of 3.2 and 3.3, we may then pick the map
$$\sigma :_\xi ,$$
such that in fact $`\sigma `$ is the restriction of an uncollapse (also called $`\sigma `$) of a countable submodel of a large enough initial segment of $`V`$ containing all objects of current interest; in particular, we want that both $`\{a_n:n<\omega \}`$ and $`\{X_n:n<\omega \}`$ are contained in $`\mathrm{ran}(\sigma )`$.
In the end, we get an embedding
$$\sigma ^{}:^{}_\xi ^{}$$
for some $`\xi ^{}\xi `$, where $`^{}`$ is the final model of $`\sigma ^1(๐ฏ)`$. Moreover, as $`๐ฏ`$ is above $`\kappa `$, $`\sigma ^1(๐ฏ)`$ is above $`\sigma ^1(\kappa )`$, which implies that
$$\sigma ^{}\sigma ^1(\kappa )=\sigma \sigma ^1(\kappa ).$$
As $`\mathrm{c}.\mathrm{p}.(F)<\kappa `$, $`\sigma ^{}\sigma ^1(F)`$ is countably complete. We may thus pick an order-preserving $`\rho `$ such that
$$\rho \mathrm{"}\sigma ^{}\sigma ^1(a_n)\sigma ^{}\sigma ^1(X_n)\mathrm{for}\mathrm{all}n<\omega .$$
But $`\sigma ^{}\sigma ^1(X_n)=X_n`$, as $`\mu <\kappa `$ and $`\sigma ^{}\sigma ^1(\kappa )=\sigma \sigma ^1(\kappa )`$, which means that
$$\rho \mathrm{"}\sigma ^{}\sigma ^1(a_n)X_n\mathrm{for}\mathrm{all}n<\omega ,$$
and $`\rho \sigma ^{}\sigma ^1`$ witnesses that $`F`$ is countably complete with respect to $`((a_n,X_n):n<\omega )`$.
$`\mathrm{}`$ (4.7)
## 5 Generalized bicephali.
We let $`C_0`$ denote the class of all limit cardinals $`\kappa `$ of $`V`$ such that $`๐ฅ_\kappa ^{K^c}`$ is universal for coiterable premice of height $`<\kappa `$. I.e., if $`\kappa C_0`$, $``$ is a premouse with $`\mathrm{OR}<\kappa `$, and $`๐ฏ`$, $`๐ฐ`$ are the iteration trees arising from the successful comparison of $`๐ฅ_\kappa ^{K^c}`$ with $``$, then $`๐^๐ฐ(0,\mathrm{}]_U=\mathrm{}`$, and $`_{\mathrm{}}^๐ฐ_{\mathrm{}}^๐ฏ`$.
###### Lemma 5.1
$`(\neg 0^{^{}})`$ $`C_0`$ is closed unbounded in $`\mathrm{OR}`$.
Proof. $`C_0`$ is trivially be seen to be closed. Suppose that $`C_0\eta `$ for some $`\eta \mathrm{OR}`$. We may then define $`F:\{\kappa :\kappa `$ is a limit cardinal $`\}\eta \mathrm{OR}`$ by $`F(\kappa )=`$ the least $`\alpha `$ such that there is a coiterable premouse $``$ of height $`\alpha <\kappa `$ and $``$ wins the coiteration against $`๐ฅ_\kappa ^{K^c}`$. By Fodor, there is an unbounded class $`A\mathrm{OR}`$ such that $`F\mathrm{"}A=\{\alpha \}`$ for some $`\alpha \mathrm{OR}`$. There are at most $`2^{\mathrm{Card}(\alpha )}`$ many premice of height $`\alpha `$, so that by 4.5 there is some cardinal $`\gamma `$ such that $`๐ฅ_\gamma ^{K^c}`$ wins the coiteration against all of them which are coiterable with $`K^c`$. This gives a contradiction!
$`\mathrm{}`$ (5.1)
###### Lemma 5.2
$`(\neg 0^{^{}})`$ Let $`\kappa C_0`$, and let $`\mathrm{}๐ฅ_\kappa ^{K^c}`$ be a $`0`$-iterable premouse. Let $`F=E_\nu ^{}\mathrm{}`$ be such that $`\nu \kappa `$ and $`\mathrm{c}.\mathrm{p}.(F)<\kappa `$. Then $`F`$ is countably complete.
Proof. Let $`\kappa `$, $``$, $`F`$, and $`\nu `$ be as in the statement of 5.2. Set $`\mu =\mathrm{c}.\mathrm{p}.(F)`$. Let $`((a_n,X_n):n<\omega )`$ be such that $`a_n[F(\mu )]^{<\omega }`$, $`X_n๐ซ([\mu ]^{\mathrm{Card}(a_n)})`$, and $`a_nF(X_n)`$ for every $`n<\omega `$. We aim to find an order-preserving $`\tau :_{n<\omega }a_n\mu `$ such that $`\tau \mathrm{"}a_nX_n`$ for every $`n<\omega `$. Let
$$\sigma :\overline{}$$
be fully elementary such that $`\sigma \mu ^++1=\mathrm{id}`$, $`\mathrm{Card}(\overline{})<\mu ^{++}`$, and $`\{a_n:n<\omega \}\mathrm{ran}(\sigma )`$. Let $`\overline{F}=\sigma ^1(F)`$ in case $`F`$ (where we then assume without loss of generality that $`F\mathrm{ran}(\sigma )`$), and let $`\overline{F}`$ be the top extender of $`\overline{}`$ otherwise (i.e., if $`F`$ is the top extender of $``$).
Let $`๐ฏ_0`$ and $`๐ฏ_1`$ be $`0`$-maximal iteration trees on the phalanx
$$((,\overline{}),\mathrm{c}.\mathrm{p}.(\sigma ))$$
and on $``$, respectively, stemming from the coiteration of that phalanx with $``$. Also let $`๐ฐ`$ and $`๐ฏ`$ be the $`0`$-maximal iteration trees arising from the comparison of $`\overline{}`$ with $``$, respectively. Let $`^{}`$ be the last model of $`๐ฏ_0`$. Then $`\overline{}T_0^{}`$ by \[9, ยง8 Lemma 1\]. But this implies that $`๐ฐ`$ only uses extenders whose critical points are less than or equal to $`\mathrm{c}.\mathrm{p}.(\sigma )`$, by 2.7.
Obviously, the coiteration of $`\overline{}`$ with $`๐ฅ_\kappa ^{K^c}`$ is an initial segment of the coiteration $`๐ฐ`$, $`๐ฏ`$. Moreover, as $`\kappa C_0`$ and $`\mathrm{Card}(\overline{})<\kappa `$, $`๐ฅ_\kappa ^{K^c}`$ wins the coiteration against $`\overline{}`$. Hence $`๐ฏ`$ actually only uses extenders from $`๐ฅ_\kappa ^{K^c}`$ and its images, and $`๐ฐ`$ does not drop.
This means that the main branch of $`๐ฐ`$ gives us an embedding
$$\pi :\overline{}_{\mathrm{\Sigma }_0}_{\mathrm{}}^๐ฐ,$$
where $`_{\mathrm{}}^๐ฐ`$ is an initial segment of $`_{\mathrm{}}^๐ฏ`$. Notice that $`\pi \mu ^++1=\mathrm{id}`$ as $`๐ฐ`$ is above $`\mathrm{c}.\mathrm{p}.(\sigma )`$. Let $`\stackrel{~}{F}=\pi (\overline{F})`$ in case $`\overline{F}\overline{}`$, and let $`\stackrel{~}{F}`$ be the top extender of $`_{\mathrm{}}^๐ฐ`$ otherwise.
We claim that $`\stackrel{~}{F}`$ is countably complete. This trivially follows from 4.7 if $`๐ฏ`$ is above $`\mu ^+`$. Otherwise let $`E_\nu ^{_i^๐ฏ}`$ be the first extender used on $`๐ฏ`$ with critical point $`\mu `$. By Claim 1 in the proof of 2.4, $`\nu `$ is then strictly greater than the index of $`\stackrel{~}{F}`$. But then 4.7 applied to $`๐ฏi`$ yields that $`\stackrel{~}{F}`$ is countably complete.
Notice that $`\pi \sigma ^1(a_n)\stackrel{~}{F}(X_n)`$ for all $`n<\omega `$. Now let $`\overline{\tau }:_{n<\omega }\pi \sigma ^1(a_n)\mu `$ be order-preserving and such that $`\overline{\tau }\mathrm{"}\pi \sigma ^1(a_n)X_n`$ for all $`n<\omega `$. Putting $`\tau =\overline{\tau }\pi \sigma ^1`$ hence gives us a witness as desired.
$`\mathrm{}`$ (5.2)
We now have to turn towards our theory of โbicephali.โ It turned out to be most convenient to let a (pre-)bicephalus be a pair of premice rather than a premouse with two top extenders.
###### Definition 5.3
An ordered pair $`๐ฉ=(๐ฉ^0,๐ฉ^1)`$ is called a generalized prebicephalus provided the following hold.
(a) $`๐ฉ^0=(J_\alpha [\stackrel{}{E}];,\stackrel{}{E},F^0)`$ is a premouse with $`F^0\mathrm{}`$,
(b) $`๐ฉ^1=(J_\beta [\stackrel{}{E}^{}];,\stackrel{}{E}^{},F^1)`$ is a premouse with $`\beta \alpha `$ and $`F^1\mathrm{}`$,
(c) $`\alpha `$ is a cardinal in $`๐ฉ^1`$ in case $`\beta >\alpha `$,
(d) $`\stackrel{}{E}^{}\alpha =\stackrel{}{E}`$, i.e., $`(J_\alpha [\stackrel{}{E}^{}\alpha ];,\stackrel{}{E}^{}\alpha )=(J_\alpha [\stackrel{}{E}];,\stackrel{}{E})`$, and
(e) $`\mathrm{c}.\mathrm{p}.(F^0)\mathrm{c}.\mathrm{p}.(F^1)<F^0(\mathrm{c}.\mathrm{p}.(F^0))`$.
In this case $`๐ฉ^0`$ is called the left part of $`๐ฉ`$, $`๐ฉ^1`$ is called the right part of $`๐ฉ`$, and $`๐ฉ`$ is called the generalized prebicephalus derived from $`๐ฉ^0`$, $`๐ฉ^1`$.
Notice that if $`\beta >\alpha `$ then $`\alpha `$ has to be a successor cardinal in $`๐ฉ^1`$.
We aim to prove that g-prebicephali trivialize, i.e., that $`๐ฉ^0=๐ฉ^1`$ provided that the generalized prebicephalus $`๐ฉ=(๐ฉ^0,๐ฉ^1)`$ meets a certain iterability criterion. Here is an immediate trivial observation:
###### Lemma 5.4
Let $`๐ฉ=(๐ฉ^0,๐ฉ^1)`$ be a generalized prebicephalus, and let $``$ be a premouse. Suppose that $`๐ฉ^0\mathrm{}`$ as well as $`๐ฉ^1\mathrm{}`$. Then $`๐ฉ^0=๐ฉ^1`$.
Proof. The hypothesis tells us that $`๐ฉ^0\mathrm{}๐ฉ^1`$. In particular $`F^0`$, the top extender of $`๐ฉ^0`$, is on the sequence of $`๐ฉ^1`$. Assume that $`๐ฉ^0๐ฉ^1`$, i.e., $`\alpha =๐ฉ^0\mathrm{OR}<๐ฉ^1\mathrm{OR}`$. Then $`F^0=E_\alpha ^{๐ฉ^1}`$, and so $`\rho _1(๐ฉ^0)<\alpha `$. But $`\alpha `$ is supposed to be a cardinal of $`๐ฉ^1`$ (cf. (c) of 5.3). Contradiction! Hence $`\alpha =๐ฉ^1\mathrm{OR}`$, and $`๐ฉ^0=๐ฉ^1`$.
$`\mathrm{}`$ (5.4)
###### Definition 5.5
Let $`๐ฉ=(๐ฉ^0,๐ฉ^1)`$ be a generalized prebicephalus. Then
$$๐ฏ=((๐ฉ_\alpha ^0,\pi _{\alpha \beta }^{๐ฉ^0},๐ฉ_\alpha ^1,\pi _{\alpha \beta }^{๐ฉ^1}:\alpha T\beta <\theta ),(E_\alpha :\alpha +1<\theta ),T)$$
is an unpadded iteration tree on $`๐ฉ`$ provided the following hold.
(a) $`T`$ is a tree order (in the sense of \[20, Def. 5.0.1\]),
(b) $`๐ฉ_0^0=๐ฉ^0`$ and $`๐ฉ_0^1=๐ฉ^1`$,
(c) for all $`\alpha <\theta `$, $`(๐ฉ_\alpha ^0,๐ฉ_\alpha ^1)`$ is a generalized prebicephalus, or else $`๐ฉ_\alpha ^0=๐ฉ_\alpha ^1`$,
(d) for all $`\alpha +1<\theta `$, $`E_\alpha \mathrm{}`$, and $`E_\alpha =E_\nu ^{๐ฉ_\alpha ^0}`$ for some $`\nu ๐ฉ_\alpha ^0\mathrm{OR}`$, or else $`E_\alpha `$ is the top extender of $`๐ฉ_\alpha ^1`$,
(e) for all $`\alpha +1<\theta `$, $`T`$-pred$`(\alpha +1)=`$ the least $`\beta \alpha `$ such that $`\mathrm{c}.\mathrm{p}.(E_\alpha )<\mathrm{min}\{E_\beta (\mathrm{c}.\mathrm{p}.(E_\beta )),๐ฉ_\beta ^0\mathrm{OR}\}`$,
(f) for all $`\alpha +1<\theta `$ and $`h=0,1`$, if $`\beta =T`$-pred$`(\alpha +1)`$ then
$$\pi _{\beta \alpha +1}^{๐ฉ^h}:๐ฅ_{\eta ^h}^{๐ฉ_\beta ^h}_{E_\alpha }๐ฉ_{\alpha +1}^h,$$
where $`\eta ^h๐ฉ_\beta ^h`$ is maximal such that $`E_\alpha `$ measures all the subsets of its critical point in $`๐ฅ_{\eta ^h}^{๐ฉ_\beta ^h}`$ (where we understand that $`\mathrm{deg}^๐ฏ(\alpha +1)=0`$ if $`๐^๐ฏ[0,\alpha +1]=\mathrm{}`$ and $`\mathrm{deg}^๐ฏ(\alpha +1)=`$ that $`n`$ such that $`\mathrm{dom}(\pi _{\beta \alpha +1}^{๐ฉ^h})`$ is $`n`$-sound if $`๐^๐ฏ[0,\alpha +1]\mathrm{}`$),
(g) if $`\alpha <\theta `$ is a limit ordinal then $`(๐ฉ_\alpha ^h,(\pi _{\beta \alpha }^{๐ฉ^h}:\beta T\alpha ))`$ is the transitive direct limit of $`(๐ฉ_\beta ^h,\pi _{\gamma \beta }^{๐ฉ^h}:\gamma T\beta <\alpha )`$, for $`h=0,1`$,
(h) for all $`\alpha <\theta `$, the set
$$๐^๐ฏ(0,\alpha ]_T=\{\beta +1T\alpha :\mathrm{dom}(\pi _{Tpred(\beta +1)\beta +1}^{๐ฉ^h})๐ฉ_{Tpred(\beta +1)}^h\}$$
is finite, for $`h=0,1`$, and
(i) for all $`\alpha T\beta <\theta `$, $`\pi _{\alpha \beta }^{๐ฉ^0}=\pi _{\alpha \beta }^{๐ฉ^1}\mathrm{dom}(\pi _{\alpha \beta }^{๐ฉ^0})`$.
By a โputativeโ iteration tree on $`๐ฉ`$ we shall mean a tree $`๐ฏ`$ of successor length $`\theta =\overline{\theta }+1`$ which is as in 5.5 except for the fact that possibly (c) fails for $`\alpha =\overline{\theta }`$ or (h) fails for $`\beta =\overline{\theta }`$. It is crucial, but straightforward, that if $`๐ฏ`$ is a putative iteration tree on $`๐ฉ`$ such that $`๐ฉ_{\overline{\theta }}^0`$ and $`๐ฉ_{\overline{\theta }}^1`$ are both transitive and (h) holds for $`๐ฏ`$, then $`๐ฏ`$ is in fact an iteration tree (to get (i) one uses (c) of 5.3).
For our purposes, we shall also consider padded iteration trees on g-prebicephali; as usual, this just means that the indexing of the models is slowed down by possible repetition of models. The reader will have no trouble with modifying 5.5 accordingly.
Let
$$๐ฏ=((๐ฉ_\alpha ^0,\pi _{\alpha \beta }^{๐ฉ^0},๐ฉ_\alpha ^1,\pi _{\alpha \beta }^{๐ฉ^1}:\alpha T\beta <\theta ),(E_\alpha :\alpha +1<\theta ),T)$$
be a(n) (putative) iteration tree on the generalized prebicephalus $`๐ฉ`$. Then weโll write $`_\beta ^๐ฏ`$ for $`(๐ฉ_\beta ^0,๐ฉ_\beta ^1)`$ if $`๐^๐ฏ(0,\beta ]_T=\mathrm{}`$, and weโll write $`_\beta ^๐ฏ`$ for $`๐ฉ_\beta ^0=๐ฉ_\beta ^1`$ if $`๐^๐ฏ(0,\beta ]_T\mathrm{}`$. We shall also use $`(_\beta ^๐ฏ)^0`$ for $`๐ฉ_\beta ^0`$, and $`(_\beta ^๐ฏ)^1`$ for $`๐ฉ_\beta ^1`$. Weโll write $`\pi _{\alpha \beta }^๐ฏ`$ for $`\pi _{\alpha \beta }^{๐ฉ^1}`$, and $`E_\beta ^๐ฏ`$ for $`E_\beta `$.
###### Definition 5.6
Let $`๐ฉ=(๐ฉ^0,๐ฉ^1)`$ be a generalized prebicephalus, and let $``$ be a premouse. We want to determine a pair $`(_0^{๐ฉ,},_1^{๐ฉ,})`$, the โleast disagreementโ of $`๐ฉ`$ and $``$.
Case 1. $`๐ฉ^0`$ and $``$ are not lined up. Let $`\nu `$ be least such that
$$E_\nu ^{๐ฉ^0}E_\nu ^{},$$
and let
$$_0^{๐ฉ,}=E_\nu ^{๐ฉ^0}\mathrm{and}_1^{๐ฉ,}=E_\nu ^{}.$$
Case 2. Not Case 1, but $`๐ฉ^1`$ and $``$ are not lined up. Then let $`_0^{๐ฉ,}`$ be the top extender of $`๐ฉ^1`$, and let $`_1^{๐ฉ,}`$ be the top extender of $`๐ฉ^0`$ (sic!).
Case 3. Neither Case 1 nor Case 2. Then we let $`(_0^{๐ฉ,},_1^{๐ฉ,})=(\mathrm{},\mathrm{})`$.
Case 2 needs a brief discussion. Weโll have that $`๐ฉ^0\mathrm{}`$, as $`๐ฉ^0`$ and $``$ are lined up but $`๐ฉ^1`$ and $``$ are not. In particular, the top extender of $`๐ฉ^0`$ appears on the sequence of $``$. But the top extender of $`๐ฉ^1`$ canโt appear on the sequence of $``$, as $`๐ฉ^1`$ and $``$ are not lined up.
The content of 5.4 may now be stated as saying that if $``$ is a premouse with $`(_0^{๐ฉ,},_1^{๐ฉ,})=(\mathrm{},\mathrm{})`$ then $`๐ฉ^0`$ or else $`๐ฉ^0=๐ฉ^1`$.
###### Definition 5.7
We call a generalized prebicephalus $`๐ฉ`$ a generalized bicephalus if it is coiterable with $`K^c`$ in the following sense.
Let $`๐ฏ`$, $`๐ฐ`$ be putative padded iteration trees of successor length $`\theta +1`$ on $`๐ฉ`$, $`K^c`$ given by the following characterization. If $`\beta <\theta `$ then
$$(E_\beta ^๐ฏ,E_\beta ^๐ฐ)=(_0^{_\beta ^๐ฏ,_\beta ^๐ฐ},_1^{_\beta ^๐ฏ,_\beta ^๐ฐ})$$
in case $`_\beta ^๐ฏ`$ is a generalized prebicephalus, and let otherwise
$$(E_\beta ^๐ฏ,E_\beta ^๐ฐ)=(E_\nu ^{_\beta ^๐ฐ},E_\nu ^{_\beta ^๐ฏ})$$
where $`\nu `$ is least with $`E_\nu ^{_\beta ^๐ฐ}E_\nu ^{_\beta ^๐ฏ}`$. Then $`๐ฏ`$ is an iteration tree on $`๐ฉ`$.
We want to emphasize that by 5.6 we have that the $`๐ฉ`$-side of the coiteration of the generalized prebicephalus $`๐ฉ`$ with $`K^c`$, being performed as in 5.7, will only use extenders which are allowed by 5.5 (d).
###### Lemma 5.8
Let $`๐ฉ`$, $`๐ฏ`$, $`๐ฐ`$ and $`\theta `$ be as in 5.7. Then $`๐ฐ`$ is normal. Moreover, for all $`\beta +1<\theta `$ do we have that the index of $`E_\beta ^๐ฏ`$ is $``$ the index of $`E_\beta ^๐ฐ`$.
Proof. The only thing to notice here is that if $`(E_\beta ^๐ฏ,E_\beta ^๐ฐ)`$ is chosen according to Case 2 in 5.6 and $`\nu `$ denotes the index of $`E_\beta ^๐ฐ`$ then $`_{\beta +1}^๐ฐ`$, $`(_{\beta +1}^๐ฏ)^0`$, and $`(_{\beta +1}^๐ฏ)^1`$ pairwise agree up to $`\nu +1`$. This holds because then $`\nu =(_\beta ^๐ฏ)^0\mathrm{OR}`$, so that either $`\nu `$ is the index of $`E_\beta ^๐ฏ`$ or else $`\nu `$ is a (successor) cardinal in $`(_\beta ^๐ฏ)^1`$. In any event,
$$E_\nu ^{(_{\beta +1}^๐ฏ)^0}=E_\nu ^{(_{\beta +1}^๐ฏ)^1}=\mathrm{}=E_\nu ^{_{\beta +1}^๐ฐ}.$$
$`\mathrm{}`$ (5.8)
###### Lemma 5.9
$`(\neg 0^{^{}})`$ Let $`๐ฉ=(๐ฉ^0,๐ฉ^1)`$ be a generalized bicephalus. Then $`๐ฉ^0=๐ฉ^1`$.
Proof. Let $`๐ฏ`$, $`๐ฐ`$ denote the (padded) iteration trees arising from the coiteration of $`๐ฉ`$ with $`K^c`$ built as in 5.7, where
$$๐ฏ=((๐ฉ_\alpha ^0,\pi _{\alpha \beta }^{๐ฉ^0},๐ฉ_\alpha ^1,\pi _{\alpha \beta }^{๐ฉ^1}:\alpha \beta <\theta ),(E_\alpha :\alpha +1<\theta ),T).$$
Claim. $`\mathrm{lh}(๐ฏ)=\mathrm{lh}(๐ฐ)<\mathrm{OR}`$.
Given this Claim, the proof of 5.9 can be completed as follows. By the proof of 4.4 we shall have that $`_{\mathrm{}}^๐ฏ`$ is a generalized prebicephalus, and $`(_{\mathrm{}}^๐ฏ)^0\mathrm{}_{\mathrm{}}^๐ฐ`$ as well as $`(_{\mathrm{}}^๐ฏ)^1\mathrm{}_{\mathrm{}}^๐ฐ`$. This implies that $`(_{\mathrm{}}^๐ฏ)^0=(_{\mathrm{}}^๐ฏ)^1`$ by 5.4. Now suppose that $`๐ฉ^0๐ฉ^1`$, and let $`w๐ฉ^0`$ be such that
$$wF^0wF^1,$$
where $`F^0`$ and $`F^1`$ are the top extenders of $`๐ฉ^0`$ and $`๐ฉ^1`$. (Notice that such $`w`$ would have to exist!) We have that $`\mathrm{dom}(\pi _0\mathrm{}^{๐ฉ^0})=๐ฉ^0`$, $`\mathrm{dom}(\pi _0\mathrm{}^{๐ฉ^1})=๐ฉ^1`$, and $`\pi _0\mathrm{}^{๐ฉ^0}=\pi _0\mathrm{}^{๐ฉ^1}๐ฉ^0`$. If we let $`\stackrel{~}{F}^0`$ and $`\stackrel{~}{F}^1`$ denote the top extenders of $`๐ฉ_{\mathrm{}}^0`$ and $`๐ฉ_{\mathrm{}}^1`$ then $`\stackrel{~}{F}^0=\stackrel{~}{F}^1`$, and thus
$$wF^0\pi _0\mathrm{}^{๐ฉ^0}(w)\stackrel{~}{F}^0\pi _0\mathrm{}^{๐ฉ^1}(w)\stackrel{~}{F}^1wF^1.$$
Contradiction!
Proof of the Claim. Let us assume that $`\mathrm{lh}(๐ฏ)=\mathrm{lh}(๐ฐ)=\mathrm{OR}+1`$. The proof of 4.4 yields that then $`\pi _0\mathrm{}^๐ฐ\mathrm{"}\mathrm{OR}\mathrm{OR}`$, so that usual arguments give some $`\lambda [0,\mathrm{OR})_T[0,\mathrm{OR})_U`$ so that
(1)
$$\mathrm{c}.\mathrm{p}.(\pi _\lambda \mathrm{}^๐ฏ)=\mathrm{c}.\mathrm{p}.(\pi _\lambda \mathrm{}^๐ฐ)=\lambda \mathrm{and}$$
(2)
$$\pi _\lambda \mathrm{}^๐ฏ๐ซ(\lambda )_\lambda ^๐ฏ=\pi _\lambda \mathrm{}^๐ฐ๐ซ(\lambda )_\lambda ^๐ฐ.$$
Let $`\alpha +1`$ be least in $`(\lambda ,\mathrm{OR}]_T`$, and let $`\beta +1`$ be least in $`(\lambda ,\mathrm{OR}]_U`$. Let $`E_\beta ^๐ฐ=E_{\nu _1}^{_\beta ^๐ฐ}`$, and let $`E_\alpha ^๐ฏ=E_{\nu _0}^{(_\alpha ^๐ฏ)^0}`$ or $`E_\alpha ^๐ฏ=E_{\nu _0}^{(_\alpha ^๐ฏ)^1}`$ (if the former is not true but the latter, then $`_\alpha ^๐ฏ`$ is a generalized prebicephalus and $`E_\alpha ^๐ฏ`$ has to be the top extender of $`(_\alpha ^๐ฏ)^1`$). Of course, $`\lambda =\mathrm{c}.\mathrm{p}.(E_\beta ^๐ฐ)=\mathrm{c}.\mathrm{p}.(E_\alpha ^๐ฏ)`$.
Case 1. $`E_\alpha ^๐ฏ=E_{\nu _0}^{(_\alpha ^๐ฏ)^0}`$.
In this case, the rest of the coiteration is beyond $`\nu _0+1`$ on both sides, i.e., all $`E_\gamma ^๐ฏ`$ and $`E_\gamma ^๐ฐ`$ for $`\gamma >\alpha `$ have index $`>\nu _0`$. Moreover, by 5.5 (e) weโll have that $`\pi _{\alpha +1\mathrm{}}^๐ฏE_\alpha ^๐ฏ(\lambda )=\mathrm{id}`$. For this reason, and because $`๐ฐ`$ is normal by 5.8, (1) and (2) give that $`E_\alpha ^๐ฏ`$ and $`E_\beta ^๐ฐ`$ are compatible.
Case 1.1. $`\nu _0<\nu _1`$.
Then $`E_\alpha ^๐ฏ๐ฅ_{\nu _1}^{_\beta ^๐ฐ}`$ by the initial segment condition for $`_\beta ^๐ฐ`$. But $`๐ฐ`$ is normal by 5.8, so that $`E_\alpha ^๐ฏ_{\mathrm{}}^๐ฐ`$, i.e., $`E_\alpha ^๐ฏ๐ฅ_{\mathrm{OR}}^_{\mathrm{}}^๐ฐ`$. On the other hand $`E_\alpha ^๐ฏ(_{\alpha +1}^๐ฏ)^0`$, which by $`\pi _{\alpha +1\mathrm{}}^๐ฏE_\alpha ^๐ฏ(\lambda )=\mathrm{id}`$ implies that $`E_\alpha ^๐ฏ(_{\mathrm{}}^๐ฏ)^0`$. But of course we have $`๐ฅ_{\mathrm{OR}}^_{\mathrm{}}^๐ฐ=๐ฅ_{\mathrm{OR}}^{(_{\mathrm{}}^๐ฏ)^0}`$. Contradiction!
Case 1.2. $`\nu _1<\nu _0`$.
We then have $`E_\beta ^๐ฐ๐ฅ_{\nu _0}^{(_\alpha ^๐ฏ)^0}`$ by the initial segment condition for $`(_\alpha ^๐ฏ)^0`$. So by $`\pi _{\alpha +1\mathrm{}}^๐ฏE_\alpha ^๐ฏ(\lambda )=\mathrm{id}`$ we know that $`E_\beta ^๐ฐ(_{\mathrm{}}^๐ฏ)^0`$. On the other hand, $`E_\beta ^๐ฐ_{\beta +1}^๐ฐ`$, which by the normality of $`๐ฐ`$ implies that $`E_\beta ^๐ฐ_{\mathrm{}}^๐ฐ`$. We thus have a contradiction as in case 1.1!
Case 1.3. $`\nu _1=\nu _0`$.
Then we get $`\alpha \beta `$ by (the proof of) 5.8; but then as $`E_\gamma ^๐ฐ`$ has index $`>\nu _0`$ for $`\gamma >\alpha `$ we must have $`\alpha =\beta `$. Now by the normality of $`๐ฐ`$ and by $`\pi _{\alpha +1\mathrm{}}^๐ฏE_\alpha ^๐ฏ(\lambda )=\mathrm{id}`$ we must have $`E_\alpha ^๐ฏ=E_\alpha ^๐ฐ`$. This clearly contradicts the choice of $`(E_\alpha ^๐ฏ,E_\alpha ^๐ฐ)`$ as $`(_0^{_\alpha ^๐ฏ,_\alpha ^๐ฐ},_1^{_\alpha ^๐ฏ,_\alpha ^๐ฐ})`$.
Case 2. $`_\alpha ^๐ฏ`$ is a generalized prebicephalus and $`E_\alpha ^๐ฏ`$ is the top extender of $`(_\alpha ^๐ฏ)^1`$.
In this case, by 5.6, $`E_\alpha ^๐ฐ`$ is the top extender of $`(_\alpha ^๐ฏ)^0`$. Let $`\nu ^{}`$ be the index of $`E_\alpha ^๐ฐ`$. So $`\nu ^{}=(_\alpha ^๐ฏ)^0\mathrm{OR}`$.
If $`\alpha <\beta `$ then by 5.3 (e), as $`_\alpha ^๐ฏ`$ is a generalized prebicephalus, and by 5.8 we would get that
$$\mathrm{c}.\mathrm{p}.(E_\alpha ^๐ฐ)\mathrm{c}.\mathrm{p}.(E_\alpha ^๐ฏ)=\lambda =\mathrm{c}.\mathrm{p}.(E_\beta ^๐ฐ)<\nu ^{}<\nu _1.$$
However $`K^c`$ is below $`0^{^{}}`$, and $`๐ฏ`$ is a normal iteration tree on $`K^c`$. We thus have a contradiction with Claim 1 in the proof of 2.4.
Thus $`\alpha \beta `$. If $`\alpha >\beta `$ then, because $`\mathrm{c}.\mathrm{p}.(E_\alpha ^๐ฐ)\mathrm{c}.\mathrm{p}.(E_\alpha ^๐ฏ)=\lambda =\mathrm{c}.\mathrm{p}.(E_\beta ^๐ฐ)`$, we would get $`\beta +1(0,\mathrm{})_U`$, which is nonsense.
We thus have that $`\alpha =\beta `$. In particular, $`\nu _1=\nu ^{}=(_\alpha ^๐ฏ)^0\mathrm{OR}`$. Hence by 5.5 (e) weโll have that
(3)
$$\pi _{\alpha +1\mathrm{}}^๐ฏ\mathrm{min}\{E_\alpha ^๐ฏ(\lambda ),\nu _1\}=\mathrm{id}.$$
Because $`๐ฐ`$ is normal by 5.8 weโll have that
(4)
$$\pi _{\alpha +1\mathrm{}}^๐ฐE_\alpha ^๐ฐ(\lambda )=\mathrm{id}.$$
Recall that $`E_\alpha ^๐ฐ`$ is the top extender of $`(_\alpha ^๐ฏ)^0`$, and $`E_\alpha ^๐ฏ`$ is the top extender of $`(_\alpha ^๐ฏ)^1`$. In particular,
(5)
$$E_\alpha ^๐ฐ(\lambda )\mathrm{min}\{E_\alpha ^๐ฏ(\lambda ),\nu _1\}.$$
But now (1) to (5) will tell us that
(6)
$$E_\alpha ^๐ฐ=E_\alpha ^๐ฏ|E_\alpha ^๐ฐ(\lambda ).$$
The initial segment condition for $`(_\alpha ^๐ฏ)^1`$ then yields $`E_\alpha ^๐ฐ=E_\alpha ^๐ฏ`$, or else $`E_\alpha ^๐ฐ(_\alpha ^๐ฏ)^1`$.
Case 2.1. $`E_\alpha ^๐ฐ(_\alpha ^๐ฏ)^1`$.
In this case, weโll have that $`E_\alpha ^๐ฐ(_{\alpha +1}^๐ฏ)^1`$, too. By (3) and (5) we have that $`\pi _{\alpha +1\mathrm{}}^๐ฏE_\alpha ^๐ฐ(\lambda )=\mathrm{id}`$, and thus $`(X,Y)E_\alpha ^๐ฐ`$ if and only if
$$\stackrel{~}{Y}[(X,\stackrel{~}{Y})\pi _{\alpha +1\mathrm{}}^๐ฏ(E_\alpha ^๐ฐ)Y=\stackrel{~}{Y}E_\alpha ^๐ฐ(\lambda )].$$
Hence $`E_\alpha ^๐ฐ(_{\mathrm{}}^๐ฏ)^1`$.
On the other hand, the normality of $`๐ฐ`$ yields that $`E_\alpha ^๐ฐ_{\mathrm{}}^๐ฐ`$. However, $`๐ฅ_{\mathrm{OR}}^_{\mathrm{}}^๐ฐ=๐ฅ_{\mathrm{OR}}^{(_{\mathrm{}}^๐ฏ)^1}`$. This is a contradiction!
Case 2.2. $`E_\alpha ^๐ฐ=E_\alpha ^๐ฏ`$.
But $`E_\alpha ^๐ฐ=E_\alpha ^๐ฏ`$ of course contradicts the choice of $`(E_\alpha ^๐ฏ,E_\alpha ^๐ฐ)`$ as $`(_0^{_\alpha ^๐ฏ,_\alpha ^๐ฐ},_1^{_\alpha ^๐ฏ,_\alpha ^๐ฐ})`$.
$`\mathrm{}`$ (Claim)
$`\mathrm{}`$ (5.9)
We want to point out that 5.9 would no longer be true if we dropped (e) in the definition 5.3 (there are easy counterexamples).
###### Lemma 5.10
$`(\neg 0^{^{}})`$ Let $`๐ฉ=(๐ฉ^0,๐ฉ^1)`$ be a generalized prebicephalus, and let $`\mu `$ be the critical point of $`๐ฉ^1`$. Let $`\kappa >\mu `$ be a cardinal in $`๐ฉ^0`$ (and hence in $`๐ฉ^1`$, too). Suppose that for $`h\{0,1\}`$ do we have the following.
For all iteration trees $`๐ฑ`$ on $`๐ฉ^h`$ with last model $`_{\mathrm{}}^๐ฑ`$ and for all $`๐ฉ^h`$-cardinals $`\rho \kappa `$ such that either
(a) $`\rho <\kappa `$, and $`๐ฑ`$ lives on $`๐ฅ_\kappa ^{๐ฉ^h}`$ and is above $`\rho `$, but $`๐ฑ`$ doesnโt use extenders which are total on $`๐ฉ^h`$, or else
(b) $`\rho =\kappa `$, and $`๐ฑ`$ is above $`\rho `$,
we have that
($``$) if $`F=E_\nu ^_{\mathrm{}}^๐ฑ\mathrm{}`$ is such that $`\nu >\rho `$ and $`\mathrm{c}.\mathrm{p}.(F)<\rho `$, then $`F`$ is countably complete.
Then $`๐ฉ`$ is a generalized bicephalus.
Notice that if $`๐ฉ`$ satisfies the assumption in the statement of 5.10 then in particular every $`E_\nu ^{๐ฉ^h}\mathrm{}`$ with $`\mathrm{c}.\mathrm{p}.(E_\nu ^{๐ฉ^h})<\kappa `$ and which is total on $`๐ฉ^h`$ is countably complete. (Let $`๐ฑ=`$ the trivial tree, and $`\rho =\mathrm{c}.\mathrm{p}.(E_\nu ^{๐ฉ^h})^{+๐ฉ^h}`$.)
Proof of 5.10. We have to show that $`๐ฉ`$ is iterable. Let $`๐ฏ`$ be a putative iteration tree on $`๐ฉ`$ arising from the comparison with $`K^c`$. By the proof of 5.9, $`๐ฏ`$ is a set. Pick an elementary embedding $`\sigma :NH_\theta `$ (for some large enough regular $`\theta `$) with $`N`$ countable and transitive, and such that $`\{๐ฏ,\kappa \}\mathrm{ran}(\sigma )`$. Let $`\overline{๐ฏ}=\sigma ^1(๐ฏ)`$, and $`\overline{\kappa }=\sigma ^1(\kappa )`$. It suffices to embed the models of $`\overline{๐ฏ}`$ into transitive structures. For notational convenience, we shall assume that $`\overline{๐ฏ}`$ is unpadded. Say,
$$\overline{๐ฏ}=(((_\alpha ^{\overline{๐ฏ}})^0,\pi _{\alpha \beta }^0,(_\alpha ^{\overline{๐ฏ}})^1,\pi _{\alpha \beta }^1:\alpha \overline{T}\beta <\mathrm{lh}(\overline{T})),(E_\alpha :\alpha +1<\mathrm{lh}(\overline{T})),\overline{T}).$$
Put $`\kappa _\alpha =\mathrm{c}.\mathrm{p}.(E_\alpha )`$ for $`\alpha +1<\mathrm{lh}(\overline{T})`$. We set $`n(0)=0`$, and for $`\alpha >0`$ we let $`n(\alpha )`$ be such that
$$\beta (\beta \overline{T}\alpha \gamma +1(\beta ,\alpha ]_{\overline{T}}\eta _{\gamma +1}^๐ฏ=\mathrm{Ult}_{n(\alpha )}(๐ฅ_\eta ^{_{Tpred(\gamma +1)}^๐ฏ};E_\gamma ^๐ฏ))).$$
Notice that for all $`\alpha <\mathrm{lh}(๐ฏ)`$ and $`h\{0,1\}`$ do we have that $`(_\alpha ^{\overline{๐ฏ}})^h`$ is $`n(\alpha )`$-sound.
We shall in fact determine some $`\vartheta \mathrm{lh}(\overline{๐ฏ})`$, and define, for both $`h\{0,1\}`$, sequences $`\stackrel{}{๐ฑ}^h=(๐ฑ_i^h:i<\vartheta )`$ of iteration trees on $`๐ฉ^h`$ such that any of the $`๐ฑ_i^h`$โs is as in (a) or (b) of 5.10; for this purpose weโll simultaneously define a sequence $`(\rho _i:i<\vartheta )`$ of ordinals. We shall also construct maps $`\sigma _\alpha ^h`$ from any $`(_\alpha ^{\overline{๐ฏ}})^h`$ into a model of some $`๐ฑ_i^h`$. For all $`i<\vartheta `$ weโll have that $`๐ฑ_i^0`$ and $`๐ฑ_i^1`$ are given by the very same sequence of extenders; in particular, $`\mathrm{lh}(๐ฑ_i^0)=\mathrm{lh}(๐ฑ_i^1)`$. We shall index the models of $`๐ฑ_i^h`$ in a non-standard way, namely, we shall start counting them with
$$\mathrm{}(i)=\underset{j<i}{}\mathrm{lh}(๐ฑ_j^h)$$
(where $``$ denotes ordinal summation), so that the models of $`๐ฑ_i^h`$ will be indexed by the elements of $`[\mathrm{}(i),\mathrm{}(i)+\mathrm{lh}(๐ฑ_i^h))`$, and weโll conveniently have that
$$\sigma _\alpha ^h:(_\alpha ^{\overline{๐ฏ}})^h_\alpha ^{๐ฑ_i^h}$$
for that $`i<\vartheta `$ such that $`\mathrm{}(i)\alpha <\mathrm{}(i)+\mathrm{lh}(๐ฑ_i^h)`$. It will be clear from the construction that always
R 0<sub>ฮฑ</sub> $`\sigma _\alpha ^0=\sigma _\alpha ^1(_\alpha ^{\overline{๐ฏ}})^0`$.
Letting $`h`$ range over $`\{0,1\}`$, we shall inductively maintain that the following requirements are met as well.
R 1<sub>ฮฑ</sub> $`\beta [0,\alpha ]_{\overline{T}}`$, if $`i`$ is maximal with $`\mathrm{}(i)\beta `$ then
$$\sigma _\beta ^h:(_\beta ^{\overline{๐ฏ}})^h_\beta ^{๐ฑ_i^h}$$
is a weak $`n(\beta )`$-embedding.
R 2<sub>ฮฑ</sub> if $`i`$ is maximal with $`\mathrm{}(i)\alpha `$ then $`\mathrm{}(i)\overline{T}\alpha `$, and for all $`\beta [0,\mathrm{}(i)]_{\overline{T}}`$ do we have that $`ji(\beta =\mathrm{}(j)_\beta ^{๐ฑ_j^h}=๐ฉ^h)`$.
R 3<sub>ฮฑ</sub> for all $`\beta +1(0,\alpha ]_{\overline{T}}`$, setting $`\beta ^{}=\overline{T}`$-pred$`(\beta +1)`$, we have that $`\sigma _\alpha ^h`$ agrees with $`\sigma _\beta ^{}^h`$ up to $`\kappa _\beta `$; i.e.,
$$๐ฅ_{\sigma _\alpha ^h(\kappa _\beta )}^{_\alpha ^{๐ฑ_i^h}}=๐ฅ_{\sigma _\alpha ^h(\kappa _\beta )}^{_\beta ^{}^{๐ฑ_j^h}},\mathrm{and}$$
$$\sigma _\alpha ^h๐ฅ_{\kappa _\beta }^{(_\alpha ^{\overline{๐ฏ}})^h}=\sigma _\beta ^h๐ฅ_{\kappa _\beta }^{(_\beta ^{}^{\overline{๐ฏ}})^h},$$
where $`i`$ is maximal with $`\mathrm{}(i)\alpha `$ and $`j`$ is maximal with $`\mathrm{}(j)\beta ^{}`$.
R 4<sub>ฮฑ</sub> if $`i`$ is maximal with $`\mathrm{}(i)\alpha `$, $`\beta \overline{T}\gamma [0,\alpha )_{\overline{T}}`$ and $`๐^{\overline{๐ฏ}}(\beta ,\gamma ]_{\overline{T}}=\mathrm{}`$ then we have:
(a) $`\mathrm{}(i)\overline{T}\beta \overline{T}\gamma \sigma _\gamma ^h\pi _{\beta \gamma }^h=\pi _{\beta \gamma }^{๐ฑ_i^h}\sigma _\beta ^h`$, and
(b) $`\beta \overline{T}\gamma \overline{T}\mathrm{}(i)\sigma _\gamma ^h\pi _{\beta \gamma }^h=\sigma _\beta ^h`$.
R 5<sub>ฮฑ</sub> if $`i`$ is maximal with $`\mathrm{}(i)\alpha `$, then for all $`ji`$ do we have that $`๐ฑ^{}=๐ฑ_j^h\mathrm{min}\{\mathrm{}(j)+\mathrm{lh}(๐ฑ_j^h),\alpha +1\}`$ is an iteration tree on $`๐ฉ^h`$ above $`\rho _j`$ as in (a) or (b) of 5.10; i.e., either
$`\rho _j<\kappa `$, and $`๐ฑ^{}`$ lives on $`๐ฅ_\kappa ^{๐ฉ^h}`$ and is above $`\rho _j`$, but $`๐ฑ^{}`$ doesnโt use extenders which are total on $`๐ฉ^h`$, or else
$`\rho _j=\kappa `$, and $`๐ฑ^{}`$ is above $`\rho _j`$.
We are now going to run our construction. In what follows we again let $`h`$ range over $`\{0,1\}`$. Put $`\mathrm{}(0)=0`$. To commence, set $`\sigma _0^h=\sigma (_0^{\overline{๐ฏ}})^h`$ (notice $`(_0^{\overline{๐ฏ}})^h=\sigma ^1(๐ฉ^h)`$), and let $`๐ฑ_0^h1`$ be trivial. It is clear that R 1<sub>0</sub> to R 5<sub>0</sub> hold.
Now suppose that we have defined everything up to $`\alpha `$, that is, suppose that for some $`i`$ are we given
$$\stackrel{}{๐ฑ}^hi=(๐ฑ_j^h:j<i),๐ฑ_i^h\alpha +1,(\sigma _\beta ^h:\beta \alpha ),(\rho _j:j<i),\mathrm{and}\rho _i\mathrm{if}\alpha >\mathrm{}(i)$$
in such a way that R 1<sub>ฮฑ</sub> through R 5<sub>ฮฑ</sub> are satisfied. Suppose that $`\alpha <\mathrm{lh}(\overline{๐ฏ})`$, and set $`\overline{F}=E_\alpha ^{\overline{๐ฏ}}`$, and $`\overline{\mu }=\mathrm{c}.\mathrm{p}.(\overline{F})`$. (If $`\alpha =\mathrm{lh}(\overline{๐ฏ})`$ weโre done with our construction.) Then either $`\overline{F}=E_\nu ^{(_\alpha ^{\overline{๐ฏ}})^0}`$ for some $`\nu (_\alpha ^{\overline{๐ฏ}})^0\mathrm{OR}`$, or else $`\overline{F}`$ is the top extender of $`(_\alpha ^{\overline{๐ฏ}})^1`$. Depending on whether the former holds, or the latter, we set $`F=\sigma _\alpha ^0(\overline{F})`$ and $`h^{}=0`$, or $`F=\sigma _\alpha ^1(\overline{F})`$ and $`h^{}=1`$. Set $`\mu ^{}=\mathrm{c}.\mathrm{p}.(F)`$.
Case 1. $`\alpha =\mathrm{}(i)\mu ^{}<\kappa `$ $`\overline{F}`$ is total on $`(_\alpha ^{\overline{๐ฏ}})^h^{}`$, or $`\alpha >\mathrm{}(i)\mu ^{}<\rho _i`$.
By R 5<sub>ฮฑ</sub> and our assumptions on $`๐ฉ`$, we get that $`F`$ is countably complete in this case. Let $`\tau :\sigma _\alpha ^h^{}\mathrm{"}\overline{F}(\overline{\mu })\mu ^{}`$ be order preserving such that for appropriate $`a`$, $`X\mathrm{ran}(\sigma _\alpha ^h^{})`$ we have that $`aF(X)\tau \mathrm{"}aX`$. We now declare $`๐ฑ_i^h=๐ฑ_i^h\alpha +1`$, i.e., $`\mathrm{}(i+1)=\alpha +1`$. If $`\alpha =\mathrm{}(i)`$ we also put $`\rho _i=0`$ (we wonโt be interested in this value).
By 5.5 and R 5<sub>ฮฑ</sub> it is easy to see that $`\overline{T}`$-pred$`(\alpha +1)\mathrm{}(i)`$, so that by R 2<sub>ฮฑ</sub> we know that $`ji\mathrm{}(j)=\overline{T}`$-pred$`(\alpha +1)`$. By R 3<sub>ฮฑ</sub>, $`\sigma _\alpha ^h`$ agrees with $`\sigma _{\mathrm{}(j)}^h`$ up to $`\kappa _\beta `$, where $`\beta +1`$ is least in $`(\mathrm{}(j),\alpha ]_{\overline{T}}`$. We have $`\kappa _\alpha <\kappa _\beta `$ by $`\neg 0^{^{}}`$ (cf. Claim 2 in the proof of 2.4), so that $`\sigma _\alpha ^h`$ agrees with $`\sigma _{\mathrm{}(j)}^h`$ up to $`\kappa _\alpha ^+`$ (calculated in $`(_\alpha ^{\overline{๐ฏ}})^h`$). By R 0<sub>ฮฑ</sub>, hence, $`\sigma _\alpha ^h^{}`$ agrees with $`\sigma _{\mathrm{}(j)}^h`$ up to $`\kappa _\alpha ^+`$.
We now define $`\sigma _{\alpha +1}^h:(_{\alpha +1}^{\overline{๐ฏ}})^h๐ฉ^h`$ by setting
$$[a,f]\sigma _{\mathrm{}(j)}^h(f)(\tau \mathrm{"}\sigma _\alpha ^h^{}(a)).$$
This is well-defined and $`\mathrm{\Sigma }_0`$-elementary, as we may reason as follows. Let $`\mathrm{\Phi }`$ be a $`\mathrm{\Sigma }_0`$ formula. Then
$$(_{\alpha +1}^{\overline{๐ฏ}})^h\mathrm{\Phi }([a_1,f_1],\mathrm{},[a_k,f_k])$$
$$(a_1,\mathrm{},a_k)\overline{F}(\{(u_1,\mathrm{},u_k):(_{\mathrm{}(j)}^{\overline{๐ฏ}})^h\mathrm{\Phi }(f_1(u_1),\mathrm{},f_k(u_k))\})$$
$`(\sigma _\alpha ^h^{}(a_1),\mathrm{},\sigma _\alpha ^h^{}(a_k))`$
$$F(\sigma _\alpha ^h^{}(\{(u_1,\mathrm{},u_k):(_{\mathrm{}(j)}^{\overline{๐ฏ}})^h\mathrm{\Phi }(f_1(u_1),\mathrm{},f_k(u_k))\})),$$
which, by the amount of agreement of $`\sigma _\alpha ^h^{}`$ with $`\sigma _{\mathrm{}(j)}^h`$ and by R 1<sub>ฮฑ</sub>, holds if and only if
$`(\sigma _\alpha ^h^{}(a_1),\mathrm{}\sigma _\alpha ^h^{}(a_k))`$
$$F(\{(u_1,\mathrm{},u_k):๐ฉ^h\mathrm{\Phi }(\sigma _{\mathrm{}(j)}^h(f_1)(u_1),\mathrm{},\sigma _{\mathrm{}(j)}^h(f_k)(u_k))\})$$
$$๐ฉ^h\mathrm{\Phi }(\sigma _{\mathrm{}(j)}^h(f_1)(\tau \mathrm{"})\sigma _\alpha ^h^{}(a_1),\mathrm{},\sigma _{\mathrm{}(j)}^h(f_k)(\tau \mathrm{"})\sigma _\alpha ^h^{}(a_k)).$$
It is straightforward to check that now R 1<sub>ฮฑ+1</sub> through R 5<sub>ฮฑ+1</sub> hold.
Case 2. $`\alpha =\mathrm{}(i)\mu ^{}\kappa `$, or $`\alpha =\mathrm{}(i)\mu ^{}<\kappa \overline{F}`$ is partial on $`(_\alpha ^{\overline{๐ฏ}})^h^{}`$, or $`\alpha >\mathrm{}(i)\mu ^{}\rho _i`$.
In this case, we start or continue copying $`\overline{๐ฏ}[\mathrm{}(i),\alpha +2)`$ onto $`๐ฉ^h`$, getting $`๐ฑ_i^h\alpha +2`$. We let $`๐ฑ_i^h`$-pred$`(\alpha +1)=\overline{๐ฏ}`$-pred$`(\alpha +1)`$, call it $`\alpha ^{}`$. Weโll have $`\mathrm{}(i)\alpha ^{}`$. We then use the shift lemma \[20, Lemma 5.2\] to get $`\pi _{\alpha ^{}\alpha +1}^{๐ฑ_i^h}`$ together with the copy map $`\sigma _{\alpha +1}^h`$. It is easy to check that R 1<sub>ฮฑ+1</sub> to R 5<sub>ฮฑ+1</sub> hold.
We put $`\rho _i=\kappa `$ if $`\alpha =\mathrm{}(i)\mu ^{}\kappa `$, and we put $`\rho _i=`$ the cardinality of $`\mathrm{c}.\mathrm{p}.(F)`$ in $`๐ฉ^h^{}`$ if $`\alpha =\mathrm{}(i)\mu ^{}<\kappa \overline{F}`$ is partial on $`(_\alpha ^{\overline{๐ฏ}})^h^{}`$.
Now let $`\lambda `$ be a limit ordinal and suppose that we have defined everything up to $`\lambda `$. To state this more precisely, we have to consider two cases.
Case 1. There are cofinally in $`\lambda `$ many $`\alpha <\lambda `$ such that there is a $`j`$ with $`\alpha =\mathrm{}(j)`$.
In this case for some $`i`$ are we given
$$\stackrel{}{๐ฑ}^hi=(๐ฑ_j^h:j<i),(\sigma _\beta ^h:\beta <\lambda ),\mathrm{and}(\rho _j:j<i)$$
in such a way that R 1<sub>ฮฑ</sub> through R 5<sub>ฮฑ</sub> are satisfied for every $`\alpha <\lambda `$. We then declare $`\mathrm{}(i)=\lambda `$, and we define
$$\sigma _\lambda ^h:(_\lambda ^{\overline{๐ฏ}})^h๐ฉ^h$$
by setting
$$\sigma _\lambda ^h(x)=\sigma _\beta ^h((\pi _{\beta \lambda }^h)^1(x)),\mathrm{where}x\mathrm{ran}(\pi _{\beta \lambda }^h).$$
This works by $`\alpha <\lambda (`$ R 2<sub>ฮฑ</sub> and R 4<sub>ฮฑ</sub> (b) $`)`$. It is easy to check that R 1<sub>ฮป</sub> to R 5<sub>ฮป</sub> hold.
Case 2. Otherwise.
In this case, there is a largest $`i`$ such that $`\mathrm{}(i)<\lambda `$, and we are given
$$\stackrel{}{๐ฑ}^hi=(๐ฑ_j^h:j<i),๐ฑ_i^h\lambda ,(\sigma _\beta ^h:\beta <\lambda ),\mathrm{and}(\rho _j:ji)$$
in such a way that R 1<sub>ฮฑ</sub> to R 5<sub>ฮฑ</sub> are satisfied for every $`\alpha <\lambda `$. We then continue with copying $`\overline{๐ฏ}[\mathrm{}(i),\lambda ]`$ onto $`๐ฉ^h`$, getting $`๐ฑ_i^h\lambda +1`$ and the copy map $`\sigma _\lambda ^h`$. This works by $`\alpha <\lambda `$ R 2<sub>ฮฑ</sub> (a). We leave the easy details to the reader. It is straighforward to check that R 1<sub>ฮป</sub> to R 5<sub>ฮป</sub> hold.
$`\mathrm{}`$ (5.10)
In practice weโll always know that the hypothesis of 5.10 is satisfied. An example is given in the proof of the following.
###### Lemma 5.11
$`(\neg 0^{^{}})`$ Let $`\kappa `$ be a limit cardinal in $`V`$ such that for all $`\lambda <\kappa `$ we have that
$$\lambda \mathrm{is}<\kappa \mathrm{strong}\mathrm{in}K^c\lambda \mathrm{is}<\mathrm{OR}\mathrm{strong}\mathrm{in}K^c,$$
and $`\kappa C_0`$. Let $`\mathrm{}๐ฅ_\kappa ^{K^c}`$ be a premouse. Suppose that $`1n<\omega `$ is such that
(a) $`\rho _n()\kappa <\rho _{n1}()`$,
(b) $``$ is $`n`$-sound above $`\kappa `$ (i.e., $``$ is $`(n1)`$-sound, and $`^{n1}`$ is generated by $`h_{}^{n1}`$ from $`\kappa \{p_{,n}\}`$), and
(c) $``$ is $`(n1)`$-iterable.
Then $`K^c`$, i.e., $``$ is an initial segment of $`K^c`$.
Proof. We first verify the following.
Claim 1. $`\rho _n()\kappa `$, and hence $``$ is $`n`$-sound.
Proof. Suppose not. Then $`\rho _n()<\kappa `$, so that $`_n()`$ has size $`<\kappa `$. Note that $`_n()`$ is $`(n1)`$-iterable and $`n`$-sound. By $`\kappa C_0`$, there is an $`(n1)`$\- maximal tree $`๐ฏ`$ on $`_n()`$ with $`๐^๐ฏ(0,\mathrm{}]_T=\mathrm{}`$ and such that $`_{\mathrm{}}^๐ฏ`$ is a non-simple iterate of $`๐ฅ_\kappa ^{K^c}`$ (it is non-simple by $`\rho _n()<\kappa `$). Moreover, the $`(n1)`$-maximal coiteration $`(\overline{๐ฐ},๐ฐ)`$ of $`_n()`$, $``$ is simple on both sides and produces some common coiterate $`Q=_{\mathrm{}}^{\overline{๐ฐ}}=_{\mathrm{}}^๐ฐ`$. Using $`\pi =\pi _0\mathrm{}^๐ฏ`$, we may copy $`\overline{๐ฐ}`$ onto $`_{\mathrm{}}^๐ฏ`$, getting an iteration $`\overline{๐ฐ}^\pi `$ of $`_{\mathrm{}}^๐ฏ`$ together with a last copy map
$$\sigma :QQ^{\prime \prime }=_{\mathrm{}}^{\overline{๐ฐ}^\pi }.$$
Now on the one hand we have
$$\sigma \pi _0\mathrm{}^๐ฐ:Q^{\prime \prime },$$
and on the other hand, we have that $`Q^{\prime \prime }`$ is a non-simple iterate of $``$ (as $`_{\mathrm{}}^๐ฏ`$ is). This contradicts the Dodd-Jensen Lemma (cf. \[20, Lemma 5.3\]).
$`\mathrm{}`$ (Claim 1)
Weโll only need that $`\rho _1()\kappa `$ in what follows.
We now use 5.2 to show that the phalanx $`๐ซ=((K^c,),\kappa )`$ is normally $`(n1)`$-iterable. For suppose $`๐ฐ`$ to be a putative normal tree on $`๐ซ`$. By 2.7, we can write
$$๐ฐ=๐ฐ_0{}_{}{}^{}๐ฐ_{1}^{}$$
where $`๐ฐ_0`$ is an iteration of $``$ above $`\kappa `$, and $`๐ฐ_1`$ is an iteration of $`K^c`$ except for the fact that the first extender, call it $`F`$, used for building $`๐ฐ_1`$ comes from the last model of $`๐ฐ_0`$, $`\mathrm{c}.\mathrm{p}.(F)<\kappa `$, and the index of $`F`$ is larger than $`\kappa `$ (possibly, $`๐ฐ_0=\mathrm{}`$, or $`๐ฐ_1=\mathrm{}`$). But by 5.2, $`F`$ is countably complete, and hence $`๐ฐ_1`$ is well-behaved by standard arguments.
Now let $`๐ฐ`$, $`๐ฏ`$ be the iteration trees arising from the comparison of $`๐ซ`$ with $`K^c`$, where we understand $`๐ฐ`$ to be $`(n1)`$-maximal. We have to verify the following.
Claim 2. $`\mathrm{root}^๐ฐ(\mathrm{})=0`$, i.e., $`_{\mathrm{}}^๐ฐ`$, the last model of $`๐ฐ`$, sits above $``$.
Before turning to its proof, let us first show that Claim 2 implies 5.11. Note that by 2.7, then, $`๐ฐ`$ is an iteration of $``$.
Suppose that $`๐ฐ`$ is non-trivial. Then $`_{\mathrm{}}^๐ฐ`$ is not sound, which by 4.5 gives that we must have $`_{\mathrm{}}^๐ฐ=_{\mathrm{}}^๐ฏ`$. This gives a standard contradiction if $`๐^๐ฐ(0,\mathrm{}]_U\mathrm{}`$ (cf. the proof of Claim 4 in the proof of \[20, Lemma 6.2\]). Hence we have to have $`๐^๐ฐ(0,\mathrm{}]_U=\mathrm{}`$. As $`๐ฏ`$ only uses extenders with indices $`\kappa `$ weโll clearly have $`\rho _\omega (_{\mathrm{}}^๐ฏ)\kappa `$. But then $`\rho _\omega (_{\mathrm{}}^๐ฐ)\kappa `$, and so
$$_\omega (_{\mathrm{}}^๐ฐ)=_n(_{\mathrm{}}^๐ฐ)=,$$
and again we get a standard contradiction (as in the proof of Claim 4 in the proof of \[20, Lemma 6.2\]).
We have shown that $`๐ฐ`$ has to be trivial; that is, $`\mathrm{}_{\mathrm{}}^๐ฏ`$. But it can then easily be verified that $`๐ฏ`$ has to be trivial, too. This implies $`K^c`$ as desired.
Proof of Claim 2. Suppose not. That is, suppose that $`\mathrm{root}^๐ฐ(\mathrm{})=1`$, i.e., that the last model of $`๐ฐ`$ sits above $`K^c`$. By 4.5 and the Dodd-Jensen Lemma applied to $`\pi _1\mathrm{}^๐ฐ`$ we easily get that $`๐^๐ฐ[1,\mathrm{}]_U=๐^๐ฏ[0,\mathrm{}]_U=\mathrm{}`$, and
$$_{\mathrm{}}^๐ฐ=_{\mathrm{}}^๐ฏ,\mathrm{call}\mathrm{it}Q.$$
In particular, we have maps $`\pi _1\mathrm{}^๐ฐ:K^cQ`$ and $`\pi _0\mathrm{}^๐ฏ:K^cQ`$. Let $`F^0=E_\alpha ^๐ฐ`$ be the first extender used on $`(1,\mathrm{}]_U`$. Set $`\mu =\mathrm{c}.\mathrm{p}.(F^0)=\mathrm{c}.\mathrm{p}.(\pi _1\mathrm{}^๐ฐ)<\kappa `$. A simple trick gives the following.
Subclaim 1. $`\mathrm{c}.\mathrm{p}.(\pi _0\mathrm{}^๐ฏ)=\mu `$.
Proof. We have to show that $`\pi _0\mathrm{}^๐ฏ\mathrm{id}`$, and that $`\mu `$ is the critical point of $`\pi _0\mathrm{}^๐ฏ`$. Let $`F^0=E_{\nu _0}^{_\alpha ^๐ฐ}`$. As $`๐ฐ`$ only uses extenders with indices $`\kappa `$, and by Claim 1 above, it is clear that $`\rho _1(๐ฅ_{\nu _0}^{_\alpha ^๐ฐ})\kappa `$. Hence 1.3 immediately gives that $`\mu `$ is $`<\kappa `$-strong in $`๐ฅ_{\nu _0}^{_\alpha ^๐ฐ}`$. But $`๐ฅ_\kappa ^{_\alpha ^๐ฐ}=๐ฅ_\kappa ^{K^c}`$, so $`\mu `$ is $`<\kappa `$-strong in $`K^c`$. By our assumption on $`\kappa `$, this means that $`\mu `$ is $`<\mathrm{OR}`$-strong in $`K^c`$. On the other hand, 2.5 gives that $`\mu `$ is not $`<\mathrm{OR}`$-strong in $`Q`$.
This already implies that $`\pi _0\mathrm{}^๐ฏ\mathrm{id}`$. Let $`F^1=E_\beta ^๐ฏ=E_{\nu _1}^{_\beta ^๐ฏ}`$ be the first extender used on $`(0,\mathrm{}]_T`$.
Let us first assume that $`\pi _0\mathrm{}^๐ฏ\mu \mathrm{id}`$, i.e., that $`\mathrm{c}.\mathrm{p}.(F^1)<\mu `$. As $`๐ฏ`$ only uses extenders with indices $`\kappa `$, we clearly have that $`\rho _1(๐ฅ_{\nu _1}^{_\beta ^๐ฏ})\kappa `$. Thus, again by using 1.3 and 2.5, we end up with getting $`\mathrm{c}.\mathrm{p}.(F^1)`$ is $`<\mathrm{OR}`$-strong in $`K^c`$, but $`\mathrm{c}.\mathrm{p}.(F^1)`$ is not $`<\mathrm{OR}`$-strong in $`Q`$. However, consider now $`\pi _1\mathrm{}^๐ฐ:K^cQ`$. As $`\pi _1\mathrm{}^๐ฐ\mathrm{c}.\mathrm{p}.(F^1)+1=\mathrm{id}`$, the fact that $`\mathrm{c}.\mathrm{p}.(F^1)`$ is $`<\mathrm{OR}`$-strong in $`K^c`$ implies that $`\mathrm{c}.\mathrm{p}.(F^1)`$ remains $`<\mathrm{OR}`$-strong in $`Q`$. Contradiction!
A symmetric argument shows that we cannot have $`\pi _0\mathrm{}^๐ฏ\mu +1=\mathrm{id}`$.
$`\mathrm{}`$ (Subclaim 1)
We could also have used the Dodd-Jensen Lemma to show $`\pi _0\mathrm{}^๐ฏ\mu =\mathrm{id}`$ in the proof of Subclaim 1. However, we chose not to do so in order to exhibit the symmetry of the argument.
Now let $`F^1=E_\beta ^๐ฏ`$ still denote the first extender used on $`(0,\mathrm{}]_T`$. Subclaim 1 says that $`\mathrm{c}.\mathrm{p}.(F^1)=\mu =\mathrm{c}.\mathrm{p}.(F^0)`$. Let, again,
$$F^0=E_{\nu _0}^{_\alpha ^๐ฐ}\mathrm{and}F^1=E_{\nu _1}^{_\beta ^๐ฏ}.$$
The rest of this proof is entirely symmetric, so that we may assume without loss of generality that $`\nu _0\nu _1`$. Both $`\nu _0`$ and $`\nu _1`$ are cardinals in $`Q`$, which implies that $`\nu _0<\nu _1\nu _0`$ is a cardinal in $`๐ฅ_{\nu _1}^{_\beta ^๐ฏ}`$. It is hence clear that we may derive from
$$๐ฅ_{\nu _0}^{_\alpha ^๐ฐ},๐ฅ_{\nu _1}^{_\beta ^๐ฏ}$$
a generalized prebicephalus, call it $`๐ฉ`$. The following is a straightforward consequence of 4.7, 5.2, and 5.10:
Subclaim 2. $`๐ฉ`$ is iterable.
Proof. It suffices to verify that $`๐ฉ`$ satisfies the hypothesis of 5.10. Specifically, we want to verify that our current $`\kappa `$ may serve as the $`\kappa `$ in the statement of 5.10. However, this is trivial by virtue of 4.7 and 5.2!
$`\mathrm{}`$ (Subclaim 2)
Now 5.9 yields that $`E_\alpha ^๐ฐ=F^0=F^1=E_\beta ^๐ฏ`$. This is a contradiction, as this canโt happen in a comparison!
$`\mathrm{}`$ (Claim 2)
$`\mathrm{}`$ (5.11)
Using 6.1 of the next section we could have shown 5.11 for all $`\kappa C_0`$.
## 6 Maximality of $`K^c`$.
###### Lemma 6.1
$`(\neg 0^{^{}})`$ (Maximality of $`K^c`$) Let $`๐ฏ`$ be a normal iteration tree on $`K^c`$ with last model $`Q=_{\mathrm{}}^๐ฏ`$. Suppose that $`(0,\mathrm{}]_T๐^๐ฏ=\mathrm{}`$. Let $`\mu <\kappa \nu `$ be cardinals of $`Q`$, and suppose that $`\pi _0\mathrm{}^๐ฏ\kappa =\mathrm{id}`$. Then there is no extender $`F`$ with critical point $`\mu `$ such that the following hold.
$``$ $`๐ฉ=(๐ฅ_\nu ^Q;F)`$ is a premouse,
$``$ if $`๐ฐ`$ is an iteration tree on $`๐ฉ`$ with last model $`๐ฉ^{}=_{\mathrm{}}^๐ฐ`$, $`(0,\mathrm{}]_U๐^๐ฐ=\mathrm{}`$, and $`\pi _0\mathrm{}^๐ฐ\kappa =\mathrm{id}`$ then the top extender of $`๐ฉ^{}`$ is countably complete, and
$``$ $`\rho _1(๐ฉ)>\mu `$.
Proof. Suppose not, and let $`๐ฏ`$, $`Q`$, $`\mu `$, $`\kappa `$, $`\nu `$, $`F`$, and $`๐ฉ`$ be as in the statement of 6.1. The proof will actually be by induction on $`\mu `$. That is, we assume that $`\mu `$ is least such that there are $`๐ฏ`$, $`Q`$, $`\kappa `$, $`\nu `$, $`F`$, and $`๐ฉ`$ as in the statement of 6.1, and we fix such $`๐ฏ`$, $`Q`$, $`\kappa `$, $`\nu `$, $`F`$, and $`๐ฉ`$.
We first aim to define the pull back of $`๐ฉ`$ via $`\pi _0\mathrm{}^๐ฏ`$. Set $`\pi =\pi _0\mathrm{}^๐ฏ`$. Let $`\lambda =F(\mu )`$, i.e., $`\lambda `$ is the largest cardinal of $`๐ฉ`$. Of course, $`\lambda `$ is an inaccessible cardinal of both $`๐ฉ`$ and $`Q`$. Let $`\overline{\lambda }=\pi ^1\mathrm{"}\lambda `$, so that $`\overline{\lambda }`$ is least with $`\pi (\overline{\lambda })\lambda `$. We define a Dodd-Jensen extender $`\overline{F}`$ by setting
$$X\overline{F}_a\pi (a)F(X)$$
for $`a[\overline{\lambda }]^{<\omega }`$ and $`X๐ซ([\mu ]^{\mathrm{Card}(a)})K^c`$. I.e., $`\overline{F}`$ is a $`(\mu ,\overline{\lambda })`$-extender over $`K^c`$ in the sense of \[20, Definition 1.0.1\]. $`\overline{F}`$ inherits the countable completeness from $`F`$. Let
$$i_{\overline{F}}:๐ฅ_{\mu ^{+K^c}}^{K^c}_{\overline{F}}\overline{}$$
be the $`\mathrm{\Sigma }_0`$ ultrapower map, and set
$$=(\overline{},i_{\overline{F}}๐ซ(\mu )K^c);$$
that is, $`\overline{}`$ is the model theoretic reduct of $``$, where the latter has the additional top extender $`i_{\overline{F}}๐ซ(\mu )K^c`$. We call $``$ the pull back of $`๐ฉ`$ via $`\pi `$. We shall also write $`\overline{F}`$ for $`i_{\overline{F}}๐ซ(\mu )K^c`$. We can define a cofinal $`\mathrm{\Sigma }_\omega `$-elementary map
$$\pi ^{}:\overline{}๐ฅ_\nu ^Q\mathrm{by}\mathrm{setting}$$
$$[a,f]_{\overline{F}}^{๐ฅ_{\mu ^{+K^c}}^{K^c}}[\pi (a),f]_F^{๐ฅ_{\mu ^{+K^c}}^{K^c}}$$
for $`a[\overline{\lambda }]^{<\omega }`$ and $`f:[\mu ]^{\mathrm{Card}(a)}๐ฅ_{\mu ^{+K^c}}^{K^c}`$, $`fK^c`$. Weโll in fact have that
$$\pi ^{}:๐ฉ$$
is $`\mathrm{\Sigma }_0`$-elementary (and hence $`\mathrm{\Sigma }_1`$-elementary, as $`\pi ^{}`$ is cofinal), because $`\pi ^{}(\overline{F}x)=F\pi ^{}(x)`$ for all $`x`$. This readily implies that $``$ is a premouse: the initial segment condition for $``$ is true as $`C_{\mathrm{OR}}^{}\mathrm{}C_{๐ฉ\mathrm{OR}}^๐ฉ\mathrm{}`$; but $`C_{๐ฉ\mathrm{OR}}^๐ฉ=\mathrm{}`$.
The following statements are easy to verify.
$``$ $`\overline{\lambda }`$ is a limit cardinal of $`K^c`$,
$``$ $`๐ฅ_{\overline{\lambda }}^{}=๐ฅ_{\overline{\lambda }}^{K^c}`$,
$``$ $`\pi ^{}๐ฅ_{\overline{\lambda }}^{}=\pi ๐ฅ_{\overline{\lambda }}^{K^c}`$, and
$``$ $`\pi ^{}(\overline{F}(\mu ))=F(\mu )=\lambda `$.
Claim 1. There is no $`\alpha \overline{\lambda }`$ such that $`E_\alpha ^{K^c}\mathrm{}`$, and $`\mathrm{c}.\mathrm{p}.(E_\alpha ^{K^c})[\mu ,\overline{\lambda })`$.
Proof. Suppose otherwise. Set $`\overline{\mu }=\mathrm{c}.\mathrm{p}.(E_\alpha ^{K^c})`$. By 1.3, then, $`\overline{\mu }`$ is $`<\overline{\lambda }`$-strong in $`K^c`$ as witnessed by $`\stackrel{}{E}^{K^c}`$. By the elementarity of $`\pi `$ we hence have $`\pi (\overline{\mu })[\mu ,\pi (\overline{\lambda }))`$ is $`<\pi (\overline{\lambda })`$-strong in $`Q`$ as witnessed by $`\stackrel{}{E}^Q`$. By the definition of $`\overline{\lambda }`$, $`\pi (\overline{\mu })<\lambda `$. So, trivially, $`\pi (\overline{\mu })`$ is $`<\lambda `$-strong in $`Q`$ as witnessed by $`\stackrel{}{E}^Q`$, and thus $`\pi (\overline{\mu })`$ is $`<\lambda `$-strong in $`๐ฉ`$ as witnessed by $`\stackrel{}{E}^๐ฉ`$. But then $`๐ฉ`$ is easily seen not to be below $`0^{^{}}`$.
$`\mathrm{}`$ (Claim 1)
Claim 2. There is no $`\alpha [\overline{\lambda },\mathrm{OR})`$ such that $`E_\alpha ^{}\mathrm{}`$, and $`\mathrm{c}.\mathrm{p}.(E_\alpha ^{})[\mu ,\overline{\lambda })`$.
Proof. Suppose otherwise. Set $`\overline{\mu }=\mathrm{c}.\mathrm{p}.(E_\alpha ^{})`$. By 1.3, then, $`\overline{\mu }`$ is $`<\overline{\lambda }`$-strong in $``$ as witnessed by $`\stackrel{}{E}^{}`$, and thus $`\overline{\mu }`$ is $`<\overline{\lambda }`$-strong in $`K^c`$ as witnessed by $`\stackrel{}{E}^{K^c}`$. This then gives a contradiction as in the proof of Claim 1.
$`\mathrm{}`$ (Claim 2)
Claim 3. $`๐ฅ_{\overline{\lambda }^+}^{}K^c`$.
Proof. This is shown by coiterating $`K^c`$ with $`๐ฅ_{\overline{\lambda }^+}^{}`$, or rather with the phalanx $`๐ซ=((K^c,๐ฅ_{\overline{\lambda }^+}^{}),\overline{\lambda })`$.
Subclaim 1. $`๐ซ`$ is iterable.
Proof. By standard arguments and 2.7, it is enough to see that if $`๐ฑ`$ is an iteration of successor length of $`๐ฅ_{\overline{\lambda }^+}^{}`$ which only uses extenders with critical point $`\overline{\lambda }`$ then any $`E_\rho ^_{\mathrm{}}^๐ฑ\mathrm{}`$ with $`\rho \overline{\lambda }`$ and $`\mathrm{c}.\mathrm{p}.(E_\rho ^_{\mathrm{}}^๐ฑ)<\overline{\lambda }`$ is countably complete. However, using $`\pi ^{}`$ we may copy $`๐ฑ`$ onto $`๐ฉ`$, getting an iteration tree $`๐ฐ`$ on $`๐ฉ`$. Let $`\pi _{\mathrm{}}^{}:_{\mathrm{}}^๐ฑ_{\mathrm{}}^๐ฐ`$ be the last copy map.
$`_{\mathrm{}}^๐ฑ`$ inherits the property expressed by Claim 2 above, that is, we must have $`\mathrm{c}.\mathrm{p}.(E_\rho ^_{\mathrm{}}^๐ฑ)<\mu `$. But then $`\mathrm{c}.\mathrm{p}.(E_{\pi _{\mathrm{}}^{}(\rho )}^_{\mathrm{}}^๐ฐ)<\mu `$, too, and hence $`E_{\pi _{\mathrm{}}^{}(\rho )}^_{\mathrm{}}^๐ฐ`$ is countably complete by 4.7. But this implies that $`E_\rho ^_{\mathrm{}}^๐ฑ`$ is countably complete.
$`\mathrm{}`$ (Subclaim 1)
Now let $`๐ฒ`$, $`๐ฑ`$ denote the iteration trees arising from the comparison of $`K^c`$ with $`๐ซ`$. By 4.5 and standard arguments it suffices to see that the last model of $`๐ฑ`$ sits above $`๐ฅ_{\overline{\lambda }^+}^{}`$. Let us suppose not. We shall derive a contradiction.
By the Dodd-Jensen lemma and 4.5 weโll then have that $`_{\mathrm{}}^๐ฒ=_{\mathrm{}}^๐ฑ`$, $`๐^๐ฒ[0,\mathrm{}]_W=\mathrm{}`$, $`๐^๐ฑ[0,\mathrm{}]_V=\mathrm{}`$, and $`\pi _0\mathrm{}^๐ฒ`$ is lexicographically $`\pi _0\mathrm{}^๐ฑ`$. In particular, we must have either $`\mathrm{c}.\mathrm{p}.(\pi _0\mathrm{}^๐ฒ)=\mathrm{c}.\mathrm{p}.(\pi _0\mathrm{}^๐ฑ)`$ or else $`\pi _0\mathrm{}^๐ฒ\mathrm{c}.\mathrm{p}.(\pi _0\mathrm{}^๐ฑ)+1=\mathrm{id}`$.
Case 1. $`\mathrm{c}.\mathrm{p}.(\pi _0\mathrm{}^๐ฒ)=\mathrm{c}.\mathrm{p}.(\pi _0\mathrm{}^๐ฑ)`$.
Let $`\alpha +1`$ be least in $`(0,\mathrm{}]_W`$, and let $`\beta +1`$ be least in $`(0,\mathrm{}]_V`$. Let
$$F^0=E_{\nu _0}^{_\alpha ^๐ฒ}\mathrm{and}F^1=E_{\nu _1}^{_\beta ^๐ฑ}.$$
The case assumption says that $`\mathrm{c}.\mathrm{p}.(F^0)=\mathrm{c}.\mathrm{p}.(F^1)`$. It is easy to see that we may derive from
$$๐ฅ_{\nu _0}^{_\alpha ^๐ฒ},๐ฅ_{\nu _1}^{_\beta ^๐ฑ}$$
a generalized prebicephalus, call it $``$.
Subclaim 2. $``$ is iterable.
Proof. It suffices to verify that $``$ satisfies the hypothesis of 5.10. Specifically, we want to verify that $`\mathrm{c}.\mathrm{p}.(\pi _0\mathrm{}^๐ฒ)^{+K^c}`$ may serve as the $`\kappa `$ in the statement of 5.10. However, this shown by combining 4.7 with arguments which should be standard by now, and with a copying argument as in the proof of Subclaim 1 above.
$`\mathrm{}`$ (Subclaim 2)
Now 5.9 yields that $`E_\alpha ^๐ฒ=F^0=F^1=E_\beta ^๐ฐ`$. This is a contradiction, as this canโt happen in a comparison!
Case 2. $`\pi _0\mathrm{}^๐ฒ\mathrm{c}.\mathrm{p}.(\pi _0\mathrm{}^๐ฑ)+1=\mathrm{id}`$.
In this case we shall use our inductive hypothesis. In fact, luckily, this is easy to do. Namely, arguments which again should be standard by now yield that we have with $`๐ฒ`$, $`_{\mathrm{}}^๐ฒ`$, $`\mathrm{c}.\mathrm{p}.(\pi _0\mathrm{}^๐ฑ)`$, $`\mathrm{sup}(\{\xi <\nu _1:\pi _0\mathrm{}^๐ฒ(\xi )=\xi \})`$, $`\nu _1`$, $`F^1`$, and $`๐ฅ_{\nu _1}^{_\beta ^๐ฑ}`$ a series of objects which are as $`๐ฏ`$, $`Q`$, $`\mu `$, $`\kappa `$, $`\nu `$, and $`๐ฉ`$ in the statement of 6.1. We leave the straightforward details to the reader. However, by Claim 2 weโll have that $`\mathrm{c}.\mathrm{p}.(\pi _0\mathrm{}^๐ฑ)<\mu `$, which contradicts the choice of $`\mu `$ being minimal!
$`\mathrm{}`$ (Claim 3)
Weโll now have to split the argument into three cases.
Case 1. $`\overline{F}(\mu )=\overline{\lambda }`$.
In this case Claim 3 immediately gives that $`\overline{}K^c`$. (Recall that $`\overline{}`$ is the model theoretic reduct of $``$ obtained by removing the top extender $`\overline{F}`$.) But now we have a contradiction with 4.2.
Case 2. $`\overline{F}(\mu )>\overline{\lambda }`$ and $`\pi (\overline{\lambda })>\lambda `$.
In this case Claim 3 gives that we may define the lift up
$$\stackrel{~}{}=\mathrm{Ult}_0(;\pi ๐ฅ_{\overline{\lambda }^+}^{}).$$
Claim 4. $`\stackrel{~}{}`$ is transitive and $`0`$-iterable.
Proof. We first note that $`\pi `$ is countably complete:
Fact. Let $`(a_n,X_n:n<\omega )`$ be such that for all $`n<\omega `$ do we have $`a_n[\mathrm{OR}]^{<\omega }`$, $`X_n๐ซ([\eta _n]^{\mathrm{Card}(a_n)})K^c`$ for some $`\eta _n`$, and $`a_n\pi (X_n)`$. Then there is an order preserving $`\tau :_{n<\omega }a_n\mathrm{OR}`$ with $`\tau \mathrm{"}a_nX_n`$ for all $`n<\omega `$.
Proof. This is a straightforward consequence of the proof of 3.2. For this purpose we may consider $`๐ฏ`$ as a tree on $`๐ฅ_\theta ^{K^c}`$, for some $`\theta `$. We may and shall assume that $`\theta `$ is large enough and regular so that we may pick $`\sigma :\overline{H}H_\theta `$ with $`\overline{H}`$ countable and transitive, and $`\{๐ฏ\}\{a_n,\pi (X_n):n<\omega \}\mathrm{ran}(\sigma )`$. Set $`\overline{๐ฏ}=\sigma ^1(๐ฏ)`$. The construction from the proof of 3.2 will then give us an embedding
$$\sigma _{\mathrm{}}:_{\mathrm{}}^{\overline{๐ฏ}}๐ฅ_\theta ^{K^c}$$
such that for all $`X`$ with $`\pi (X)\mathrm{ran}(\sigma )`$ do we have that
$$\sigma _{\mathrm{}}\sigma ^1\pi (X)=X.$$
We hence have in
$$\tau =\sigma _{\mathrm{}}\sigma ^1\underset{n<\omega }{}a_n$$
a function as desired.
$`\mathrm{}`$ (Fact)
This implies Claim 4 by standard arguments.
$`\mathrm{}`$ (Claim 4)
Notice that $`\stackrel{~}{}`$ is a premouse (rather than a protomouse, as the top extender of $``$ has critical point $`\mu `$ and $`\pi \mu ^+=\mathrm{id}`$). Moreover, $`\nu `$ (the height of $`๐ฉ`$) is a cardinal in $`\stackrel{~}{}`$ because $`\pi (\overline{\lambda })>\lambda `$ and so $`\pi (\overline{\lambda })>\nu `$, and
$$=(๐ฉ,\stackrel{~}{})$$
is a generalized prebicephalus. Notice that we canโt have that $`๐ฉ=\stackrel{~}{}`$, just because $`\nu =๐ฉ\mathrm{OR}<\stackrel{~}{}`$. Thus, by 5.9, in order to reach a contradiction it suffices to verify the following.
Claim 5. $``$ is iterable.
Proof. We shall apply 5.10. Specifically, we plan on letting the current $`\kappa `$ play the rรดle of the $`\kappa `$ in the statement of 5.10. We have to show that any iteration tree $`๐ฑ`$ on $`๐ฉ`$, or on $`\stackrel{~}{}`$, meets the hypothesis in the statement of 5.10. However, notice that this is clear for $`๐ฉ`$ by 4.7 and our assumptions on $`๐ฉ`$. We are hence left with having to verify the following.
Subclaim 3. Let $`๐ฑ`$ be an iteration tree on $`\stackrel{~}{}`$ with last model $`_{\mathrm{}}^๐ฑ`$ such that for some $`\stackrel{~}{}`$-cardinal $`\rho \kappa `$ we have that either
(a) $`\rho <\kappa `$, and $`๐ฑ`$ lives on $`๐ฅ_\kappa ^\stackrel{~}{}`$ and is above $`\rho `$, but $`๐ฑ`$ doesnโt use extenders which are total on $`\stackrel{~}{}`$, or else
(b) $`\rho =\kappa `$, and $`๐ฑ`$ is above $`\rho `$,
we have that
($``$) if $`F=E_\nu ^_{\mathrm{}}^๐ฑ\mathrm{}`$ is such that $`\nu >\rho `$ and $`\mathrm{c}.\mathrm{p}.(F)<\rho `$, then $`F`$ is countably complete.
Proof. Let $`๐ฑ`$ be as in the statement of this Subclaim, with some $`\rho \kappa `$. Let $`F=E_\nu ^_{\mathrm{}}^๐ฑ\mathrm{}`$ be such that $`\nu >\rho `$ and $`\mathrm{c}.\mathrm{p}.(F)<\rho `$, and let $`(a_n,X_n:n<\omega )`$ be such that $`a_nF(X_n)`$ for all $`n<\omega `$. Let $`\sigma :\overline{H}H_\theta `$ for some large enough regular $`\theta `$ be such that $`\overline{H}`$ is countable and transitive, and $`\{๐ฑ,F\}\{a_n,X_n:n<\omega \}\mathrm{ran}(\sigma )`$. Set $`\overline{๐ฑ}=\sigma ^1(๐ฑ)`$. By the above Fact (in the proof of Claim 4) and standard arguments there is some
$$\overline{\sigma }:_0^{\overline{๐ฑ}}=\sigma ^1(\stackrel{~}{}).$$
Recall that we also have a cofinal map
$$\pi ^{}:_{\mathrm{\Sigma }_1}๐ฉ.$$
We may hence copy $`\overline{๐ฑ}`$ onto $`๐ฉ`$ using $`\pi ^{}\overline{\sigma }`$, getting an iteration tree $`๐ฑ^{}`$ on $`๐ฉ`$, together with a last copy map
$$\sigma _{\mathrm{}}:_{\mathrm{}}^{\overline{๐ฑ}}_{\mathrm{}}^๐ฑ^{}.$$
Now by 4.7 and the assumptions on $`๐ฉ`$ weโll have that $`\sigma _{\mathrm{}}\sigma ^1(F)`$ is countably complete. We may hence pick some order preserving
$$\tau :\underset{n<\omega }{}\sigma _{\mathrm{}}\sigma ^1(a_n)\mathrm{c}.\mathrm{p}.(\sigma _{\mathrm{}}\sigma ^1(F))\mathrm{with}$$
$$\tau \sigma _{\mathrm{}}\sigma ^1\mathrm{"}a_n\sigma _{\mathrm{}}\sigma ^1(X_n)\mathrm{for}\mathrm{all}n<\omega .$$
However,
$$\sigma _{\mathrm{}}\sigma ^1(X_n)=X_n,$$
as $`\mathrm{c}.\mathrm{p}.(F)<\rho `$, and $`๐ฑ`$ is above $`\rho `$. We therefore have
$$\tau \sigma _{\mathrm{}}\sigma ^1\mathrm{"}a_nX_n\mathrm{for}\mathrm{all}n<\omega ,$$
and hence we have in $`\tau \sigma _{\mathrm{}}\sigma ^1`$ a function as desired.
$`\mathrm{}`$ (Subclaim 3)
$`\mathrm{}`$ (Claim 5)
Case 3. $`\overline{F}(\mu )>\overline{\lambda }`$ and $`\pi (\overline{\lambda })=\lambda `$.
Set $`\stackrel{~}{\lambda }=\mathrm{sup}\pi \mathrm{"}\overline{\lambda }`$. It is easy to see that $`\pi `$ canโt be continuous at $`\overline{\lambda }`$ in this case, i.e., that
$``$ $`\stackrel{~}{\lambda }<\lambda `$,
because otherwise we would have to have $`\pi ^{}(\overline{\lambda })\lambda `$, whereas $`\pi ^{}(\overline{F}(\mu ))=\lambda `$. Moreover, as $`\lambda `$ is inaccessible in $`Q`$, we must have that $`\pi ^1(\lambda )=\overline{\lambda }`$ is inaccessible in $`K^c`$. This implies, because $`\pi `$ is an iteration map which is discontinuous at $`\overline{\lambda }`$, that $`\overline{\lambda }`$ is measurable in $`K^c`$, and that in fact we can write
$$๐ฏ=๐ฏ_0{}_{}{}^{}๐ฏ_{1}^{}$$
where $`\pi _0\mathrm{}^{๐ฏ_0}`$ is continuous at $`\overline{\lambda }`$, $`\mathrm{c}.\mathrm{p}.(\pi _0\mathrm{}^{๐ฏ_1})=\pi _0\mathrm{}^{๐ฏ_0}(\overline{\lambda })=\stackrel{~}{\lambda }`$, and $`\pi =\pi _0\mathrm{}^{๐ฏ_1}\pi _0\mathrm{}^{๐ฏ_0}`$. (Let $`\rho `$ be the least $`\overline{\rho }[0,\mathrm{})_T`$ with $`\mathrm{c}.\mathrm{p}.(\pi _{\overline{\rho }\mathrm{}}^๐ฏ)=\pi _{0\overline{\rho }}^๐ฏ(\overline{\lambda })`$. Then $`\pi _0\mathrm{}^{๐ฏ_0}=\pi _{0\rho }^๐ฏ`$, and $`\pi _0\mathrm{}^{๐ฏ_1}=\pi _\rho \mathrm{}^๐ฏ`$.) Let us write $`\pi _0=\pi _0\mathrm{}^{๐ฏ_0}`$ and $`\pi _1=\pi _0\mathrm{}^{๐ฏ_1}`$. We have that $`\pi _0(\overline{\lambda })=\stackrel{~}{\lambda }`$, and $`\pi _1(\stackrel{~}{\lambda })=\lambda `$.
Claim 6. $`๐ฅ_{\overline{\lambda }^+}^{}=๐ฅ_{\overline{\lambda }^{+K^c}}^{K^c}`$.
Proof. We already know that $`๐ฅ_{\overline{\lambda }^+}^{}\mathrm{}๐ฅ_{\overline{\lambda }^{+K^c}}^{K^c}`$. It is easy to see that $`_{\mathrm{}}^{๐ฏ_0}`$ and $`Q`$ agree up to $`\stackrel{~}{\lambda }^{+_{\mathrm{}}^{๐ฏ_0}}=\stackrel{~}{\lambda }^{+Q}`$ and that
(1)
$$\pi _0๐ฅ_{\overline{\lambda }^{+K^c}}^{K^c}:๐ฅ_{\overline{\lambda }^{+K^c}}^{K^c}๐ฅ_{\stackrel{~}{\lambda }^{+Q}}^Q$$
is exactly the lift up of $`๐ฅ_{\overline{\lambda }^{+K^c}}^{K^c}`$ by $`\pi _0๐ฅ_{\overline{\lambda }}^{K^c}=\pi ๐ฅ_{\overline{\lambda }}^{K^c}`$, which in turn is given by letting $`๐ฏ_0`$ act on $`๐ฅ_{\overline{\lambda }^{+K^c}}^{K^c}`$.
Weโll now use the interpolation technique. Let
$$\sigma _0:=\mathrm{Ult}_0(;\pi ๐ฅ_{\overline{\lambda }}^{K^c})$$
be the lift up of $``$ by $`\pi ๐ฅ_{\overline{\lambda }}^{K^c}`$. Notice that $`\rho _1()\stackrel{~}{\lambda }`$ and $``$ is $`\stackrel{~}{\lambda }`$-sound, because every element of $``$ can be written in the form $`[a,i_{\overline{F}}(f)(b)]`$ where $`a[\stackrel{~}{\lambda }]^{<\omega }`$, $`b[\overline{\lambda }]^{<\omega }`$, and $`f:[\mu ]^{\mathrm{Card}(b)}๐ฅ_{\mu ^{+K^c}}^{K^c}`$ with $`f๐ฅ_{\mu ^{+K^c}}^{K^c}`$. Of course, $`\sigma _0(\overline{\lambda })=\stackrel{~}{\lambda }`$.
We may define
$$\sigma _1:๐ฉ$$
by letting
$$[a,f]_{\pi ๐ฅ_{\overline{\lambda }}^{K^c}}^{}\pi ^{}(f)(a)$$
for $`a[\stackrel{~}{\lambda }]^{<\omega }`$, and $`f:[\eta ]^{\mathrm{Card}(a)}`$ for some $`\eta <\overline{\lambda }`$ with $`f`$ and $`a\pi (\mathrm{dom}(f))`$. It is straightforward to see that $`\sigma _1`$ is $`\mathrm{\Sigma }_0`$-elementary, that $`\sigma _1\stackrel{~}{\lambda }=\mathrm{id}`$, and that
$$\pi ^{}=\sigma _1\sigma _0.$$
Of course $`\stackrel{~}{\lambda }`$ is inaccessible in both $``$ and $`๐ฉ`$.
Clearly, $`๐ฅ_{\stackrel{~}{\delta }^{+๐ฉ}}^๐ฉ=๐ฅ_{\stackrel{~}{\delta }^{+๐ฌ}}^๐ฌ`$. Let us assume that we also have
(2)
$$๐ฅ_{\stackrel{~}{\delta }^+}^{}=๐ฅ_{\stackrel{~}{\delta }^{+๐ฌ}}^๐ฌ.$$
Then weโll have that
(3)
$$\sigma _0๐ฅ_{\overline{\lambda }^+}^{}:๐ฅ_{\overline{\lambda }^+}^{}๐ฅ_{\stackrel{~}{\lambda }^+}^{}=๐ฅ_{\stackrel{~}{\lambda }^{+Q}}^Q$$
is the map obtained by taking the ultrapower of $`๐ฅ_{\overline{\lambda }^+}^{}`$ by the long extender $`\pi ๐ฅ_{\overline{\lambda }}^{K^c}`$, which of course is exactly what is obtained when we let $`๐ฏ_0`$ act on $`๐ฅ_{\overline{\lambda }^+}^{}`$. Hence both $`\pi _0๐ฅ_{\overline{\lambda }^{+K^c}}^{K^c}`$ and $`\sigma _0๐ฅ_{\overline{\lambda }^+}^{}`$ are maps induced by $`๐ฏ_0`$, and they have the same target model, and $`๐ฅ_{\overline{\lambda }^+}^{}\mathrm{}๐ฅ_{\overline{\lambda }^{+K^c}}^{K^c}`$. This implies that we must have $`๐ฅ_{\overline{\lambda }^+}^{}=๐ฅ_{\overline{\lambda }^{+K^c}}^{K^c}`$.
In order to finish the proof of Claim 6 it hence suffices to verify (2). We first prove:
Subclaim 4. There is no $`E_\alpha ^๐ฉ\mathrm{}`$ with $`\mathrm{c}.\mathrm{p}.(E_\alpha ^๐ฉ)[\mu ,\stackrel{~}{\lambda })`$ and $`\alpha (\stackrel{~}{\lambda },\nu )`$.
Proof. Otherwise by 1.3 we get that $`\overline{\mu }=\mathrm{c}.\mathrm{p}.(E_\alpha ^๐ฉ)`$ is $`<\stackrel{~}{\lambda }`$-strong in $`๐ฉ`$ as witnessed by $`\stackrel{}{E}^๐ฉ`$. Using $`\pi _0`$, we then get that there is some $`\overline{\mu }[\mu ,\overline{\lambda })`$ which is $`<\overline{\lambda }`$-strong in $`K^c`$ as witnessed by $`\stackrel{}{E}^{K^c}`$, which is a contradiction as in the proof of Claim 2 above.
$`\mathrm{}`$ (Subclaim 4)
Now (2) is shown by applying the condensation lemma \[9, ยง8 Lemma 4\] to $`\sigma _1`$. Recall that $`\rho _1()\stackrel{~}{\lambda }`$, $``$ is $`\stackrel{~}{\lambda }`$-sound, and $`\sigma _1\stackrel{~}{\lambda }=\mathrm{id}`$. We may assume that $`\sigma _1\mathrm{id}`$ and $`๐ฉ`$, as otherwise (2) is trivial. Let $`\eta =\mathrm{c}.\mathrm{p}.(\sigma _1)\stackrel{~}{\lambda }`$. Well, if (b) or (c) in the conclusion of \[9, ยง8 Lemma 4\] were to hold then (in much the same way as in the proof of 1.1) weโd get that there is some $`E_\alpha ^๐ฉ\mathrm{}`$ with $`\mathrm{c}.\mathrm{p}.(E_\alpha ^๐ฉ)=\mu `$ and $`\alpha (\stackrel{~}{\lambda },\nu )`$, contradicting Subclaim 4. We hence must have that (a) in the conclusion of \[9, ยง8 Lemma 4\] holds, i.e., that $``$ is the $`\eta `$-core of $`๐ฉ`$, and that $`\sigma _1`$ is the core map.
We are now finally going to use our assumption that $`\rho _1(๐ฉ)>\mu `$ which we make in the statement of 6.1. Let $`๐ฑ`$, $`๐ฑ^{}`$ denote the iteration trees arising from the comparison of $``$ with $`๐ฉ`$. We know that $`(0,\mathrm{}]_V๐^๐ฑ=(0,\mathrm{}]_V^{}๐^๐ฑ^{}=\mathrm{}`$, and that $`_{\mathrm{}}^๐ฑ=_{\mathrm{}}^๐ฑ^{}`$, because $``$ is the $`\eta `$-core of $`๐ฉ`$ and $`\sigma _1`$ is the core map. We also know that $`\pi _0\mathrm{}^๐ฑ\eta =\mathrm{id}`$: otherwise we can consider the coiteration of $`((๐ฉ,),\eta )`$ with $`๐ฉ`$ and by 2.7 weโd get a contradiction as in the solidity proof (cf. \[9, ยง7\]).
Suppose that $`\pi _0\mathrm{}^๐ฑ^{}\stackrel{~}{\lambda }\mathrm{id}`$. Weโd then have that $`\mathrm{c}.\mathrm{p}.(\pi _0\mathrm{}^๐ฑ^{})\mu `$, because by Subclaim 4 above $`๐ฉ`$ does not have any extenders $`E_\alpha ^๐ฉ\mathrm{}`$ with $`\mathrm{c}.\mathrm{p}.(E_\alpha ^๐ฉ)(\mu ,\stackrel{~}{\lambda })`$ and $`\alpha >\stackrel{~}{\lambda }`$, a fact which is inherited by iterations of $`๐ฉ`$ above $`\stackrel{~}{\lambda }`$. But then $`\mathrm{c}.\mathrm{p}.(\pi _0\mathrm{}^๐ฑ^{})\mu `$ and $`\rho _1(๐ฉ)>\mu `$ imply that
$$\rho _1(_{\mathrm{}}^๐ฑ^{})>\pi _0\mathrm{}^๐ฑ^{}(\mu )\pi _0\mathrm{}^๐ฑ^{}(\mathrm{c}.\mathrm{p}.(\pi _0\mathrm{}^๐ฑ^{}))>\stackrel{~}{\lambda }.$$
On the other hand, as $`\pi _0\mathrm{}^๐ฑ\stackrel{~}{\lambda }=\mathrm{id}`$, we must have
$$\rho _1(_{\mathrm{}}^๐ฑ)\stackrel{~}{\lambda }.$$
This is a contradiction!
Hence we must have that $`\pi _0\mathrm{}^๐ฑ^{}\stackrel{~}{\lambda }=\mathrm{id}`$, that is, the coiteration of $``$ with $`๐ฉ`$ is above $`\stackrel{~}{\lambda }`$ on both sides, $`(0,\mathrm{}]_V๐^๐ฑ=(0,\mathrm{}]_V^{}๐^{๐ฑ\mathrm{`}}=\mathrm{}`$, and $`_{\mathrm{}}^๐ฑ=_{\mathrm{}}^๐ฑ^{}`$. This yields (2).
$`\mathrm{}`$ (Claim 6)
Given Claim 6, we can now easily finish the proof of 6.1. As $`\pi `$ is an iteration map, we must have that $`\pi `$ is continuous at $`\overline{\lambda }^{+K^c}`$. By Claim 6, hence, $`\pi \mathrm{"}\overline{\lambda }^+`$ is cofinal in $`\pi (\overline{\lambda }^{+K^c})=\nu `$. But then, if we define the lift up
$$\stackrel{~}{}=\mathrm{Ult}_0(;\pi ๐ฅ_{\overline{\lambda }^+}^{})$$
as in Case 2 above, then we shall again have that $`\nu `$ is a cardinal in $`\stackrel{~}{}`$. We may then continue and derive a contradiction exactly as in Case 2 above.
$`\mathrm{}`$ (6.1)
It can be shown that 6.1 also holds when the assumption that $`\rho _1(๐ฉ)>\mu `$ is dropped.
## 7 Weak covering for $`K^c`$.
Recall that a cardinal $`\kappa `$ is called countably closed if $`\overline{\kappa }^\mathrm{}_0<\kappa `$ for all $`\overline{\kappa }<\kappa `$. We let $`B_0`$ denote the class of all countably closed singular cardinals $`\kappa `$ such that $`\kappa C_0`$ and for all $`\mu <\kappa `$ do we have that if $`\mu `$ is $`<\kappa `$-strong in $`K^c`$ then $`\mu `$ is $`<\mathrm{OR}`$-strong in $`K^c`$.
It is easy to see (using 5.1) that $`B_0`$ is a stationary class. The significance of $`B_0`$ is due to the following fact, the โweak covering lemma for $`K^c`$.โ<sup>5</sup><sup>5</sup>5We want to remark that 7.1 is not the strongest covering lemma which is provable for $`K^c`$. On the other hand, 7.1 will suffice for the construction of $`K`$.
###### Lemma 7.1
($`\neg 0^{^{}}`$) Let $`\kappa B_0`$. Then $`\kappa ^{+K^c}=\kappa ^+`$.
Proof. Fix $`\kappa B_0`$, and set $`\lambda =\kappa ^{+K^c}`$. Let us assume that $`\lambda <\kappa ^+`$. We shall eventually derive a contradiction.
Fix $`\mathrm{\Omega }>\kappa ^+`$, a regular cardinal. Notice that for $`W=(K^c)^{H_\mathrm{\Omega }}`$ we have that $`๐ฅ_\mathrm{\Omega }^{K^c}=๐ฅ_\mathrm{\Omega }^W`$. We may pick an elementary embedding
$$\pi :NH_\mathrm{\Omega }$$
such that $`{}_{}{}^{\omega }NN`$, $`\mathrm{Card}(N)<\kappa `$, $`\{\kappa ,\lambda \}\mathrm{ran}(\pi )`$, and $`\mathrm{ran}(\pi )\lambda `$ is cofinal in $`\lambda `$. We may and shall assume that
$$\mathrm{Card}(N)=\mathrm{cf}(\lambda )^\mathrm{}_0,$$
and that $`N(\mathrm{Card}(N))^+`$ is transitive.
Set $`\overline{\kappa }=\pi ^1(\kappa )`$, $`\overline{\lambda }=\pi ^1(\lambda )`$, and $`\overline{K}=๐ฅ_{\overline{\lambda }}^{(K^c)^N}`$. Let $`\delta =\mathrm{c}.\mathrm{p}.(\pi )<\overline{\kappa }`$. Notice that $`\delta =N(\mathrm{Card}(N))^+`$, $`\mathrm{Card}(\delta )=\mathrm{Card}(N)`$, and $`\pi (\delta )=\mathrm{Card}(N)^+=\delta ^{+V}`$.We can make an immediate observation.
Claim 1. $`(๐ซ(\delta )K^c)N\mathrm{}`$.
Proof. Suppose otherwise. Set $`\stackrel{~}{F}=\pi ๐ซ(\delta )K^c`$. $`\stackrel{~}{F}`$ is countably complete by $`{}_{}{}^{\omega }NN`$. It is easy to see that we have a contradiction with 4.2.
$`\mathrm{}`$ (Claim 1)
Let $`\eta `$ be least such that $`(๐ซ(\eta )K^c)N\mathrm{}`$. Either $`\eta =\delta `$ and $`\delta `$ is inaccessible in $`\overline{K}`$, or else $`\delta =\eta ^{+\overline{K}}`$ (and then $`\pi (\delta )=\eta ^{+K^c}`$).
Let $`๐ฐ`$ and $`๐ฏ`$ denote the iteration trees arising from the comparison of $`\overline{K}`$ with $`K^c`$. By 4.5, weโll have that $`_{\mathrm{}}^๐ฏ_{\mathrm{}}^๐ฐ`$, and $`๐^๐ฐ[0,\mathrm{}]_U=\mathrm{}`$. In fact, the Dodd-Jensen Lemma can easily be used to see that $`๐ฏ`$ actually only uses extenders from $`๐ฅ_\mathrm{\Omega }^{K^c}`$ and its images.
We now aim to prove:
Claim 2. $`๐ฐ`$ is trivial, i.e., $`_{\mathrm{}}^๐ฐ=\overline{K}`$.
Proof. Suppose that $`๐ฐ`$ is non-trivial, and let $`F=E_\alpha ^๐ฐ`$ be the first extender used on $`[0,\mathrm{})_U`$. Recall that $`D^๐ฐ[0,\mathrm{})_U=\mathrm{}`$. Let $`\mu `$ be the critical point of $`F`$. Notice that by 2.4 $`๐ฐ\alpha `$ only uses extenders with critical point $`>\mu ^{+\overline{K}}`$.
Subclaim 1. $`F`$ is countably complete.
Proof. Let $`((a_n,X_n):n<\omega )`$ be such that $`a_nF(X_n)`$ for all $`n<\omega `$. Pick
$$\sigma :\overline{H}H_\vartheta $$
such that $`\vartheta `$ is regular and large enough, $`\overline{H}`$ is transitive, $`\mathrm{Card}(\overline{H})=\mathrm{}_0`$, and $`\{๐ฐ\alpha ,F\}\{a_n,X_n:n<\omega \}\mathrm{ran}(\sigma )`$. Let $`\overline{๐ฐ}=\sigma ^1(๐ฐ\alpha )`$, and $`\overline{}=\sigma ^1(\overline{K})`$. Notice that $`\overline{๐ฐ}`$, $`\overline{}`$, as well as $`\sigma \overline{}`$ are all elements of $`N`$, by $`{}_{}{}^{\omega }NN`$.
Let us copy $`\overline{๐ฐ}`$ onto $`\overline{K}`$, using $`\sigma \overline{}`$. The entire copying construction takes place within $`N`$, and it gives an iteration tree $`๐ฐ^{}`$ on $`\overline{K}`$ together with copying maps $`\sigma _i:_i^{\overline{๐ฐ}}_i^๐ฐ^{}`$ (where $`\sigma _0=\sigma \overline{}`$).
As $`๐ฐ\alpha `$ only uses extenders with critical point $`>\mu ^{+\overline{K}}`$, we have that $`๐ฐ^{}`$ only uses extenders with critical point $`>\mu ^{+\overline{K}}`$, too. This implies that $`\sigma _0`$ and $`\sigma _{\mathrm{}}:_{\mathrm{}}^{\overline{๐ฐ}}_{\mathrm{}}^๐ฐ^{}`$ are such that $`\sigma _{\mathrm{}}\sigma _0^1(X_n)=X_n`$ for all $`n<\omega `$. Moreover, by 4.7 applied inside $`N`$ together with $`{}_{}{}^{\omega }NN`$, we have that $`F^{}=\sigma _{\mathrm{}}\sigma _0^1(F)`$ is (really) countably complete. Let $`a_n^{}=\sigma _{\mathrm{}}\sigma _0^1(a_n)`$ for $`n<\omega `$, and let $`\tau :_{n<\omega }a_n^{}\mu `$ be order preserving and such that $`\tau \mathrm{"}a_n^{}X_n`$ for all $`n<\omega `$. We then have in
$$\tau \sigma _{\mathrm{}}\sigma _0^1\underset{n<\omega }{}a_n$$
a function as desired.
$`\mathrm{}`$ (Subclaim 1)
Now notice that \[9, ยง8 Lemma 1\] (a corollary to the solidity proof) and 2.6 imply that $`\mu \eta `$. Set $`\mu _0=\eta `$, and let $`(\mu _i:0<i<\gamma )`$ enumerate the successor cardinals of $`\overline{K}`$ in the half-open interval $`(\eta ,\mu ^{+\overline{K}}]`$ (if $`\mu =\eta `$ then $`\gamma =2`$; otherwise, weโll have $`\gamma =\mu +1`$). For every $`i<\gamma `$ let $`\tau (i)`$ be the least $`\tau `$ such that $`๐ฅ_{\mu _i}^{_\tau ^๐ฏ}=๐ฅ_{\mu _i}^{\overline{K}}`$, and let $`๐ซ_i`$ be the longest initial segment of $`_{\tau (i)}^๐ฏ`$ which has the same bounded subsets of $`\mu _i`$ as $`\overline{K}`$ has.
Subclaim 2. $`๐ซ_0=K^c`$, and $`i>0\rho _\omega (๐ซ_i)<\mu _i`$.
Proof. It is enough to verify that $`๐ซ_i`$ is a set-sized premouse for $`i>0`$. Let $`i>0`$, and assume that $`๐ซ_i`$ is a weasel. (In particular, $`๐ซ_i=_{\tau (i)}^๐ฏ`$.) Notice that $`\tau (i)>0`$, because $`(๐ซ(\eta )K^c)\overline{K}\mathrm{}`$. Let $`G=E_j^๐ฏ=E_\nu ^{๐ซ_j}`$ be the first extender used on $`[0,\tau (i))_T`$. $`G`$ is total on $`K^c`$, and countably complete by 4.7. Moreover, of course, $`\nu <\mu _i`$. Also, $`\mathrm{c}.\mathrm{p}.(G)<\eta `$. So $`\mathrm{ran}(G)\overline{K}`$ and $`\stackrel{~}{G}=\pi G`$ is well-defined. We have that $`G(\mu )`$ is a cardinal of $`\overline{K}`$, and hence $`\stackrel{~}{G}(\mu )`$ is a cardinal of $`K^c`$, using the elementarity of $`\pi `$. But as $`G`$ is countably complete, $`\stackrel{~}{G}`$ is countably complete, too, by the countable completeness of $`\pi `$ (i.e., by $`{}_{}{}^{\omega }NN`$). We thus have a contradiction with 4.2.
$`\mathrm{}`$ (Subclaim 2)
Now set $`๐ซ_\gamma =_\alpha ^๐ฐ`$. We shall be interested in the phalanx
$$\stackrel{}{๐ซ}=((๐ซ_i:i<\gamma +1),(\mu _i:i<\gamma )).$$
For $`0<i<\gamma `$ let $`\mu _i^{}`$ denote the cardinal predecessor of $`\mu _i`$ in $`๐ซ_i`$. Standard arguments give that for such $`i`$ is $`๐ซ_i`$ sound above $`\mu _i^{}`$, i.e., if $`k<\omega `$ is such that $`\rho _{k+1}(๐ซ_i)<\mu _i\rho _k(๐ซ_i)`$ then $`๐ซ_i`$ is $`k`$-sound and $`(๐ซ_i)^k`$ is the hull of $`\mu _i^{}\{p_{(๐ซ_i)^k,1}\}`$ generated by $`h_{๐ซ_i}^k`$.
Subclaim 3. $`\stackrel{}{๐ซ}`$ is coiterable with $`K^c`$.
Proof. Let us coiterate $`K^c`$ with $`\stackrel{}{๐ซ}`$ using $`\omega `$-maximal trees, getting iteration trees $`๐ฑ`$ and $`๐ฒ`$. Subclaim 3 says that all models of $`๐ฒ`$ are transitive, and that we therefore get comparable $`_{\mathrm{}}^๐ฑ`$ and $`_{\mathrm{}}^๐ฒ`$.
Let $`\mathrm{}<\mathrm{OR}`$ denote that ordinal such that either $`_{\mathrm{}}^๐ฒ`$ is ill-founded, or else $`_{\mathrm{}}^๐ฑ`$ and $`_{\mathrm{}}^๐ฒ`$ are comparable. We aim to show that it is the latter alternative which holds. Notice that $`๐ฑ\alpha +1=๐ฏ\alpha +1`$, that $`๐ฒ\alpha `$ is trivial, and that $`E_\alpha ^๐ฒ=F`$ which is applied to $`๐ซ_{\gamma 1}`$.
Well, by 2.7, there is $`n<\omega `$ such that $`๐ฒ`$ can be written as
$$๐ฒ_0{}_{}{}^{}๐ฒ_{1}^{}{}_{}{}^{}\mathrm{}{}_{}{}^{}๐ฒ_{n}^{}$$
where each $`๐ฒ_k`$ is an iteration of some $`๐ซ_{i(k)}`$ (with $`k^{}>ki(k^{})<i(k)`$) except for the fact that the very first extender used in $`๐ฒ_k`$ is equal to $`F`$ if $`k=0`$, or else is taken from the last model of $`๐ฒ_{k1}`$ if $`k>0`$. Let $`F_k`$ denote the first extender used in $`๐ฒ_k`$, and let $`\kappa _k=\mathrm{c}.\mathrm{p}.(F_k)`$. Notice that $`k^{}>k\kappa _k^{}<\kappa _k`$, and that $`๐ฒ_k`$ is an iteration which uses only extenders with critical point $`\kappa _k`$ (by the rules for iterating a phalanx).
Let us pick an elementary embedding
$$\sigma :\overline{H}H_\theta $$
where $`\theta `$ is regular and large enough, $`\overline{H}`$ is countable and transitive, and $`๐ฒ\mathrm{ran}(\sigma )`$. Set $`\stackrel{}{๐ฌ}=\sigma ^1(\stackrel{}{๐ซ})`$, $`\overline{๐ฒ}=\sigma ^1(๐ฒ)`$, $`\overline{\kappa }_k=\sigma ^1(\kappa _k)`$, and $`\overline{F}=\sigma ^1(F)`$. Then $`\overline{๐ฒ}`$ is an iteration of $`\stackrel{}{๐ฌ}`$ which can be written as
$$\overline{๐ฒ}_0{}_{}{}^{}\overline{๐ฒ}_{1}^{}{}_{}{}^{}\mathrm{}{}_{}{}^{}\overline{๐ฒ}_{n}^{}$$
where $`\overline{๐ฒ}_k=\sigma ^1(๐ฒ_k)`$.
Now as $`F=F_0`$ is countably complete by Subclaim 1, we may pick some $`\tau :(F(\mu )\mathrm{ran}(\sigma ))\kappa _0=\mu `$ order preserving such that $`aF(X)\tau \mathrm{"}aX`$ for appropriate $`a`$, $`X\mathrm{ran}(\sigma )`$. Set $`\overline{\gamma }=\sigma ^1(\gamma )`$. Then
$$_{\overline{\gamma }+1}^{\overline{๐ฒ}}=\mathrm{Ult}_k(\sigma ^1(๐ซ_{\gamma 1});\overline{F}),$$
where $`\rho _{k+1}(๐ซ_{\gamma 1})\mu <\rho _k(๐ซ_{\gamma 1})`$, is the first model of $`\overline{๐ฒ}`$ above the list of models $`\stackrel{}{๐ฌ}`$. We may define an embedding
$$\overline{\sigma }:_{\overline{\gamma }+1}^{\overline{๐ฒ}}=\mathrm{Ult}_k(\sigma ^1(๐ซ_{\gamma 1});\overline{F})๐ซ_{\gamma 1},$$
as being the extension (via the upward extensions of embeddings lemma) of the map
$$\overline{\overline{\sigma }}:(_{\overline{\gamma }+1}^{\overline{๐ฒ}})^k(๐ซ_{\gamma 1})^k,$$
where $`\overline{\overline{\sigma }}`$ is defined by
$$[a,f]\sigma (f)(\tau \mathrm{"}\sigma (a)).$$
Let $`\overline{๐ฒ}^{}`$ be that iteration of the phalanx $`\sigma ^1(\stackrel{}{๐ซ}\gamma )^{}_{\overline{\gamma }+1}^{\overline{๐ฒ}}`$ which uses exactly the same extenders (in the same order) as $`\overline{๐ฒ}`$ does, except for the very first one. In particular, $`\overline{๐ฒ}^{}`$ and $`\overline{๐ฒ}`$ have exactly the same models, except that $`\overline{๐ฒ}^{}`$ misses $`_{\overline{\gamma }}^{\overline{๐ฒ}}`$. Notice that $`\stackrel{}{๐ซ}\gamma =(๐ซ_i:i<\gamma )`$ is iterable, as it is generated by an iteration of $`K^c`$. It hence suffices to copy the iteration $`\overline{๐ฒ}^{}`$ onto $`\stackrel{}{๐ซ}\gamma `$, using the maps $`\sigma `$ and $`\overline{\sigma }`$.
However, we have that $`\overline{\sigma }`$ agrees with $`\sigma \sigma ^1(๐ซ_{\gamma 1})`$ up to $`\sigma ^1(\kappa _1)^+`$ (calculated in $`\sigma ^1(๐ซ_{\gamma 1})`$). Hence the standard copying construction goes through.
$`\mathrm{}`$ (Subclaim 3)
Let still $`๐ฑ`$, $`๐ฒ`$ denote the trees coming from the coiteration of $`K^c`$ with $`\stackrel{}{๐ซ}`$. It is now easy to see that we have to have that $`0=\mathrm{root}^๐ฒ(\mathrm{})`$, i.e., that the final model of $`๐ฒ`$ sits above $`๐ซ_0=K^c`$: otherwise $`_{\mathrm{}}^๐ฒ`$ would be non-sound which would give an immediate contradiction to 4.5. By applying the Dodd-Jensen lemma to $`\pi _0\mathrm{}^๐ฒ`$ otherwise, we get that $`๐ฑ`$ has to be simple along its main branch, and $`_{\mathrm{}}^๐ฑ\mathrm{}_{\mathrm{}}^๐ฒ`$. But 4.5 implies that $`๐ฒ`$ has to be simple along its main branch, and $`_{\mathrm{}}^๐ฒ\mathrm{}_{\mathrm{}}^๐ฑ`$.
Set $`Q=_{\mathrm{}}^๐ฑ=_{\mathrm{}}^๐ฒ`$. Let $`G`$ be the first extender applied on the main branch of $`W`$, and set $`\overline{\mu }=\mathrm{c}.\mathrm{p}.(G)`$. As the final model of $`๐ฒ`$ is above $`K^c`$, we know that $`\overline{\mu }<\eta `$. The Dodd-Jensen lemma tells us that $`\pi _0\mathrm{}^๐ฑ\pi _0\mathrm{}^๐ฒ`$ lexicographically; in particular, $`\pi _0\mathrm{}^๐ฑ\overline{\mu }=\mathrm{id}`$.
Case 1. $`\mathrm{c}.\mathrm{p}.(\pi _0\mathrm{}^๐ฑ)=\overline{\mu }`$.
Let $`\beta _0+1`$ be least in $`(0,\mathrm{}]_V`$, and let $`\beta _1+1`$ be least in $`(0,\mathrm{}]_W`$. We may in this case derive from
$$^0=๐ฅ_{\nu _{\beta _0}}^{_{\beta _0}^๐ฑ},^1=๐ฅ_{\nu _{\beta _1}}^{_{\beta _1}^๐ฒ}$$
a generalized prebicephalus, call it $`๐ฉ`$.
Subclaim 4. $`๐ฉ`$ is iterable.
Proof. This is shown by applying 5.10. Specifically, we plan on letting $`\overline{\mu }^{+K^c}`$ play the rรดle of $`\kappa `$ in the statement of 5.10. Set $`\overline{\nu }=\overline{\mu }^{+K^c}`$. As $`\overline{\nu }`$ is a cardinal of $`K^c`$, $`^0\mathrm{}๐ฅ_{\overline{\nu }}^{K^c}`$, and $`^1\mathrm{}๐ฅ_{\overline{\nu }}^{K^c}`$, by virtue of 4.7 it suffices to verify the following two statements in order to prove Subclaim 4.
A$`_0`$ if $`\stackrel{~}{}`$ is an iterate of $`^0`$ above $`\overline{\nu }`$, and if $`H=E_\nu ^\stackrel{~}{}\mathrm{}`$ is an extender with $`\mathrm{c}.\mathrm{p}.(H)<\overline{\nu }`$ and $`\nu >\overline{\nu }`$ then $`H`$ is countably complete.
A$`_1`$ if $`\stackrel{~}{}`$ is an iterate of $`^1`$ above $`\overline{\nu }`$, and if $`H=E_\nu ^\stackrel{~}{}\mathrm{}`$ is an extender with $`\mathrm{c}.\mathrm{p}.(H)<\overline{\nu }`$ and $`\nu >\overline{\nu }`$ then $`H`$ is countably complete.
Now by $`\neg 0^{^{}}`$, $`^0`$ is easily seen to be an iterate of $`K^c`$ above $`\overline{\nu }`$. Hence if $`\stackrel{~}{}`$ and $`H`$ are as in A<sub>0</sub> then $`\stackrel{~}{}`$ is an iterate of $`K^c`$ above $`\overline{\nu }`$, too, and $`H`$ is countably complete by 4.7. It thus remains to verify A<sub>1</sub>.
Fix $`\stackrel{~}{}`$ and $`H=E_\nu ^\stackrel{~}{}`$ as in A<sub>1</sub>, and let $`(a_n,X_n:n<\omega )`$ be such that $`a_nH(X_n)`$ for every $`n<\omega `$. Let $`๐ฒ^{}`$ be the iteration of $`^1`$ which gives $`\stackrel{~}{}`$. Note that $`^1`$ in turn is given by an iteration, call it $`๐ฒ^+`$, of the phalanx $`\stackrel{}{๐ซ}`$. Recall that $`\stackrel{}{๐ซ}=(๐ซ_i:i<\gamma +1)`$, and that $`F`$ is the first extender used in $`๐ฒ`$. Let $`\stackrel{}{๐ซ}^{}`$ denote the phalanx $`(๐ซ_i:i<\gamma )^{}\mathrm{Ult}(๐ซ_{\gamma 1};F)`$ (with the same exchange ordinals as $`\stackrel{}{๐ซ}`$). Then, trivially, $`^1`$ is given by that iteration, call it $`๐ฒ^{}`$, of $`\stackrel{}{๐ซ}^{}`$ which uses exactly the same extenders (in the same order) as $`๐ฒ^+`$ does, except for the very first one, $`F`$.
Let us pick an elementary embedding
$$\sigma :\overline{H}H_\theta $$
where $`\theta `$ is regular and large enough, $`\overline{H}`$ is countable and transitive, and
$$\{\stackrel{}{๐ซ}^{},๐ฒ^{},^1,๐ฒ^{},H\}\{a_n,X_n:n<\omega \}\mathrm{ran}(\sigma ).$$
Set $`\stackrel{}{๐ฌ}=\sigma ^1(\stackrel{}{๐ซ}^{})`$, $`\overline{๐ฒ}=\sigma ^1(๐ฒ^{})`$, $`\overline{}=\sigma ^1(^1)`$, and $`\overline{๐ฒ}^{}=\sigma ^1(๐ฒ^{})`$. Exactly as in the proof of Subclaim 2 above, we may first find a map $`\overline{\sigma }`$ re-embedding the last model of $`\stackrel{}{๐ฌ}`$ into $`๐ซ_{\gamma 1}`$ (using the countable completeness of $`F`$), and we may then use the agreement between $`\overline{\sigma }`$ and $`\sigma `$ to copy $`\overline{๐ฒ}`$ onto $`\stackrel{}{๐ซ}\gamma `$, which gives an iteration $`\overline{๐ฒ}^c`$ of $`\stackrel{}{๐ซ}\gamma `$ together with the copy maps. Let
$$\overline{\sigma }_{\mathrm{}}:_{\mathrm{}}^{\overline{๐ฒ}}_{\mathrm{}}^{\overline{๐ฒ}^c}$$
be the copy map from the last model of $`\overline{๐ฒ}`$ to the last model of $`\overline{๐ฒ}^c`$.
Now $`\overline{๐ฒ}^{}`$ is an iteration of (a truncation of) $`_{\mathrm{}}^{\overline{๐ฒ}}`$, and we may hence continue with copying $`\overline{๐ฒ}^{}`$ onto $`_{\mathrm{}}^{\overline{๐ฒ}^c}`$, using $`\overline{\sigma }_{\mathrm{}}`$, which gives an iteration $`\overline{๐ฒ}^{cc}`$ of $`_{\mathrm{}}^{\overline{๐ฒ}^c}`$ together with the copy maps. Let
$$\sigma _{\mathrm{}}:_{\mathrm{}}^{\overline{๐ฒ}^{}}_{\mathrm{}}^{\overline{๐ฒ}^{cc}}$$
be the copy map from the last model of $`\overline{๐ฒ}^{}`$ to the last model of $`\overline{๐ฒ}^{cc}`$.
It is now easy to verify that $`\sigma `$ and $`\sigma _{\mathrm{}}`$ agree up to $`\sigma ^1(๐ฅ_{\overline{\nu }}^{K^c})`$. Moreover, as $`\stackrel{}{๐ซ}\gamma `$ is given by a normal iteration of $`K^c`$, weโll have that $`\overline{๐ฒ}^c`$ is in fact an iteration of $`K^c`$ above $`\overline{\nu }`$. But then
$$\overline{๐ฒ}^c{}_{}{}^{}\overline{๐ฒ}_{}^{cc}$$
is an iteration of $`K^c`$ above $`\overline{\nu }`$, too. By 4.7, we thus know that $`\sigma _{\mathrm{}}\sigma ^1(H)`$ is countably complete. We may hence pick $`\tau :_{n<\omega }\sigma _{\mathrm{}}\sigma ^1(a_n)\mathrm{c}.\mathrm{p}.(H)`$ such that $`\tau \mathrm{"}\sigma _{\mathrm{}}\sigma ^1(a_n)X_n`$ for every $`n<\omega `$. We have found a function as desired!
$`\mathrm{}`$ (Subclaim 4)
But now by 5.9 we have that Subclaim 4 implies
$$E_{\beta _0}^๐ฑ=E_{\beta _1}^๐ฒ.$$
This is a contradiction!
Case 2. $`\pi _0\mathrm{}^๐ฑ\overline{\mu }+1=\mathrm{id}`$.
In this case, we get a contradiction with 6.1. That 6.1 is applicable here follows from arguments exactly as in Case 1 above. We leave it to the reader to chase through the obvious details.
$`\mathrm{}`$ (Claim 2)
Claim 3. $`๐ฏ`$ only uses extenders with critical point greater than or equal to $`\eta `$.
Proof. This follows from the argument which was presented in the proof of Subclaim 2 above.
$`\mathrm{}`$ (Claim 3)
The combination of Claims 1, 2, and 3 now give an initial segment, $``$, of $`_{\mathrm{}}^๐ฏ`$ with $`\overline{K}`$, $`\overline{\lambda }`$ is a cardinal in $``$, $``$ is $`\overline{\kappa }`$-sound, and $`\rho _\omega ()\overline{\kappa }`$. Let $`n<\omega `$ be such that $`\rho _{n+1}()\overline{\kappa }<\rho _n()`$. Let
$$\stackrel{~}{}=\mathrm{Ult}_n(;\pi \overline{K}).$$
Then either $`\stackrel{~}{}`$ is an $`n`$-iterable premouse with $`\stackrel{~}{}๐ฅ_\lambda ^{K^c}`$, $`\rho _{n+1}(\stackrel{~}{})\kappa `$, and $`\stackrel{~}{}`$ is sound above $`\kappa `$, or else $`n=0`$, $``$ has a top extender, $`\overline{F}`$, $`\pi \mathrm{"}\mathrm{c}.\mathrm{p}.(\overline{F})^+`$ is not cofinal in $`\pi (\mathrm{c}.\mathrm{p}.(\overline{F})^+)`$, and $`\stackrel{~}{}`$ is a protomouse (see \[19, ยง2.3\]). In the latter case, let $`\overline{F}^{}`$ be the top extender of $`\stackrel{~}{}`$, and let $`๐ฉ=๐ฅ_\eta ^{K^c}`$ be the longest initial segment of $`K^c`$ which has only subsets of $`\mathrm{c}.\mathrm{p}.(\overline{F}^{})`$ which are measured by $`\overline{F}^{}`$. Let
$$\stackrel{~}{๐ฉ}=\mathrm{Ult}_m(๐ฉ;\overline{F}^{}),$$
where $`\rho _{m+1}(๐ฉ)\mathrm{c}.\mathrm{p}.(\overline{F}^{})<\rho _m(๐ฉ)`$. Using the countable completeness of $`\pi `$, by standard arguments weโll have that $`\stackrel{~}{๐ฉ}`$ is an $`m`$-iterable premouse with $`\stackrel{~}{๐ฉ}๐ฅ_\lambda ^{K^c}`$, $`\rho _\omega (\stackrel{~}{๐ฉ})\kappa `$, and $`\stackrel{~}{๐ฉ}`$ is sound above $`\kappa `$.
But now, finally, 5.11 gives a contradiction!
$`\mathrm{}`$ (7.1)
## 8 Beavers and the existence of K.
In this section we shall isolate $`K`$, the true core model below $`0^{^{}}`$. We closely follow \[27, ยง5\]; however, of course, some revisions are necessary, as works in the theory โ$`\mathrm{๐น๐ฅ๐ข}+\mathrm{\Omega }`$ is measurable.โ
We first need a concept of thick classes, which is originally due to . We commence with a simple observation (which is standard).
###### Lemma 8.1
$`(\neg 0^{^{}})`$ Let $`W`$ be an iterable weasel. Let $`S`$ be a stationary class such that for all $`\beta S`$ we have that $`\beta `$ is a strong limit cardinal, $`\mathrm{cf}^W(\beta )`$ is not measurable in $`W`$, and $`\beta ^{+W}=\beta ^+`$. Then $`W`$ is universal.
Proof. Suppose not. Then there is a (set- or class-sized) premouse $``$ such that if $`๐ฏ`$, $`๐ฐ`$ are the iteration trees arising from the comparison of $``$ with $`W`$ we have that: $`\mathrm{lh}(๐ฏ)=\mathrm{lh}(๐ฐ)=\mathrm{OR}+1`$, there is no drop on $`[0,\mathrm{OR}]_U`$, $`\pi _0\mathrm{}^๐ฐ\mathrm{"}\mathrm{OR}\mathrm{OR}`$, and if $`\alpha `$ is largest in $`\{0\}(๐^๐ฏ(0,\mathrm{})_T)`$ then $`\pi _\alpha \mathrm{}^๐ฏ\mathrm{"}\mathrm{OR}\mathrm{OR}`$. (We have $`\mathrm{}=\mathrm{OR}`$ here.) There are club many $`\beta [0,\mathrm{})_U`$ such that $`\pi _{0\beta }^๐ฐ\mathrm{"}\beta \beta `$. Also, there are $`\gamma [\alpha ,\mathrm{})_T`$ and $`\kappa _\gamma ^๐ฏ`$ such that for club many $`\beta [\gamma ,\mathrm{})_T`$ do we have that $`\pi _{\gamma \beta }^๐ฏ(\kappa )=\beta `$ and $`\beta `$ is the critical point of $`\pi _\beta \mathrm{}^๐ฏ`$.
Now pick $`\beta S(_\gamma ^๐ฏ\mathrm{OR})`$ with $`\pi _{0\beta }^๐ฐ\mathrm{"}\beta \beta `$, $`\pi _{\gamma \beta }^๐ฏ(\kappa )=\beta `$, and $`\mathrm{c}.\mathrm{p}.(\pi _\beta \mathrm{}^๐ฏ)=\beta `$. As $`\beta `$ is a strong limit cardinal and $`\mathrm{cf}^W(\beta )`$ is not measurable in $`W`$, weโll have that $`\pi _0\mathrm{}^๐ฐ(\beta )=\beta `$, which implies that $`\beta ^{+_{\mathrm{}}^๐ฐ}=\beta ^+`$. On the other hand, weโll have that $`\pi _{\gamma \beta }^๐ฏ\mathrm{"}\kappa ^{+_\gamma ^๐ฏ}`$ is cofinal in $`\beta ^{+_\beta ^๐ฏ}=\beta ^{+_{\mathrm{}}^๐ฏ}`$, so that $`\beta ^{+_{\mathrm{}}^๐ฏ}<\beta ^+`$. Hence $`\beta ^{+_{\mathrm{}}^๐ฏ}<\beta ^+=\beta ^{+_{\mathrm{}}^๐ฐ}`$. This is a contradiction!
$`\mathrm{}`$ (8.1)
###### Definition 8.2
Let $`W`$ be a weasel, let $`S\mathrm{OR}`$, and let $`\tau \mathrm{OR}`$. A class $`\mathrm{\Gamma }\mathrm{OR}`$ is called $`S,\tau `$-thick in $`W`$ provided the following clauses hold.
(i) $`S`$ is stationary in $`\mathrm{OR}`$,
(ii) for all but nonstationary many $`\beta S`$ do we have that
(a) $`\beta ^{+W}=\beta ^+`$,
(b) $`\beta `$ is a strong limit cardinal,
(c) $`\mathrm{cf}^W(\xi )`$ is not measurable in $`W`$ as witnessed by $`\stackrel{}{E}^W`$ (i.e., there is no $`E_\nu ^W\mathrm{}`$ with $`\mathrm{c}.\mathrm{p}.(E_\nu ^W)=\mathrm{cf}^W(\xi )`$ and $`E_\nu ^W`$ is total on $`W`$) for all $`\xi \mathrm{\Gamma }[\beta ,\beta ^+)`$, and
(d) $`\mathrm{\Gamma }(\beta ,\beta ^+)`$ is unbounded in $`\beta ^+`$, and for all regular $`\theta >\tau `$ we have that $`\mathrm{\Gamma }(\beta ,\beta ^+)`$ is $`\theta `$-closed, and $`\beta \mathrm{\Gamma }`$.
A class $`\mathrm{\Gamma }\mathrm{OR}`$ is called $`S`$-thick in $`W`$ if there is some $`\tau `$ such that $`\mathrm{\Gamma }`$ is $`S,\tau `$-thick in $`W`$.
By the proof of 8.1, if $`\mathrm{\Gamma }\mathrm{OR}`$ is $`S`$-thick in $`W`$ for some $`S\mathrm{OR}`$ then $`W`$ is universal.
The following four lemmas are easy to prove (cf. \[27, Lemmas 3.9 to 3.11\]).
###### Lemma 8.3
Let $`W`$ be a weasel, let $`S\mathrm{OR}`$ be a class, and let $`(\mathrm{\Gamma }_i:i<\theta )`$ be such that $`\mathrm{\Gamma }_i`$ is $`S`$-thick in $`W`$ for all $`i<\theta `$. Then $`_{i<\theta }\mathrm{\Gamma }_i`$ is $`S`$-thick in $`W`$.
###### Lemma 8.4
Let $`\overline{W}`$ and $`W`$ be weasels, and let $`S\mathrm{OR}`$ be a class. Let $`\pi :\overline{W}W`$ be an elementary embedding such that there is some $`\mathrm{\Gamma }\mathrm{ran}(\pi )`$ which is $`S`$-thick in $`W`$. Then $`\mathrm{\Gamma }\{\xi :\pi (\xi )=\xi \}`$ is $`S`$-thick in both $`\overline{W}`$ and $`W`$.
###### Lemma 8.5
Let $`W`$ be a weasel, let $`S\mathrm{OR}`$ be a class, and let $`\mathrm{\Gamma }`$ be $`S`$-thick in $`W`$. Let $`F`$ be a total extender on $`W`$ such that $`\mathrm{Ult}(W;F)`$ is transitive. Let $`i_F`$ denote the ultrapower map. Then $`\mathrm{\Gamma }\{\xi :i_F(\xi )=\xi \}`$ is $`S`$-thick in both $`W`$ and $`\mathrm{Ult}(W;F)`$.
###### Lemma 8.6
Let $`๐ฏ`$ be an iteration tree on the weasel $`W`$, let $`S\mathrm{OR}`$ be a class, and let $`\mathrm{\Gamma }`$ be $`S`$-thick in $`W`$. Let $`\alpha \mathrm{lh}(๐ฏ)`$ (possibly, $`\alpha =\mathrm{OR}`$) be such that $`๐^๐ฏ(0,\alpha ]_T=\mathrm{}`$ and $`\pi _{0\alpha }^๐ฏ\mathrm{"}\mathrm{OR}\mathrm{OR}`$. Then $`\mathrm{\Gamma }\{\xi :\pi _{0\alpha }^๐ฏ(\xi )=\xi \}`$ is $`S`$-thick in both $`W`$ and $`_\alpha ^๐ฏ`$.
In order to isolate $`K`$ we now have to work with some $`K_\mathrm{\Gamma }^c`$ for a small $`\mathrm{\Gamma }`$ rather than with $`K^c`$, just because there need not be $`S`$ and $`\mathrm{\Gamma }^{}`$ such that $`\mathrm{\Gamma }^{}`$ is $`S`$-thick in $`K^c`$. Let us consider $`K_{\{\omega _1\}}^c`$. We get versions of everything we proved in sections 4 through 7 about $`K^c`$ also for $`K_{\{\omega _1\}}^c`$. (When quoting a result from one of these sections in what follows, what weโll in fact intend to quote is the corresponding version for $`K_{\{\omega _1\}}^c`$.)
In particular, there is a stationary class $`B_0^{\{\omega _1\}}`$ (defined in a similar fashion as the old $`B_0`$ was defined for $`K^c`$) with $`\kappa ^{+K_{\{\omega _1\}}^c}=\kappa ^+`$ for all $`\kappa B_0^{\{\omega _1\}}`$. With this meaning of $`B_0^{\{\omega _1\}}`$ in mind we now let $`A_0`$ denote the class of all elements of $`B_0^{\{\omega _1\}}`$ which are strong limit cardinals of cofinality $`>\omega _1`$. I.e., $`A_0`$ consists of all singular strong limit cardinals $`\kappa `$ of cofinality $`>\omega _1`$ such that $`๐ฅ_\kappa ^{K_{\{\omega _1\}}^c}`$ is universal for coiterable premice of height $`<\kappa `$ and such that if $`\mu `$ is $`<\kappa `$-strong in $`K_{\{\omega _1\}}^c`$ then $`\mu `$ is $`<\mathrm{OR}`$-strong in $`K_{\{\omega _1\}}^c`$. Of course, $`A_0`$ is a stationary subclass of $`B_0^{\{\omega _1\}}`$.
The following is now immediate from 7.1.
###### Lemma 8.7
$`(\neg 0^{^{}})`$ Let $`\mathrm{\Gamma }`$ be the class of all ordinals of cofinality $`>\omega _1`$. Then $`\mathrm{\Gamma }`$ is $`A_0`$-thick in $`K_{\{\omega _1\}}^c`$.
Proof. As for (ii) (c) in 8.2, notice that if $`\mathrm{cf}^{K_{\{\omega _1\}}^c}(\xi )`$ were measurable in $`K_{\{\omega _1\}}^c`$ as witnessed by $`\stackrel{}{E}^{K_{\{\omega _1\}}^c}`$ then we would have to have $`\mathrm{cf}^V(\xi )=\mathrm{cf}^V(\mathrm{cf}^{K_{\{\omega _1\}}^c}(\xi ))=\omega _1`$, because if $`E_\nu ^{K_{\{\omega _1\}}^c}\mathrm{}`$ is total on $`K_{\{\omega _1\}}^c`$ then its critical point has cofinality $`=\omega _1`$. Contradiction!
$`\mathrm{}`$ (8.7)
Let $`W`$ be a weasel, and let $`\stackrel{}{\alpha }\mathrm{OR}`$. For the following purposes we shall denote by $`x=\tau ^W[\stackrel{}{\alpha }]`$ the fact that there is some formula $`\mathrm{\Phi }`$ such that $`x`$ is the unique $`\overline{x}`$ with $`W\mathrm{\Phi }(\overline{x},\stackrel{}{\alpha })`$, and $`\tau `$ is the term given by $`\mathrm{\Phi }`$. For $`\mathrm{\Gamma }\mathrm{OR}`$ we shall let $`xH^W(\mathrm{\Gamma })`$ mean that $`x=\tau ^W[\stackrel{}{\alpha }]`$ for some term $`\tau `$ and $`\stackrel{}{\alpha }\mathrm{\Gamma }`$.
###### Definition 8.8
Let $`W`$ be a weasel, and let $`S\mathrm{OR}`$. Let $`\alpha \mathrm{OR}`$. We say that $`W`$ has the $`S`$-hull property at $`\alpha `$ just in case that for all $`S`$-thick $`\mathrm{\Gamma }`$ do we have that $`๐ซ(\alpha )W`$ is a subset of the transitive collapse of $`H^W(\alpha \mathrm{\Gamma })`$.
###### Lemma 8.9
$`(\neg 0^{^{}})`$ There is some club $`C\mathrm{OR}`$ such that $`K_{\{\omega _1\}}^c`$ has the $`A_0`$-hull property at any $`\kappa C`$.
Proof. Let $`C`$ be the class of all limit cardinals $`\kappa `$ (of $`V`$) such that for any $`\mu <\kappa `$ we have that if $`\mu `$ is $`<\kappa `$-strong in $`K_{\{\omega _1\}}^c`$ then $`\mu `$ is $`<\mathrm{OR}`$-strong in $`K_{\{\omega _1\}}^c`$. Let $`\kappa C`$. We aim to show that $`K_{\{\omega _1\}}^c`$ has the $`A_0`$-hull property at $`\kappa `$.
Let $`\mathrm{\Gamma }`$ be $`A_0`$-thick in $`K_{\{\omega _1\}}^c`$, and let $`\sigma :WK_{\{\omega _1\}}^c`$ be elementary with $`W`$ transitive, and $`\mathrm{ran}(\sigma )=H^{K_{\{\omega _1\}}^c}(\kappa \mathrm{\Gamma })`$. Notice that $`W`$ is universal. It suffices to prove that the coiteration of $`W`$, $`K_{\{\omega _1\}}^c`$ is above $`\kappa `$ on both sides. We first show the following easy
Claim. Let $``$ be an iterate of $`W`$ above $`\kappa `$, and let $`F=E_\nu ^{}\mathrm{}`$ be such that $`\nu >\kappa `$ and $`\mathrm{c}.\mathrm{p}.(F)<\kappa `$. Then $`F`$ is countably complete.
Proof. Let $`=_{\mathrm{}}^๐ฏ`$ where $`๐ฏ`$ is the iteration tree on $`W`$ giving $``$. Using $`\sigma :WK_{\{\omega _1\}}^c`$ we may copy $`๐ฏ`$ onto $`K_{\{\omega _1\}}^c`$, getting an iteration tree $`๐ฐ`$ on $`K_{\{\omega _1\}}^c`$ together with a last copy map $`\sigma _{\mathrm{}}:_{\mathrm{}}^๐ฐ`$. As $`๐ฏ`$ is above $`\kappa `$, $`๐ฐ`$ is above $`\kappa `$, too. Hence $`\sigma _{\mathrm{}}(F)`$ is countably complete by 4.7. This clearly implies that $`F`$ is countably complete, because $`\sigma _{\mathrm{}}๐ซ(\mathrm{c}.\mathrm{p}.(F))W=\mathrm{id}`$.
$`\mathrm{}`$ (Claim)
Now let $`๐ฏ`$, $`๐ฐ`$ denote the iteration trees arising from the comparison of $`W`$ with $`K_{\{\omega _1\}}^c`$. If $`\pi _0\mathrm{}^๐ฏ\kappa \mathrm{id}`$ or $`\pi _0\mathrm{}^๐ฐ\kappa \mathrm{id}`$ then exactly as in the proof of 5.11 Claim 2 Subclaim 1 weโll get that $`\mathrm{c}.\mathrm{p}.(\pi _0\mathrm{}^๐ฏ)=\mathrm{c}.\mathrm{p}.(\pi _0\mathrm{}^๐ฐ)<\kappa `$. (Notice that any $`\mu <\kappa `$ is $`<\mathrm{OR}`$-strong in $`W`$ if and only if it is $`<\mathrm{OR}`$-strong in $`K_{\{\omega _1\}}^c`$!) Hence if $`\alpha +1`$ is least in $`(0,\mathrm{}]_T`$ and $`\beta +1`$ is least in $`(0,\mathrm{}]_U`$, and if $`E_\alpha ^๐ฏ=E_{\nu _0}^{_\alpha ^๐ฏ}`$ and $`E_\beta ^๐ฐ=E_{\nu _1}^{_\beta ^๐ฐ}`$, then we may derive from
$$๐ฅ_{\nu _0}^{_\alpha ^๐ฏ},๐ฅ_{\nu _1}^{_\beta ^๐ฐ}$$
a generalized prebicephalus, call it $`๐ฉ`$.
But using the Claim above, $`๐ฉ`$ is iterable by 5.10 (where we let the current $`\kappa `$ play the rรดle of the $`\kappa `$ in the statement of 5.10). This gives a contradiction as in the proof of 5.11.
$`\mathrm{}`$ (8.9)
###### Lemma 8.10
Let $`W`$ be a universal weasel, and let $`\mathrm{\Gamma }`$ be $`S`$-thick in $`W`$. Then $`W`$ has the hull property at all $`\kappa `$ such that $`\kappa `$ is $`<\mathrm{OR}`$-strong in $`W`$.
Proof. This is shown by induction on $`\kappa `$. Let $`W`$ be a universal weasel, and let $`\mathrm{\Gamma }`$ be $`S`$-thick in $`W`$. Let
$$\sigma :\overline{W}H^W(\kappa \mathrm{\Gamma })W.$$
Let $`๐ฐ`$, $`๐ฏ`$ denote the iteration trees arising from the comparison of $`W`$ with $`\overline{W}`$. It suffices to prove that $`\pi _0\mathrm{}^๐ฐ\kappa =\pi _0\mathrm{}^๐ฏ\kappa =\mathrm{id}`$. Suppose not. Then $`\mu =\mathrm{min}\{\mathrm{c}.\mathrm{p}.(\pi _0\mathrm{}^๐ฐ),`$ $`\mathrm{c}.\mathrm{p}.(\pi _0\mathrm{}^๐ฏ)\}`$ is $`<\kappa `$-strong in $`W`$, and hence $`<\mathrm{OR}`$-strong in both $`W`$ and $`\overline{W}`$. By an argument as in the proof of 5.11 Claim 2 Subclaim 1 we then get that in fact $`\mu =\mathrm{c}.\mathrm{p}.(\pi _0\mathrm{}^๐ฐ)=\mathrm{c}.\mathrm{p}.(\pi _0\mathrm{}^๐ฏ)`$. Let $`\alpha +1`$ be least in $`(0,\mathrm{}]_U`$, and let $`\beta +1`$ be least in $`(0,\mathrm{}]_T`$. By inductive hypothesis both $`W`$ and $`\overline{W}`$ have the hull-property at $`\mu `$. Set $`Q=_{\mathrm{}}^๐ฐ=_{\mathrm{}}^๐ฏ`$.
By 8.3, 8.4, and 8.6, $`\mathrm{\Gamma }^{}=\{\xi \mathrm{\Gamma }:\sigma (\xi )=\pi _0\mathrm{}^๐ฐ(\xi )=\pi _0\mathrm{}^๐ฏ(\xi )=\xi \}`$ is $`A_0`$-thick in $`W`$, $`\overline{W}`$, and $`Q`$. Let $`F=E_\alpha ^๐ฐ`$, and $`F^{}=E_\beta ^๐ฏ`$. Let $`a[F(\mu )F^{}(\mu )]^{<\omega }`$, and $`X๐ซ([\mu ]^{\mathrm{Card}(a)})W=๐ซ([\mu ]^{\mathrm{Card}(a)})\overline{W}`$. Let $`X=\tau ^W[\stackrel{}{\xi }][\mu ]^{\mathrm{Card}(a)}`$ for $`\stackrel{}{\xi }\mathrm{\Gamma }^{}`$. Then $`X=\tau ^{\overline{W}}[\stackrel{}{\xi }][\mu ]^{\mathrm{Card}(a)}`$, too. Moreover, $`aF(X)`$ if and only if $`a\pi _0\mathrm{}^๐ฐ(X)=\tau ^Q[\stackrel{}{\xi }]\pi _0\mathrm{}^๐ฐ([\mu ]^{\mathrm{Card}(a)})`$ if and only if $`a\tau ^Q[\stackrel{}{\xi }]\pi _0\mathrm{}^๐ฏ([\mu ]^{\mathrm{Card}(a)})=\pi _0\mathrm{}^๐ฏ(X)`$ if and only if $`aF^{}(X)`$. Hence $`F`$ and $`F^{}`$ are compatible. Contradiction!
$`\mathrm{}`$ (8.10)
Let $`W`$ be a weasel, and let $`S\mathrm{OR}`$. We shall let $`x\mathrm{Def}(W,S)`$ mean that whenever $`\mathrm{\Gamma }`$ is $`S`$-thick in $`W`$ then $`xH^W(\mathrm{\Gamma })`$.
If $`W`$, $`W^{}`$ are both (universal) weasels, and $`S\mathrm{OR}`$, such that there is some $`\mathrm{\Gamma }\mathrm{OR}`$ which is $`S`$-thick in $`W`$ as well as in $`W^{}`$ then it is easy to verify that $`\mathrm{Def}(W,S)`$ and $`\mathrm{Def}(W^{},S)`$ have the same transitive collapse (cf. \[27, Corollary 5.7\]. Hence for $`S`$ a stationary class we denote by $`K(S)`$ the transitive collapse of $`\mathrm{Def}(W,S)`$ for any (some) universal $`W`$ such that there is some $`\mathrm{\Gamma }\mathrm{OR}`$ which is $`S`$-thick in $`W`$.
The proof of the following lemma 8.13 needs some care in order to not to go beyond the bounds of predicative class theory. We shall need the following auxiliary concept, which was introduced in \[4, Definition 2.7\] (and generalizes a concept of ; see also ).
###### Definition 8.11
Let $`=(J_\alpha [\stackrel{}{E}];,\stackrel{}{E},F)`$ be a premouse with top extender $`F\mathrm{}`$. Set $`\kappa =\mathrm{c}.\mathrm{p}.(F)`$. Then $``$ is called a beaver provided there is a universal weasel $`W`$ with the following properties.
(a) $`๐ฅ_\lambda ^W=๐ฅ_\lambda ^{}`$ where $`\lambda =\kappa ^{+W}=\kappa ^+`$,
(b) $`W`$ has the definability property at all $`\mu <\kappa `$ such that $`\mu `$ is $`<\kappa `$-strong in $`W`$, and
(c) $`\mathrm{Ult}_0(W;F)`$ is $`0`$-iterable.
It is straightforward to check that the proof of \[4, Lemma 2.9\] works below $`0^{^{}}`$; this observation establishes:
###### Lemma 8.12
$`(\neg 0^{^{}})`$ Suppose that $`=(J_\alpha [\stackrel{}{E}];,\stackrel{}{E},F)`$ and $`^{}=(J_\alpha [\stackrel{}{E}];,\stackrel{}{E},F^{})`$ are beavers. Then $`F=F^{}`$.
###### Lemma 8.13
$`(\neg 0^{^{}})`$ $`K(A_0)`$ is a weasel.
Proof. We aim to show that $`\mathrm{Def}(K_{\{\omega _1\}}^c,A_0)`$ is unbounded in $`\mathrm{OR}`$. Let $`\mathrm{\Gamma }`$ be the class of all ordinals of cofinality $`>\omega _1`$. By 8.7, $`\mathrm{\Gamma }`$ is $`A_0`$-thick in $`K_{\{\omega _1\}}^c`$.
Claim 1. There is a sequence $`(\mathrm{\Gamma }_\kappa :\kappa \mathrm{OR})`$ of classes such that
(A) $`\mathrm{\Gamma }_0=\mathrm{\Gamma }`$,
(B) $`\mathrm{\Gamma }_\kappa \mathrm{\Gamma }_\kappa ^{}`$ for $`\kappa \kappa ^{}`$,
(C) if there is some $`\overline{\mathrm{\Gamma }}\mathrm{\Gamma }_\kappa `$ which is $`A_0`$-thick in $`K_{\{\omega _1\}}^c`$ and such that $`\kappa H^{K_{\{\omega _1\}}^c}(\overline{\mathrm{\Gamma }})`$ then $`\kappa H^{K_{\{\omega _1\}}^c}(\mathrm{\Gamma }_{\kappa +1})`$, for all ordinals $`\kappa `$, and
(D) $`\mathrm{\Gamma }_\lambda =_{\kappa <\lambda }\mathrm{\Gamma }_\kappa `$ for all limit ordinals $`\lambda `$.
Proof. Of course, Claim 1 is supposed to say that the class $`\{(\kappa ,x):x\mathrm{\Gamma }_\kappa \}`$ exists. The point of Claim 1 is that we have to be able to define such a class. For this purpose, we need a refinement of the argument in \[7, Appendix II\]. Let us first indicate how we aim to choose $`\mathrm{\Gamma }_{\kappa +1}`$, given $`\mathrm{\Gamma }_\kappa `$.
Suppose that $`\mathrm{\Xi }`$ is $`A_0`$-thick in $`K_{\{\omega _1\}}^c`$, and $`\beta <\alpha \beta \mathrm{Def}(K_{\{\omega _1\}}^c,A_0)\beta H^{K_{\{\omega _1\}}^c}(\mathrm{\Xi })`$. We want to find some canonical $`\mathrm{\Xi }^{}`$ which is $`A_0`$-thick in $`K_{\{\omega _1\}}^c`$, and $`\alpha \mathrm{Def}(K_{\{\omega _1\}}^c,A_0)\alpha H^{K_{\{\omega _1\}}^c}(\mathrm{\Xi }^{})`$. Let us suppose without loss of generality that $`\alpha H^{K_{\{\omega _1\}}^c}(\mathrm{\Xi })\mathrm{Def}(K_{\{\omega _1\}}^c,A_0)`$. Let $`\mathrm{\Xi }_0\mathrm{\Xi }`$ be $`A_0`$-thick in $`K_{\{\omega _1\}}^c`$ such that $`\alpha H^{K_{\{\omega _1\}}^c}(\mathrm{\Xi }_0)`$. Let
$$\sigma _0:\overline{W}_0H^{K_{\{\omega _1\}}^c}(\mathrm{\Xi }_0)K_{\{\omega _1\}}^c,\mathrm{and}$$
$$\sigma :\overline{W}H^{K_{\{\omega _1\}}^c}(\mathrm{\Xi })K_{\{\omega _1\}}^c.$$
Then $`\sigma ^1(\alpha )`$ is the critical point of $`\sigma ^1\sigma _0`$, and both $`\overline{W}_0`$ and $`\overline{W}`$ have the definability property at all $`\mu <\sigma ^1(\alpha )`$. Set $`\kappa =\sigma ^1(\alpha )`$. The proof of \[4, Lemma 1.3\] then works below $`0^{^{}}`$ and shows that $`\kappa ^{+\overline{W}_0}=\kappa ^{+\overline{W}}`$. Let us write $`\lambda =\kappa ^{+\overline{W}}`$. Set $`F=\sigma ^1\sigma _0๐ซ(\kappa )\overline{W}_0`$, and $`\stackrel{~}{\lambda }=\mathrm{sup}\sigma ^1\sigma _0\mathrm{"}\lambda `$. Now consider
$$=(๐ฅ_{\stackrel{~}{\lambda }}^{\overline{W}},F).$$
Let us suppose without loss of generality that $``$ is a premouse (if not, then the initial segment condition fails for $`F`$, and we may replace $``$ by a premouse obtained as in the proof of 4.2). Moreover, as $`\mathrm{Ult}_0(\overline{W}_0;F)`$ can be embedded into $`\overline{W}`$, we immediately get that in fact $``$ is a beaver. On the other hand, by 8.12, there is at most one $`\overline{F}`$ such that $`(๐ฅ_{\stackrel{~}{\lambda }}^{\overline{W}},\overline{F})`$ is a beaver.
By \[4, Lemma 1.2\] we now know that $`\mathrm{Ult}_0(\overline{W};F)`$ is $`0`$-iterable, too. Let
$$\pi :\overline{W}_F\stackrel{~}{W}.$$
By coiterating $`\overline{W}`$, $`\stackrel{~}{W}`$ we get a common coiterate $`Q`$ together with iteration maps $`\pi _{\overline{W},Q}`$ and $`\pi _{\stackrel{~}{W},Q}`$. Set
$$\mathrm{\Xi }^{}=\{\xi \mathrm{\Xi }:\sigma (\xi )=\pi _{\overline{W},Q}(\xi )=\pi _{\stackrel{~}{W},Q}\pi (\xi )=\xi \}.$$
Then $`\mathrm{\Xi }^{}`$ is $`A_0`$-thick in $`K_{\{\omega _1\}}^c`$ by 8.3, 8.4, 8.5, and 8.6. Notice that $`xH^{K_{\{\omega _1\}}^c}(\mathrm{\Xi }^{})`$ of course implies that $`\pi _{\overline{W},Q}\sigma ^1(x)\mathrm{ran}(\pi _{\stackrel{~}{W},Q}\pi )`$.
Subclaim 1. $`\alpha H^{K_{\{\omega _1\}}^c}(\mathrm{\Xi }^{})`$.
Proof. Suppose otherwise, and let $`\alpha =\tau ^{K_{\{\omega _1\}}^c}[\stackrel{}{\xi }]`$ where $`\stackrel{}{\xi }\mathrm{\Xi }^{}`$. Then $`\kappa =\tau ^{\overline{W}}[\stackrel{}{\xi }]`$. Notice that both $`\overline{W}`$ and $`\stackrel{~}{W}`$ have the definability property at all $`\mu <\kappa `$. By \[4, Corollary 1.5\], hence, $`\pi _{\overline{W},Q}\kappa =\pi _{\stackrel{~}{W},Q}\kappa =\mathrm{id}`$.
Case 1. $`\pi _{\overline{W},Q}\kappa +1=\mathrm{id}`$.
In this case we get that $`\kappa =\tau ^Q[\stackrel{}{\xi }]\mathrm{ran}(\pi _{\stackrel{~}{W},Q}\pi )`$. However, $`\kappa `$ is the critical point of $`\pi `$, and $`\pi _{\stackrel{~}{W},Q}\kappa =\mathrm{id}`$. Contradiction!
Case 2. $`\kappa `$ is the critical point of $`\pi _{\overline{W},Q}`$.
Because $``$ is a premouse, we have that $`๐ฅ_{\stackrel{~}{\lambda }}^{\overline{W}}=๐ฅ_{\stackrel{~}{\lambda }}^{\stackrel{~}{W}}`$. By $`\neg 0^{^{}}`$ (i.e., by 2.5), and because $`\pi _{\stackrel{~}{W},Q}\kappa =\mathrm{id}`$, we hence get that in fact $`\pi _{\stackrel{~}{W},Q}F(\kappa )=\mathrm{id}`$. By \[4, Lemma 1.3\] both $`\overline{W}`$ and $`\stackrel{~}{W}`$ have the hull property at $`\kappa `$. Let $`F^{}`$ be the first extender used on the main branch giving $`\pi _{\overline{W},Q}`$.
Fix $`a[F(\kappa )F^{}(\kappa )]^{<\omega }`$, and $`X๐ซ([\kappa ]^{\mathrm{Card}(a)})\overline{W}=๐ซ([\kappa ]^{\mathrm{Card}(a)})\stackrel{~}{W}`$. Pick $`\overline{\tau }`$ and $`\stackrel{}{\zeta }\mathrm{\Xi }^{}`$ such that $`X=\overline{\tau }^{\overline{W}}[\stackrel{}{\zeta }][\kappa ]^{\mathrm{Card}(a)}`$. Then we get that $`aF(X)`$ if and only if $`a\pi (X)`$ if and only if $`a\pi _{\stackrel{~}{W},Q}\pi (X)`$ (as $`\pi _{\stackrel{~}{W},Q}F(\kappa )=\mathrm{id}`$) if and only if $`a\pi _{\overline{W},Q}(X)`$ (as $`\pi _{\stackrel{~}{W},Q}\pi (X)=\overline{\tau }^Q[\stackrel{}{\zeta }][\pi _{\stackrel{~}{W},Q}\pi (\kappa )]^{\mathrm{Card}(a)}`$ and $`\pi _{\overline{W},Q}(X)=\overline{\tau }^Q[\stackrel{}{\zeta }][\pi _{\overline{W},Q}(\kappa )]^{\mathrm{Card}(a)}`$) if and only if $`aF^{}(X)`$. Hence $`F`$ and $`F^{}`$ are compatible.
But we can say more. Using the initial segment for $``$ or for the the model where $`F^{}`$ comes from we can easily deduce that in fact $`F^{}=F`$. However, then, $`\rho _1(๐ฅ_{\stackrel{~}{\lambda }}^{\overline{W}})<F(\kappa )`$. On the other hand, $`F(\kappa )`$ is a cardinal in $`\overline{W}`$, by how $`F`$ was obtained. Contradiction!
$`\mathrm{}`$ (Subclaim 1)
We are now going to define $`\{(\kappa ,x):x\mathrm{\Gamma }_\kappa \}`$ in such a way that $`\mathrm{\Gamma }_\kappa =\mathrm{\Xi }\mathrm{\Gamma }_{\kappa +1}=\mathrm{\Xi }^{}`$ for all ordinals $`\kappa `$, where $`\mathrm{\Xi }\mathrm{\Xi }^{}`$ is as above.
Let $`A_0^{}`$ denote the class of limit points of $`A_0`$. We closely follow \[7, Appendix II\]. Weโll construct $`\mathrm{\Gamma }_i^\delta `$, $`Y_i^\delta `$ for all $`i\mathrm{OR}`$ and certain $`\delta A_0^{}`$, by recursion on $`i\mathrm{OR}`$. We shall inductively maintain that the following statements are true (whenever the objects referred to are defined).
(1)
$$\mathrm{\Gamma }_i^\delta \delta ,\mathrm{and}\mathrm{\Gamma }_i^\delta Y_i^\delta _{\mathrm{\Sigma }_1}๐ฅ_\delta ^{K_{\{\omega _1\}}^c},$$
(2)
$$ij\mathrm{\Gamma }_i^\delta \mathrm{\Gamma }_j^\delta Y_i^\delta Y_j^\delta ,$$
(3)
$$\delta ^{}\delta \mathrm{\Gamma }_i^\delta =\mathrm{\Gamma }_i^\delta ^{}\delta Y_i^\delta =Y_i^\delta ^{}๐ฅ_\delta ^{K_{\{\omega _1\}}^c}$$
To commence, we let $`\mathrm{\Gamma }_0^\delta `$ be all ordinals $`<\delta `$ of cofinality $`\omega `$, and we let $`Y_0^\delta =๐ฅ_\delta ^{K_{\{\omega _1\}}^c}`$, for all $`\delta A_0^{}`$. If $`\lambda `$ is a limit ordinal then we let $`\mathrm{\Gamma }_\lambda ^\delta =_{i<\lambda }\mathrm{\Gamma }_i^\delta `$ and $`Y_\lambda ^\delta =_{i<\lambda }Y_i^\delta `$, for all $`\delta `$ such that $`\mathrm{\Gamma }_i^\delta `$ and $`Y_i^\delta `$ are defined whenever $`i<\lambda `$ (otherwise $`\mathrm{\Gamma }_\lambda ^\delta `$ and $`Y_\lambda ^\delta `$ will be undefined). Notice that (1) will then be true for $`\lambda `$, as any $`๐ฅ_\delta ^{K_{\{\omega _1\}}^c}`$ has a $`\mathrm{\Sigma }_1`$-definable $`\mathrm{\Sigma }_1`$ Skolem function.
Now suppose that $`\mathrm{\Gamma }_i^\delta `$ and $`Y_i^\delta `$ are defined. If $`iY_i^\delta `$ then we put $`\mathrm{\Gamma }_{i+1}^\delta =\mathrm{\Gamma }_i^\delta `$ and $`Y_{i+1}^\delta =Y_i^\delta `$. Suppose now that $`iY_i^\delta `$. Consider
$$\sigma :W=W_i^\delta Y_i^\delta _{\mathrm{\Sigma }_1}๐ฅ_\delta ^{K_{\{\omega _1\}}^c}.$$
Suppose that $`W\delta =\delta `$ (otherwise $`\mathrm{\Gamma }_{i+1}^\delta `$ and $`Y_{i+1}^\delta `$ will be undefined). If there are no $`\alpha <\delta `$ and $`F`$ such that
$$=(๐ฅ_\alpha ^W,F)$$
is a beaver with $`\mathrm{c}.\mathrm{p}.(F)=\sigma ^1(i)`$ then we let $`\mathrm{\Gamma }_{i+1}^\delta =\mathrm{\Gamma }_i^\delta `$ and $`Y_{i+1}^\delta =Y_i^\delta `$. Otherwise let $`\alpha <\lambda `$ be least such that $`F`$ is the unique (by 8.12) $`F`$ so that $``$ as above is a beaver. Let
$$i_F:W_F\stackrel{~}{W}=\stackrel{~}{W}_i^\delta =\mathrm{Ult}_0(W;F).$$
Notice that $`\stackrel{~}{W}`$ must be $`0`$-iterable (as $``$ is a beaver), and $`\stackrel{~}{W}\mathrm{OR}=\delta `$. Let $`๐ฐ`$ and $`๐ฏ`$ denote the iteration trees arising from the coiteration of $`W`$ with $`\stackrel{~}{W}`$. Suppose that:
$``$ $`๐^๐ฐ(0,\mathrm{}]_U=๐^๐ฏ(0,\mathrm{}]_T=\mathrm{}`$, and
$``$ $`_{\mathrm{}}^๐ฐ\mathrm{OR}=_{\mathrm{}}^๐ฏ\mathrm{OR}=\delta `$.
(Otherwise $`\mathrm{\Gamma }_{i+1}^\delta `$ and $`Y_{i+1}^\delta `$ will be undefined.) Now put
$$\mathrm{\Gamma }_{i+1}^\delta =\{\xi \mathrm{\Gamma }_i^\delta :\xi =\sigma (\xi )=\pi _0\mathrm{}^๐ฐ(\xi )=\pi _0\mathrm{}^๐ฏi_F(\xi )\},\mathrm{and}$$
$$Y_{i+1}^\delta =\{\sigma (x):xW\pi _0\mathrm{}^๐ฐ(x)\mathrm{ran}(\pi _0\mathrm{}^๐ฏi_F)\}.$$
This finishes the recursive definition of $`\mathrm{\Gamma }_i^\delta `$ and $`Y_i^\delta `$. Now set $`D_i=\{\delta :\mathrm{\Gamma }_i^\delta `$ and $`Y_i^\delta `$ are defined$`\}`$, and let
$$\mathrm{\Gamma }_i=\underset{\delta D_i}{}\mathrm{\Gamma }_i^\delta ,\mathrm{and}$$
$$Y_i=\underset{\delta D_i}{}Y_i^\delta .$$
Besides (1), (2), (3), we also want to verify, inductively, that:
(4)
$$i\mathrm{OR}\eta (D_i\eta =A_0^{}\eta ),$$
(5)
$$H^{K_{\{\omega _1\}}^c}(\mathrm{\Gamma }_i)Y_iK_{\{\omega _1\}}^c,$$
(6)
$$(\mathrm{\Gamma }_i:i\mathrm{OR})\mathrm{is}\mathrm{as}\mathrm{in}\mathrm{the}\mathrm{statement}\mathrm{of}\mathrm{Claim}1.$$
It is now straightforward to see that in order to inductively prove that (1) through (6) hold for all $`i\mathrm{OR}`$ it suffices to show the following (cf. \[7, Appendix II\]).
###### Proposition 8.14
Suppose that (1) through (6) hold for some $`i\mathrm{OR}`$. Let $`\delta A_0^{}`$ be large enough. Suppose that $`Y_{i+1}^\delta Y_i^\delta `$, and let $`=(๐ฅ_\alpha ^{W_i^\delta },F)`$ be the beaver used to define $`\mathrm{\Gamma }_{i+1}^\delta `$ and $`Y_{i+1}^\delta `$. Let $`๐ฐ`$, $`๐ฏ`$ denote the coiteration of $`W_i^\delta `$ with $`\stackrel{~}{W}_i^\delta =\mathrm{Ult}_0(W_i^\delta ;F)`$. Let $`\theta =\mathrm{lh}(๐ฐ)=\mathrm{lh}(๐ฏ)`$. Let
$$\sigma ^{}:W_iY_iK_{\{\omega _1\}}^c.$$
Let $`๐ฐ^{}`$, $`๐ฏ^{}`$ denote the coiteration of $`W_i`$ with $`\mathrm{Ult}_0(W_i;F)`$. Then $`\delta D_{i+1}`$, $`\theta [0,\mathrm{}]_U^{}[0,\mathrm{}]_T^{}`$, and $`\pi _\theta \mathrm{}^๐ฐ^{}\delta =\pi _\theta \mathrm{}^๐ฏ^{}\delta =\mathrm{id}.`$
It is clear that $`๐ฐ`$ and $`๐ฏ`$ are initial segments of $`๐ฐ^{}`$ and $`๐ฏ^{}`$, respectively. What 8.14 says is that the rest of the coiteration is above $`\delta `$ (and that $`\delta D_{i+1}`$).
Proof of 8.14. It is easy to see that $`\delta D_{i+1}`$. Suppose that $`\theta [0,\mathrm{}]_U^{}[0,\mathrm{}]_T^{}`$. Let $`\alpha +1`$ be least in $`(0,\mathrm{}]_U^{}\theta `$, and let $`\alpha ^{}=U^{}`$-$`\mathrm{pred}(\alpha +1)`$. Let $`\beta +1`$ be least in $`(0,\mathrm{}]_T^{}\theta `$, and let $`\beta ^{}=T^{}`$-$`\mathrm{pred}(\alpha +1)`$. Then $`\mu =\mathrm{min}\{\mathrm{c}.\mathrm{p}.(\pi _\alpha ^{}\mathrm{}^๐ฐ^{}),`$ $`\mathrm{c}.\mathrm{p}.(\pi _\beta ^{}\mathrm{}^๐ฏ^{})\}<\delta `$. As $`\delta A_0^{}`$, we in fact get essentially as in the proof of 5.11 Claim 2 Subclaim 1 that $`\mu =\mathrm{c}.\mathrm{p}.(\pi _\alpha ^{}\mathrm{}^๐ฐ^{})=\mathrm{c}.\mathrm{p}.(\pi _\beta ^{}\mathrm{}^๐ฏ^{})`$. But then 8.10 gives a standard contradiction (see the proof of 8.10)! This yields 8.14.
Now the previous construction together with 8.12 as well as the set theoretical definability of โbeaver-hoodโ (see \[4, ยง2\]) finishes the proof of Claim 1.
$`\mathrm{}`$ (Claim 1)
Now let us assume that $`\mathrm{Def}(K_{\{\omega _1\}}^c,A_0)`$ is bounded, $`\beta =\mathrm{sup}(\mathrm{Def}(K_{\{\omega _1\}}^c,A_0))`$, say. We aim to derive a contradiction. Fix $`(\mathrm{\Gamma }_\kappa :\kappa \mathrm{OR})`$ as given by Claim 1.
Let $`b_\kappa `$ denote the least ordinal in $`H^{K_{\{\omega _1\}}^c}(\mathrm{\Gamma }_\kappa )\beta `$ for $`\kappa \beta `$ (hence $`b_\kappa \kappa `$). By further thinning out the $`\mathrm{\Gamma }_\kappa `$โs if necessary we may assume without loss of generality that $`b_\kappa <b_{\kappa +1}`$ for $`\kappa \beta `$.
Claim 2. There is some $`\nu >\beta `$, a limit of $`b_\kappa `$โs, such that $`\nu H^{K_{\{\omega _1\}}^c}(\nu \mathrm{\Gamma }_{\nu +1})`$.
Given Claim 2, the proof of 8.13 can be completed as follows (cf. \[27, p. 38\]). Fix $`\nu `$ as in Claim 2. Let $`\nu =\tau ^{K_{\{\omega _1\}}^c}[\stackrel{}{\alpha },\stackrel{}{\rho }]`$ where $`\stackrel{}{\alpha }<\nu `$ and $`\stackrel{}{\rho }\mathrm{\Gamma }_{\nu +1}`$. Then $`\stackrel{}{\alpha }<b_\kappa `$ for some $`b_\kappa <\nu `$. Hence
$$K_{\{\omega _1\}}^c\stackrel{}{\alpha }<b_\kappa (b_\kappa <\tau ^{K_{\{\omega _1\}}^c}[\stackrel{}{\alpha },\stackrel{}{\rho }]<b_{\nu +1}).$$
But $`b_\kappa `$, $`\stackrel{}{\rho }`$, $`b_{\nu +1}\mathrm{\Gamma }_\kappa `$, and $`H^{K_{\{\omega _1\}}^c}(\mathrm{\Gamma }_\kappa )K_{\{\omega _1\}}^c`$, so that there is some $`\stackrel{}{\alpha }^{}b_\kappa H^{K_{\{\omega _1\}}^c}(\mathrm{\Gamma }_\kappa )\beta `$ with
$$b_\kappa <\tau ^{K_{\{\omega _1\}}^c}[\stackrel{}{\alpha }^{},\stackrel{}{\rho }]<b_{\nu +1}.$$
By $`\stackrel{}{\alpha }^{}\beta H^{K_{\{\omega _1\}}^c}(\mathrm{\Gamma }_\kappa )`$, we have that $`\stackrel{}{\alpha }^{}\mathrm{Def}(K_{\{\omega _1\}}^c,A_0)`$. Hence $`\tau ^{K_{\{\omega _1\}}^c}[\stackrel{}{\alpha }^{},\stackrel{}{\rho }]H^{K_{\{\omega _1\}}^c}(\mathrm{\Gamma }_{\nu +1})`$ and $`\beta b_\kappa \tau ^{K_{\{\omega _1\}}^c}[\stackrel{}{\alpha }^{},\stackrel{}{\rho }]<b_{\nu +1}`$. This contradicts the definition of $`b_{\nu +1}`$!
Proof of Claim 2. Let $`C`$ denote the class of all limit points of $`\{b_\kappa :\beta \kappa <\mathrm{OR}\}`$. Of course, $`C`$ is club in $`\mathrm{OR}`$.
Let us suppose Claim 2 to be false, so that for all $`\kappa C`$ do we have
$$\sigma _\kappa :K_{\{\omega _1\}}^{c}{}_{\kappa }{}^{}H^{K_{\{\omega _1\}}^c}(\kappa \mathrm{\Gamma }_{\kappa +1})K_{\{\omega _1\}}^c$$
with $`\mathrm{c}.\mathrm{p}.(\sigma _\kappa )=\kappa `$ and $`\sigma _\kappa (\kappa )b_{\kappa +1}`$. Let $`F_\kappa =\sigma _\kappa ๐ซ(\kappa )K_{\{\omega _1\}}^{c}{}_{\kappa }{}^{}`$ for $`\kappa C`$.
The following is due to John Steel and is included here with his permission.
Subclaim 2 (Steel). The class $`\{\kappa C:F_\kappa `$ is countably complete $`\}`$ contains a club.
Proof. Suppose not, so that
$$S_0=\{\kappa C:F_\kappa \mathrm{is}\mathrm{not}\mathrm{countably}\mathrm{complete}\}$$
is stationary. We may then pick $`((a_\kappa ^n,X_\kappa ^n):n<\omega \kappa C)`$ such that $`a_\kappa ^nF_\kappa (X_\kappa ^n)`$, but there is no order preserving $`\tau :_{n<\omega }a_\kappa ^n\kappa `$ with $`\tau \mathrm{"}a_\kappa ^nX_\kappa ^n`$ for any $`\kappa C`$.
We have to define a function $`G:S_0V`$, saying how the $`a_\kappa ^n`$โs sit inside $`_na_\kappa ^n`$. Let $`\tau _\kappa :_na_\kappa ^n\mathrm{otp}(_na_\kappa ^n)<\omega _1`$ be order preserving, and let $`G(\kappa )=(\mathrm{otp}(_na_\kappa ^n),(\tau _\kappa \mathrm{"}a_\kappa ^n:n<\omega ))`$.
Now by 8.9, there is some stationary $`S_1S_0`$ such that $`K_{\{\omega _1\}}^c`$ has the hull property at every $`\kappa S_1`$. Hence for any $`n<\omega `$ and $`\kappa S`$ may we pick a term $`\tau _\kappa ^n`$ and $`\stackrel{}{\alpha }_\kappa ^n<\kappa `$ and $`\stackrel{}{\gamma }_\kappa ^n\mathrm{\Gamma }_{\kappa +1}`$ such that
$$\sigma _\kappa (X_\kappa ^n)=(\tau _\kappa ^n)^{K_{\{\omega _1\}}^c}[\stackrel{}{\alpha }_\kappa ^n,\stackrel{}{\gamma }_\kappa ^n].$$
We now consider
$$F(\kappa )=(G(\kappa ),(\tau _\kappa ^n,\stackrel{}{\alpha }_\kappa ^n,((\tau _\kappa ^n)^{K_{\{\omega _1\}}^c}[\stackrel{}{\alpha },\stackrel{}{\gamma }_\kappa ^n]\beta :\stackrel{}{\alpha }<\beta ):n<\omega )),$$
being essentially a regressive function on $`S`$. By Fodor, there is an unbounded $`DS`$ such that $`F`$ is constant on $`D`$. Fix $`\kappa <\kappa ^{}D`$ with $`\sigma _\kappa (\kappa )b_{\kappa +1}\kappa ^{}`$.
Notice first that we have an order preserving map
$$\tau :_na_\kappa ^{}^n_na_\kappa ^n$$
with $`\tau \mathrm{"}a_\kappa ^{}^n=a_\kappa ^n`$ for all $`n<\omega `$, due to the fact that $`G(\kappa )=G(\kappa ^{})`$. Hence we would have a contradiction to the choice of $`(a_\kappa ^{}^n,X_\kappa ^{}^n:n<\omega )`$ if we were able to show that $`a_\kappa ^nX_\kappa ^{}^n`$ for all $`n<\omega `$.
In order to do this, as $`a_\kappa ^n<\sigma _\kappa (\kappa )b_{\kappa +1}`$ and $`a_\kappa ^n\sigma _\kappa (X_\kappa ^n)`$ by the choice of $`(a_\kappa ^n,X_\kappa ^n:n<\omega )`$ (and because $`X_\kappa ^{}^nb_{\kappa +1}=\sigma _\kappa ^{}(X_\kappa ^{}^n)b_{\kappa +1}`$ by $`\kappa ^{}b_{\kappa +1}`$), it suffices to establish that
$$\sigma _\kappa (X_\kappa ^n)b_{\kappa +1}=\sigma _\kappa ^{}(X_\kappa ^{}^n)b_{\kappa +1}.$$
We have that $`\tau _\kappa ^n=\tau _\kappa ^{}^n=\tau `$ and $`\stackrel{}{\alpha }_\kappa ^n=\stackrel{}{\alpha }_\kappa ^{}^n=\stackrel{}{\alpha }`$ for some $`\tau `$, $`\stackrel{}{\alpha }`$, and $`\sigma _\kappa (A_\kappa ^n)=\tau ^{K_{\{\omega _1\}}^c}(\stackrel{}{\alpha },\stackrel{}{\gamma }_\kappa ^n)`$ and $`\sigma _\kappa ^{}(X_\kappa ^{}^n)=\tau ^{K_{\{\omega _1\}}^c}(\stackrel{}{\alpha },\stackrel{}{\gamma }_\kappa ^{}^n)`$. So if $`\sigma _\kappa (A_\kappa ^n)b_{\kappa +1}\sigma _\kappa ^{}(X_\kappa ^{}^n)b_{\kappa +1}`$ then
$$K_{\{\omega _1\}}^c\stackrel{}{\alpha }<b_{\kappa +1}\tau ^{K_{\{\omega _1\}}^c}[\stackrel{}{\alpha },\stackrel{}{\gamma }_\kappa ^n]\tau ^{K_{\{\omega _1\}}^c}[\stackrel{}{\alpha },\stackrel{}{\gamma }_\kappa ^{}^n].$$
As $`b_{\kappa +1}`$, $`\stackrel{}{\gamma }_\kappa ^n`$, $`\stackrel{}{\gamma }_\kappa ^{}^n\mathrm{\Gamma }_{\kappa +1}`$, there is some witness $`\stackrel{}{\alpha }^{}h^{K_{\{\omega _1\}}^c}(\mathrm{\Gamma }_{\kappa +1})`$. But $`\stackrel{}{\alpha }^{}h^{K_{\{\omega _1\}}^c}(\mathrm{\Gamma }_{\kappa +1})b_{\kappa +1}`$, so that $`\stackrel{}{\alpha }^{}\mathrm{Def}(K_{\{\omega _1\}}^c,A_0)\beta `$.
Now we have that
$$K_{\{\omega _1\}}^cd<b_{\kappa +1}(d\tau ^{K_{\{\omega _1\}}^c}[\stackrel{}{\alpha }^{},\stackrel{}{\gamma }_\kappa ^n]d\tau ^{K_{\{\omega _1\}}^c}[\stackrel{}{\alpha }^{},\stackrel{}{\gamma }_\kappa ^{}^n]).$$
Here we have that $`\stackrel{}{\alpha }^{}`$, $`\stackrel{}{\gamma }_\kappa ^n`$, $`\stackrel{}{\gamma }_\kappa ^{}^n\mathrm{\Gamma }_{\kappa +1}`$, so that there is some witness $`dh^{K_{\{\omega _1\}}^c}(\mathrm{\Gamma }_{\kappa +1})`$. But $`d<b_{\kappa +1}`$ and $`b_{\kappa +1}h^{K_{\{\omega _1\}}^c}(\mathrm{\Gamma }_{\kappa +1})=\mathrm{Def}(K_{\{\omega _1\}}^c,A_0)\beta `$, so that we may conclude that $`d<\beta `$. I.e., $`\tau ^{K_{\{\omega _1\}}^c}[\stackrel{}{\alpha }^{},\stackrel{}{\gamma }_\kappa ^n]\beta \tau ^{K_{\{\omega _1\}}^c}[\stackrel{}{\alpha }^{},\stackrel{}{\gamma }_\kappa ^{}^n]\beta `$.
However, as $`\stackrel{}{\alpha }^{}<\beta `$, $`\kappa <\kappa ^{}D`$ gives us that $`\tau ^{K_{\{\omega _1\}}^c}[\stackrel{}{\alpha }^{},\stackrel{}{\gamma }_\kappa ^n]\beta =\tau ^{K_{\{\omega _1\}}^c}[\stackrel{}{\alpha }^{},\stackrel{}{\gamma }_\kappa ^{}^n]\beta `$. Contradiction!
$`\mathrm{}`$ (Subclaim 2)
We may hence pick $`\kappa C`$ such that $`F_\kappa `$ is countably complete and $`K_{\{\omega _1\}}^c`$ has the hull property at $`\kappa `$. We then have an immediate contradiction with 4.2.
$`\mathrm{}`$ (Claim 2)
$`\mathrm{}`$ (8.13)
###### Definition 8.15
$`(\neg 0^{^{}})`$ We shall write $`K`$ for $`K(A_0)`$. $`K`$ is called the core model below $`0^{^{}}`$.
###### Lemma 8.16
$`(\neg 0^{^{}})`$ $`K`$ is full, i.e., whenever $`F`$ is a countably complete extender such that $`(๐ฅ_\alpha ^K,F)`$ is a premouse then $`F=E_\alpha ^K`$.
Proof. Let $`W`$ be such that $`\alpha \mathrm{Def}(W,A_0)`$. As $`F`$ is countably complete, $`\mathrm{Ult}(W,F)`$ is iterable. In particular, $`๐ฅ_\alpha ^W=๐ฅ_\alpha ^K`$. Then \[4, Lemma 1.6\] gives the desired result.
$`\mathrm{}`$
Combined with 4.4 this gives at once:
###### Corollary 8.17
$`(\neg 0^{^{}})`$ $`K`$ is universal.
It is now clear that the proofs in and give the following.
###### Theorem 8.18
$`(\neg 0^{^{}})`$ Let $`\beta \omega _2`$. Then $`\mathrm{cf}^V(\beta ^{+K})\mathrm{Card}^V(\beta )`$.
As in \[24, Theorem 1.1\], weโll also have the following.
###### Theorem 8.19
$`(\neg 0^{^{}})`$ Let $`1\alpha \omega _1`$, and suppose that
$$๐ฅ_{\omega _1}^K\mathrm{there}\mathrm{are}<\omega \alpha \mathrm{many}\mathrm{strong}\mathrm{cardinals}.$$
Then the set of reals coding $`๐ฅ_{\omega _1}^K`$ is an element of $`J_{1+\alpha }()`$.
For the projective level, \[4, Theorems 3.4 and 3.6\] (which are due to the present author) give the following refinements.
###### Theorem 8.20
$`(\neg 0^{^{}})`$ Let $`n<\omega `$. Suppose that
$$๐ฅ_{\omega _1}^K\mathrm{there}\mathrm{are}\mathrm{exactly}n\mathrm{strong}\mathrm{cardinals}.$$
If $`\omega _1`$ is inaccessible in $`K`$ then $`๐ฅ_{\omega _1}^K`$ is (lightface) $`\mathrm{\Delta }_{n+5}^1`$ in the codes; and if $`\omega _1`$ is a successor cardinal in $`K`$ then $`๐ฅ_{\omega _1}^K`$ is $`\mathrm{\Delta }_{n+4}^1(x)`$ in the codes for some $`x`$ coding an initial segment of $`๐ฅ_{\omega _1}^K`$.
Jensen has shown in \[8, ยง5.3 Lemma 5\] that if $`0^{\mathrm{}}`$ does not exist then every universal weasel is an iterate of $`K`$ (via a normal iteration). His proof in fact straightforwardly generalizes to the situation where $`K^c`$ does not have a cardinal which is strong up to a measurable cardinal, i.e., a measurable cardinal $`\kappa `$ such that
$$๐ฅ_\kappa ^{K^c}\mathrm{"}\mathrm{there}\mathrm{is}\mathrm{a}\mathrm{strong}\mathrm{cardinal}.\mathrm{"}$$
On the other hand, we are now going to construct an example of a universal weasel not being an iterate of $`K`$, assuming that $`K^c`$ does have a โstrong up to a measurable.โ The example is due to John Steel and is included here with his permission. (Recall that \[27, p. 86\] had already shown indirectly that such a weasel has to exist if $`K\mathrm{HC}`$ is not $`\mathrm{\Delta }_5^1`$.)
###### Lemma 8.21
(Steel, $`\neg 0^{^{}}`$ ) Suppose that $`\mu <\kappa `$ are such that $`\kappa `$ is measurable in $`K`$ and $`๐ฅ_\kappa ^K`$$`\mu `$ is a strong cardinal.โ Then if $`H`$ is $`Col(\omega ,\kappa ^{+K})`$-generic over $`K`$, inside $`K[H]`$ there exists a universal weasel $`W`$ which is not an iterate of $`K`$. In fact, $`W`$ may be chosen such that $`W(๐ฅ_{\kappa ^{+K}}^K,F)`$ for some extender $`F`$.
Moreover, if $`\mu `$ is the only $`\overline{\mu }`$ with $`๐ฅ_\kappa ^K`$$`\overline{\mu }`$ is a strong cardinal,โ then no universal $`W^{}(๐ฅ_{\kappa ^{+K}}^K,F)`$ has the definability property at $`\mu `$.
Proof. First fix $`gK[H]`$ being $`Col(\omega ,\mu ^{++K})`$-generic over $`K`$. We shall use a theorem of Woodin (unpublished) which tells us that in $`K[g]`$ there is a tree $`T_3`$ projecting to a universal $`\mathrm{\Pi }_3^1`$-set of reals in all further small extensions of $`K[g]`$, where โsmallโ means that they are obtained by forcing with some $`P๐ฅ_\alpha ^K[g]`$ where $`\alpha =E_\beta ^K(\mu )`$ for some $`E_\beta ^K`$ with critical point $`\mu `$. This implies that in $`K[g]`$ there is a tree $`T`$ projecting to the set of all real codes for mice $``$ such that $`K^c()`$ exists, where $`T`$ works in all such extensions; the reason is that this is a $`\mathrm{\Pi }_3^1`$-set of reals (cf. \[4, Corollary 2.18 (a)\]).
Let $`E=E_\nu ^K`$ witness that $`\kappa `$ is measurable in $`K`$, and let $`\pi :K_EM`$ with $`M`$ being transitive. We may extend $`\pi `$ to $`\stackrel{~}{\pi }:K[g]M[g]`$, and by elementarity, $`\stackrel{~}{\pi }(T)M[g]`$ is a tree projecting to the set of all real codes for mice $``$ such that $`K^c()`$ exists, with $`\stackrel{~}{\pi }(T)`$ working in all further small extensions of $`M[g]`$.
Now let $`G`$ be $`Col(\omega ,\nu )`$-generic over $`K[g]`$, so that $`G`$ is $`Col(\omega ,\nu )`$-generic over $`M[g]`$, too. Notice that $`\stackrel{~}{\pi }(T)`$ still โworksโ in $`M[g][G]`$. In $`M[g][G]`$, let $`x`$ code $`๐ฅ_\nu ^M=๐ฅ_\nu ^K`$. We may build a tree $`T^{}M[g][G]`$ searching for a pair $`(y,f)`$ such that $`(xy,f)[\stackrel{~}{\pi }(T)]`$. Here, by $`ab`$ we mean a canonical code for $`(๐ฉ,F)`$ obtained from $`(a,b)`$ where $`a`$ codes $`๐ฉ`$ and $`b`$ codes $`F`$.
We claim that $`T^{}`$ is ill-founded (in $`K[g][G]`$, and hence in $`M[g][G]`$). Let $`T_n`$ denote $`T`$ up to the $`n^{\mathrm{th}}`$ level. Well, if $`yK[g][G]`$ codes $`E_\nu ^K`$ then $`xyp[T]`$, and hence if $`fK[g][G]`$ is such that $`(xy,f)[T]`$ then
$$n(xyn,fn)T_n,\mathrm{hence}$$
$$n(xyn,\pi (fn))T_n,$$
so that $`(y,_n\pi (fn))T^{}`$.
Thus in $`M[g]`$, the following holds true:
$$||_{Col(\omega ,\nu )}\mathrm{`}\mathrm{`}\mathrm{there}\mathrm{is}\mathrm{an}\mathrm{extender}F\mathrm{such}\mathrm{that}K^c((๐ฅ_\nu ^M,F))\mathrm{exists},^{\prime \prime }$$
and by elementarity of $`\stackrel{~}{\pi }`$, in $`K[g]`$ we have that
$$||_{Col(\omega ,\kappa ^+)}\mathrm{`}\mathrm{`}\mathrm{there}\mathrm{is}\mathrm{an}\mathrm{extender}F\mathrm{such}\mathrm{that}K^c((๐ฅ_{\kappa ^+}^K,F))\mathrm{exists}.^{\prime \prime }$$
Now let $`\overline{H}K[H]`$ be $`Col(\omega ,\kappa ^+)`$-generic over $`K[g]`$. By what we have shown, in $`K[g][\overline{H}]`$, and hence in $`K[H]`$ we may pick some $`F`$ such that $`(๐ฅ_{\kappa ^+}^K,F)`$ is a premouse and $`W=K^c((๐ฅ_{\kappa ^+}^K,F))`$ exists. Notice that $`\rho _1((๐ฅ_{\kappa ^+}^K,F))<\kappa `$.
It is now easy to see that $`W`$ cannot be an iterate of $`K`$, as $`\kappa ^{+K}`$ is a cardinal in every such iterate, whereas it is not a cardinal in $`W`$.
Now suppose that the universal weasel $`W^{}(๐ฅ_{\kappa ^+}^K,F)`$ would have the definability property at $`\mu `$, and that $`\mu `$ is the only $`\overline{\mu }`$ with $`๐ฅ_\kappa ^K`$$`\overline{\mu }`$ is a strong cardinal.โ Then the coiteration of $`W^{}`$ and $`K`$ would be above $`\kappa `$ on both sides by \[4, Corollary 1.5\]. Weโd thus get that $`\kappa ^{+W^{}}=\kappa ^{+K}`$. Contradiction!
$`\mathrm{}`$ (8.21)
## 9 An application.
By classical results, Projective Determinacy ($`\mathrm{๐ฏ๐ฃ}`$, for short) implies that every projective set of reals is Lebesgue measurable and has the property of Baire, and that every projective subset of $`\times `$ has a uniformizing function with a projective graph (cf. for example \[21, 6A.16\], \[21, 6A.18\], and \[21, 6C.5\]). In , Woodin asked whether $`\mathrm{๐ฏ๐ฃ}`$ is actually equivalent with the fact that every projective set of reals is Lebesgue measurable and has the property of Baire, and that every projective subset of $`\times `$ has a uniformizing function with a projective graph. This question became the $`12^{\mathrm{th}}`$ Delfino problem (see ), and it was widely believed to have an affirmative answer, โ a belief which was supported by a theorem of Woodinโs according to which at least $`๐ท`$$`{}_{\mathrm{๐}}{}^{}{}_{}{}^{\mathrm{๐}}`$ determinacy does follow from these โanalyticโ consequences of $`\mathrm{๐ฃ๐}`$.
However, Steel in the fall of 1997 proved that the answer to Woodinโs question is negative. This left open the question of the exact consistency strength of the assumption in the statement of the $`\mathrm{๐๐}^{\mathrm{๐ญ๐ก}}`$ Delfino problem, that is, of above-mentioned โanalyticโ consequences of $`\mathrm{๐ฃ๐}`$. Our application, which yields the following theorem 9.1, determines that strength in terms of large cardinal assumptions.<sup>6</sup><sup>6</sup>6It was actually the wish to complete the proof of this theorem which motivated our work reported in the previous sections.
###### Theorem 9.1
The following theories are equiconsistent.
(1) $`\mathrm{๐ญ๐๐}\mathbf{+}`$ โevery projective set of reals is Lebesgue measurable and has the property of Baireโ $`\mathbf{+}`$ โevery projective subset of $`\mathbf{\times }`$ has a uniformizing function with a projective graph,โ
(2) $`\mathrm{๐ญ๐๐}\mathbf{+}`$ โevery projective set of reals is Lebesgue measurable and has the property of Baireโ $`\mathbf{+}`$ โthe pointclass consisting of the projective sets of reals has the scale property,โ and
(3) $`\mathrm{๐ญ๐๐}\mathbf{+}`$ โthere are infinitely many cardinals $`๐_\mathrm{๐}\mathbf{<}๐_\mathrm{๐}\mathbf{<}๐_\mathrm{๐}\mathbf{<}\mathbf{}`$ each of which is $`\mathbf{(}\mathrm{๐ฌ๐ฎ๐ฉ}_{๐ง\mathbf{<}๐}๐_๐ง\mathbf{)}^\mathbf{+}`$-strong, i.e., for every $`๐ง\mathbf{<}๐`$ and for every $`๐\mathbf{}\mathrm{๐ฌ๐ฎ๐ฉ}_{๐ง\mathbf{<}๐}๐_๐ง`$ there is an elementary embedding $`๐\mathbf{:}๐\mathbf{}๐`$ with $`๐`$ being transitive, $`๐\mathbf{.}๐ฉ\mathbf{.}\mathbf{(}๐\mathbf{)}\mathbf{=}๐_๐ง`$, and $`๐\mathbf{}๐`$.โ
Proof. Con(3) $`\mathbf{}`$ Con(2) was shown by Steel (building on work of Woodin, cf. ). (2) $`\mathbf{}`$ (1) is classical (cf. for example \[21, 4E.3\]). Con(1) $`\mathbf{}`$ Con(3) was shown in (building on Woodinโs and the authorโs work on the complexity of $`๐ฒ\mathbf{}\mathrm{๐๐}`$, cf. 8.19 and 8.20 above; cf. \[4, Theorem 4.2\]) with $`\mathrm{๐ญ๐๐}`$ replaced by $`\mathrm{๐ญ๐๐}\mathbf{+}`$$`^\mathbf{\#}`$ existsโ (the authors of used the core model theory of which needs the existence of a measurable cardinal, or some substitute for it). By our work done here in the previous sections, $`^\mathbf{\#}`$ is now no longer needed for running the arguments of .
$`\mathbf{}`$ (9.1)
The reader will have no problems with finding the appropriate reformulations of \[24, Theorems 1.1, 1.4, and 1.5\] in light of our work done here.
We finally remark that the following variants of the $`\mathrm{๐๐}^{\mathrm{๐ญ๐ก}}`$ Delfino problem remain wide open.
Question. Suppose that every projective set of reals is Lebesgue measurable and has the property of Baire. Suppose that either
(a) for each $`๐\mathbf{<}๐`$ does every $`๐ท`$$`{}_{\mathrm{๐}๐\mathbf{+}\mathrm{๐}}{}^{}{}_{}{}^{\mathrm{๐}}`$ subset of $`\mathbf{\times }`$ have a uniformizing function with a graph in $`๐ท`$$`{}_{\mathrm{๐}๐\mathbf{+}\mathrm{๐}}{}^{}{}_{}{}^{\mathrm{๐}}`$, or else
(b) every lightface projective subset of $`\mathbf{\times }`$ has a uniformizing function with a graph which is lightface projective.
Does then Projective Determinacy hold?
Steel has shown (cf. \[27, Corollary 7.14\]) that variant (a) implies $`๐ซ`$$`{}_{\mathrm{๐}}{}^{}{}_{}{}^{\mathrm{๐}}`$ determinacy. It might be reasonable to conjecture that variant (a) in fact implies $`\mathrm{๐ฃ๐}`$. On the other hand it is conceivable that a refinement of coding arguments as being developed in will establish that variant (b) does not imply $`๐ซ`$$`{}_{\mathrm{๐}}{}^{}{}_{}{}^{\mathrm{๐}}`$ determinacy. |
warning/0002/cond-mat0002049.html | ar5iv | text | # Second Topological Moment โจ๐ยฒโฉ of Two Closed Entangled Polymers
## Abstract
We calculate exactly by field theoretical techniques the second topological moment $`m^2`$ of entanglement of two closed polymers $`P_1`$ and $`P_2`$. This result is used to estimate approximately the mean square average of the linking number of a polymer $`P_1`$ in solution with other polymers.
1. Consider two closed polymers $`P_1`$ and $`P_2`$ which statistically can be linked with each other any number of times $`m=0,1,2,\mathrm{}`$. The situation is illustrated in Fig. 1 for $`m=2`$.
An important physical quantity is the probability distribution of the linking number $`m`$ as a function of the lengths of $`P_1`$ and $`P_2`$. As a first step towards finding it we calculate, in this note, an exact expression for the second moment of the distribution, $`m^2`$.
An approximate result for this quantity was obtained before in Ref. on the basis of a a mean-field method, considering the density of bond vectors of $`P_2`$ as Gaussian random variables. Such methods are usually quite accurate when a large number of polymers is involved . As an unpleasant feature, however, they they introduce a dependence on the source of Gaussian noise, and modify the critical behavior of the system, whereas topological interactions are not expected to do that . Our note goes therefore an important step beyond this approximation. It treats the two-polymer problem exactly, and contains an application to the topological entanglement in diluted solutions. The relevance of the two-polymer systems to such systems was emphasized in . Focusing attention upon a particular molecule, $`P_1`$, one may imagine all others to form a single long effective molecule $`P_2`$.
2. Let $`G_m(๐ฑ_1,๐ฑ_2;L_1,L_2)`$ be the configurational probability to find the polymer $`P_1`$ of length $`L_1`$ with fixed coinciding end points at $`๐ฑ_1`$ and the polymer $`P_2`$ of length $`L_2`$ with fixed coinciding end points at $`๐ฑ_2`$, topologically entangled with a Gaussian linking number $`m`$.
The second moment $`m^2`$ is defined by the ratio of integrals
$`m^2={\displaystyle \frac{d^3x_1d^3x_2_{\mathrm{}}^+\mathrm{}๐mm^2G_m(๐ฑ_1,๐ฑ_2;L_1,L_2)}{d^3x_1d^3x_2_{\mathrm{}}^+\mathrm{}๐mG_m(๐ฑ_1,๐ฑ_2;L_1,L_2)}}`$ (1)
The denominator in (1) plays the role of a partition function:
$`Z{\displaystyle d^3x_1d^3x_2_{\mathrm{}}^+\mathrm{}๐mG_m(๐ฑ_1,๐ฑ_2;L_1,L_2)}.`$ (2)
Due to the translational invariance of the system, the probabilities depend only on the differences between the end point coordinates:
$`G_m(๐ฑ_1,๐ฑ_2;L_1,L_2)=G_m(๐ฑ_1๐ฑ_2;L_1,L_2).`$ (3)
Thus, after a shift of variables, the spatial double integrals in (1) can be rewritten as
$`{\displaystyle d^3x_1d^3x_2G_m(๐ฑ_1๐ฑ_2;L_1,L_2)}=V{\displaystyle d^3xG_m(๐ฑ;L_1,L_2)},`$ (4)
where $`V`$ denotes the total volume of the system.
3. The most efficient way of describing the statistical fluctuations of the polymers $`P_1`$ and $`P_2`$ is by two complex polymer fields $`\psi _1^{a_1}(๐ฑ_1)`$ and $`\psi _2^{a_2}(๐ฑ_2)`$ with $`n_1`$ and $`n_2`$ replicas $`(a_1=1,\mathrm{},n_1,a_2=1,\mathrm{},n_2)`$. At the end we shall take $`n_1,n_20`$ to ensure that these fields describe only one polymer each .
For these fields we define an auxiliary probability $`G_\lambda (\stackrel{}{๐ฑ}_1,\stackrel{}{๐ฑ}_2;\stackrel{}{z})`$ to find the polymer $`P_1`$ with open ends at $`๐ฑ_1,๐ฑ_1^{}`$ and the polymer $`P_2`$ with open ends at $`๐ฑ_2,๐ฑ_2^{}`$. The double vectors $`\stackrel{}{๐ฑ}_1(๐ฑ_1,๐ฑ_1^{})`$ and $`\stackrel{}{๐ฑ}_2(๐ฑ_2,๐ฑ_2^{})`$ collect initial and final endpoints of the two polymers $`P_1`$ and $`P_2`$. Here we follow the approach of Edwards , in which one starts with open polymers with fixed ends. The case of closed polymers, where $`m`$ becomes a true topological number and it is thus relevant in the present context, is recovered in the limit of coinciding extrema. We notice that in this way one introduces in the configurational probability an artificial dependence on the fixed points $`๐ฑ_1`$ and $`๐ฑ_2`$. In physical situations, however, the fluctuations of the polymers are entirely free. For this reason we have averaged in (1) over all possible fixed points by means of the integrations in $`d^3๐ฑ_1d^3๐ฑ_2`$.
The auxiliary probability $`G_\lambda (\stackrel{}{๐ฑ}_1,\stackrel{}{๐ฑ}_2;\stackrel{}{z})`$ is given by a functional integral
$`G_\lambda (\stackrel{}{๐ฑ}_1,\stackrel{}{๐ฑ}_2;\stackrel{}{z})=\underset{n_1,n_20}{lim}{\displaystyle ๐(\mathrm{fields})}`$ (5)
$`\times \psi _1^{a_1}(๐ฑ_1)\psi _1^{a_1}(๐ฑ_1^{})\psi _2^{a_2}(๐ฑ_2)\psi _2^{a_2}(๐ฑ_2^{})e^๐,`$ (6)
where $`๐(\text{fields})`$ indicates the measure of functional integration and $`a_1,a_2`$ are now fixed replica indices. $`๐=๐_{\mathrm{CS}}+๐_{\mathrm{pol}}`$ is the action governing the fluctuations. It consists of a polymer action
$`๐_{\mathrm{pol}}={\displaystyle \underset{i=1}{\overset{2}{}}}{\displaystyle d^3๐ฑ\left[|\overline{๐}^i\mathrm{\Psi }_i|^2+m_i^2|\mathrm{\Psi }_i|^2\right]}.`$ (7)
and a Chern-Simons action to describe the linking number $`m`$
$`๐_{\mathrm{CS}}=i\kappa {\displaystyle d^3x\epsilon _{\mu \nu \rho }A_1^\mu _\nu A_2^\rho },`$ (8)
In Eq. (8) we have omitted a gauge fixing term, which enforces the Lorentz gauge. The effects of self-entanglement and of the so-called excluded-volume interactions are ignored. The Chern-Simons fields are coupled to the polymer fields by the covariant derivatives $`๐^i=\mathbf{}+i\gamma _i๐^i,`$ with the coupling constants $`\gamma _{1,2}`$ given by $`\gamma _1=\kappa ,\gamma _2=\lambda .`$ The square masses of the polymer fields are given by $`m_i^2=2Mz_i`$, where $`M=3/a`$, with $`a`$ being the length of the polymer links, and $`z_i`$ the chemical potentials of the polymers, measured in units of the temperature. The chemical potentials are conjugate variables to the length parameters $`L_1`$ and $`L_2`$, respectively. The symbols $`\mathrm{\Psi }_i`$ collect the replicas $`\psi _i^{a_i}`$ of the two polymer fields. Let us note that in the topological Landau-Ginzburg model (6) the Chern-Simons fields do not change the critical behavior of the system, as expected.
The parameter $`\lambda `$ is conjugate to the linking number $`m`$. We can therefore calculate the desired probability $`G_m(\stackrel{}{๐ฑ}_1,\stackrel{}{๐ฑ}_2;L_1,L_2)`$ from the auxiliary one $`G_\lambda (\stackrel{}{๐ฑ}_1,\stackrel{}{๐ฑ}_2;\stackrel{}{z})`$ by Laplace integrals over $`\stackrel{}{z}=(z_1,z_1)`$ and an inverse Fourier transformation over $`\lambda `$.
4. Let us use the polymer field theory to calculate the partition function (2). It is given by the integral over the auxiliary probabilities
$`Z`$ $`=`$ $`{\displaystyle d^3x_1d^3x_2\underset{\genfrac{}{}{0pt}{}{๐ฑ_1^{}๐ฑ_1}{๐ฑ_2^{}๐ฑ_2}}{lim}_{ci\mathrm{}}^{c+\mathrm{}}\frac{Mdz_1}{2\pi i}\frac{Mdz_2}{2\pi i}e^{z_1L_1+z_2L_2}}`$ (10)
$`\times {\displaystyle _{\mathrm{}}^+\mathrm{}}dm{\displaystyle _{\mathrm{}}^+\mathrm{}}d\lambda e^{im\lambda }G_\lambda (\stackrel{}{๐ฑ}_1,\stackrel{}{๐ฑ}_2;\stackrel{}{z}).`$
The integration over $`m`$ is trivial and gives $`2\pi \delta (\lambda )`$, enforcing $`\lambda =0`$, so that
$`Z`$ $`=`$ $`{\displaystyle d^3x_1d^3x_2\underset{\genfrac{}{}{0pt}{}{๐ฑ_1๐ฑ_1^{}}{๐ฑ_2^{}๐ฑ_2}}{lim}_{ci\mathrm{}}^{c+i\mathrm{}}\frac{Mdz_1}{2\pi i}\frac{Mdz_2}{2\pi i}e^{z_1L_1+z_2L_2}}`$ (12)
$`\times G_{\lambda =0}(\stackrel{}{๐ฑ}_1,\stackrel{}{๐ฑ}_2;\stackrel{}{z}).`$
To calculate $`G_{\lambda =0}(\stackrel{}{๐ฑ}_1,\stackrel{}{๐ฑ}_2;\stackrel{}{z})`$ we observe that the action $`๐`$ is quadratic in $`\lambda `$. A trivial calculation gives
$`G_{\lambda =0}(\stackrel{}{๐ฑ}_1,\stackrel{}{๐ฑ}_2;\stackrel{}{z})={\displaystyle ๐(\mathrm{fields})e^{๐_0}}`$ (13)
$`\times \psi _1^{a_1}(๐ฑ_1)\psi _1^{a_1}(๐ฑ_1^{})\psi _2^{a_2}(๐ฑ_2)\psi _2^{a_2}(๐ฑ^{})`$ (14)
where
$`๐_0`$ $``$ $`๐_{\mathrm{CS}}+{\displaystyle d^3๐ฑ\left[|๐_1\mathrm{\Psi }_1|^2+|\mathbf{}\mathrm{\Psi }_2|^2+\underset{i=1}{\overset{2}{}}m_i^2|\mathrm{\Psi }_i|^2\right]},`$ (15)
From Eq. (15) it is clear that $`G_{\lambda =0}(\stackrel{}{๐ฑ}_1,\stackrel{}{๐ฑ}_2;\stackrel{}{z})`$ is the product of the configurational probabilites of two free polymers. In fact, the fields $`\mathrm{\Psi }_2,\mathrm{\Psi }_2^{}`$ are free, whereas the fields $`\mathrm{\Psi }_1,\mathrm{\Psi }_1^{}`$ are apparently not free because of the couplings with the Chern-Simons fields through the covariant derivative $`๐^1`$. This is, however, an illusion: Integrating out $`A_2^\mu `$ in (14), we find the flatness condition: $`\epsilon ^{\mu \nu \rho }_\nu A_\mu ^i=0.`$ On a flat space with vanishing boundary conditions at infinity this implies $`A_1^\mu =0`$. As a consequence, the functional integral (14) factorizes
$`G_{\lambda =0}(\stackrel{}{๐ฑ}_1,\stackrel{}{๐ฑ}_2;\stackrel{}{z})=G_0(๐ฑ_1๐ฑ_1^{};z_1)G_0(๐ฑ_2๐ฑ_2^{};z_2),`$ (16)
where $`G_0(๐ฑ_i๐ฑ_i^{};z_i)`$ are the free correlation functions of the polymer fields
$`G_0(๐ฑ_i๐ฑ_i^{};L_i)`$ $`=`$ $`{\displaystyle _{ci\mathrm{}}^{c+i\mathrm{}}}{\displaystyle \frac{Mdz_i}{2\pi i}}e^{z_iL_i}G_0(๐ฑ_i๐ฑ_i^{};z_i)`$ (18)
$`={\displaystyle \frac{1}{2}}\left({\displaystyle \frac{M}{4\pi L_i}}\right)^{3/2}e^{M(๐ฑ_i๐ฑ_i^{})^2/2L_i}.`$
Thus we obtain for (12) the integral
$`Z=2\pi {\displaystyle d^3x_1d^3x_2\underset{\genfrac{}{}{0pt}{}{๐ฑ_1^{}๐ฑ_1}{๐ฑ_2^{}๐ฑ_2}}{lim}G_0(๐ฑ_1๐ฑ_1^{};L_1)G_0(๐ฑ_2๐ฑ_2^{};L_2)},`$ (19)
yielding the partition function
$`Z={\displaystyle \frac{2\pi M^3V^2}{(4\pi )^3}}(L_1L_2)^{3/2}.`$ (20)
It is important to realize that in Eq. (10) the limits of coinciding end points $`๐ฑ_i^{}๐ฑ_i`$ and the inverse Laplace transformations do not commute unless a proper renormalization scheme is chosen to eliminate the divergences caused by the insertion of the composite operators $`|\psi _i^{a_i}(๐ฑ)|^2`$ and $`|\mathrm{\Psi }_i(๐ฑ)|^2`$.
5. Let us now turn to the numerator in Eq. (1). Exploiting the identity $`m^2e^{im\lambda }=^2e^{im\lambda }/\lambda ^2`$, and performing two partial integrations in $`\lambda `$, the same technique used above to evaluate the partition function $`Z`$ yields
$`N`$ $`=`$ $`\kappa ^2{\displaystyle d^3x_1d^3x_2\underset{\genfrac{}{}{0pt}{}{n_10}{n_20}}{lim}_{ci\mathrm{}}^{c+i\mathrm{}}\frac{Mdz_1}{2\pi i}\frac{Mdz_2}{2\pi i}e^{z_1L_1+z_2L_2}}`$ (21)
$`\times `$ $`{\displaystyle ๐(\text{fields})\mathrm{exp}(๐_0)|\psi _1^{a_1}(๐ฑ_1)|^2|\psi _2^{a_2}(๐ฑ_2)|^2}`$ (22)
$`\times `$ $`\left[\left({\displaystyle d^3x๐_1\mathrm{\Psi }_1^{}\mathbf{}\mathrm{\Psi }_1}\right)^2+{\displaystyle \frac{1}{2}}{\displaystyle d^3x๐_1^2|\mathrm{\Psi }_1|^2}\right]`$ (23)
$`\times `$ $`\left[\left({\displaystyle d^3x๐_2\mathrm{\Psi }_2^{}\mathbf{}\mathrm{\Psi }_2}\right)^2+{\displaystyle \frac{1}{2}}{\displaystyle d^3x๐_2^2|\mathrm{\Psi }_2|^2}\right].`$ (24)
where $`๐_0`$ has been defined in (15). In the above equation we have taken the limits of coinciding endpoints inside the Laplace integral over $`z_1,z_2`$. This will be justified later on the grounds that the potentially dangerous Feynman diagrams containing the insertions of operations like $`|\mathrm{\Psi }_i|^2`$ vanish in the limit $`n_1,n_20`$. The functional integral in Eq. (24) can be calculated exactly by diagrammatic methods since only four diagrams shown in Fig. (2) contribute.
Only the first diagram in Fig. 2 is divergent from the loop integral formed by two correlation functions of the vector field. This infinity may be absorbed in the four-$`\mathrm{\Psi }`$ interaction accounting for the excluded volume effect which we do not consider at the moment. No divergence arises from the insertion of the composite fields $`|\psi _i^{a_i}(๐ฑ_i)|^2`$.
6. In this section we evaluate the first term appearing in the right hand side of Eq. (24):
$`N_1`$ $`=`$ $`{\displaystyle \frac{\kappa ^2}{4}}\underset{\genfrac{}{}{0pt}{}{n_10}{n_20}}{lim}{\displaystyle _{ci\mathrm{}}^{c+i\mathrm{}}}{\displaystyle \frac{Mdz_1}{2\pi i}}{\displaystyle \frac{Mdz_2}{2\pi i}}e^{z_1L_1+z_2L_2}`$ (27)
$`{\displaystyle d^3x_1d^3x_2d^3x_1^{}d^3x_2^{}}`$
$`|\psi _1^{a_1}(๐ฑ_1)|^2|\psi _2^{a_2}(๐ฑ_2)|^2\left(|\mathrm{\Psi }_1|^2๐_1^2\right)_{๐ฑ_1^{}}\left(|\mathrm{\Psi }_2|^2๐_2^2\right)_{๐ฑ_2^{}}.`$
There is an ultraviolet-divergence which must be regularized. This is done by cutting the spatial integrals off at the persistence length $`\xi `$ over which a polymer is stiff. This contains the stiffness caused by the excluded-volume effects. To be rigorous, we define the integral (27) on a lattice with spacing $`\xi `$.
Replacing the expectation values by the Wick contractions corresponding to the first diagram in Fig. 2, we obtain
$`N_1={\displaystyle \frac{V}{4\pi }}{\displaystyle \frac{M^4}{(4\pi )^6}}(L_1L_2)^{\frac{1}{2}}{\displaystyle _0^1}๐s\left[(1s)s\right]^{\frac{3}{2}}{\displaystyle d^3xe^{\frac{M๐ฑ^2}{2s(1s)}}}`$ (28)
$`\times `$ $`{\displaystyle _0^1}๐t\left[(1t)t\right]^{\frac{3}{2}}{\displaystyle d^3ye^{\frac{M๐ฒ^2}{2t(1t)}}d^3x_1^{\prime \prime }\frac{1}{|๐ฑ_1^{\prime \prime }|^4}}.`$ (29)
The variables $`๐ฑ`$ and $`๐ฒ`$ have been rescaled with respect to the original ones in order to extract the behavior of $`N_1`$ in $`L_1`$ and $`L_2`$. As a consequence, the lattices where $`๐ฑ`$ and $`๐ฒ`$ are defined have now spacings $`\xi /\sqrt{L_1}`$ and $`\xi /\sqrt{L_2}`$ respectively. The $`๐ฑ,๐ฒ`$ integrals may be explicitly computed by analytical methods in the physical limit $`L_1,L_2>>\xi `$, in which the above spacings become small. This has a physical explanation. Indeed, if the polymer lengths are much larger than the persistence length, the effects due to the finite monomer size become negligible and can be ignored.
Finally, it is possible to approximate the integral in $`๐ฑ_1^{\prime \prime }`$ with an integral over a continuous variable $`\rho `$ and a cutoff in the ultraviolet region: $`d^3x_1^{\prime \prime }1/|๐ฑ_1^{\prime \prime }|^44\pi ^2_\xi ^{\mathrm{}}๐\rho /\rho ^2.`$ After these approximations, we obtain
$`N_1=V\sqrt{\pi }{\displaystyle \frac{M}{(4\pi )^3}}(L_1L_2)^{1/2}\xi ^1.`$ (30)
7. For the second diagram in Fig. 2 we have to calculate
$`N_2=\kappa ^2\underset{\genfrac{}{}{0pt}{}{n_10}{n_20}}{lim}{\displaystyle _{ci\mathrm{}}^{c+i\mathrm{}}}{\displaystyle \frac{Mdz_1}{2\pi i}}{\displaystyle \frac{Mdz_2}{2\pi i}}e^{z_1L_1+z_2L_2}`$ (31)
$`\times {\displaystyle }d^3x_1d^3x_2{\displaystyle }d^3x_1^{}d^3x_1^{\prime \prime }d^3x_2^{}`$ (32)
$`\times |\psi _1^{a_1}(๐ฑ_1)|^2|\psi _2^{a_2}(๐ฑ_2)|^2(๐_1\mathrm{\Psi }_1^{}\mathbf{}\mathrm{\Psi }_1)_{๐ฑ_1^{}}`$ (33)
$`\times (๐_1\mathrm{\Psi }_1^{}\mathbf{}\mathrm{\Psi }_1)_{๐ฑ_1^{\prime \prime }}(๐_2^2|\mathrm{\Psi }_2|^2)_{๐ฑ_2^{}}.`$ (34)
The above amplitude has no ultraviolet divergence, so that no regularization is required. The Wick contractions pictured in the second Feynman diagrams of Fig. 2 lead to the integral
$`N_2=4\sqrt{2}VL_2^{1/2}L_1^1{\displaystyle \frac{M^3}{\pi ^6}}{\displaystyle _0^1}๐t{\displaystyle _0^t}๐t^{}C(t,t^{}),`$ (35)
where $`C(t,t^{})`$ is a function independent of $`L_1`$ and $`L_2`$:
$`C(t,t^{})=\left[(1t)t^{}(tt^{})\right]^{3/2}{\displaystyle d^3xd^3yd^3ze^{\frac{M(๐ฒ๐ฑ)^2}{2(1t)}}}`$ (36)
$`\times \left(_๐ฒ^\nu e^{M๐ฒ^2/2t^{}}\right)\left(_๐ฑ^\mu e^{M๐ฑ^2/2(tt^{})}\right)P_{\mu \nu }(๐ฑ,๐ฒ,๐ฑ),`$ (37)
with $`P_{\mu \nu }(๐ฑ,,๐ฒ,๐ฑ)[\delta _{\mu \nu }๐ณ(๐ณ+๐ฑ)(z+x)_\mu z_\nu ]/(|๐ณ|^3|๐ณ+๐ฑ|^3)`$. As in the previous section, the variables $`๐ฑ,๐ฒ,๐ณ`$ have been rescaled with respect to the original ones in order to extract the behavior in $`L_1`$. Again, if $`L_1,L_2>>\xi `$ the analytical evaluation of $`C(t,t^{})`$ becomes possible, leading to
$`N_2`$ $`=`$ $`{\displaystyle \frac{VL_2^{1/2}L_1^1}{(2\pi )^6}}M^{3/2}4K,`$ (38)
where $`K`$ is the constant $`\frac{1}{6}B(\frac{3}{2},\frac{1}{2})+\frac{1}{2}B(\frac{5}{2},\frac{1}{2})B(\frac{7}{2},\frac{1}{2})+\frac{1}{3}B(\frac{9}{2},\frac{1}{2})=19\pi /3840.154,`$ and $`B(a,b)`$ is the Beta function. For large $`L_1\mathrm{}`$, this diagram gives a negligible contribution with respect to $`N_1`$.
The third diagram in Fig. 2 give the same as the second, except that $`L_1`$ and $`L_2`$ are interchanged: $`N_3=N_2|_{L_1L_2}.`$
8. The fourth Feynman diagram in Fig. 2 has no ultraviolet divergence. As before, it can be exactly evaluated apart from the lattice integrations. However, the behavior of the related Feynman integral $`N_4`$ can be easily estimated in the following limits:
1. $`L_11;L_1L_2`$, where $`N_4L_1^1,`$
2. $`L_21;L_2L_1`$ where $`N_4L_2^1,`$
3. $`L_1,L_21,L_2/L_1=\text{finite},`$ where $`N_4L_1^{3/2}`$.
Moreover, if the lengths of the polymers are considerably larger than the persistence length, $`N_4`$ can be computed in a closed form:
$`N_4`$ $``$ $`{\displaystyle \frac{128V}{\pi ^5}}{\displaystyle \frac{M}{\pi ^{3/2}}}(L_1L_2)^{1/2}`$ (39)
$`\times `$ $`{\displaystyle _0^1}๐s{\displaystyle _0^1}๐t(1s)(1t)(st)^{1/2}`$ (40)
$`\times `$ $`\left[L_1t(1s)+L_2(1t)s\right]^{1/2}.`$ (41)
It is simple to check that this expression has exactly the above behaviors.
9. Collecting all contributions we obtain the result for the second topological moment $`m^2=(N_1+N_2+N_3+N_4)/Z,`$ with $`N_1,N_2,N_3,N_4,Z`$ given by Eqs. (20), (30), (38), and (41). In all formulas, we have assumed that the volume $`V`$ of the system is much larger than the size of the volume occupied by a single polymer, i.e., $`VL_1^3`$.
To discuss the physical content of the above expression for $`m^2`$, we consider a number $`N_p`$ of polymers $`p_1\mathrm{}p_{N_p}`$ with lengths $`l_1\mathrm{}l_{N_p}`$ in an uniform solution. We introduce the polymer concentration $`\rho =/V`$ as the average mass density of the polymers per unit volume, where $``$ is the total mass of the polymers $`=_{k=1}^{N_p}m_al_k/a`$ and $`m_a`$ is the mass of a single monomer of length $`a`$. Thus $`l_k/a`$ is the number of monomers in the polymer $`p_k`$. The polymer $`P_1`$ is singled out as anyone of the polymers $`p_k`$, say $`p_{\overline{k}}`$, of length $`L_1=l_{\overline{k}}`$. The remaining ones are replaced by a long effective polymer $`P_2`$ of length $`L_2=_{k\overline{k}}l_k`$. From the above relations we may also write
$`L_2{\displaystyle \frac{aV\rho }{m_a}}.`$ (42)
In this way, the length of the effective molecule $`P_2`$ is expressed in terms of physical parameters. Keeping only the leading terms for $`V1`$, we find for the average square number of intersections $`m^2_{sol}`$ formed by $`P_1`$ with the other polymers the approximate result
$`m^2_{sol}{\displaystyle \frac{N_1+N_2}{Z}},`$ (43)
which, in turn, has the approximate form
$`m^2_{sol}={\displaystyle \frac{a\rho }{m_a}}\left[{\displaystyle \frac{\xi ^1L_1}{2\pi ^{1/2}M^2}}{\displaystyle \frac{2KL_1^{1/2}}{\pi ^4M^{3/2}}}\right],`$ (44)
with $`K`$ as defined after Eq. (38). This is the announced final result. Since the persistence length is of the same order of the monomer length $`a`$ and $`Ma^1`$, $`m^2`$ is positive for large $`L_1`$ as it should.
10. In conclusion, we have found an exact field theoretic formula for the second topological moment of two polymers. Only the final integrations over the spatial variables in the Feynman diagrams of Fig. 2 were done approximately. These were defined on a lattice related to the finite monomer size. Our Chern-Simons-based theory is free of the shortcomings of previous mean-field procedures. Our formula for $`m^2`$ has been applied to the realistic case of long flexible polymers in a solution. When the polymer lengths become large, the Feynman integrals in $`N_1,\mathrm{},N_4`$ can be evaluated analytically. In this way we have been able to derive the result (44) for the average square number of intersections formed by a polymer $`P_1`$ with all the others. This calculation is exact in the long-polymer limit. The corrections to (44) are suppressed by further inverse square roots of the polymer lengths.
To leading order in $`L_1`$, our result (44) agrees with that of , but our exact subleading correction go beyond the approximation of . Note that there is no direct comparison of our result with that of , since there the polymer $`P_2`$ was considered as a fixed obstacle causing a dependence on the choice of the configuration of $`P_2`$.
Finally, let us emphasize the absence of infrared divergences in the topological field theory (6) in the limit of vanishing masses $`m_1,m_2=0`$. As a consequence, the second topological moment does not diverge in the limit of large $`L_1`$ if $`m^2`$ is calculated from (6) for polymers passing through two fixed points $`๐ฑ_1,๐ฑ_2`$. This indicates a much stronger reduction of the configurational fluctuations by topological constraints than one might have anticipated. |
warning/0002/math0002217.html | ar5iv | text | # Bassโs Work in Ring Theory and Projective Modules
(Date: March 23, 1999.)
## Abstract.
The early papers of Hyman Bass in the late 50s and the early 60s leading up to his pioneering work in algebraic $`K`$-theory have played an important and very special role in ring theory and the theory of projective (and injective) modules. In this article, we give a general survey of Bassโs fundamental contributions in this early period of his work, and explain how much this work has influenced and shaped the thinking of subsequent researchers in the area.
The work on this paper was supported in part by a grant from NSA
Contents
ยง0. Introduction
Part I: Projective (and Torsionfree) Modules
ยง1. Big Projectives
ยง2. Stable Structure of Projective Modules
ยง3. Work Related to Serreโs Conjecture
ยง4. Rings with Binary Generated Ideals: Bass Rings
Part II: Ring Theory
ยง5. Semiperfect Rings as Generalizations of Semiprimary Rings
ยง6. Perfect Rings and Restricted DCC
ยง7. Perfect Rings and Representation Theory
ยง8. Stable Range of Rings
ยง9. Rings of Stable Range One
References
ยง0. Introduction
It gives me great pleasure to have this opportunity to write about Professor Hyman Bassโs work in ring theory and projective modules. This was work done by a young Hyman in the early 60s when he was a junior faculty member at Columbia. The now-classical paper on the homological generalization of semiprimary rings \[B<sub>1</sub>\], an outgrowth of his 1959 Chicago thesis written under the direction of Irving Kaplansky, was followed in quick succession by a brilliant series of papers \[B<sub>2</sub>โB<sub>7</sub>\] on various aspects of the structure of projective modules (decomposition, extendibility, and freeness) and injective modules (injective dimensions, Gorenstein rings, and Bass rings). The 1962 announcement of a homotopy theory of projective modules with Schanuel \[BS\] and the subsequent 1964 announcement \[B<sub>8</sub>\] on the stable structure of the general linear group over an arbitrary ring culminated in his famous IHES paper \[B<sub>9</sub>\] which, perhaps more than any other work in that era, marked the monumental creation of the new mathematical subject of algebraic $`K`$-theory. The usual way to put it would be to say, I guess, that all of this work โshowed the future masterโ; but in fact, by this time, Hyman Bass had already proved himself to be a true master of the great art of algebra โ and a mathematician of extraordinary creativity and insight.
Since Chuck Weibel \[We\] will be touching upon Bassโs work in algebraic $`K`$-theory and Craig Huneke \[Hu\] will be reporting on his work on injective modules and Gorenstein rings, I shall focus the discussion in this article mainly on Bassโs contributions to ring theory and the structure of projective modules. Starting with the latter, we first survey, in ยง1 and ยง2 below, Bassโs work on big projectives and on his generalizations and extensions of Serreโs results on the stable structure of projective modules. In ยง3, we move on to Bassโs work surrounding the theme of Serreโs Conjecture on f.g.<sup>1</sup><sup>1</sup>1Throughout this paper, โf.g.โ will be used as an abbreviation for โfinitely generatedโ. projective modules over polynomial rings. The last section of Part I concludes with a survey on Bassโs work on the decomposition of torsionfree modules and on commutative rings with binary generated ideals, which resulted in the important notion of Bass rings and Bass orders. In this section (ยง4), Gorenstein rings play a substantial role, but our exposition is designed to overlap only peripherally with that of Hunekeโs article.
In Part II of this article, we return to Bassโs maiden work on the homological generalizations of semiprimary rings, and explain the historic context of this work as well as the role it played in the subsequent development of noncommutative ring theory. Guided by the ideas of stability in the homotopy theory of vector bundles, Bass also single-handedly invented the notion of stable range of a ring $`R`$, which he successfully applied to the study of the stable structure of the infinite general linear group over $`R`$. (This work may be thought of as the โ$`K_1`$-analogueโ of the stable theory of projective modules reported in ยง2.) Our survey concludes with a discussion (ยงยง8-9) of Bassโs work on the stable range, with a special emphasis on the case of stable range one which in turn has deep significance on the arithmetic of rings, and on questions concerning the cancellation and substitution of modules with respect to direct sum decompositions.
Throughout this exposition, our aim is not only to survey Bassโs work in ring theory and projective modules, but also to point out how much this work has influenced and shaped the thinking of subsequent researchers in the area. It will be seen that, in a number of new lines of investigation in algebra in the last few decades, it was the decisive pioneering steps of Bass that broke open the new path. Although Bassโs work in modules and rings spanned only the decade of the 60s, its impact has been enormous indeed, and will certainly continue to be felt as ring theory moves into the next century.
Writing this article is for me a very pleasant task, though undertaking such a task inevitably involves somewhat of a nostalgic trip down the memory lane. Since it is perhaps not out of place in these Proceedings to talk about oneโs connections to Hyman, I will indulge myself in a few personal reminiscences below.
As a beginning graduate student at Columbia in the mid-60s, I was more than a little awed by the power and fame of the senior faculty: Professors Lipman Bers, Samuel Eilenberg, Ellis Kolchin, Masatake Kuranishi, Serge Lang, Edgar Lorch, and Paul A. Smith, among others, from whom I took my first and second year graduate courses.<sup>2</sup><sup>2</sup>2The students were nothing short of stellar either! My graduate peers included Mike Engber, Audun Holme, Fred Gardiner, and Irwin Kra; Winfried Scharlau was a visiting student from Germany, and Alexander Mikhalev was an exchange scholar from U.S.S.R.; Bill Haboush, Tony Bak, Mike Stein, Allen Altman, Bruce Bennett, Spencer Bloch, Bob and Jane Gilman were a couple of years behind; Julius Shaneson, Sylvain Cappell and Ethan Akin were โhot-shotโ Columbia undergraduates taking graduate courses with us. Each person in this list is now a Professor of Mathematics. Many of them have been chairs, provosts, and deans. After a few skirmishes with functional analysis (taught to me by Professor Lorch and Visiting Professor B. Sz. Nagy), I fell under the spell of Sammy Eilenberg, and positioned myself to become a student of his, hoping to study with him category theory and homological algebra. Under Sammyโs guidance, I wrote my first paper \[La<sub>1</sub>\], which he kindly communicated to the Proceedings of the National Academy of Sciences (of which he was a member). But then for the year 1966-67, Professor Eilenberg had his Sabbatical coming, and he was to go off to Paris to spend the entire academic year. To go with him to Paris would have been the โcoolโ thing to do for a mid-career graduate student except, alas, I found myself both financially and linguistically poorly equipped to make the trip. So Professor Eilenberg said to me: โWhy donโt you stay here, and study with Bass while I am gone?โ This was how I became a student of Bass! Hyman was then a young Assistant Professor, who, just a few years ago, was brought from Chicago to Columbia by Sammy himself.
Hyman became an Associate Professor around that time, and in another year he was moved up to Full Professor. Even I could see he must be working on something hot! So I xeroxed all of his papers,<sup>3</sup><sup>3</sup>3Yes, there were already xerox machines in the mid-60s, though stencils and mimeographed copies were still not entirely out of fashion. and started poring over them assiduously. Hymanโs $`K`$-theory paper \[B<sub>9</sub>\] had just come out in the โBlue Journalโ; this and his earlier papers in ring theory and projective modules eventually became a staple of my graduate education. The most challenging open problem in homological algebra in those days was Serreโs Conjecture (ibid.); many a graduate student in algebra from that era had no doubt tried his/her hand at it โ and I was no exception. Hyman, who had generalized Seshadriโs solution of the Conjecture in two variables, was the natural leader for this small circle of aspiring graduate students. I still recall that, one time, one of us thought he had a brilliant idea to solve Serreโs Conjecture. A few of us excitedly met with Hyman in an impromptu seminar to go over the โideaโ; but of course Hyman quickly found the hole. Although we never got anywhere with our fledgling efforts, my fascination with Serreโs Conjecture continued, and culminated in the writing of my 1978 Springer Lecture Notes \[La<sub>2</sub>\], two years after the Conjecture was fully solved independently by Suslin and Quillen. Hymanโs great influence was evident from cover to cover of my modest book.
I spent my last graduate year at Columbia in 1966-67. In that year, Hyman taught what was most probably the first graduate course ever given in algebraic $`K`$-theory in the US. Pavaman Murthy had come from the Tata Institute to do postdoctoral work with Hyman, and was in the audience. (In his modest Morningside Heights apartment, Pavaman served what he claimed to be the best coffee in Manhattan. It was free, so I had no reason to disagree; we became good friends.) Some of the lecture notes taken by Murthy, Charles Small and me eventually evolved into Hymanโs famed tome \[Ba<sub>10</sub>\]. At that time, Hyman was busy at work with Milnor and Serre on the Congruence Subgroup Problem for the special linear groups and the symplectic groups. I got lucky and proved several little things about Mennicke symbols and $`SK_1`$ of abelian group rings, which earned me a few attributions in \[BMS\]. Needless to say, I was proud to be mentioned in a paper by such distinguished authors. I got lucky in some other fronts too, and in May, 1967 completed a thesis in algebraic $`K`$-theory under Hyman dealing with Artinโs Induction Theorem and induction techniques for Grothendieck groups and Whitehead groups of finite groups. In those days, algebraic $`K`$-theory meant only $`K_0`$ and $`K_1`$; even Milnorโs $`K_2`$ had not been defined yet.
As it happened, I was Hymanโs first Ph.D. student. I have always viewed this as a special honor, and I am sure that this fact has helped me a great deal professionally. Now the list of Ph.D. students of Hyman is 25(?) strong. The longevity of Hyman as a thesis advisor is rather strikingly illustrated by the fact that at least several of my mathematical brothers and sisters in this list were not even born yet when I completed my Ph.D. degree at Columbia.
In supervising my work, Hyman never tried to tell me what I should work on. Rather, he let me find my own path, and instilled in me the needed confidence to grow into a research mathematician. What I learned from Hyman was not just mathematics, but how to do mathematics, and, what is perhaps even more important, how to conduct myself as a mathematician. All of this he taught me in the best way โ by his own example. For this, and for the many other favors he has rendered over the years, I shall always be grateful.
Part I: Projective (and Torsionfree) Modules
ยง1. Big Projectives
While the projective modules occurring in number theory, representation theory and algebraic geometry are mostly f.g. ones, non-f.g. projective modules do arise naturally over various kinds of rings, for instance, rings of continuous functions. Thus, the quest for information about non-f.g. projective modules is not a frivolous one.
The first significant result in the study of general (that is, not necessarily f.g.) projective modules was found by Kaplansky in 1958. In his seminal paper \[K<sub>2</sub>\], Kaplansky proved that, for any ring $`R`$, any projective $`R`$-module is always a direct sum of countably generated (projective) modules. This result has two interesting consequences. First, any indecomposable projective module over any ring is countably generated; and second, if any countably generated projective module over some ring $`R`$ is free, then any projective $`R`$-module is also free. The case of a local ring $`R`$ provides a particularly striking illustration for the power of the second statement: in this case, Kaplansky used an ingenious argument (reminiscent of the proof of Nakayamaโs Lemma) to show that any countably generated projective $`R`$-module is free, from which it then follows that any projective $`R`$-module is free.
Inspired by Kaplanskyโs result, Bass took up the study of โbigโ projective modules in \[B<sub>5</sub>\]. The overall theme of \[B<sub>5</sub>\] is that, under certain fairly mild conditions on a ring $`R`$, โbigโ projective $`R`$-modules are necessarily free. To formulate this precisely, Bass introduced the notion of a uniformly $`\mathrm{}`$-big module: for an infinite cardinal $`\mathrm{}`$, an $`R`$-module $`P`$ is uniformly $`\mathrm{}`$-big if $`P`$ can be generated by $`\mathrm{}`$ elements, and for any ideal $`IR`$, the $`R/I`$-module $`P/IP`$ cannot be generated by fewer than $`\mathrm{}`$ elements. For instance, the free module of rank $`\mathrm{}`$ over a nonzero ring $`R`$ is uniformly $`\mathrm{}`$-big. Bassโs first result on big projectives is the following converse of this statement:
(1.1) Theorem. If $`R/\text{rad}(R)`$ is left noetherian, then, for any infinite cardinal $`\mathrm{}`$, any uniformly $`\mathrm{}`$-big projective left $`R`$-module is free.
Bassโs proof of this result is modeled upon Kaplanskyโs proof of his main theorem in \[K<sub>2</sub>\]: one makes a reduction to the crucial case when $`\mathrm{}`$ is the countable infinite cardinal $`\mathrm{}_0`$, and in this case one achieves the desired goal by a clever juggling with infinite matrices.
For concrete applications of (1.1), one needs to find situations where we can say, for instance, that all non-f.g. projective modules are uniformly big for some infinite cardinal. Bass showed that this is the case when, say, $`R`$ is a commutative noetherian ring with only trivial idempotents. Thus, one has the following
(1.2) Corollary. If $`R`$ is a commutative noetherian ring with only trivial idempotents, then any non-f.g. projective $`R`$-module is free.
For instance, if $`R=k[x_1,\mathrm{},x_n]`$ where $`k`$ is a field, Serre asked in 1955 (see ยง2) if every f.g. projective $`R`$-module $`P`$ is free. The result above would give the freeness of $`P`$ if $`P`$ was not f.g.
The general results of Bass (such as (1.1) and (1.2)) describing the behavior of big projectives have remained essentially unsurpassed to this date. In \[B<sub>5</sub>\], Bass remarked that these results seemed to indicate that big projective modules โinvite little interest.โ We can say, today, that this is perhaps not quite true. In a recent paper \[LR\], Levy and Robson have determined the structure of all infinitely generated projective modules over (noncommutative) hereditary noetherian prime rings. Their results showed that, over such rings, there may exist non-f.g. projective modules which are not free and which have rather interesting structures.
Two other results of Bass on projective modules over any ring have become folklore in the subject, so it is fitting to end this section by recalling them. The first of these says that:
(1.3) Theorem. Any nonzero projective module $`P`$ over any ring $`R`$ has a maximal submodule. More precisely, rad$`(R)P=\text{rad}(P)P`$, where rad$`(P)`$ denotes the intersection of all maximal submodules of $`P`$.
This result first appeared in \[B<sub>1</sub>: p. 474\], and has been used by many authors since. The second result, also involving the Jacobson radical of a ring $`R`$, states the following:
(1.4) Theorem. Let $`I`$ be any ideal of a ring $`R`$ such that $`I\text{rad}(R)`$, and $`P,Q`$ be f.g. projective left $`R`$-modules. Then
$`(1.5)`$
$$PQasR\text{-}modulesP/IPQ/IQasR/I\text{-}modules.$$
In particular, $`P`$ is free as an $`R`$-module iff $`P/IP`$ is free as an $`R/I`$-module.
This result, which made its first appearance in the literature as Lemma 2.4 in \[B<sub>2</sub>\], is now in virtually every textbook which treats the subject of projective modules, and is well-known to any student knowledgeable about the subject of homological algebra. We should point out that, with Bassโs notion of left $`T`$-nilpotency on ($`1`$-sided) ideals (introduced later in ยง6), (1.4) can also be given the following โtransfiniteโ formulation, for any pair of projective modules $`P,Q`$:
(1.6) Theorem. If $`I`$ is any left $`T`$-nilpotent left ideal in a ring $`R`$, then $`(1.5)`$ holds for any (not necessarily f.g.) projective left $`R`$-modules $`P,Q`$.
For a proof of this, see \[La<sub>3</sub>: (23.17)\].
ยง2. Stable Structure of Projective Modules
In his epoch-making papers \[S<sub>1</sub>, S<sub>2</sub>\], Serre established an analogy between projective modules (in algebra) and vector bundles (in topology). This important analogy, which was further promulgated by the work of Swan \[Sw<sub>1</sub>\], enables one to establish the rudiments of a dictionary to translate the language of projective modules into that of vector bundles. Now by the late 1950s, the topology of vector bundles was already rather well-developed; the Serre-Swan analogy mentioned above made it possible, therefore, for algebraists to โpredictโ (if not prove) theorems in the theory of projective modules based on their knowledge of topological results in the theory of vector bundles.
One of the best known elementary facts about vector bundles is the following. If $`[X,Y]`$ denotes the set of homotopy classes of maps from $`X`$ to $`Y`$, and BO$`(r)`$ denotes the classifying space of the orthogonal group O$`(r)`$, then
(2.1) Theorem. For any connected finite CW complex $`X`$ of dimension $`d`$, the natural map
$$i_r:[X,\text{BO}(r)][X,\text{BO}(r+1)]$$
is surjective for $`rd`$, and injective (and hence bijective) for $`rd+1`$.
Now by the classification theorem of vector bundles, $`[X,\text{BO}(r)]`$ represents the set of equivalence classes of real $`r`$-plane bundles over $`X`$, and, with this interpretation, the map $`i_r`$ in (2.1) above is given by adding a trivial line bundle. To make the transfer into algebra, we replace $`X`$ by a commutative noetherian ring $`R`$ with only trivial idempotents (so that the Zariski prime spectrum Spec$`(R)`$ is a connected space), and replace $`[X,\text{BO}(r)]`$ by the set $`๐ซ_r(R)`$ of isomorphism classes of f.g. projective $`R`$-modules of rank $`r`$. The $`i_r`$ in (2.1) is then to be replaced by the map given by โadding a copy of $`R`$โ. With such a transfer in place, Serre broke new ground by coming up with commutative algebra techniques to prove the following remarkable analogue of the surjectivity part of (2.1).
(2.2) Theorem. (\[S<sub>2</sub>: Thรฉorรจme 1\]) In the above setting, assume that the maximal ideal spectrum max$`(R)`$ (with the Zariski topology) is a space of dimension<sup>4</sup><sup>4</sup>4The dimension of a topological space $`X`$ is defined to be the supremum of the codimensions of its nonempty closed sets. Here, the codimension of an irreducible closed set $`F`$ is the supremum of the lengths of finite chains of irreducible closed sets โaboveโ $`F`$, and the codimension of an arbitrary closed set $`C`$ is the infimum of the codimensions of all irreducible closed sets $`FC`$. $`d`$. Then the map
$`(2.3)`$
$$i_r:๐ซ_r(R)๐ซ_{r+1}(R)$$
is surjective if $`rd`$. In other words, any f.g. projective $`R`$-module of rank $`p>d`$ is of the form $`R^{pd}P_0`$ for some (projective) module $`P_0`$.
In this setting (and under the same hypotheses), the question whether the map $`i_r`$ in (2.3) is injective for $`rd+1`$ (in analogy to the second part of (2.1)) begs to be asked. This amounts to the following cancellation question for f.g. projective $`R`$-modules $`P,Q`$ and $`M`$: if $`MPMQ`$ and $`P`$ has rank $`>d`$, does it follow that $`PQ`$? It would be difficult to imagine that Serre was not aware of this very natural question in 1957-58, but anyhow, Serre did not pursue it in \[Se<sub>1</sub>\].
In his two papers \[B<sub>2</sub>\] and \[B<sub>9</sub>\] (see also the research announcement \[BS\] written jointly with S. Schanuel), Bass not only answered this question in the affirmative, but also relaxed some of the assumptions in (2.2), and extended all results to a noncommutative setting. Working now with a commutative ring $`A`$ and a module-finite $`A`$-algebra $`R`$, Bass considered $`R`$-modules $`P`$ which need not be projective or f.g. He defined such a module to be of $`f`$-rank $`r`$ if, at every maximal ideal $`๐ช\text{max}(A)`$, the $`R_๐ช`$-module $`P_๐ช`$ contains an $`R_๐ช`$-free direct summand of rank $`r`$.
(2.4) Theorem. In the above (noncommutative) setting, assume that max$`(A)`$ is a noetherian space of dimension $`d`$.
(1) (Splitting) If an $`R`$-module $`P`$ is a direct summand of a direct sum of finitely presented $`R`$-modules and has f-rank $`>d`$, then $`P`$ has a direct summand isomorphic to $`R`$.
(2) (Cancellation) Let $`P,Q`$ be $`R`$-modules such that $`P`$ has a projective direct summand of f-rank $`>d`$. Then, for any f.g. projective $`R`$-module $`M`$,
$`(2.5)`$
$$MPMQPQ.$$
If we now define $`๐ซ_r(R)`$ to be the set of isomorphism classes of f.g. projective $`R`$-modules of $`f`$-rank $`r`$, clearly (1) and (2) above imply that:
(2.6) Corollary. The map $`i_r:๐ซ_r(R)๐ซ_{r+1}(R)`$ (defined by adding $`R`$) is surjective for $`rd`$, and injective (and hence bijective) for $`rd+1`$.
Just as (2.1) is quantitatively the best result for bundles in topology, both (2.2) and (2.5), (2.6) are quantitatively the best for f.g. projective modules in algebra. Weโll mention the usual examples to substantiate this statement in the commutative and noncommutative cases below.
In the commutative case, let $`R`$ be the real coordinate ring of the sphere $`S^2`$, so $`R=[x,y,z]`$, with the relation $`x^2+y^2+z^2=1`$. This is a noetherian domain of (Krull) dimension 2. Let $`P=\text{ker}(\phi )`$ where $`\phi `$ is the $`R`$-epimorphism $`R^3R`$ given by $`\phi (e_1)=x,\phi (e_2)=y`$, and $`\phi (e_3)=z`$. We have $`R^3RP`$, so $`P`$ is a f.g. projective $`R`$-module of rank 2. Here, $`P`$ corresponds to the tangent bundle of $`S^2`$. Since the tangent bundle is known to be indecomposable, $`P`$ is also indecomposable (and in particular $`PR^2`$), so the class of $`P`$ in $`๐ซ_2(R)`$ is not in the image of the map $`i_1`$ in (2.3). This shows that the condition $`rd`$ for the surjective part in (2.6) is the best possible. On the other hand, the fact that $`RPRR^2`$ and $`PR^2`$ shows that the condition $`rd+1`$ for the injective part in (2.6) is also the best possible.
In the noncommutative setting, a very simple case for application is that of a group ring $`R=\pi `$, where $`\pi `$ is a finite group. Here we take $`A`$ to be $``$ in (2.4), and set $`d=1`$. It follows from (2.4)(1) that any f.g. projective $`R`$-module is a direct sum of rank 1 projective modules. But, according to a famous example of Swan \[Sw<sub>2</sub>\], if $`\pi `$ is the generalized quaternion group of order 32, there exists a nonfree rank 1 f.g. projective $`R`$-module $`P`$ such that $`RPRR`$. Since $`PR`$, we see again that, in the conclusion in (2.4)(2), the condition rank$`(P)>d`$ cannot be further weakened. In Swanโs construction, the module $`P`$ was, in fact, chosen such that, for a suitable maximal order $`S`$ containing $`R=\pi `$, $`S_RPS`$. Therefore, by tensoring up to $`S`$, we obtain a similar example (of non-cancellation) over a left and right hereditary module-finite algebra over the one-dimensional ring $``$.
Of course, the bounds $`rd`$ and $`rd+1`$ are just general bounds for surjectivity and injectivity to hold, respectively, in (2.6). For specific rings, there can be stronger results. For instance, in contrast with Swanโs example mentioned above, if $`\pi _0`$ is the (ordinary) quaternion group of order 8, then full cancellation holds for f.g. projective modules over the group ring $`\pi _0`$, according to a result of J. Martinet \[Ma\].
Bassโs fundamental results (2.4) and (2.6) in the early 60s provided the general framework for much of the subsequent investigations on the Splitting and Cancellation Problems for projective (and more general) modules over various classes of rings. Without giving any details, let us just mention, in this direction, the work of Chase, Mohan Kumar, Murthy, Nori, Sridharan, Suslin, Swan, Towber, Wiegand, and others on the splitting and cancellation of f.g. projective modules over affine algebras. A survey of some aspects of this appears in Murthyโs article \[Mu\] in this volume. On the noncommutative side, there are the important cancellation results of Jacobinski, Guralnick, Levy, Roฤญter, Swan and many others for lattices over group rings and orders in separable algebras.
Now it may be said that a cancellation result such as (2.4)(2) is not quite truly in the noncommutative spirit. In this theorem, the ring $`R`$ in question is a module-finite algebra over a commutative ring $`A`$; orders in finite-dimensional separable algebras are also of the same nature. Such rings are simply not sufficiently representative of a general noncommutative ring. Later developments show that there are indeed some โtruly noncommutativeโ cancellation theorems, where the cancellation of a module $`M`$ in (2.5) depends on the โstable rangeโ of the endomorphism ring of $`M`$ and on the structure of the module $`P`$. Such results, which in essence generalize (2.4)(2), were first found by Warfield \[Wa<sub>2</sub>\] in 1980 (and in part by Evans in 1973). We shall return to formulate these more general cancellation results after we introduce the notion of stable range of (noncommutative) rings in ยง8.
ยง3. Work Related to Serreโs Conjecture
In the Serre-Swan analogy between vector bundles and projective modules, the counterpart of the affine $`n`$-space $`k^n`$ over a field $`k`$ is the polynomial ring in $`n`$ variables $`R=k[x_1,\mathrm{},x_n]`$, so vector bundles over $`k^n`$ correspond to f.g. projective modules over $`R`$. From the viewpoint of topology, the (real) affine $`n`$-space is contractible, so the vector bundles over it are all trivial. This led Serre to ask the question, in his famous FAC paper \[S<sub>1</sub>\], whether every f.g. projective module over the polynomial ring $`R=k[x_1,\mathrm{},x_n]`$ is free (for any field $`k`$). An affirmative answer to this question seemed so plausible and convincing to the mathematical public that, almost from the very beginning, it became known under the misnomer of โSerreโs Conjectureโ. This โconjectureโ is clearly true when $`n=1`$, since in this case $`R=k[x_1]`$ is a PID, and f.g. projective modules over a PID are well known to be free. In the general case, it is known (essentially from Hilbertโs Syzygy Theorem) that any f.g. projective module $`P`$ over $`R=k[x_1,\mathrm{},x_n]`$ is stably free, so Serreโs Conjecture boils down to a cancellation statement: that $`PR`$ free should imply $`P`$ free.<sup>5</sup><sup>5</sup>5In this form, the Conjecture is capable of a completely elementary statement meaningful to any student with a high school background in algebra: whenever $`f_1g_1+\mathrm{}+f_rg_r=1R`$, there is an $`r\times r`$ matrix of determinant $`1`$ over $`R`$ with first row $`(f_1,\mathrm{},f_r)`$. See \[La<sub>2</sub>: (I.4.5)\] for details. By Bassโs Theorem (2.6)(2), this would follow, for instance, if the rank of $`P`$ is at least $`n+1`$.
In the 1960s, Serreโs Conjecture became one of the premier open problems in algebra. The fact that the Conjecture was prompted by a natural analogy with vector bundle theory gave it a certain sense of inevitability; on the other hand, the fact that the conjecture can be stated so directly and in such completely elementary terms made it enticing to all. Many algebraists in the 1960s, junior and senior alike, must have tried their hands at solving this famous conjecture. Bassโs interest in the structure of projective modules, evident from his first two papers \[B<sub>1</sub>, B<sub>2</sub>\], naturally steered him in this direction.
In 1958, Seshadri \[Se<sub>1</sub>\] confirmed Serreโs Conjecture in the case of two variables, proving, as he put it, the โtriviality of vector bundles over the affine space $`K^2`$โ. As a matter of fact, Seshadri showed more generally that any f.g. projective module $`P`$ over $`A[t]`$ โcomes fromโ $`A`$ if $`A`$ is a PID \[Se<sub>1</sub>\], or the coordinate ring of a nonsingular affine curve over an algebraically closed field \[Se<sub>2</sub>\]. These results led Bass to consider the same problem over $`R=A[t]`$ when $`A`$ is a Dedekind domain. The desired goal in this case would be the same as Seshadriโs โ to prove that any f.g. projective $`P`$ over $`R`$ โcomes fromโ $`A`$, that is, $`PR_AP_0`$ for some (necessarily f.g. projective) module $`P_0`$ over $`A`$.
In the early 60s, Bass succeeded in extending Seshadriโs argument, and proved the following result in \[B<sub>4</sub>: (2.4)\].
(3.1) Theorem. Let $`A`$ be a Dedekind ring, $`R=A[t]`$, and $`P`$ be a f.g. $`R`$-module such that, for any prime ideal $`๐ญA`$, $`P/๐ญP`$ is a torsionfree $`(A/๐ญ)[t]`$-module. Then $`P`$ is extended from a f.g. projective $`A`$-module; in particular, $`P`$ is projective.
This implies, in particular, that any f.g. projective $`R`$-module is extended from $`A`$. As it turned out, the same result was obtained independently by Serre at about the same time; see \[S<sub>3</sub>\]. In retrospect, this result is perhaps most appropriately called the Seshadri-Bass-Serre Theorem.
In 1964, Bass took up what may be considered a noncommutative version of Serreโs Problem. Perhaps not too surprisingly, the noncommutative case turned out to be more tractable. Generalizing the work of P. M. Cohn, Bass succeeded in proving that, if $`\pi `$ is a free group (resp. a free monoid), then for any principal ideal domain $`A`$, any f.g. projective module over the group ring (resp. monoid ring) $`A\pi `$ is free. This result appeared in Bassโs paper \[B<sub>7</sub>\], in the first volume of the Journal of Algebra. In the case when $`\pi `$ is a free monoid on one generator, it retrieves, of course, Seshadriโs theorem in \[Se<sub>1</sub>\].
Although Bass did not publish further results on Serreโs Conjecture after the 1960s, his keen interest in it continued well into the 70s. When Suslin and Vaserstein began to make significant progress on the conjecture in the early 70s, the former Soviet Union was still quite isolated mathematically from Europe and from the U.S. In order to make their latest findings known to the West, Suslin and Vaserstein could only communicate them by letter to Bass. I still remember vividly the AMS Annual Meeting in San Francisco in 1974, in which Bass gave an โimpromptuโ lecture on the most recent Suslin-Vaserstein results on Serreโs Conjecture โ to a room-full of people eager to find out how close Serreโs Conjecture had come to being solved. One of these โFrom Russia, with Loveโ results proclaimed the freeness of projective $`k[x_1,\mathrm{},x_n]`$-modules โof rank $`1+n/2`$โ. This surely looked wonderful, but a high-school algebra question unwittingly came up: when Suslin wrote โ$`1+n/2`$โ in his letter to Bass, did he mean $`\mathrm{\hspace{0.17em}1}+\frac{n}{2}`$, or could he have meant the more ambitious $`(1+n)/2`$? It was anybodyโs guess $`\mathrm{}`$. (As it turned out, Suslin did mean $`\mathrm{\hspace{0.17em}1}+\frac{n}{2}`$, as he, perhaps, should.) Later in June that year, in Paris, Bass was to give a similar lecture on the status of Serreโs Conjecture in the โSรฉminaire Bourbakiโ. The write-up of these survey lectures subsequently appeared in Bassโs article \[B<sub>12</sub>\], with the charming title โLibรฉration des modules projectifs $`\mathrm{}`$โ.
Serreโs Conjecture stood open for over twenty years, and was finally proved in 1976, completely independently and almost simultaneously, by D. Quillen \[Qu\] and A. Suslin \[Su\]. (For a detailed exposition on this, see \[La<sub>2</sub>\].) As is often the case in mathematics, however, the solution of one important conjecture was only to be followed by the formulation of a new, more powerful, conjecture. After 1976, Serreโs Conjecture was generalized into the so-called Bass-Quillen Conjecture, which states the following:
(3.2)<sub>d</sub> If $`A`$ is a commutative regular ring<sup>6</sup><sup>6</sup>6A commutative noetherian ring $`A`$ is said to be regular if $`A_๐ญ`$ is a regular local ring for any prime ideal $`๐ญA`$. of Krull dimension $`d<\mathrm{}`$, then any f.g. projective $`A[x_1,\mathrm{},x_n]`$-module is extended from $`A`$.
When $`A`$ is a field, of course, this gives back the original Serre Conjecture (now the Quillen-Suslin Theorem). When $`A`$ is a Dedekind ring and $`n=1`$, $`(3.2)_1`$ is the Theorem of Seshadri, Bass, and Serre. By a powerful general technique known as โQuillen Inductionโ, which is gleaned from Quillenโs solution of Serreโs Conjecture (see \[La<sub>2</sub>: p. 139\]), one can reduce the proof of $`(3.2)_d`$ to a demonstration of the following special case of it:
(3.3)<sub>d</sub> If $`A`$ is a regular local ring of Krull dimension $`d`$, then any f.g. projective $`A[x_1]`$-module is free.
The best result on $`(3.2)_d`$ and $`(3.3)_d`$ known to me is that they are both true when $`d2`$, or when $`A`$ is a formal power series ring over a field; see \[La<sub>2</sub>: p. 138\]. On several occasions, I have heard research announcements claiming the general truth of $`(3.2)_d`$ (for all $`n`$ and all $`d`$), but so far I have not seen any published proofs. Thus, it appears that โSerreโs Conjectureโ first raised in the 1950s is still very much alive today: it has simply undergone a mathematical metamorphosis and has now become the even more challenging โBass-Quillen Conjectureโ. Given this, we can say with a reasonable amount of certainly that the work of Bass on Serreโs Conjecture and its generalizations will continue to have its impact on the mathematics of the next century.
ยง4. Rings with Binary Generated Ideals: Bass Rings
One way to try to prove Serreโs Conjecture over $`R=k[x_1,\mathrm{},x_n]`$ ($`k`$ a field) would be to show that any f.g. projective $`R`$-module is isomorphic to a direct sum of ideals in $`R`$ (and then use the fact that Pic$`(R)=\{1\}`$ for the unique factorization domain $`R`$). The conclusion of Bassโs Theorem (3.1) has a rather similar flavor: for any Dedekind ring $`A`$, this result says that any f.g. module $`P`$ over $`R=A[t]`$ satisfying the torsionfree hypothesis in that theorem is extended from a f.g. projective $`A`$-module $`P_0`$. Over the Dedekind ring $`A`$, $`P_0`$ is isomorphic to a direct sum of ideals, so $`P`$ is likewise isomorphic to a direct sum of ideals in $`R`$. Considerations such as this led Bass to the following general question on the decomposition of torsionfree modules:
(4.1) When is it true that any f.g. torsionfree module over a (commutative) noetherian domain $`R`$ is isomorphic to a direct sum of ideals?
Equivalently, when is it true that any f.g. indecomposable torsionfree $`R`$-module has rank 1 (that is, isomorphic to an ideal of $`R`$) ?
As it turned out, this interesting question led Bass to a fruitful program of research. In his paper \[B<sub>4</sub>\], Bass not only proved the result (3.1), but also obtained a criterion for the decomposability of all f.g. torsionfree modules into rank one modules over a noetherian domain $`R`$, under a mild assumption on the integral closure $`\stackrel{~}{R}`$ of $`R`$. The main theorem (1.7) in \[B<sub>4</sub>\] gives the following definitive result.
(4.2) Theorem. For any commutative noetherian domain $`R`$ such that $`\stackrel{~}{R}`$ is f.g. as an $`R`$-module, the following two conditions are equivalent:
(1) Any f.g. torsionfree $`R`$-module is isomorphic to a direct sum of ideals ;
(2) Any ideal in $`R`$ can be generated by two elements.
In Bassโs proof of this theorem, the hypothesis on $`\stackrel{~}{R}`$ is needed only for the implication $`(2)(1)`$.<sup>7</sup><sup>7</sup>7In other words, $`(1)(2)`$ holds for any commutative noetherian domain $`R`$. This hypothesis is, of course, a very natural one from the viewpoint of algebraic geometry. As a matter of fact, Bassโs proof (in \[B<sub>4</sub>\]) for $`(2)(1)`$ in the above theorem is based on a rather subtle induction on the length of the $`R`$-module $`\stackrel{~}{R}/R`$. The fact that this module has finite length is a consequence of the first conclusion in the following result describing some of the key properties of commutative domains with binary generated ideals.
(4.3) Theorem. If a commutative domain $`R`$ has the property that any ideal in $`R`$ can be generated by two elements, then:
(1) $`R`$ has Krull dimension $`1`$;
(2) every ideal in $`R`$ is reflexive<sup>8</sup><sup>8</sup>8An $`R`$-module $`P`$ is said to be reflexive if the natural map from $`P`$ to its double dual $`P^{}`$ is an isomorphism. as an $`R`$-module .
Here, the first conclusion, (1), goes back to I. S. Cohen. In fact, Cohen has proved already in 1949 that the conclusion $`K`$-dim $`R1`$ will follow if every ideal of the domain $`R`$ can be generated by $`k`$ elements for a fixed integer $`k`$ \[C: p. 37, Cor. 1\]. Once we have $`K`$-dim $`R1`$, the conclusion (2) follows from Lemma (1.6) of \[B<sub>4</sub>\].
Of course, the driving force behind all of these conditions in (4.2) and (4.3) is the basic example of a Dedekind ring $`R`$. For such a ring, the properties (1), (2) in (4.2) and (4.3) are well-known to every student of abstract algebra. The case of Dedekind rings $`R`$ is precisely the integrally closed case of these results, that is, when $`\stackrel{~}{R}=R`$. This is, in fact, the beginning case for Bassโs proof for $`(2)(1)`$ in (4.2) by induction on $`\text{length}_R(\stackrel{~}{R}/R)`$. How about the case of non-Dedekind rings? It behooves us to recall the two basic examples mentioned by Bass in \[B<sub>4</sub>: p. 324\]:
* Let $`R`$ be the ring $`\{a+2bi:a,b\}`$, as a subring of the Dedekind ring of Gaussian integers $`S`$. Here, $`\stackrel{~}{R}=S`$, and $`\text{length}_R(\stackrel{~}{R}/R)=1`$. There are only two isomorphism types of nonzero ideals in $`R`$, namely, $`R`$ and $`\mathrm{\hspace{0.17em}2}R+2iR`$, so indeed every $`R`$-ideal is binary generated.
* For another typical example, let $`I`$ be the additive submonoid of the non-negative integers generated by $`2`$ and an odd integer $`n3`$, and let $`R`$ be the subring of $`S=[[t]]`$ consisting of power series $`_{i0}a_it^i`$ with $`a_i=0`$ for $`iI`$. Then $`\stackrel{~}{R}=S`$, $`\text{length}_R(\stackrel{~}{R}/R)=n2`$, and again it is easy to check that every ideal of $`R`$ is binary generated.
So far we have quoted two main results ((4.2) and (4.3)) from \[B<sub>4</sub>\], which arose from the consideration of torsionfree modules over commutative noetherian domains, and have the common theme of binary generated ideals in such domains. For Bass, there was another important motivation for these results ! In fact, (4.3) is very much a part of Bassโs research program studying (not necessarily commutative) noetherian rings of finite injective dimensions over themselves. This is a very important program he started in \[B<sub>3</sub>\], and continued in \[B<sub>6</sub>\]. Let us now explain the connections.
Classically, among the noetherian rings, those that are (say, left) self-injective are the so-called quasi-Frobenius (QF) rings. These rings of self-injective dimension zero have been extensively studied in the ring theory literature. In generalization of this, Bass sought characterizations of noetherian rings $`R`$ with $`\text{id}({}_{R}{}^{}R)<\mathrm{}`$, where โidโ stands for the injective dimension (of a module). The case of $`\text{id}({}_{R}{}^{}R)1`$ was successfully characterized by Jans \[Ja\] and Bass \[B<sub>3</sub>: (3.3)\], as follows.
(4.4) Theorem. For any (left and right) noetherian ring $`R`$, the following are equivalent:
(1) $`\text{id}({}_{R}{}^{}R)1`$;
(2) every f.g. torsionless right $`R`$-module<sup>9</sup><sup>9</sup>9Yes, there is a change of side in this statement ! Torsionless modules will be defined in the paragraph following the statement of the Theorem. is reflexive ;
(3) If $`{}_{R}{}^{}P`$ is f.g. projective and $`{}_{R}{}^{}N`$ is f.g. torsionless, then any short exact sequence $`\mathrm{\hspace{0.17em}0}PMN0`$ splits.
Here, since $`R`$ is not a domain any more, we do not have a natural notion of torsionfree modules. In their place, Bass introduced the notion of a torsionless module: an $`R`$-module $`P`$ is said to be torsionless if the natural map $`\theta _P`$ from $`P`$ to its double-dual $`P^{}`$ is a monomorphism (that is, for any nonzero $`pP`$, $`f(p)0`$ for some $`fP^{}`$). If $`R`$ happens to be a commutative domain, a torsionless $`R`$-module $`P`$ is easily seen to be torsionfree; the converse does not hold in general, but does hold if $`P`$ is f.g.
(We should note in passing that, in stating the above theorem (4.4) in \[B<sub>3</sub>\], Bass had another โequivalentโ condition: โEvery right ideal in $`R`$ is reflexive.โ However, while this condition is implied by those in (4.4), Bassโs argument for the converse contained a gap, as was later acknowledged in \[B<sub>6</sub>: p.12\]. In (6.2) of \[B<sub>6</sub>\], this extra condition is restored in the commutative case, in the form โ$`R`$ is Cohen-Macaulay, and every ideal in $`R`$ is reflexive.โ)
Let us mention two well-known classes of rings that satisfy the conditions in (4.4). First, consider any noetherian left hereditary ring $`R`$. For $`E=E({}_{R}{}^{}R)`$ (the injective hull of $`{}_{R}{}^{}R`$), we have an exact sequence
$$0REE/R0.$$
Here, the quotient module $`E/R`$ must be injective, by \[La<sub>5</sub>: (3.22)\]. Therefore, we have $`\text{id}({}_{R}{}^{}R)1`$. Another interesting (and important) class of noetherian rings $`R`$ satisfying $`\text{id}({}_{R}{}^{}R)1`$ is given by the group rings $`kG`$, where $`k`$ is any Dedekind ring, and $`G`$ is any finite group. These rings are โone-step awayโ from the quasi-Frobenius rings, in that, if $`a`$ is any nonzero element in $`k`$, the quotient $`kG/(a)(k/(a))[G]`$ is a well-known example of a QF ring (see \[La<sub>5</sub>: Exer. (15.14)\]).
For a commutative noetherian local ring $`R`$, the condition $`\text{id}({}_{R}{}^{}R)<\mathrm{}`$ turns out to be one of several equivalent conditions defining a (local) Gorenstein ring, and for such a ring $`R`$, we have in fact $`\text{id}(R)=K`$-dim$`R`$. This homological characterization of a local Gorenstein ring is close in spirit to the usual homological characterization of regular local rings: recall that, by the theorem of Auslander-Buchsbaum and Serre, regular local rings are exactly those noetherian local rings $`R`$ for which we have $`\text{gl.dim}(R)<\mathrm{}`$, or equivalently, $`\text{gl.dim}(R)=K`$-dim$`R`$.
In the case of $`K`$-dim $`R=0`$, the local Gorenstein rings are precisely the local QF rings. In the case of $`K`$-dim $`R=1`$, the earliest manifestation of local Gorenstein rings was in the form of localizations of plane curves, and more generally, complete intersection curves, as was noted by Apรฉry, Samuel, Gorenstein himself, and Rosenlicht. We refer the reader to Craig Hunekeโs article in this volume for a thorough survey on the history of Gorenstein rings. In \[Hu\], Huneke traced the Gorenstein ring notion from the work of the above-named authors to that of Grothendieck and Serre, who defined local Gorenstein rings in the context of duality theory, via the use of dualizing sheaves. After Serre observed the connection to rings with finite self-injective dimensions, Bass wrote the famous โUbiquityโ paper \[B<sub>6</sub>\] in 1963 to put the whole theory of Gorenstein rings on a firm footing. One of the basic things he did in this splendid paper was to give a โglobalโ definition for Gorenstein rings: he called a general commutative noetherian ring $`R`$ Gorenstein if all localizations of $`R`$ at prime ideals are local Gorenstein rings. If $`K\text{-dim}R`$ happens to be finite, this was shown to be equivalent to $`\text{id}(R)<\mathrm{}`$, and again, in this case, $`\text{id}(R)=K`$-dim$`R`$. Various other characterizations for Gorenstein rings (e.g. in terms of primary decompositions of ideals, multiplicities, etc.) are given in the โFundamental Theoremโ in ยง1 of \[B<sub>6</sub>\].
To bring the notion of Gorenstein rings to bear on the question of decompositions of modules considered in (4.2), Bass redid and generalized this result in ยง7 of \[B<sub>6</sub>\]. To develop a theory suitable for applications to, say, integral representation theory, Bass considered now a commutative reduced noetherian ring $`R`$ of Krull dimension $`1`$, with a total ring of quotients $`K`$. It is no longer assumed that $`R`$ is a domain, but we retain the reasonable assumption that $`\stackrel{~}{R}`$, the integral closure of $`R`$ in $`K`$, is a f.g. $`R`$-module. (Of course, $`\stackrel{~}{R}`$ is just a finite direct product of Dedekind rings.) The f.g. torsionless $`R`$-modules can be seen to be exactly those f.g. $`R`$-modules $`M`$ for which the natural map $`MK_RM`$ is an injection. These $`R`$-modules may be called $`R`$-lattices, in analogy with the terminology used in integral representation theory. With this setting in place, Bass considered the following three conditions on $`R`$:
(1) Every $`R`$-ideal is generated by two elements.
(2) Any ring between $`R`$ and $`\stackrel{~}{R}`$ is a Gorenstein ring.
(3) Every indecomposable $`R`$-lattice is isomorphic to an $`R`$-ideal.
Bassโs main result in ยง7 of \[B<sub>6</sub>\] is that $`(1)(2)(3)`$ (and that (3) fails to imply (1) and (2) โonly in a situation that can be analyzed completelyโ.<sup>10</sup><sup>10</sup>10In analyzing the situations in which (3) fails to imply (1) and (2), however, Bass seemed to have overlooked certain cases. A more complete analysis when $`R`$ is local was given later by Nazarova and Roฤญter in \[NR\]; for the general case, see the paper of Haefner and Levy \[HL\]. The main difference between (3) and (1) lies in the fact that (1) is a local property (as Bass had shown in \[Ba<sub>6</sub>: (7.4)\]), while (3) is not a local property (see \[Gr: ยง2\]). This result constitutes an expansion and simplification of (4.2). The novel feature about this result is the emergence of the โhereditarily Gorensteinโ condition (2). In the post-1963 literature, a ring $`R`$ in the above setting satisfying this condition (2) (or equivalently (1)) has been, quite justifiably, called a Bass ring.
Later, various authors have obtained new characterizations for Bass rings. For instance, if $`R`$ and $`\stackrel{~}{R}`$ are as above, then the Bass ring conditions (1), (2) are further shown to be equivalent to each of the following:
(4) (Greither \[Gr\]) $`\stackrel{~}{R}`$ is binary generated as an $`R`$-module.
(5) (Levy-Wiegand \[LW\]) $`\stackrel{~}{R}/R`$ is a cyclic $`R`$-module.
(6) (Wiegand \[Wi\]) Every faithful $`R`$-lattice has a direct summand isomorphic to a faithful ideal.
Yet other characterizations are obtained by Handelman in \[Ha<sub>2</sub>\]. Not to be outdone by Bassโs โubiquityโ title, Handelman called his own paper โPropinquity of one-dimensional Gorenstein ringsโ.<sup>11</sup><sup>11</sup>11Propinquity$`:=`$ the state of being near in space or in time. Handelman used this term to refer to the fact that the rings between a Bass ring $`R`$ and its integral closure $`\stackrel{~}{R}`$ are rather โcloseโ to one another. In the 80s, Bass was one of the communicating editors for the Journal of Pure and Applied Algebra; it was perhaps not a coincidence, therefore, that the papers of Handelman, Greither, and Levy-Wiegand all appeared in that same journal !
In classical commutative algebra, there is a well-known Steinitz-Chevalley theory for f.g. modules over Dedekind domains. In studying Bass rings, a natural topic to investigate is therefore the classification problem for $`R`$-lattices. This problem has been successfully tackled by Levy and Wiegand. In \[LW\], they show, in generalization of the Steinitz-Chevalley theory, that any given $`R`$-lattice $`M`$ over a Bass ring $`R`$ is determined by its genus,<sup>12</sup><sup>12</sup>12Recall that two $`R`$-modules $`M`$ and $`N`$ are said to be in the same genus if $`M_๐ชN_๐ช`$ for every maximal ideal $`๐ชR`$. together with the class of a faithful ideal $`c\mathrm{}(M)R`$ associated with $`M`$. Levy and Wiegand have also obtained very interesting cancellation theorems for projective lattices over Bass rings, and later, Levy \[Le\] even extended these results from $`R`$-lattices to general f.g. $`R`$-modules over a specific class of Bass rings.
Of course, the most classical examples of Bass rings are the Dedekind rings (and their finite direct products). The next class of (non-integrally closed) examples are the quadratic orders, that is, $``$-orders in a quadratic number field. Any nonzero ideal in such an order $`R`$ is isomorphic to $`^2`$, and is therefore binary generated; therefore, $`R`$ is a Bass order. For yet another class of examples, consider the integral group ring $`R=G`$ for a finite abelian group $`G`$. As we have indicated before, $`R`$ is always Gorenstein, but it need not be Bass. Bassโs result, in this case, shows that (1), (2), (3) above are in fact equivalent, and, by the results of Dade and Heller-Reiner, they amount to the fact that the order of $`G`$ is square-free (see also \[Gr: Th. 8.1\]). In other words, $`|G|`$ being square-free is the necessary and sufficient condition for the integral group ring $`G`$ to be a Bass ring.
To any perceptive reader of the two papers \[B<sub>3</sub>\], \[B<sub>6</sub>\], there should be little doubt that one of the objectives Bass had in mind for his theory of Gorenstein rings and the decompositions of torsionfree modules was the potential applications to integral representation theory. Bassโs explicit mention (in \[B<sub>6</sub>: ยง7\]) of the $`R=G`$ example in the paragraph above was a clear indication of his vision in this direction. As it turned out, the task of carrying out the program of applying Gorenstein (and Bass) rings to integral representation theory was to fall on the shoulders of researchers on representation modules in the Russian School.
Of course, in integral representation theory, the rings to be considered are no longer commutative; but in a sense they are sufficiently close to commutative rings. To put ourselves into this new setting, we start with a Dedekind ring $`R`$ with quotient field $`K`$, and consider an $`R`$-order $`\mathrm{\Lambda }`$ in a finite-dimensional separable $`K`$-algebra $`A`$. In this setting, we shall be exclusively concerned with (say, left) $`\mathrm{\Lambda }`$-lattices; that is, f.g. left $`\mathrm{\Lambda }`$-modules $`L`$ which embed into the $`A`$-module $`K_RL`$ by the natural map). Following Curtis and Reiner \[CR: ยง37\], we say that $`\mathrm{\Lambda }`$ is a Gorenstein order if the left regular module $`{}_{\mathrm{\Lambda }}{}^{}\mathrm{\Lambda }`$ has the following โweakly injectiveโ property: for any $`\mathrm{\Lambda }`$-lattices $`M`$ and $`N`$, any $`\mathrm{\Lambda }`$-exact sequence
$$\mathrm{\hspace{0.17em}0}\mathrm{\Lambda }MN0$$
splits.<sup>13</sup><sup>13</sup>13This is actually just an equivalent way to say that $`\text{id}({}_{\mathrm{\Lambda }}{}^{}\mathrm{\Lambda })1`$: one can see this, essentially, by applying $`(1)(3)`$ in (4.4). It can also be seen that $`\mathrm{\Lambda }`$ being a Gorenstein order is a left/right symmetric property \[CR: (37.8)\]; hence the omission of any reference to side. We then define Bass orders by a hereditary property: $`\mathrm{\Lambda }`$ is said to be a Bass order if every $`R`$-order in $`A`$ containing $`\mathrm{\Lambda }`$ is a Gorenstein order.
Comparing these definitions with the traditionally well known ones for hereditary orders and maximal orders, we can easily verify the following hierarchy:
$$\{\text{Maximal\hspace{0.33em}Orders}\}\{\text{Hereditary\hspace{0.33em}Orders}\}\{\text{Bass\hspace{0.33em}Orders}\}\{\text{Gorenstein\hspace{0.33em}Orders}\},$$
where, as indicated, each inclusion is proper. In parallel to the properties (1), (2) and (3) in the commutative case, one can now consider the following three properties of a given $`R`$-order $`\mathrm{\Lambda }`$:
(a) Every left ideal of $`\mathrm{\Lambda }`$ is generated by two elements.
(b) $`\mathrm{\Lambda }`$ is a Bass order.
(c) Every indecomposable $`\mathrm{\Lambda }`$-lattice is isomorphic to a left ideal of $`\mathrm{\Lambda }`$.
Shortly after the appearance of \[B<sub>6</sub>\], Russian workers in integral representation theory mounted an ambitious program to try to determine the exact relationships between the three properties above in the setting of noncommutative $`R`$-orders. The definitive results were obtained around 1966-67. In \[Ro\], Roฤญter showed that $`\text{(a)}\text{(b)}`$, and in \[DKR\], Drozd, Kirichenko, and Roฤญter showed that $`\text{(b)}\text{(c)}`$. These results are the best possible, since in general, (b) does not imply (a), nor does (c) imply (b). A detailed exposition on the proofs of the implications $`\text{(a)}\text{(b)}\text{(c)}`$ can be found in ยง37 of the book of Curtis and Reiner \[CR\], on which our present discussion is based.
The results of Roฤญter, Drozd and Kirichenko are quite deep, involving rather subtle analysis of the decompositions of $`\mathrm{\Lambda }`$-lattices. But it was clearly the paradigm of the results of Bass in the commutative case that had guided the Russians in their work in this phase of integral representation theory. Subsequently, Bass orders were used as a fundamental tool in the work of Drozd-Roฤญter and Drozd-Kirichenko in their approaches to the characterization of orders of finite representation type (the corresponding work on group rings over rings of algebraic integers having been completed earlier by H. Jacobinski).
In another direction, we should mention that the idea of Gorenstein rings has also found recent applications in noncommutative algebraic geometry: a notion of (noncommutative) โAuslander-Gorenstein ringsโ has been introduced and studied by K. Ajitabh, S. P. Smith, and J. J. Zhang (see \[ASZ\]).
Today, in no small measure due to the influence of Bassโs paper, Gorenstein rings have lived up to their โubiquityโ billing, and are widely used in number theory, arithmetic and algebraic geometry, commutative (and noncommutative) algebra, theory of invariants, and combinatorics. Even the work of Andrew Wiles \[W\] on elliptic curves and modular forms leading to his spectacular proof of Fermatโs Last Theorem made use of Gorenstein rings at several crucial points: in the Appendix, in Ch. 2 (ยง1), and then in Ch. 3 on the estimates for the Selmer group. Recall that it was exactly Wilesโs earlier attempts to use Euler systems for obtaining the upper bounds on the Selmer group that had led to the โfatal flawโ in his first proof of FLT announced in Cambridge on June 23, 1993. In 1994, finally realizing that the Euler system approach was irreparable, Wiles returned to his original approach in estimating the Selmer group using ideas from Iwasawa Theory. Basically, in \[W: Ch. 3\], Wiles needed to show that certain minimal Hecke rings, which are known to be Gorenstein rings, are indeed complete intersections. This was eventually accomplished jointly with Richard Taylor in \[TW\]. The role played by commutative algebra (and Gorenstein rings in particular) in Wilesโs paper was described vividly in the following words of his \[W: p. 451\]:
> โThe turning point in this and indeed in the whole proof came in the Spring of 1991. In searching for a clue from commutative algebra I had been particularly struck some years earlier by a paper of Kunz. I had already needed to verify that the Hecke rings were Gorenstein in order to compute the congruences developed in Chapter 2. $`\mathrm{}`$ Kunzโs paper suggested the use of an invariant (the $`\eta `$-invariant $`\mathrm{}`$) which I saw could be used to test for isomorphisms between Gorenstein rings. A different invariant (the $`๐ญ/๐ญ^2`$-invariant $`\mathrm{}`$) I had already observed could be used to test for isomorphisms between complete intersections. $`\mathrm{}`$ Not long afterwards I realized that, unlikely though it seemed at first, the equality of these invariants was actually a criterion for a Gorenstein ring to be a complete intersection.โ
For related literature, see also Lenstraโs paper \[Len\] on complete intersections and Gorenstein rings. In this paper, Lenstra sharpened Wilesโs criterion (in the Appendix of \[W\]) for a finite $`๐ช`$-free local Gorenstein algebra $`T`$ over a complete discrete valuation ring $`๐ช`$ to be a complete intersection. (Lenstra was able to remove the Gorenstein assumption on $`T`$.)
Part II: Ring Theory
ยง5. Semiperfect Rings as Generalizations of Semiprimary Rings
In Part II of this paper, we come to Bassโs work in noncommutative ring theory. As we have mentioned in the Introduction, Bassโs maiden work \[B<sub>1</sub>\], developed from his Chicago thesis in 1959, is a ring-theoretic paper dealing with the homological generalization of semiprimary rings. This work turned out to be one of the most influential ring theory papers written in that period, as can be partly gauged from the following fact. In L. Smallโs collected reviews \[Sm\] of ring theory papers published in Math. Reviews in 1940-79, an average paper got at most a few cross citations from other reviews, but Bassโs paper \[B<sub>1</sub>\] managed to pull as many as 29! Usually, a reviewer would only cite a paper in order to indicate the source of a crucial topic or an important idea; the fact that, in its twenty years of existence, \[B<sub>1</sub>\] drew as many 29 cross citations from other reviews was almost without parallel in \[Sm\].
What makes \[B<sub>1</sub>\] a masterpiece was the fact that it wove together many of the themes in ring theory and homological algebra that were being developed at that time. On the ring theory side, these themes include: the Krull-Schmidt Theorem (Azumaya version), chain conditions (suitably restricted), maximal and minimal submodules (existence questions), the Jacobson radical (nilpotency questions and lifting of idempotents), and Nakayamaโs Lemma (for general, not necessarily f.g., modules). On the side of homological algebra (a pretty new subject in 1959), the themes include: projective, injective, and flat modules, projective covers, โExtโ and โTorโ functors, and all kinds of homological dimensions. Quoting from Bassโs comments on โTheorem Pโ in \[B<sub>1</sub>\], โthis result provides one of those gratifying instances in which several ostensibly diverse notions were shown to be intimately related.โ Indeed, it seems clear in retrospect that it was this remarkable bridge-building role played by the various results in \[B<sub>1</sub>\] which helped secure it a permanent place in the ring theory literature. Only a few years after the appearance of \[B<sub>1</sub>\], the key ingredients of the paper (the theory of perfect and semiperfect rings) were incorporated into a standard textbook \[L\], Lambekโs โLectures on Rings and Modulesโ, ca. 1966. Today, perfect and semiperfect rings continue to be used extensively in a wide variety of ring-theoretic settings.
In this and the next two sections, weโll give a report on Bassโs paper \[B<sub>1</sub>\] and its impact on noncommutative ring theory. In order to keep the size of these sections within bounds, however, we shall only survey below the first part of \[B<sub>1</sub>\] on perfect and semiperfect rings, and will not try to cover its second part on the finitistic homological dimensions of rings.
For the readerโs convenience, we first recall a couple of basic definitions. In noncommutative ring theory, a ring $`R`$ is said to be semiprimary if its Jacobson radical rad$`(R)`$ is nilpotent, and the quotient ring $`R/\text{rad}(R)`$ is (artinian) semisimple. Rings of this type were well known in classical ring theory, and had been studied in part as a viable generalization of one-sided artinian rings by K. Asano, G. Azumaya and T. Nakayama, among others. With the advent of the new style of algebra in the 1950s, homological properties of such rings also attracted attention, and had been explored, for instance, in some of the papers in the โNagoya seriesโ (ca. 1955-56) on the homological dimensions of modules and rings. Bass had the genius to recognize that the artinian condition on $`R/\text{rad}(R)`$ is important in its own right, and later (in \[B<sub>9</sub>: p. 504\]) defined a ring $`R`$ to be semilocal if it satisfies this condition. In the commutative case, this condition amounts to the finiteness of the number of maximal ideals in $`R`$, so Bassโs definition of semilocal rings agrees with the usual one in commutative algebra. For noncommutative rings, however, โsemilocalโ is no longer equivalent to the finiteness of the number of (one-sided or two-sided) maximal ideals. Ring theorists now know that Bass had chosen the right definition (as well as the right name!) for an important class of rings.
Like semiprimary rings, the semiperfect rings and left (right) perfect rings introduced in \[B<sub>1</sub>\] are special cases of semilocal rings. Bass was led to these rings partly by the classical idea of an injective hull due to Eckmann and Schรถpf. In general, an injective hull of a module $`M`$ is an injective module $`I`$ containing $`M`$ as a large (or essential) submodule, in the sense that:
$$\text{submodule}SI,SM=0S=0.$$
Eckmann and Schรถpf showed that (over any ring) an injective hull for $`M`$ always exists, and is unique up to an isomorphism over $`M`$. In \[Ei\], Eilenberg initiated a notion of a minimal epimorphism from a projective module to $`M`$, and used this notion to study minimal resolutions, homological dimensions and syzygies. Bass observed that, by slightly changing Eilenbergโs definition, one gets a precise dual of the notion of an injective hull: according to Bass, a projective cover of a module $`M`$ is a projective module $`P`$ with an epimorphism $`f:PM`$ such that ker$`(f)`$ is small (or superfluous) in $`P`$ in the sense that:
$$\text{submodule}SP,S+\text{ker}(f)=PS=P.$$
Such a projective cover for $`M`$ is easily seen to be unique; the only problem is that it may not exist. For instance, in the category of $``$-modules, the only objects with projective covers (in the above sense) are the free abelian groups.
The task of studying the existence of projective covers in module categories was taken up by Bass \[B<sub>1</sub>\], who defined a ring $`R`$ to be left perfect (resp. left semiperfect) if every left $`R`$-module (resp. cyclic left $`R`$-module) has a projective cover.<sup>14</sup><sup>14</sup>14The adjective โperfectโ is due to Eilenberg, who used the term โperfect categoriesโ in \[Ei\] for categories of modules possessing projective covers. The main work to be done was that of characterizing such rings in terms of other interesting conditions.
In this section, we shall focus on the semiperfect case. Here, Bassโs major result is the following.
(5.1) Theorem. For any ring $`R`$, the following are equivalent:
(1) $`R`$ is left semiperfect;
(2) Every f.g. left $`R`$-module has a projective cover;<sup>15</sup><sup>15</sup>15It is worth noting that, according to later work of F. Sandomierski (\[Sa<sub>1</sub>\], \[Mue\]; see also \[La<sub>4</sub>: (24.3)\]), this condition (2) is also equivalent to the following: $`(2)^{}`$ Every simple right $`R`$-module has a projective cover. In a same vein, Sandomierski proved that a ring $`R`$ is right perfect iff every every semisimple right $`R`$-module has a projective cover.
(3) $`R`$ is semilocal, and idempotents in $`R/\text{rad}(R)`$ can be lifted to idempotents in $`R`$.
In case $`R`$ is a commutative ring, the above conditions are also equivalent to:
(4) $`R`$ is a finite direct product of commutative local rings.
Note that $`(1)(3)`$ shows, in particular, the somewhat surprising left/right symmetry of semiperfect rings.<sup>16</sup><sup>16</sup>16In fact, the left-right symmetric characterization of semiperfect rings in (5.1)(3) is often used as their definition in the literature: see, for instance, \[La<sub>3</sub>: p. 346\]. The second condition in (3) was already quite well known in ring theory at that time, and had been studied by Kaplansky, Jacobson, and Zelinsky, among others. For instance, if rad$`(R)`$ is a nil ideal, this condition is always satisfied (see \[L: p. 72, Prop. 1\]). Thus, semiperfect rings include all semiprimary rings (which in turn include all one-sided artinian rings). It is in this sense that semiperfect rings (and the 1-sided perfect rings to be discussed later) are homological generalizations of the classically well known semiprimary rings.
The equivalence $`(1)(4)`$ in Theorem 5.1 shows that, for commutative rings, the notion of semiperfect rings has essentially nothing new to add to the existing theory. But the interesting case is that of noncommutative semiperfect rings. Since a major source of noncommutative rings is the class of endomorphism rings of modules, it is significant to ask when such endomorphism rings are semiperfect. The answer to this question is contained in the following result, which is a remarkable extension of the well-known classical theorem that the endomorphism ring of a module of finite length is always semiprimary.
(5.2) Theorem. Let $`M`$ be a right module over a ring $`S`$. Then the endomorphism ring $`\text{End}_S(M)`$ is semiperfect iff $`M`$ has a finite Azumaya decomposition, that is, a decomposition $`M=M_1\mathrm{}M_n`$ such that each endomorphism ring $`\text{End}_S(M_i)`$ is local.
Although this result was not explicitly stated in \[Ba<sub>1</sub>\], it can be proved easily using the techniques of that paper. (An explicit proof can be found in \[La<sub>3</sub>: (23.6)\].) In some sense, (5.2) shows the โubiquityโ of semiperfect rings. For instance, for any injective module $`M_S`$, if $`M`$ has finite uniform dimension, then $`\text{End}_S(M)`$ is semiperfect (and conversely). As a special case of (5.2), we also see that a ring $`R`$ is semiperfect iff the regular module $`R_R`$ has a finite Azumaya decomposition; or, in terms of idempotents, iff there is a decomposition $`\mathrm{\hspace{0.17em}1}=e_1+\mathrm{}+e_n`$ where the $`e_i`$โs are mutually orthogonal idempotents with each $`e_iRe_i`$ a local ring. The existence of such a decomposition makes it possible to generalize a considerable amount of the elementary theory of artinian rings to a semiperfect ring $`R`$, including, for instance: the classification of f.g. projective $`R`$-modules, the construction of projective covers for f.g. $`R`$-modules, the definition of a Cartan matrix, block decomposition and basic ring for $`R`$, etc. These are parts of the foundational material for the theory of artinian rings, first developed (by Brauer, Osima and others) in the context of group algebras of finite groups for applications to the theory of modular representations. It is gratifying to see that a large part of this well-known classical theory can be carried over verbatim to semiperfect rings.
The ultimate justification for the introduction of the class of semiperfect rings lies in the fact that, besides semiprimary rings, there are also many other natural classes of rings which turn out to be semiperfect. Let us make a list of some such classes below.
* If $`N_S`$ is any module of finite uniform dimension over a ring $`S`$ and $`M`$ is its injective hull, then the endomorphism ring End$`(M_R)`$ is semiperfect. (This is essentially a rehash of the remark on injective modules we made in the last paragraph.)
* If a semilocal ring $`R`$ is 1-sided self-injective, then $`R`$ is semiperfect; see \[La<sub>5</sub>: (13.4)\].
* Certain semigroup rings called Kupisch rings by H. Sato are semiperfect; see \[Sat: (5.1)\].
* Any right serial ring (that is, a ring $`R`$ such that $`R_R`$ is a direct sum of uniserial modules) is a semiperfect ring; see, e.g., \[F<sub>2</sub>: p. 81\].
* It is well-known to researchers in duality theory that if a ring $`R`$ admits a Morita duality into some other ring $`S`$, then $`R`$ and $`S`$ must be semiperfect rings. Indeed, if $`U`$ is an $`(S,R)`$-bimodule defining the duality, then there is a duality between the Serre subcategories of $`U`$-reflexive right $`R`$-modules and left $`S`$-modules. Since (left) $`S`$-modules have injective hulls, it is not difficult to check from the above duality that f.g. (right) $`R`$-modules have projective covers. Thus, $`R`$, and hence also $`S`$, must be semiperfect rings. This basic observation is due to Barbara Osofsky. As a special case, it follows that any cogenerator ring<sup>17</sup><sup>17</sup>17A ring $`R`$ is called a cogenerator ring if $`R`$ is a cogenerator both as a left and as a right $`R`$-module. A cogenerator ring is also known as a Morita ring in the literature. is a semiperfect ring.
* More generally, Sandomierski \[Sa<sub>2</sub>: p. 335\] has shown that right linearly compact rings are semiperfect, and Azumaya \[Az<sub>1</sub>\], Osofsky \[Os\] and others have shown that right PF (pseudo-Frobenius) rings<sup>18</sup><sup>18</sup>18A ring $`R`$ is called a right PF ring if any faithful right $`R`$-module is a generator. For a list of many other characterizations of such rings, see \[La<sub>5</sub>: (19.26)\]. The cogenerator rings in Footnote (17) are precisely the 2-sided PF-rings. are also semiperfect.
The notion of semiperfect rings has also been generalized in several directions. The following are two of them.
* First, prompted by Bassโs characterization (5.1)(2) for a semiperfect ring, ring-theorists have come up with a slightly more general notion: a ring $`R`$ is said to be $`F`$-semiperfect if every finitely presented left $`R`$-module has a projective cover. This terminology is due to Oberst and Schneider \[OS\]: โ$`F`$โ here stands for โfiniteโ. Clearly, every semiperfect ring is $`F`$-semiperfect. In parallel to (5.1)(3), there is the following characterization of an $`F`$-semiperfect ring: $`R`$ is $`F`$-semiperfect iff it is semiregular; that is, idempotents of $`R/\text{rad}(R)`$ can be lifted to $`R`$, and the epimorphic image $`R/\text{rad}(R)`$ of $`R`$ is a von Neumann regular ring (instead of a semisimple ring). (See, e.g. \[OS: (1.2)\], and \[Ni<sub>3</sub>, Ni<sub>4</sub>\].) There is a large supply of such rings. In fact, if $`M`$ is any quasi-injective module,<sup>19</sup><sup>19</sup>19A module $`M`$ is said to be quasi-injective if, for any submodule $`NM`$, any homomorphism from $`N`$ to $`M`$ can be extended to an endomorphism of $`M`$. Needless to say, quasi-injective modules include all injective ones. then the endomorphism ring $`\text{End}_R(M)`$ is always a semiregular ring. This fact, first proved by Faith and Utumi \[FU\], is one of the underpinnings for the Findlay-Lambek-Utumi theory of maximal rings of quotients; a self-contained proof of it can be found in \[La<sub>5</sub>: (13.1)\].
* In the theory of direct sum decompositions of modules, there is an important class of rings called โexchange ringsโ that are formally christened by Warfield in \[Wa<sub>1</sub>\]. (A ring $`R`$ is said to be an exchange ring if $`R_R`$ satisfies the exchange property introduced in the work of Crawley and Jรณnsson \[CJ\]. For more details on (and characterizations of) such rings, see \[Ni<sub>4</sub>\].) The following hierarchy shows exactly how semiregular rings and exchange rings compare with semiperfect rings:
$$\{\text{Semiperfect\hspace{0.33em}Rings}\}\{\text{Semiregular\hspace{0.33em}Rings}\}\{\text{Exchange\hspace{0.33em}Rings}\},$$
where the second inclusion was shown by Warfield (\[Wa<sub>1</sub>: Th. 3\]; see also \[Ni<sub>4</sub>: Cor. (2.3)\]). Here, both inclusions are proper, as indicated. It turns out that the main difference between semiperfect rings and exchange rings lies in a finiteness condition. In fact, by combining Nicholsonโs results \[Ni<sub>1</sub>: (4.3)\] and \[Ni<sub>4</sub>: (1.9)\], one sees that semiperfect rings are just the exchange rings which do not have infinite sets of nonzero orthogonal idempotents; this is stated as (8.4C) in Faithโs recent book \[Fa<sub>3</sub>\]. Note that this statement is much easier to prove if we replace the word โexchange ringsโ by โsemiregular ringsโ, since it is well-known that, whenever idempotents can be lifted from $`R/\text{rad}(R)`$ to $`R`$, any countable set of nonzero orthogonal idempotents in $`R/\text{rad}(R)`$ can be lifted to a similar set in $`R`$ \[La<sub>3</sub>: (21.25\], and that a von Neumann regular ring is semisimple if it has no infinite sets of nonzero orthogonal idempotents (\[Go<sub>1</sub>: (2.16)\], \[La<sub>5</sub>: Exer. 6.29\]).
Exchange rings are worthy of study in ring theory since they generalize semiperfect rings, and they form a fairly broad class of rings. Some known facts about semiperfect rings turn out to be true for exchange rings; the proofs of them are sometimes clearer when they are expressed in the context of exchange rings. For instance, Mรผllerโs well-known result \[Mue\] that any projective (right) module over a semiperfect ring $`R`$ is isomorphic to a direct sum $`_ie_iR`$ (where $`e_i=e_i^2R`$) generalizes to, and is quite easy to prove over, any exchange ring $`R`$; see \[Wa<sub>1</sub>: Th. 1\]. Apply this to a local ring $`R`$ and youโll retrieve Kaplanskyโs classical result (mentioned in ยง1) that any projective $`R`$-module is free.
Finally, we should mention that the notion of semiperfect rings has also been successfully extended to a module-theoretic setting: a module $`M`$ over a ring $`R`$ is said to be semiperfect if every quotient of $`M`$ has a projective cover. (Thus, $`R`$ is semiperfect iff the module $`{}_{R}{}^{}R`$ is.) The theory of semiperfect modules was initiated in the projective case in Mares \[M\] and Kasch-Mares \[KM\], and has been studied further by Nicholson \[Ni<sub>2</sub>, Ni<sub>3</sub>\], and Azumaya \[Az<sub>2</sub>, Az<sub>4</sub>, Az<sub>5</sub>\], among others. For a detailed treatment of semiperfect modules in the general case, see the textbook of Kasch \[Ka\]. Ever since \[B<sub>1</sub>\] appeared in 1960, semiperfect rings (and modules) have been further studied and extensively utilized in many papers in ring theory. To reflect this trend, I have devoted the last chapter of my ring theory graduate text \[La<sub>3</sub>\] to an introductory exposition on the theory of (perfect and) semiperfect rings.
ยง6. Perfect Rings and Restricted DCC
In this section, we come to Bassโs remarkable characterizations of left perfect rings; recall that these are, by definition, rings all of whose left modules have projective covers. Since left perfect rings are obviously semiperfect, we expect that one of the characterizations should be a strengthening of the condition (3) in (5.1); this is given by the condition (2) below. The other conditions will be commented upon later.
(6.1) Theorem. For any ring $`R`$ with Jacobson radical $`J=\text{rad}(R)`$, the following are equivalent:
(1) $`R`$ is left perfect;
(2) $`R`$ is semilocal, and $`J`$ is left $`T`$-nilpotent, that is, for any $`a_1,a_2,\mathrm{}J`$, $`a_1a_2\mathrm{}a_n=0`$ for some $`n`$;
(3) $`R`$ is semilocal, and every nonzero left $`R`$-module has a maximal submodule;
(4) Every flat left $`R`$-module is projective;
(5) $`R`$ satisfies DCC on principal right ideals;
(6) Any right $`R`$-module satisfies DCC on cyclic submodules;
(7) $`R`$ is semilocal, and every nonzero right $`R`$-module has a minimal submodule;
$`(7)^{}`$ $`R`$ has no infinite sets of nonzero orthogonal idempotents, and every nonzero right $`R`$-module has a minimal submodule.
It is worth pointing out that left perfect rings are interesting mostly when they are not semiprime. In fact, as soon as a left perfect ring $`R`$ is semiprime, then $`R`$ is semisimple; see, e.g. \[La<sub>3</sub>: (10.24)\].
According to Bass, the โ$`T`$โ in โ$`T`$-nilpotencyโ in (2) stands for โtransfiniteโ. Note that this condition in (2) comes somewhere between $`J`$ being nilpotent and $`J`$ being nil.<sup>20</sup><sup>20</sup>20In particular, the characterization (2) implies immediately that semiprimary rings are both right and left perfect. In general, however, right perfect and left perfect are not equivalent conditions, as was already pointed out by Bass. Similar nil-ness conditions have appeared before in the work of Levitzki on the Levitzki radical, but Bass was the first one to realize the role of $`T`$-nilpotency in the study of modules and restricted chain conditions on 1-sided ideals. As it turned out, the successful use of the $`T`$-nilpotency condition played a pivotal role in the proof of (6.1). Results such as (1.6) showed further the efficacy of the $`T`$-nilpotency conditions.
We shall now comment on the other conditions in Theorem (6.1). Along with (1), the condition (4) is of a homological nature; it can be slightly rephrased as follows.
$`(4)^{}`$ Every left $`R`$-module has the same flat as projective dimension.
Since, in general, flat modules are precisely direct limits of projective modules (by the theorem of Lazard and Govorov), (4) can also be stated in the following form:
$`(4)^{\prime \prime }`$ Direct limits of projective left $`R`$-modules are projective.
Bassโs proof for $`(4)^{\prime \prime }(5)`$ is based on a very ingenious analysis of the left $`R`$-module with generators $`x_1,x_2,\mathrm{},`$ and relations $`x_n=a_nx_{n+1}`$ for all $`n`$, where $`a_1,a_2,\mathrm{}`$ are given elements in $`R`$.
Note that the first four conditions in (6.1) ((1)โ(4)) are conditions on left $`R`$-modules, while the last four conditions ((5)โ$`(7)^{}`$) are on right ones ! This switch from left modules to right modules, albeit not new for Bass (see Footnote (9)), is in fact one of the inherent peculiar features of his Theorem (6.1). Unfortunately, because of this unusual switch of sides, Theorem (6.1) is often misquoted in the literature, sometimes even in authoritative sources: see, for instance, \[KS: Thm. 2, p. 57\]. In a couple of standard textbooks, Bassโs left (resp. right) $`T`$-nilpotency condition was renamed the โright (resp. left) vanishing conditionโ; in another textbook, the author simply switched Bassโs definitions of left and right $`T`$-nilpotency โ and did so without even alerting the reader to the difference! All of this evidently further compounded the confusion. We hereby urge all future authors to exercise restraint in changing existing definitions, and to check their statements very carefully when quoting Bassโs Theorem (6.1).
While clearly (2) and (5) in (6.1) are purely ring-theoretic conditions, Bassโs proof of $`(2)(5)`$ was routed through the homological condition (4) (or its equivalent versions $`(4)^{}`$, $`(4)^{\prime \prime }`$). This raised the challenging question whether it is possible to give a direct proof for $`(2)(5)`$ without using any homological algebra.<sup>21</sup><sup>21</sup>21The reverse implication $`(5)(2)`$ can indeed be done directly without invoking homological tools; see, for instance, the exposition in \[La<sub>3</sub>: p. 369\]. Such a proof was eventually found by Rentschler in \[Re\].
Let us now discuss in detail the interesting condition (5), which is a natural weakening of the usual artinian condition on right ideals. Here, Kaplanskyโs ideas played a key role. In his work with Arens on the topological representations of algebras in the late 40s, Kaplansky \[K<sub>1</sub>\] was motivated to consider the condition stipulating the stabilization of all chains of the form
$$aRa^2Ra^3R\mathrm{}$$
for all elements $`a`$ in a ring $`R`$. This condition was shown to be left/right symmetric many years later by Dischinger, and defines the class of strongly $`\pi `$-regular rings. We will not try to justify the somewhat clumsy terminology here, for this can be done only by making a digression into some other definitions not essential for our purposes. Suffice it to say that, in the case of commutative rings $`R`$, strong $`\pi `$-regularity amounts to $`R`$ having Krull dimension $`0`$. And, as examples of noncommutative strongly $`\pi `$-regular rings, Kaplansky mentioned the class of algebraic algebras over a field.
In an appendix to \[K<sub>1</sub>\], Kaplansky pointed out that it would also be natural to consider the DCC for all principal right ideals in a ring, that is, the condition (5) in Theorem 6.1. Kaplansky noted that the Russian mathematician Gertschikoff had considered rings satisfying this condition as far back as 1940, and obtained characterizations of such rings that are without nonzero nilpotent elements.<sup>22</sup><sup>22</sup>22In fact, Gertschikoffโs characterization worked more generally for rings possibly without an identity element. In retrospect, this result of Gertschikoff was certainly a harbinger for the equivalence of the conditions (2) and (5) in Bassโs Theorem 6.1. By a rather strange coincidence, 1959-61 turned out to be the years in which the minimum condition (5) for principal right ideals was destined to blossom: during this period, along with the publication of \[B<sub>1</sub>\] came two papers of Faith \[Fa<sub>1</sub>\] and three papers of Szรกsz \[Sz\] on the same topic. Bass invented the term โleft perfectโ for his rings (after Eilenberg), while Faith used the English acronym โMP-ringโ, and Szรกsz introduced the German acronym โMHR-ring.โ Carl Faith told me that, in his first meeting with Bass in the early 60s, they compared notes on their respective works on rings with DCC on 1-sided ideals, but found surprisingly that what they did with these rings had almost nothing in common!
Going beyond Bassโs paper \[B<sub>1</sub>\], we should mention several significant results obtained later by other authors, which were directly inspired by Theorem 6.1. Prior to the publication of \[B<sub>1</sub>\], Chase had proved in his Chicago thesis that, over any semiprimary ring $`R`$, DCC holds for f.g. submodules of any (left or right) $`R`$-module.<sup>23</sup><sup>23</sup>23For a self-contained proof of this, see Faithโs paper \[Fa<sub>2</sub>: p. 189\]. Prompted by this, Bass asked if, in any left perfect ring, f.g. right ideals satisfy the DCC. This was affirmed in 1969 by Bjรถrk (\[Bj<sub>1</sub>\], see also \[La<sub>4</sub>: (23.3)\]) who proved the amazing result that, whenever a module (over any ring) satisfies DCC on cyclic submodules, it also satisfies the DCC on f.g. submodules. In particular, this extends Chaseโs result to left perfect rings, showing that in the condition (6) in Theorem 6.1, the word โcyclicโ can be replaced by โf.g.โ. One year later, Jonah \[Jo\] added โinfinitely manyโ more equivalent conditions to the list in (6.1):
$`(0)_n`$ Any left $`R`$-module satisfies ACC on its $`n`$-generated submodules.
(Here, $`n`$ is any natural number.) Remarkably, these new characterization of left perfect rings are in terms of ascending chain conditions; also, we are now back full circle to left $`R`$-modules! Note that, in particular, Jonahโs result implies โRings with the minimum condition for principal right ideals have the maximum condition for principal left ideals.โ This was, in fact, the title of Jonahโs paper \[Jo\].
To continue our discussion on left perfect rings, we take another look at the list of equivalent conditions in (6.1). Since the condition (7) is equivalent to the ostensibly weaker condition $`(7)^{}`$, an obvious question one can ask is whether (3) can likewise be weakened to:
$`(3)^{}`$ $`R`$ has no infinite sets of nonzero orthogonal idempotents, and every nonzero left $`R`$-module has a maximal submodule.
Note that this condition is in a way โdualโ to the condition $`(7)^{}`$ (just as (3) is dual to $`(7)`$), so it seems tempting to add $`(3)^{}`$ to the list of equivalent conditions in (6.1). An open problem proposed in \[B<sub>1</sub>\] is, indeed, whether $`(3)^{}`$ is a characterization for left-perfectness of a ring $`R`$. Some authors have referred to an affirmative answer to this as โBassโs Conjectureโ, although in fact Bass merely raised the question (see \[B<sub>1</sub>: p. 471\]). The answer to this question is possibly a bit surprising: it is โyesโ in the commutative case, and โnoโ in the general case, as is shown by Koifman \[Ko\], and partly by Cozzens, Hamsher, and Renault. To give an idea of how these conclusions were obtained, we proceed as follows.
In view of (3) and $`(3)^{}`$, it is of interest to isolate the condition: โany nonzero left $`R`$-module has a maximal submoduleโ. Let us say that $`R`$ is a left-max ring if this condition is satisfied. It is not hard to see that any left-max ring has a left $`T`$-nilpotent Jacobson radical \[La<sub>4</sub>: (24.6)\]. In the commutative case, one has the following characterization of a (left)-max ring, due to Hamsher, Renault and Koifman; see \[La<sub>4</sub>: (24.9)\].
(6.2) Theorem. A commutative ring $`R`$ is (left)-max iff $`\text{rad}(R)`$ is $`T`$-nilpotent and $`R/\text{rad}(R)`$ is a von Neumann regular ring.
It follows, in particular, from this result that a commutative max-ring $`R`$ is semiregular. If, in addition, $`R`$ has no infinite sets of nonzero orthogonal idempotents, then $`R/\text{rad}(R)`$ must be semisimple by an argument we gave in ยง5. Since $`\text{rad}(R)`$ is $`T`$-nilpotent, $`R`$ is perfect, which answers Bassโs question affirmatively for commutative rings.
To treat Bassโs question in the general case, we use the notion of a left $`V`$-ring โ a ring whose simple left modules are injective. A well-known characterization for such a ring $`R`$ is that every left $`R`$-module has a zero radical (\[La<sub>5</sub>: (3.75)\]); in particular, $`R`$ must be a left-max ring. Using differential algebra, Cozzens \[Co\] constructed remarkable examples of left $`V`$-domains that are not division rings; such domains clearly satisfy $`(3)^{}`$ above, but are not left perfect rings, thus answering Bassโs question in the negative. Moreover, Cozzensโs rings are simple, principal left/right ideal domains. The emergence of rings of this type served as an important turning point for the theory of simple noetherian rings; see, e.g. the book of Cozzens and Faith \[CF\].
Many other characterizations of left perfect rings are known. Without any attempt at completeness, let us mention a few below.
* P. A. Griffith, B. Zimmermann-Huisgen and others have characterized left perfect rings in terms of conditions similar to (6.1)(4). In the homological theory of modules, an $`R`$-module $`P`$ is said to be locally projective if every $`R`$-epimorphism $`f:QP`$ is locally split (in the sense that, for every $`pP`$, there exists $`g\text{Hom}_R(P,Q)`$ such that $`fg(p)=p`$). It is easy to see that such a module $`P`$ must be flat, so the locally projective modules form a class between the class of projective modules and that of flat modules. In \[Zi<sub>2</sub>\], Zimmermann-Huisgen showed that a ring $`R`$ is left perfect iff every locally projective left $`R`$-module is projective. This is, therefore, a variant of the criterion (4) for left perfect rings.
* In Theorem 10 loc. cit., Zimmermann-Huisgen has also characterized left perfect rings in terms of $`\mathrm{}_1`$-separable modules; these are modules $`P`$ with the property that any countable subset of $`P`$ is contained in a countably generated direct summand of $`P`$.
* In the work of Harada-Ishii \[HI\], Yamagata \[Ya\], Zimmermann-Huisgen and Zimmermann \[ZZ<sub>1</sub>\], some characterizations of left perfect rings are given in terms of the exchange property of projective left $`R`$-modules.
* It seems to be a folklore result that a ring $`R`$ is left perfect (resp. semiperfect) iff every left $`R`$-module $`P`$ (resp. f.g. left $`R`$-module $`P`$) is โsupplementedโ, that is, every submodule of $`P`$ has an addition complement. (See, e.g. \[Wis: (42.6), (43.9)\].) Recently, Keskin \[Ke\] has further extended these characterizations by using the more general notion of $``$-supplemented modules.
* Optimally, characterizations for left perfect rings should be generalizable to characterizations for the endomorphism ring of a given module to be left perfect. Some results in this direction can be found in \[AF<sub>2</sub>: ยง29\], \[Wis: (43.10)\], and \[Az<sub>3</sub>\].
In closing, we should mention the fact that much of Bassโs work on left perfect rings can be extended to functor categories. In general, if A is a small additive category, we may view A as a generalization of a ring (it is a โring with several objectsโ), and even more significantly, we may view Add(A, Ab) (the category of additive functors from A into the category of abelian groups) as a generalization of the category of modules over a ring. In view of this, it is not surprising that various module-theoretic notions can be generalized to notions concerning functors in Add(A, Ab). As it turns out, after defining flat functors, projective covers of functors, and DCC for f.g. subfunctors, etc., one can completely โtransferโ Bassโs Theorem (6.1) into a theorem on the functor category Add(A, Ab). This leads then to the notion of a left (resp. right) perfect additive category A. For a full account of all this (as well as various other relations between perfect rings and model-theoretic algebra), see the book of Jensen and Lenzing \[JL\].
Now all of this is not just generalization for generalizationโs sake! For instance, one can apply it back to the case when A is a certain additive subcategory of a module category over some ring $`R`$. By doing so, one is sometimes able to make interesting connections, and even prove nontrivial results. For instance, if A is the category of all finitely presented modules over a ring $`R`$, it turns out that the study of the perfectness of A leads to various insights about (and characterizations of) the so-called pure semisimple rings. For a detailed formulation of this, see \[JL: Thm. B.14\]. As a matter of fact, we shall return to the theme of pure semisimple rings in the next section in the context of Auslanderโs representation theory of artinian rings.
Needless to say, none of the work discussed above would have been possible without the pioneering effort of Bass in \[B<sub>1</sub>\].
ยง7. Perfect Rings and Representation Theory
In mathematics, a good notion or a good theorem often has a way of finding surprising connections to things to which it might have seemed unrelated at first. As it turned out, Bassโs notion of perfect rings and his various results on them provide such an example. The unexpected connections are to the representation theory of artinian rings, and the study of the general decomposition theory of modules into direct sums. In this section, weโll give a short account on some of these connections; our discussion will culminate in an open question in noncommutative ring theory which has remained unanswered to this date. I thank N. V. Dung for suggesting that I include such a discussion here, and for explaining to me the main results and references in this area of research.
To set the stage for our discussion, we first recall the following basic notion in representation theory. A ring $`R`$ is said to be of finite representation type (FRT) if it is left artinian and there are only finitely many (isomorphism types of) f.g. indecomposable left $`R`$-modules. According to a result of Eisenbud and Griffith \[EG\], this notion is left/right symmetric, so we are justified in suppressing the word โleftโ or โrightโ in referring to rings of FRT. The study of these rings, in the form of finite-dimensional algebras over fields, goes back a long way to Brauer, Thrall, Kasch, Kneser, Kupisch, and others. It turns out that rings of FRT have a rather subtle relationship to Bassโs perfect rings, which weโll now try to explain.
We start with a notion introduced by Anderson and Fuller \[AF<sub>1</sub>\] in the decomposition theory of modules. A decomposition of an $`R`$-module $`M`$ into $`_{iI}M_i`$ ($`M_i0`$) is said to โcomplement direct summandsโ if, for each direct summand $`N`$ of $`M`$, there exists a subset $`JI`$ such that $`M=N_{jJ}M_j`$. (Note that if such a decomposition for $`M`$ exists, the summands $`M_i`$ are necessarily indecomposable.) Using this terminology and going beyond \[B<sub>1</sub>\], Anderson and Fuller obtained in \[AF<sub>1</sub>\] (ca. 1972) the following new characterization of left perfect (and semiperfect) rings:
(7.1) Theorem. A ring $`R`$ is left perfect (resp. semiperfect) iff every projective left $`R`$-module (resp. f.g. projective left $`R`$-module) has a direct sum decomposition that complements direct summands, iff the free left $`R`$-module $`RR\mathrm{}`$ (resp. $`RR`$) has a direct sum decomposition that complements direct summands.
In view of this theorem, it is natural to ask when does every left $`R`$-module admit a decomposition that complements direct summands. It turns out that the answer to this question also involves left perfect rings, albeit in different way. Given a ring $`R`$, let $`{}_{R}{}^{}U`$ be the direct sum $`_iU_i`$, where the $`U_i`$โs consist of one isomorphic copy of each finitely presented left $`R`$-module. Now let $`E`$ be the subring (without identity) of $`\text{End}_R(U)`$ consisting of endomorphisms $`f`$ such that $`U_if=0`$ for almost all $`i`$; $`E`$ is called the โleft functor ringโ of $`R`$ (for reasons that we shall not elaborate on here). In 1976, building upon the results of Auslander and Harada, Fuller \[Ful\] proved the following remarkable result (see also Simsonโs paper \[Si<sub>1</sub>\]).
(7.2) Theorem. For any ring $`R`$, the following are equivalent:
(1) Every left $`R`$-module has a decomposition that complements direct summands;
(2) Every left $`R`$-module is a direct sum of f.g. submodules;
(3) The left functor ring $`E`$ associated to $`R`$ (defined above) is left perfect.
According to a theorem of Chase, these conditions (specifically (2)) imply that $`R`$ is a left artinian ring.
In (3) above, of course, weโll need to use the notion of a left perfect ring without an identity. This is not a big problem; in fact, with essentially the same definition of left perfectness, Harada \[H\] has shown that much of Theorem 6.1 can be proved for rings without $`1`$ but โwith enough idempotentsโ (see also \[Wis: ยง49\]). The condition (2) in (7.2) expresses a very desirable module-theoretic property (arbitrary left modules can be โconstructedโ from f.g. ones) that has also been studied by Auslander \[Au<sub>2</sub>\], Gruson and Jensen \[GJ\] (and many others in the commutative case). In \[GJ\], it was shown that (2) is equivalent to $`R`$ having โpure left global dimension zeroโ, that is, if every pure short exact sequence of left $`R`$-modules splits.<sup>24</sup><sup>24</sup>24A short exact sequence of left $`R`$-modules is said to be pure if it remains exact upon tensoring by any right $`R`$-module. In 1977, Simson \[Si<sub>2</sub>\] introduced the shorter term โleft pure semisimpleโ for rings $`R`$ with this property, and later, expanding on the work of Chase and Warfield, Zimmermann-Huisgen \[Zi<sub>2</sub>\] showed that the property (2) is also equivalent to:
(4) Every left $`R`$-module is a direct sum of indecomposable submodules.
Therefore, any of (1), (2), (3), (4) is a characterization for the left pure semisimplicity of a ring $`R`$.
Now finally, we come to the connection between pure semisimplicity and rings of FRT. This is given by the following theorem of Fuller and Reiten \[FR\] (ca. 1975), which completed the earlier work of Auslander, Ringel and Tachikawa:
(7.3) Theorem. A ring $`R`$ is of FRT iff it is both left and right pure semisimple; that is, iff every left and every right $`R`$-module is a direct sum of f.g. submodules.
The major question that remains in this area of study is whether left pure semisimplicity is equivalent to right pure semisimplicity. An affirmative answer to this question is known as the โPure Semisimplicity Conjectureโ (PSC) in ring theory. If this Conjecture holds up, then the sufficiency part of the theorem above would say that left pure semisimple rings are of FRT. As a positive evidence for this, we mention the โsparsityโ result of Prest \[Pr\], Zimmermann-Huisgen and Zimmermann \[ZZ<sub>2</sub>\], which states that, over a left pure semisimple ring $`R`$, there exist only finitely many (isomorphism types of) indecomposable left $`R`$-modules of any given composition length. This means $`R`$ comes indeed โreasonably closeโ to being of finite representation type. Also, it is worth noting that, according to results of Simson \[Si<sub>3</sub>\] and Herzog \[Her\], the general form of PSC would follow as soon as one could prove that a left pure semisimple ring is necessarily right artinian. For more perspectives and recent results on PSC, see \[Si<sub>4</sub>\] and the literature referenced therein.
For connections with the work of Bass, we note that, with the help of left functor rings, PSC can actually be formulated entirely in terms of the notion of perfect rings. We need the following additional fact, which can be gleaned from the work of Auslander \[Au<sub>1</sub>\] and Fuller \[Ful\] (see also Wisbauerโs book \[Wis: (54.3)\]).
(7.4) Proposition. A ring $`R`$ is left and right pure semisimple iff its left functor ring $`E`$ (defined in the paragraph preceding (7.2)) is both left and right perfect.
Granted this and (7.2), the issue of whether a left pure semisimple ring $`R`$ needs to be right pure semisimple boils down to testing whether the perfectness of $`E`$ on the left would imply its perfectness on the right. In an earlier footnote (Footnote (20)), we have mentioned Bassโs observation that, in general, a left perfect ring need not be right perfect. The above considerations show, however, that PSC reduces to a โleft perfect implies right perfectโ statement, for the special class of left functor rings.
As we have mentioned at the end of ยง6, Bassโs ideas of using projective covers and the condition โflat $``$ projectiveโ for module categories have been successfully carried over to categories of functors. Such a generalization has proved to be quite fruitful for the representation theory of artinian rings; see, inter alia, the work of Auslander \[Au<sub>1</sub>\] and Simson \[Si<sub>1</sub>\]. Bassโs left and right $`T`$-nilpotency conditions for ideals perhaps also inspired Auslanderโs definition in \[Au<sub>1</sub>\] of noetherian and conoetherian conditions for families of homomorphisms between modules. In fact, for a left artinian ring $`R`$, the left and right $`T`$-nilpotency conditions on the Jacobson radical of the left functor ring $`E`$ are precisely equivalent to the noetherian and conoetherian conditions for families of homomorphisms between f.g. indecomposable left $`R`$-modules. From these and the preceding discussions, we see that Bassโs early work in \[B<sub>1</sub>\] has certainly played a substantial and rather interesting role in the later development of the representation theory of artinian rings.
ยง8. Stable Range of Rings
The notion of stable range for rings is another great invention of Bass that has proved to be of lasting importance in algebra and ring theory. This notion was originally introduced by Bass (in \[B<sub>9</sub>: p. 14\]) for the study of stabilization questions in algebraic $`K`$-theory, or at least for working with the functor $`K_1`$. The definition of stable range goes as follows.
(8.1) Definition. We say that a positive integer $`n`$ is in the stable range of a ring $`R`$ (or, more informally, that $`R`$ โhas stable range $`n`$โ) if, whenever $`a_1R+\mathrm{}+a_{m+1}R=R`$ for $`mn`$ (where all $`a_iR`$), there exist elements $`x_1,\mathrm{},x_mR`$ such that
$$(a_1+a_{m+1}x_1)R+\mathrm{}+(a_m+a_{m+1}x_m)R=R.$$
It is straightforward to see that, to verify the condition for stable range $`n`$ given above, it is sufficient to do so in the critical case $`m=n`$. (It follows that if $`n`$ is in the stable range for $`R`$, then so is any larger integer.) Also, although it appears that we should have referred to the above condition as $`R`$ having โrightโ stable range $`n`$ (since the definition was based on the use of right ideals), it has been shown later by Vaserstein \[V<sub>1</sub>: Th. 2\] (and also by Warfield \[Wa<sub>2</sub>: Th. (1.6)\]) that โright stable range $`n`$โ and โleft stable range $`n`$โ are actually equivalent conditions. For this reason, we shall suppress any reference to side in referring to the stable range conditions defined in (8.1).
For readers familiar with Bassโs big book \[B<sub>10</sub>\], we should point out that there is a small discrepancy between the stable range notation used here and that used in \[B<sub>10</sub>\]. What we called โstable range $`n`$โ above corresponds to what Bass called SR<sub>n+1</sub> in his book. We believe the usage in (8.1) is now the standard one.
Bassโs study of the stable range of rings was again motivated by the basic ideas of stability in the homotopy theory of vector bundles. His main results may be summarized in the fundamental theorem (8.2) below, which is to be thought of as the โ$`K_1`$-analogueโ of Corollary (2.6). Here, $`\text{GL}_n(R)`$ denotes the $`n\times n`$ general linear group over the ring $`R`$, and $`\text{E}_n(R)`$ denotes its subgroup generated by the $`n\times n`$ elementary matrices $`I_n+ae_{ij}`$ ($`ij,aR`$); $`K_1(R)`$ denotes the Whitehead group $`\text{GL}(R)/E(R)`$ in algebraic $`K`$-theory, where $`\text{GL}(R)`$ and E$`(R)`$ are, respectively, the direct limits of the groups $`\text{GL}_r(R)`$โs and $`\text{E}_r(R)`$, taken over all positive integers $`r`$.
(8.2) Theorem. (See \[B<sub>9</sub>: (4.2), (11.1)\] and \[B<sub>10</sub>: pp. 239-240\].)
(1) If a ring $`R`$ has stable range $`n`$, then $`\text{GL}_n(R)K_1(R)`$ is a surjective homomorphism. Moreover, for any $`rn+1`$, $`\text{E}_r(R)`$ is normal in $`\text{GL}_r(R)`$, and $`\text{GL}_r(R)/\text{E}_r(R)K_1(R)`$ is an isomorphism.
(2) Let $`A`$ be a commutative ring whose maximal ideal spectrum $`\text{max}(A)`$ is a noetherian space of dimension $`d`$. Then, any module-finite $`A`$-algebra $`R`$ has stable range $`d+1`$. In particular, the conclusions of (1) apply to such an algebra $`R`$ for $`n=d+1`$.
The second part of this theorem has been generalized to noncommutative noetherian rings $`R`$ by Stafford; see \[St\]. In the case of commutative rings $`R`$, there has also been work done toward the removal of the noetherian hypothesis on the maximal ideal spectrum. For instance, Heitmann \[Hei\] has shown that, if $`R`$ is a commutative ring of Krull dimension $`d`$, then $`R`$ has stable range $`d+2`$ (and $`d+1`$ if $`R`$ is a domain).
As one would perhaps expect, the case of stable range $`1`$ has a special significance, not shared by the general case of stable range $`n`$. It is true, for instance, that having stable range $`1`$ is a Morita invariant property of a ring, while, for $`n>1`$, having stable range $`n`$ is not a Morita invariant property \[V<sub>1</sub>, V<sub>2</sub>, Wa<sub>2</sub>\]. The study of stable range has led to many new and unexpected results in the arithmetic of rings and the structure theory of modules. In the rest of this section (and the next), we shall give a survey of some of the interesting mathematics which resulted from this study, and which certainly would not have been possible without the pioneering work of Bass. We begin with the following useful observation that apparently first appeared in an unpublished note of Kaplansky \[K<sub>3</sub>\].
(8.3) Proposition. If a ring $`R`$ has stable range $`1`$, then $`R`$ is Dedekind-finite; that is, $`uv=1R`$ implies that $`vu=1R`$.
Proof (following \[K<sub>3</sub>\]). Say $`uv=1R`$. Since $`vR+(1vu)R=R`$, there exists $`tR`$ such that $`w:=v+(1vu)t`$ has a right inverse. Left-multiplying this equation by $`u`$, we get $`uw=1`$. Therefore, $`w`$ has also a left inverse. It follows that $`w,u\text{U}(R)`$ (the group of units of $`R`$), and hence $`vu=1`$. QED
For an alternative approach to this Proposition, see the beginning of ยง8 below. An immediate consequence of the Proposition is the following somewhat sharper formulation for the condition of stable range $`1`$.
(8.4) Corollary. A ring $`R`$ has stable range $`\mathrm{\hspace{0.17em}1}`$ iff, whenever $`aR+bR=R`$, there exists $`xR`$ such that $`a+bx\text{U}(R)`$ (the group of units of $`R`$).
We recorded this corollary explicitly since the characterization for stable range 1 contained herein is often used as its definition in papers in the literature dealing with the stable range of rings.
In \[B<sub>9</sub>: (6.5)\], Bass proved the following basic result:
(8.5) Theorem. Any semilocal ring $`R`$ has stable range $`1`$.
For a modern reader, this is a very natural result. One can first check (without too much difficulty) that a semisimple ring has stable range $`1`$, and then deduce (8.5) from the observation that, in general, $`R`$ has $`n`$ in its stable range if $`R/\text{rad}(R)`$ does. (For more details, see \[La<sub>3</sub>: pp. 313-314\].) Conceptually, one can think of (8.5) as the โ$`0`$-dimensional caseโ of part (2) of Theorem 8.2; in fact, in the case of commutative rings, $`R`$ semilocal means that $`\text{max}(R)`$ is finite, and therefore a 0-dimensional space.
Bassโs study of stable range was quickly picked up by Estes and Ohm, who obtained various results in 1967 on stable range in the commutative case. The case of local rings shows already that a (commutative) ring $`R`$ may have arbitrary Krull dimension, but still have $`1`$ in its stable range. In \[EO\], Estes and Ohm constructed commutative Bรฉzout domains of arbitrary dimension with $`1`$ in their stable range. For other interesting examples of rings of stable range 1 and 2, see, for instance, the work of Heinzer \[He\], Jensen \[Je\], Menal-Moncasi \[MM\], and Warfield \[Wa<sub>2</sub>\].
Warfieldโs paper \[Wa<sub>2</sub>\] is a very important contribution to the study of the stable range of noncommutative rings. In this paper, he studied the endomorphism ring $`E:=\text{End}_R(M)`$ of a right module $`M`$ over an arbitrary ring $`R`$, and sought characterizations for such a ring $`E`$ to have stable range $`n`$ (for a given integer $`n`$) . Since Warfieldโs results lead to very interesting noncommutative cancellation theorems harkening back to Bassโs results in ยง2, we shall give a quick exposition on the gist of \[Wa<sub>2</sub>\] below. The first result of Warfield characterizes the stable range of $`E:=\text{End}_R(M)`$ in terms of a certain โsubstitution propertyโ of $`M`$, as follows.
(8.6) Theorem. \[Wa<sub>2</sub>: (1.6)\] Let $`M`$ be any right $`R`$-module with endomorphism ring $`E:=\text{End}_R(M)`$ (operating on the left on $`M`$). Then $`E`$ has (right) stable range $`n`$ iff $`M`$ has the following โ$`n`$-Substitution Propertyโ:
$`(S)_n`$ For any split epimorphism $`\pi `$ from $`T=M^nX`$ to $`M`$, there exists a splitting $`\phi `$ such that $`T=\phi (M)YX`$, where $`YM^n`$.
The proof of (8.6) consists of a fairly straightforward manipulation of the definition of stable range $`n`$.
(8.7) Remarks. (1) Note that if $`(S)_n`$ holds, then in the notation there we have
$`(8.8)`$
$$\text{ker}(\pi )YX,\text{and}M^nYM.$$
The first follows since both sides of the equation are direct complements of $`\phi (M)`$ in $`T`$; the second follows similarly since both $`M^n`$ and $`Y\phi (M)`$ are direct complements of $`X`$ in $`T`$ (and $`\phi (M)M`$).
(2) Since the $`n`$-Substitution Property $`(S)_n`$ on $`M`$ depends only on $`E`$ by (8.6), it follows that if $`N`$ is any module over any other ring with endomorphism ring isomorphic to $`E`$, then $`M`$ satisfies $`(S)_n`$ iff $`N`$ does.
To see the effect of (8.6) (and (8.7)) on module cancellations, we formulate the following slight improvement on Warfieldโs result in \[Wa<sub>2</sub>: (1.3)\].
(8.9) Theorem. Let $`M`$ be any right $`R`$-module with endomorphism ring $`E:=\text{End}_R(M)`$ having stable range $`n1`$. For any right $`R`$-modules $`P`$ and $`Q`$, the following are equivalent:
(1) For some module $`X`$, we have $`P=M^{n1}X`$ and $`MPMQ`$;
(2) For some modules $`X`$ and $`Y`$, we have
$`(8.10)`$
$$PM^{n1}X,QYX,\text{and}MYM^n.$$
Proof. If (2) holds, then from (8.10):
$$MPMM^{n1}XM^nXMYXMQ.$$
Conversely, if (1) holds, let $`T:=MP=MM^{n1}X`$. Since $`TMQ`$ has a split epimorphism $`\pi `$ into $`M`$ with $`\text{ker}(\pi )Q`$, the property $`(S)_n`$ for $`M`$ implies the existence of a splitting $`\phi `$ such that $`T=\phi (M)YX`$, where $`YM^n`$. By (8.8), we have $`MYM^n`$, and $`Q\text{ker}(\pi )YX`$, as required in (2). (Note that the condition $`YM^n`$ from $`(S)_n`$ is never used in all of the arguments above.) QED
Note that, if we could have the implication
$`(8.11)`$
$$MYM^nYM^{n1},$$
then in (2) above, we would have been able to conclude that $`PQ`$. While (8.11) does hold sometimes,<sup>25</sup><sup>25</sup>25By a principle of A. Dress \[La<sub>5</sub>: (18.59)\], (8.11) holds iff $`EWE^n`$ (as right $`E`$-modules) implies $`WE^{n1}`$. Thus, (8.11) will hold if, say, stably free right modules over $`E`$ are free and $`E`$ has invariant basis number. it cannot be counted on in the most general situation. To get a good cancellation theorem out of (8.9), we can, instead, impose a slightly stronger assumption on the module $`P`$, as in the following result.
(8.12) Warfieldโs Cancellation Theorem. \[Wa<sub>2</sub>: (1.2)\] Let $`M`$ be any right $`R`$-module with endomorphism ring $`E:=\text{End}_R(M)`$ having stable range $`n1`$. If $`P`$ contains a direct summand $`M^n`$, then for any $`R`$-module $`Q`$,
$$MPMQPQ.$$
Proof. Write $`P=M^nW=M^{n1}X`$, where $`X:=MW`$. If $`MPMQ`$, weโll have the ismorphisms in (8.10) for some $`Y`$, and hence
$$QYXYMWM^nW=P,$$
as desired. QED
In the special case when $`E=\text{End}_R(M)`$ has stable range 1, the theorem above holds even without any assumption on $`P`$ (or on $`Q`$); we shall come back to this point in ยง9.
In comparison with Bassโs Cancellation Theorem (2.4)(2), the advantage of Warfieldโs (8.12) lies in the fact that there is neither a dimension assumption nor a noetherian assumption in its statement. Thus, (8.12) may be regarded as a โtruly noncommutativeโ cancellation result. To see that (8.12) essentially retrieves the cancellability of a f.g. projective module $`M`$ in the setting of (2.4)(2), we may first replace $`M`$ there by $`R^m`$ for some $`m`$, and then reduce to the case when $`M=R`$. In this case, $`E=\text{End}_R(M)R`$. If $`R`$ is module-finite over a commutative ring whose maximum spectrum is a noetherian space of dimension $`d`$, then by (8.2)(2), $`R`$ has stable range $`d+1`$. So, as long as $`P`$ has a free direct summand of rank $`d+1`$, (8.12) enables us to cancel off $`M=R`$ in (2.5).
In order to apply (8.12), one needs to know about the stable range of the endomorphism rings of modules. In this direction, Warfield has further generalized Bassโs result (2.4)(2) by taking the module-finite $`A`$-algebra $`R`$ there and considering finitely presented right modules $`M`$ over $`R`$. In \[Wa<sub>2</sub>: (3.4)\], Warfield showed that, if max$`(A)`$ is noetherian of dimension $`d`$, then $`\text{End}_R(M)`$ has $`d+1`$ in its stable range. Thus, one can apply (8.12) with $`n=d+1`$. Since $`M_R`$ need no longer be a projective module, this leads to cancellation results well beyond the reach of (2.4)(2). Warfieldโs proof involved some new techniques, since Bassโs methods did not apply to endomorphism rings.
ยง9. Rings of Stable Range One
To conclude our exposition, we shall consider in this section the case of rings of stable range 1. The crucial result here is that of Bass (Theorem 8.5), which states that semilocal rings have stable range 1. Again, as it turned out, this single result served as the fountain-head of many beautiful ideas to come; the survey in this section will show how much this result has stimulated subsequent research on the stable range of rings.
For the balance of this section, we shall only consider the stable range 1 case. First let us make the following useful observation.
(9.1) Remark. Any module $`M_R`$ satisfying the 1-Substitution Property $`(S)_1`$ in (8.6) is necessarily Dedekind-finite; that is, $`M`$ is not isomorphic to a proper direct summand of itself.
To see this, suppose there is an isomorphism $`\pi :MWM`$. Applying $`(S)_1`$ to $`\pi `$, we have a splitting $`\phi `$ for $`\pi `$ (necessarily $`\pi ^1`$!) such that $`MW=\phi (M)YW`$ for some $`Y`$. But $`\phi (M)=T`$, so $`W`$ (as well as $`Y`$) must be zero, proving (9.1).
Note that (9.1) gives us another view of the fact (8.3) that a ring $`R`$ of stable range 1 is Dedekind-finite. In fact, if $`R`$ has (right) stable range 1, then, viewing $`R`$ as $`\text{End}(R_R)`$, (8.6) shows that $`R_R`$ satisfies $`(S)_1`$, and hence (9.1) shows that $`R_R`$ is Dedekind-finite. But this is the same as saying that its endomorphism ring $`R`$ is Dedekind-finite (as a ring). This is a more conceptual, if somewhat longer, proof of (8.3).
Equipped now with the information in (9.1), let us take another look at the condition $`(S)_1`$ on $`M`$:
$`(S)_1`$ For any split epimorphism $`\pi `$ from $`T=MX`$ to $`M`$, there exists a splitting $`\phi `$ such that $`T=\phi (M)YX`$, where $`YM`$.
Here, we have $`MYM`$ (as observed in (8.8)), so if $`(S)_1`$ holds for $`M`$, (9.1) implies that $`Y=0`$, and hence the conclusion of $`(S)_1`$ simplifies to $`T=\phi (M)X`$. Thus, $`(S)_1`$ simply says that, if $`X`$ and $`\text{ker}(\pi )`$ both have direct complements isomorphic to $`M`$ in a module $`T`$, then they have a common direct complement in $`T`$. We can therefore restate (8.6) in the case $`n=1`$ as follows.
(9.2) Theorem. The endomorphism ring $`E:=\text{End}_R(M)`$ of an $`R`$-module $`M_R`$ has stable range $`1`$ iff $`M`$ has the following โSubstitution Propertyโ: Whenever a right $`R`$-module $`T`$ has direct decompositions $`T=M_iP_i`$ for $`i=1,2`$ where $`M_1M_2M`$, there exists a submodule $`CT`$ such that $`T=CP_i`$ for $`i=1,2`$.
Note that, in the above notation, if the submodule $`C`$ exists, we have, in particular, $`P_1T/CP_2`$. Therefore, we have the following consequences of (9.2).
(9.3) Corollary. (1) If an $`R`$-module $`M_R`$ has an endomorphism ring with $`1`$ in its stable range, then $`M`$ is โcancellableโ; that is, for any right $`R`$-modules $`P,Q`$,
$$MPMQPQ.$$
(2) If a ring $`R`$ has stable range $`1`$, then $`R_R`$ is cancellable, and hence so is any f.g. projective right (respectively, left) $`R`$-module.
The result (9.3) was a cancellation theorem obtained earlier (ca. 1973) by E. G. Evans in \[Ev\]. We are reporting the results in the reverse chronological order here only because Warfieldโs result (9.2) is more general than Evansโs, so it is logically more convenient to state (9.2) first and deduce (9.3) as its corollary.
To put things in the right historical perspective, we should also point out that what we called the โSubstitution Propertyโ in (9.2) had been considered as early as 1971 by L. Fuchs. In \[Fu\], apparently unaware of Bassโs work on stable range, Fuchs obtained a different characterization of the Substitution Property, and proved the result (9.2) in the case $`M_R=R`$ (and some other cases), falling just short of proving (9.2) for any $`R`$-module $`M_R`$. Fuchs also considered the case of von Neumann regular rings $`R`$ and sought conditions for $`R`$ to have the Substitution Property (that is, to have stable range $`1`$). He obtained the following result (\[Fu: Cor. 1 and Th. 4\]), which was also proved independently by Kaplansky \[K<sub>3</sub>: Th. 3\].
(9.4) Theorem. A von Neumann regular ring $`R`$ has stable range $`1`$ iff, for $`x,yR`$, $`xRyRR/xRR/yR`$.
Since any principal right ideal $`xR`$ in a von Neumann regular ring $`R`$ is a direct summand of the right module $`R_R`$ (and $`R/xR`$ is isomorphic to any direct complement of $`xR`$), the latter condition in Theorem 9.4 amounts to an โInternal Cancellation Propertyโ of the module $`R_R`$; that is, isomorphic direct summands in $`R_R`$ have isomorphic direct complements. Now by subsequent results of Ehrlich and Handelman on regular rings $`R`$ (\[Eh\], \[Ha<sub>1</sub>\]; see also \[La<sub>4</sub>: p. 242\]), the module $`R_R`$ has this property iff the ring $`R`$ is unit-regular, in the sense that, for any $`aR`$, there exists a unit $`uR`$ such that $`a=aua`$. Therefore, (9.4) amounts to the fact that a von Neumann regular ring $`R`$ has stable range $`1`$ iff $`R`$ is unit-regular. This statement seemed to have appeared explicitly for the first time in Henriksenโs paper \[Hen: Prop. 8\].
While the above results on unit-regular rings are concerned mainly with the cancellation of f.g. projective modules, we should mention at least one case where (9.3)(1) applies more generally to all f.g. modules. This is a result due to Menal \[Me\], which deals with a certain subclass of unit-regular rings.
(9.5) Theorem. Let $`R`$ be a von Neumann regular ring all of whose primitive factor rings are artinian. Then for any f.g. projective right $`R`$-module $`M`$, $`\text{End}_R(M)`$ has stable range $`1`$. In particular, for any $`R`$-modules $`P`$ and $`Q`$, $`MPMQ`$ implies $`PQ`$.
Now let us come back to Bassโs theorem (8.5). Combining this theorem with Evansโs Cancellation Theorem (9.3), we see that any module with a semilocal endomorphism ring is cancellable. Another recent result of Facchini, Herbera, Levy and Vรกmos \[FH: (2.1)\] showed that, if $`P`$, $`Q`$ are modules with semilocal endomorphism rings, then they have the โ$`n`$-cancellation propertyโ (for any integer $`n1`$); that is,
$$P^nQ^nPQ.$$
(For a detailed survey on $`n`$-cancellation, see \[La<sub>6</sub>\].) These remarkable results give us strong motivation for finding classes of modules which have semilocal endomorphism rings. Now the problem of describing the endomorphism rings of specific kinds of modules has had a time-honored history, starting with the very famous Schurโs Lemma:
(9.6) If $`M`$ is a simple $`R`$-module, $`\text{End}_R(M)`$ is a division (and hence semilocal) ring.
(9.7) Also well-known is the following classical generalization of Schurโs Lemma: if $`M_R`$ is a module of finite length, then $`\text{End}_R(M)`$ is a semiprimary (and hence semilocal) ring; see, for instance, \[La<sub>4</sub>: Ex. (21.24)\].
(9.8) Any strongly indecomposable module has, by definition, a local (and hence semilocal) endomorphism ring. Examples include, for instance, all indecomposable modules of finite length, and all injective (in fact all quasi-injective or pure-injective) indecomposable modules; see, e.g. \[La<sub>5</sub>: (3.52), Exer. (6.32)\] and \[F<sub>2</sub>: (2.27)\].
(9.9) Camps and Dicks \[CD\] showed in 1993 that any artinian module has a semilocal endomorphism ring.
(9.10) In 1995, Herbera and Shamsuddin \[HS\] generalized the Camps-Dicks result above by showing that, if a module $`M`$ has finite uniform dimension and co-uniform dimension,<sup>26</sup><sup>26</sup>26A module $`M`$ is said to have finite uniform dimension if there is a bound on the numbers $`n`$ for which $`M`$ contains a direct sum of $`n`$ nonzero modules; dually, $`M`$ is said to have finite co-uniform dimension if there is a bound on the numbers $`n`$ for which $`M`$ has a quotient that is a direct sum of $`n`$ nonzero modules. then it has a semilocal endomorphism ring. This class includes all linearly compact modules<sup>27</sup><sup>27</sup>27An $`R`$-module $`M`$ is said to be linearly compact if, for any submodules $`N_iM`$ and any elements $`m_iM`$, any system of congruences $`\{xm_i(\text{mod}N_i)\}`$ that is finitely solvable is solvable., and therefore all artinian modules. Thus, the Herbera-Shamsuddin result implies that linearly compact modules have semilocal endomorphism rings; this affirms an earlier conjecture of Faith. In the commutative case, a stronger conclusion is possible. In this case, any linearly compact module $`M_R`$ is algebraically compact (i.e. purely injective) according to Jensen and Lenzing \[JL: p. 289\]. Therefore, combining the Herbera-Shamsuddin result above with \[JL: (7.5)\], we see that $`M`$ has a semiperfect endomorphism ring. This provides a nice connection back to the material of ยง5.
(9.11) A nonzero module is said to be uniserial if its submodules form a chain under inclusion. Such a module has clearly uniform dimension and co-uniform dimension both equal to $`1`$. Thus, if $`M:=U_1\mathrm{}U_n`$ where the $`U_i`$โs are uniserial, then the โserial moduleโ $`M`$ has uniform and co-uniform dimension $`n`$, and hence it has a semilocal endomorphism ring by (9.10). For more information on the structure of $`\text{End}(U_i)`$, $`\text{End}(M)`$, and its applications to a (weak) Krull-Schmidt Theorem for serial modules, see Facchiniโs paper \[F<sub>1</sub>\] and his recent book \[F<sub>2</sub>\].
In view of the remarks made before (9.6), any modules of the type listed in (9.6)โ(9.11) above have both the cancellation property and the $`n`$-cancellation property.
New results on stable range $`1`$ are still being discovered today. Let us just mention a couple of more recent works. In \[Ca\], Canfell studied the stable range $`1`$ condition from the viewpoint of completing diagrams of modules using automorphisms, and made connections to the notions of epi-projective and mono-injective modules. In \[Ar\], inspired by earlier results of Menal and Goodearl-Menal \[GM\], Ara proved the remarkable result that
(9.12) Theorem. Any strongly $`\pi `$-regular ring (see ยง6) has stable range $`1`$.
In the commutative case, this amounts to the fact that any ring of Krull dimension $`0`$ has stable range 1, which is close in spirit to part (2) of Bassโs Theorem (8.2). In the general case, Araโs result is rather deep, and depends heavily on the use of noncommutative techniques. Combined with the earlier result of Armendariz, Fisher and Snider \[AFS\], (9.12) implies that any module with the โFitting decomposition propertyโ is cancellable: see (5.4) in \[La<sub>6</sub>\].
In closing, we should also mention some interesting variations of the stable range one condition, due to Goodearl and Menal. In \[GM\], Goodearl and Menal considered the following two conditions for a ring $`E`$:
(A) $`a,bE,aE+bE=Eu\text{U}(E)`$ such that $`a+bu\text{U}(E)`$; and
(B) $`a,bE,u\text{U}(E)`$ such that $`au\text{and}bu^1\text{U}(E).`$
They showed that $`(B)(A)`$, and obviously $`(A)`$ stable range $`1`$. They called condition (A) (right) โunit 1-stable rangeโ; the unitary analogue of it has been exploited for $`C^{}`$-algebras. Another variation, also due to Goodearl, is a weakening of stable range 1 condition. Extending Bassโs definition, Goodearl defined a ring $`E`$ to have the power-substitution property if, whenever $`aE+bE=E`$, there exists a positive integer $`n`$ and a matrix $`X๐_n(E)`$ such that $`aI_n+bX`$ is a unit in $`๐_n(E)`$. (This property is also known to be left-right symmetric.) Clearly, if $`E`$ has stable range 1, then it has the power-substitution property. The converse of this is, however, not true. For instance, it can be shown that the ring $``$ has the power-substitution property, but it certainly does not have stable range 1. The raison dโรชtre for the power-substitution property lies in the following important result of Goodearl on โpower-cancellationโ \[Go<sub>2</sub>: Cor. 4\]:
(9.13) Theorem. Let $`M,P,Q`$ be right modules over an arbitrary ring $`R`$ such that $`MPMQ`$. If $`E:=\text{End}_R(M)`$ has the power-substitution property, then $`P^nQ^n`$ for some $`n1`$.
This result of Goodearl is, of course, an extension of (9.3)(1). In fact, in the case when the endomorphism ring $`E`$ has stable range $`1`$, the proof of Goodearlโs result boils down to that of Evans, and gives the conclusion $`PQ`$. The power of (9.13) stems from the fact that many different types of rings happen to have the power-substitution property; a good list of such rings is given in my survey \[La<sub>6</sub>: p. 34\]. In ยง5 of this survey, there is also a summary of the many interesting results, due to Goodearl, Guralnick, and Levy-Wiegand, on the power-cancellation exponent $`n`$ occurring in the statement of (9.13).
Acknowledgment. The author heartily thanks Larry Levy for a critical reading of a penultimate version of this paper. His many thoughtful comments and suggestions have led to various improvements in my exposition. |
warning/0002/hep-ph0002073.html | ar5iv | text | # References
Neutrino Tests of General and Special Relativity
C. N. Leung<sup>1</sup><sup>1</sup>1Talk presented at the Nufactโ99 Workshop, Lyon, France, July 5-9, 1999.
Department of Physics and Astronomy, University of Delaware
Newark, DE 19716, U.S.A.
Abstract
We review the status of testing the principle of equivalence and Lorentz invariance from atmospheric and solar neutrino experiments.
Let us begin by asking the question: do neutrinos obey the principle of equivalence? The answer must of course come from experiments. For ordinary matter, the most sensitive test of the equivalence principle comes from torsion balance experiments which measure the gravitational acceleration of macroscopic bodies. The current best limit is $`10^{12}`$, i.e., all macroscopic bodies experience the same gravitational acceleration to an accuracy of one part in $`10^{12}`$. Surely torsion balance experiments are not applicable to neutrinos. So the question is: what types of experiments are suitable for testing if neutrinos obey the principle of equivalence?
Suppose the neutrinos violate the principle of equivalence. A consequence will be that different neutrino types (i.e., different neutrino gravitational eigenstates) will couple to gravity with a slightly different strength. Suppose the neutrino weak interaction eigenstates are not the same as their gravitational interaction eigenstates, but are linear superpositions of them. Then, as a neutrino of definite flavour (e.g., $`\nu _\mu `$) travels through a gravitational field, the gravitational components of the flavour neutrino will develop different dynamical phases, which will result in neutrino flavour oscillations. In other words, neutrino oscillation experiments provide a laboratory to test the equivalence principle for neutrinos.
To parametrize the probabilities for neutrino oscillations arising from a violation of the equivalence principle (VEP), we follow the formalism in Ref. where the general equation governing neutrino flavour evolutions is found to be
$`i{\displaystyle \frac{d}{dt}}\left(\begin{array}{c}\nu _e\\ \nu _\mu \end{array}\right)`$ $`=`$ $`\{{\displaystyle \frac{\mathrm{\Delta }m^2}{4E}}\left[\begin{array}{cc}\mathrm{cos}2\theta _M& \mathrm{sin}2\theta _M\\ \mathrm{sin}2\theta _M& \mathrm{cos}2\theta _M\end{array}\right]`$ (5)
$`+`$ $`E|\varphi (r)|\mathrm{\Delta }f\left[\begin{array}{cc}\mathrm{cos}2\theta _G& e^{i\beta }\mathrm{sin}2\theta _G\\ e^{i\beta }\mathrm{sin}2\theta _G& \mathrm{cos}2\theta _G\end{array}\right]`$ (8)
$`+`$ $`{\displaystyle \frac{G_FN_e}{\sqrt{2}}}\left[\begin{array}{cc}1& 0\\ 0& 1\end{array}\right]\left\}\right(\begin{array}{c}\nu _e\\ \nu _\mu \end{array})_.`$ (13)
We have assumed two neutrino flavours ($`\nu _e`$ and $`\nu _\mu `$) for simplicity. The first and third line on the right-hand side of Eq.(13) are familiar from the usual studies of neutrino oscillations. The first line describes flavour mixing arising from nondegenerate neutrino vacuum masses. Here $`E`$ denotes the neutrinoโs energy, $`\mathrm{\Delta }m^2m_2^2m_1^2`$ is the neutrino mass-squared difference, and $`\theta _M`$ is the mixing angle between the weak interaction eigenstates and the mass eigenstates of the neutrinos. The third line describes the matter effects on neutrino flavour evolutions. Here $`G_F`$ is the Fermi constant and $`N_e`$ is the electron density in matter. To study neutrino oscillations in vacuum, simply set $`N_e=0`$. The second line on the right-hand side of Eq.(13) describes neutrino flavour mixing arising from VEP. Here $`\varphi (r)`$, which satisfies the boundary condition $`\varphi (r\mathrm{})0`$, is the gravitational potential through which the neutrinos propagate, $`\theta _G`$ is the mixing angle between the weak interaction eigenstates and the gravitational interaction eigenstates of the neutrinos, $`\beta `$ is a phase factor which comes about because we have mixing among three generally distinct sets of neutrino eigenstates: weak, gravitational, and mass. Finally, $`\mathrm{\Delta }f`$ is a measure of VEP defined as $`\mathrm{\Delta }ff_2f_1`$, where $`f_{1,2}`$ are parameters introduced in the linearized neutrino-gravity interaction Lagrangian density,
$$_{\mathrm{int}}=\frac{if_j}{4}\sqrt{8\pi G_N}h^{\alpha \beta }[\overline{\nu _j}\gamma _\alpha _\beta \nu _j(_\alpha \overline{\nu _j})\gamma _\beta \nu _j],j=1,2$$
(14)
to gauge the deviation from universal coupling to gravity. Here $`\nu _{1,2}`$ denote the neutrino gravitational eigenstates, $`h^{\alpha \beta }`$ is the background gravitational field, and $`G_N`$ is the Newton gravitational constant. Einsteinโs theory of general relativity predicts $`f_1=f_2=1`$. VEP occurs when $`f_1f_2`$.
It can be concluded from Eq.(13) that
1. VEP leads to neutrino flavour oscillations even if the neutrinos are massless or degenerate (in this case the phase factor $`\beta `$ can be eliminated by redefinition of the neutrino fields);
2. neutrino oscillation tests of the principle of equivalence involve two parameters: $`\mathrm{\Delta }f`$ and $`\theta _G`$; it is also necessary to know the value of the local gravitational potential;
3. VEP oscillations have a different energy dependence from oscillations due to nondegenerate neutrino masses.
The dependence of VEP oscillations on $`\varphi (r)`$ is a consequence of violating the principle of equivalence. It is also a source of uncertainty. Table 1 in Ref. shows that our local value of $`|\varphi |`$ spans from $`6\times 10^{10}`$ for the Earthโs gravitational potential to about $`3\times 10^5`$ for the gravitational potential due to our local supercluster. Its value may be even bigger if the contributions from more distant sources can be estimated. Since $`\varphi `$ appears to be dominated by distant sources, its value varies little over a typical neutrino path length for earthbound experiments and may thus be treated as a constant. It is then convenient to treat $`|\varphi |\mathrm{\Delta }f`$ as a single parameter in phenomenological analyses.
In the constant $`\varphi `$ approximation, the VEP flavour transition probability has the familiar form ($`\mathrm{\Delta }m^2`$ is assumed to be zero here):
$$P(\nu _e\nu _\mu )=\mathrm{sin}^2(2\theta _G)\mathrm{sin}^2\left(\frac{\pi L}{\lambda _G}\right)_,$$
(15)
where $`L`$ is the neutrino path length and $`\lambda _G`$ is the oscillation length given by
$$\lambda _G=\frac{\pi }{E|\varphi |\mathrm{\Delta }f}=6.2\mathrm{km}\left(\frac{10^{19}}{|\varphi |\mathrm{\Delta }f}\right)\left(\frac{1\mathrm{GeV}}{E}\right)_.$$
(16)
If we compare this with the well-known vacuum oscillation probability due to a neutrino mass difference, we see that the two cases are related by the formal connection,
$$\frac{\mathrm{\Delta }m^2}{4E}E|\varphi |\mathrm{\Delta }f\mathrm{and}\theta _M\theta _G.$$
(17)
This connection can also be gleaned from Eq.(13). In contrast to the case of flavour oscillations induced by a neutrino mass difference where the oscillation length increases with the neutrinoโs energy, VEP oscillations are characterized by an oscillation length that decreases with increasing neutrino energy. The two mechanisms may therefore be distinguished by measuring the neutrino energy spectrum. Note also that VEP oscillations will be more prominent in experiments with a large $`EL`$, i.e., high energy neutrinos and/or long path lengths, which is ideal for the neutrino factory experiments discussed in this workshop.
In a certain class of string theories, VEP may arise through interactions with a massless dilaton field. This type of equivalence principle violation may also lead to neutrino oscillations, but with an energy dependence which is the same as in oscillations due to nondegenerate neutrino masses. It is therefore difficult to distinguish this type of VEP oscillations from the mass mixing oscillations, although one can still use data from oscillation experiments to constrain the relevant neutrino-dilaton couplings. We shall not consider further VEP oscillations from string theory, but will focus on the VEP oscillations from Eq.(13), which share the same distinctive energy dependence with neutrino oscillations arising from a possible breakdown of Lorentz invariance.
If Lorentz invariance is violated, different (massless) neutrino species may have different maximum attainable speeds which are close to but not necessarily equal to $`c`$. Neutrino oscillations can occur if the neutrino flavour eigenstates are linear superpositions of their velocity eigenstates, defined to be the energy eigenstates at infinite momentum, with the probability (for two-neutrino mixing)
$$P(\nu _e\nu _\mu )=\mathrm{sin}^2(2\theta _v)\mathrm{sin}^2\left(\frac{\pi L}{\lambda _v}\right)_,$$
(18)
where $`L`$ is the neutrino path length, $`\theta _v`$ is the mixing angle between the weak interaction eigenstates and the velocity eigenstates of the neutrinos, and $`\lambda _v`$ is the oscillation length given by
$$\lambda _v=\frac{2\pi }{E\mathrm{\Delta }v}=1.24\mathrm{km}\left(\frac{10^{18}}{\mathrm{\Delta }v}\right)\left(\frac{1\mathrm{GeV}}{E}\right)_.$$
(19)
Here $`\mathrm{\Delta }v=v_2v_1`$ is the difference between the speeds of the two neutrino velocity eigenstates. Comparing Eqs.(18) and (19) with Eqs.(15) and (16), the similarities between the VLI (violation of Lorentz invariance) oscillations and VEP oscillations (for constant $`\varphi `$) are obvious.
Eqs.(16) and (19) indicate the sensitivity we can expect for VEP and VLI tests from neutrino oscillation experiments. Assuming maximal mixing, there will be appreciable flavour transitions if $`\pi L/\lambda _{G,v}O(1)`$. For atmospheric neutrinos, $`E(0.110^3)`$ GeV and $`L(2010^4)`$ km, hence $`|\varphi |\mathrm{\Delta }f`$ (or $`\mathrm{\Delta }v/2`$) in the range $`(10^{26}10^{19})`$ can be probed. For solar neutrinos, $`E(0.110)`$ MeV and, for vacuum oscillations, $`L10^8`$ km, hence $`|\varphi |\mathrm{\Delta }f`$ (or $`\mathrm{\Delta }v/2`$) in the range $`(10^{25}10^{23})`$ can be probed. For MSW solution to the solar neutrino problem, the solar gravitational potential, which is about $`1.7\times 10^5`$ at the solar core and decreases to about $`2\times 10^6`$ at the surface of the Sun, plays a significant role and the constant $`\varphi `$ approximation is not as good here. Nevertheless we can obtain an order-of-magnitude estimate for the sensitivity limit : $`|\varphi |\mathrm{\Delta }f`$ (or $`\mathrm{\Delta }v/2`$) $`(10^{23}10^{22})`$, by letting $`L=7\times 10^5`$ km, the mean solar radius (see Fig. 2 in Ref. for a more accurate estimate). For long-baseline experiments such as those envisaged for future neutrino factories, $`E(10\mathrm{\hspace{0.17em}100})`$ GeV and $`L(\mathrm{few}\times 10^2\mathrm{few}\times 10^3)`$ km, hence a sensitivity of $`(10^{25}\mathrm{\hspace{0.17em}10}^{23})`$ for $`|\varphi |\mathrm{\Delta }f`$ (or $`\mathrm{\Delta }v/2`$) can be reached. These estimates demonstrate that neutrino oscillations can provide a very sensitive test for the fundamental principles of general and special relativity.
We now confront VEP and VLI with experiments. Because of the limited time, we shall consider only atmospheric and solar neutrino experiments.
Positive evidence for $`\nu _\mu \nu _\tau `$ transitions has been established by the Super-Kamiokande (SK) Collaboration. An analysis of the 535 days of SK data on sub-GeV and multi-GeV events found that VEP and VLI oscillations were consistent with the data, although the $`\chi ^2`$ fit was not as good as the fit for mass mixing oscillations. However, later analyses that included more recent SK data on upward-going muon events, which corresponded to atmospheric neutrinos with higher energies: $`E`$ up to $`10^3`$ GeV, found that VEP and VLI oscillations were not compatible with the data. In particular, Fogli et al., did a very thorough analysis of the VEP and VLI mechanisms with the SK data. They allowed a power-law energy dependence for the oscillation length: $`\lambda E^n`$, and found that, at 90% C.L., $`n=0.9\pm 0.4`$, which identified the mass mixing mechanism as the mechanism for the observed $`\nu _\mu \nu _\tau `$ transitions. VEP and VLI oscillations, which correspond to $`n=1`$, are clearly incompatible with the data. The possibility of both mass mixing and VEP (or VLI) mechanisms contributing to the flavour transitions, as described in Eq.(13), was also studied in Ref.. It was found that, for a wide range of parameter choices, including the VEP (or VLI) contributions did not improve the fit for the mass mixing mechanism, which allowed the authors to obtain the 90% C.L. upper bounds for the $`\nu _\mu `$-$`\nu _\tau `$ sector:
$$|\mathrm{\Delta }f|<10^{19},|\mathrm{\Delta }v|<6\times 10^{24},$$
(20)
independent of the values of the corresponding mixing angles. The limit on $`\mathrm{\Delta }f`$ is obtained by assuming $`|\varphi |=3\times 10^5`$, the contribution from the local supercluster. This limit on VEP is seven orders of magnitude more stringent than the best limit from torsion balance experiments. The limit on $`\mathrm{\Delta }v`$ is also the most stringent limit on VLI to-date. Note finally that these bounds are consistent with the range of sensitivity discussed above and are within about an order of magnitude of the sensitivity limit for long-baseline experiments.
As constraining as the SK atmospheric neutrino data are, they only tell us about VEP and VLI in the $`\nu _\mu `$-$`\nu _\tau `$ sector. To study flavour transitions involving $`\nu _e`$, we turn to solar neutrino experiments. The observed solar neutrino deficit may be a result of long-wavelength neutrino vacuum oscillations or matter-enhanced flavour transitions (MSW effect) in the Sun. Based on the Standard Solar Model predictions in Ref., it was found in Ref. that only matter-enhanced VEP oscillations were compatible with the available solar neutrino data then. However, using the improved Standard Solar Model predictions and the latest solar neutrino data, and treating the <sup>8</sup>B neutrino flux as a free parameter, a recent study finds that long-wavelength VEP oscillations can also account for the solar neutrino deficit. The authors of Ref. examined constraints from the observed solar neutrino rates and from the SK spectral data. They found that, for $`\nu _e\nu _\tau `$ transitions, the allowed parameter space lay in the region
$$0.65\mathrm{sin}^2(2\theta _G)1.0\mathrm{and}3.3\times 10^{20}<|\mathrm{\Delta }f|3.3\times 10^{18}$$
(21)
for VEP and
$$0.65\mathrm{sin}^2(2\theta _v)1.0\mathrm{and}2.0\times 10^{24}<|\mathrm{\Delta }v|2.0\times 10^{22}$$
(22)
for VLI. We have again used $`|\varphi |=3\times 10^5`$ to obtain the limits on $`\mathrm{\Delta }f`$ in (21). For $`\nu _e\nu _\mu `$ transitions, constraints from the recent accelerator neutrino data from the CCFR Collaboration reduce the allowed range for $`|\mathrm{\Delta }f|`$ and $`|\mathrm{\Delta }v|`$ to $`3.3\times 10^{20}<|\mathrm{\Delta }f|<6.7\times 10^{19}`$ and $`2.0\times 10^{24}<|\mathrm{\Delta }v|<4.0\times 10^{23}`$, for the same range of values for the corresponding mixing angle.
The VEP and VLI mechanisms also admit a MSW solution to the solar neutrino problem. One generally finds two disjointed allowed regions, one for small mixing angles and one for large mixings. Figure 1, which is due to P. I. Krastev, shows the 99% confidence level allowed regions for $`|\mathrm{\Delta }f|`$ and $`\mathrm{sin}^2(2\theta _G)`$ derived from comparing all available data on the solar neutrino rates, including the 825 days of SK data, with the Standard Solar Model predictions in Ref.. Figure 1 is obtained by using only the solar gravitational potential for $`\varphi (r)`$ in Eq.(13) to determine the neutrino flavour evolution. For the $`\nu _e\nu _\mu `$ channel, the allowed regions in Figure 1 are completely ruled out by the constraints from the CCFR data. For the $`\nu _e\nu _\tau `$ channel, only the large mixing angle region is incompatible with the CCFR data<sup>2</sup><sup>2</sup>2The authors of Ref. claim that a three-neutrino mixing analysis rules out the MSW allowed regions completely, for both $`\nu _e\nu _\mu `$ and $`\nu _e\nu _\tau `$ transitions.. The remaining small mixing region may be constrained by the SK data on the recoil-electron energy spectrum, although an analysis based on earlier (504 days) SK data did not find such a constraint.
In conclusion, in addition to helping determine the masses and mixing angles of neutrinos, neutrino oscillation experiments constitute a very useful tool for testing the principle of equivalence and Lorentz invariance. Current solar and atmospheric neutrino data already provide very stringent limits. For the $`\nu _\mu `$-$`\nu _\tau `$ sector, atmospheric neutrino data imply that the principle of equivalence cannot be violated by more than one part in $`10^{19}`$ and Lorentz invariance cannot be violated by more than $`6\times 10^{24}`$. Solar neutrino data provide limits for both the $`\nu _e`$-$`\nu _\mu `$ and $`\nu _e`$-$`\nu _\tau `$ sectors. The exact limits depend on the assumption of either MSW oscillations or long-wavelength vacuum oscillations, with the long-wavelength solution imposing the stronger bounds. Future long-baseline experiments can help improve these limits with sufficiently long (thousands of kilometers) baselines and sufficiently high energy (hundreds of GeV) neutrinos. A neutrino factory will be a big access for reaching these goals.
Acknowledgement
This work was supported in part by the U.S. Department of Energy under grant DE-FG02-84ER40163. |
warning/0002/quant-ph0002094.html | ar5iv | text | # Completely Positive Quantum Dissipation
## Derivation and structure of the master equation.
Let $`H=H_0+H_\text{m}+V`$ be the Hamiltonian of the whole confined system in second quantization, $`H_0=_hE_ha_h^{}a_h`$ describing the particle (either a fermion or a boson in the state $`u_h`$), $`H_\text{m}`$ matter and $`V`$ their mutual interaction. We intend to describe a single particle, so that we consider for the total system the statistical operator $`\varrho =_{kh}a_k^{}\varrho ^\text{m}a_h\varrho _{kh}`$, where $`\varrho ^\text{m}`$ describes matter, $`\varrho _{kh}`$ is a positive matrix with trace one and we have $`a_h\varrho ^\text{m}=0`$, $`\varrho ^\text{m}a_h^{}=0`$, $`h`$. In order to consider the subdynamics of the microsystem we exploit the following reduction formula, which connects the expectations of operators in the total Fock space $`_F`$ with those of operators in the one particle Hilbert space $`^{(1)}`$
$`\text{Tr}__F\left(A\varrho \right)={\displaystyle \underset{h,k}{}}A_{hk}\varrho _{kh}=\text{Tr}_{^{(1)}}\left(\widehat{๐ }\widehat{\varrho }\right)`$
where $`A`$ has the typical structure $`A=_{h,k}a_h^{}A_{hk}a_k=_{h,k}a_h^{}h|\widehat{๐ }|ka_k`$. We intend to work on a time scale $`\tau `$ much longer than microphysical collision time, being only interested in the slow dynamics of the particle, and we shall therefore approximate the time derivative with the following coarse-grained one:
$`{\displaystyle \frac{\mathrm{\Delta }_\tau \varrho _{kh}(t)}{\tau }}`$ $`=`$ $`{\displaystyle \frac{1}{\tau }}\left[\varrho _{kh}(t+\tau )\varrho _{kh}(t)\right]`$
$`=`$ $`{\displaystyle \frac{1}{\tau }}\left[\text{Tr}__F\left(a_h^{}a_ke^{\frac{i}{\mathrm{}}H\tau }\varrho (t)e^{\frac{i}{\mathrm{}}H\tau }\right)\varrho _{kh}(t)\right].`$
Exploiting the cyclic invariance of the trace we are led to consider H-picture operators $`a_h^{}(\tau )a_k(\tau )`$, to be evaluated on the given time scale using a superoperator formalism, so that e.g. $`=\frac{i}{\mathrm{}}[H,]`$, using the integral representation
$`a_k(\tau )=e^\tau a_k={\displaystyle _{i\mathrm{}+\eta }^{+i\mathrm{}+\eta }}{\displaystyle \frac{dz}{2\pi i}}e^{z\tau }\left(z\right)^1a_k.`$
Introducing the superoperator $`๐ฏ(z)๐ฑ+๐ฑ\left(z\right)^1๐ฑ,`$ which is the analog of the T-matrix, we have
$`\left(z\right)^1=\left(z_0\right)^1+\left(z_0\right)^1๐ฏ(z)\left(z_0\right)^1`$
so that in the considered structure $`a_h^{}(\tau )a_k(\tau )`$, bilinear in the field operators, the emergence of a typically incoherent term having the CP structure $`K_\alpha \varrho K_\alpha ^{}`$ bilinear in the T-matrix naturally appears, thus confirming recent phenomenological approaches . Using the fact that $`[H,_ha_h^{}a_h]=0`$ (the interaction potential being bilinear in the particle field operators) the restriction of $`๐ฏ(z)a_k`$ to the case of a single particle may be generally written $`i\mathrm{}๐ฏ(z)a_k=_hT{}_{h}{}^{k}\left(i\mathrm{}z\right)a_h,`$ where $`T{}_{h}{}^{k}\left(z\right)`$ is an operator in the Fock space of the macrosystem only. This matrix, which according to the introduction of the time scale should exhibit a slow energy dependence, plays a central role, accounting for the peculiarities of the interaction between the particle and the considered medium. The master equation we finally obtain has the following Lindblad form
$$\frac{d\widehat{\varrho }}{d\tau }=\frac{i}{\mathrm{}}[\widehat{๐ง},\widehat{\varrho }]+\frac{1}{\mathrm{}}\underset{\lambda ,\xi }{}\left[\widehat{๐ซ}_{\lambda \xi }\widehat{\varrho }\widehat{๐ซ}{}_{\lambda \xi }{}^{}{\scriptscriptstyle \frac{1}{2}}\{\widehat{๐ซ}{}_{\lambda \xi }{}^{}\widehat{๐ซ}{}_{\lambda \xi }{}^{},\widehat{\varrho }\}\right],$$
(1)
where $`\widehat{๐ง}=\widehat{๐ง}_0+\widehat{๐ต}`$ and $`\widehat{๐ต}=\frac{1}{2}(\widehat{๐ฐ}+\widehat{๐ฐ}{}_{}{}^{})`$ with
$`k|\widehat{๐ฐ}|h`$ $`=`$ $`\text{Tr}__F\left[T{}_{h}{}^{k}(E_k+i\epsilon )\varrho ^\text{m}\right]`$
$`k|\widehat{๐ซ}_{\lambda \xi }|h`$ $`=`$ $`\sqrt{2\epsilon \pi _\xi }{\displaystyle \frac{\lambda |T{}_{h}{}^{k}(E_k+i\epsilon )|\xi }{E_k+E_\lambda E_hE_\xi i\epsilon }},`$
being $`\varrho ^\text{m}=_\xi \pi _\xi |\xi \xi |`$ the statistical operator describing matter at equilibrium and $`|\lambda ,|\xi `$ denoting eigenvectors of $`H_\text{m}`$ with eigenvalues $`E_\lambda ,E_\xi `$, while $`|k,|h`$ denote eigenvectors of $`H_0`$ with eigenvalues $`E_k,E_h`$. A detailed derivation of this master equation is given in . In considering many-particle systems important corrections due to statistics of identical particles appear. This will not be the case here since we consider a single particle distinguished from its environment through the Brownian limit in which the ratio between the masses is very far from unit. A generalization of the formalism to cope with dilute many-particle systems, in which statistic effects have been accounted for, has been considered in Ref. , aiming at understanding the transition between quantum and classical .
## Completely positive quantum Brownian motion.
Following we will make the general Ansatz $`T{}_{h}{}^{k}\left(z\right)=d^3๐d^3๐\psi ^{}(๐)u_k^{}(๐)t(z,๐๐)u_h(๐)\psi (๐)`$, where we have supposed translation invariance in the interaction kernel, and $`\psi ^{}`$, $`\psi `$ denote field operators for the macrosystem. Introducing creation and destruction operators $`b^{}`$, $`b`$ in the Fock space of the macrosystem, we may write $`T{}_{h}{}^{k}\left(z\right)=_{\eta \mu }b_\eta ^{}T_{k\eta h\mu }(z)b_\mu `$. Being interested in local dissipation effects we may safely suppose that at least far away from the boundaries the system is homogeneous so as to use as quantum numbers momentum eigenvalues, thus obtaining $`T_{k\eta h\mu }(z)=\delta _{p_\eta +p_k,p_h+p_\mu }\stackrel{~}{t}(z,\left|๐_\mu ๐_\eta \right|)`$, depending on the Fourier transform of the interaction kernel, and therefore
$$T{}_{h}{}^{k}(z)=\underset{\eta \mu }{}\delta _{p_\eta +p_k,p_h+p_\mu }b_\eta ^{}\stackrel{~}{t}(z,|๐_\mu ๐_\eta |)b_\mu .$$
(2)
We may now insert this expression in (1) to evaluate the different contributions, starting from the last, typically incoherent term. In doing this we consider the medium as composed of free gas particles, so that the energy eigenstates $`|\lambda `$, $`|\xi `$ of $`\varrho ^\mathrm{m}`$ may be obtained by the repeated action of $`b_l^{}`$ on the vacuum, and we can write $`|\lambda =|\{n_l^\lambda \}`$, $`l`$ labeling the different momenta. We therefore simply have $`\lambda |b_\eta ^{}b_\mu |\lambda =\delta _{\eta ,\mu }n_\mu ^\lambda `$ and taking the slow energy dependence of $`\stackrel{~}{t}`$ into account the contributions with $`\lambda =\xi `$ cancel out in the master equation and we need only consider the primed sum for $`\lambda \xi `$. Using, for $`\lambda \xi `$, $`\lambda |b_\eta ^{}b_\mu |\xi =\left(_{\nu \mu ,\eta }\delta _{n_\nu ^\lambda ,n_\nu ^\xi }\right)\delta _{(n_\eta ^\lambda 1),n_\eta ^\xi }\delta _{n_\mu ^\lambda ,(n_\mu ^\xi 1)}(1\delta _{\eta ,\mu })\sqrt{n_\mu ^\xi }\sqrt{n_\eta ^\lambda }`$ we come to, setting $`๐ธ_{\mu \eta }๐_\mu ๐_\eta `$
$`{\displaystyle \frac{1}{\mathrm{}}}{\displaystyle \underset{\lambda ,\xi }{}}\widehat{๐ซ}_{\lambda \xi }\widehat{\varrho }\widehat{๐ซ}{}_{\lambda \xi }{}^{}={\displaystyle \frac{2\epsilon }{\mathrm{}}}{\displaystyle \underset{pp^{}}{}}{\displaystyle \underset{\eta \mu }{}}^{}`$ $`n_\mu (1\pm n_\eta ){\displaystyle \frac{\stackrel{~}{t}(\left[๐+๐ธ_{\mu \eta }\right]^2/2M+i\epsilon ,Q_{\mu \eta })}{\frac{๐_\mu ^2}{2m}\frac{(๐_\mu ๐ธ_{\mu \eta })^2}{2m}+\frac{๐^2}{2M}\frac{(๐+๐ธ_{\mu \eta })^2}{2M}+i\epsilon }}`$
$`\times e^{\frac{i}{\mathrm{}}๐ธ_{\mu \eta }\widehat{๐}}|๐๐|\widehat{\varrho }|๐^{}๐^{}|e^{\frac{i}{\mathrm{}}๐ธ_{\mu \eta }\widehat{๐}}{\displaystyle \frac{\stackrel{~}{t}^{}(\left[๐^{}+๐ธ_{\mu \eta }\right]^2/2M+i\epsilon ,Q_{\mu \eta })}{\frac{๐_\mu ^2}{2m}\frac{(๐_\mu ๐ธ_{\mu \eta })^2}{2m}+\frac{๐_{}^{}{}_{}{}^{2}}{2M}\frac{(๐^{}+๐ธ_{\mu \eta })^2}{2M}i\epsilon }},`$
where $`M`$ denotes the mass of the Brownian particle, whose position operator is $`\widehat{๐
}`$, while $`m`$ is the mass of the gas particles. The Brownian particle is immersed in a non degenerate gas, so that $`n_\mu (1\pm n_\eta )=n_\mu (1\pm n_\eta )n_\mu `$. Considering now the quasi-diagonality of the density matrix, linked to its slow variability, we substitute in the T-matrix and the denominators $`๐`$, $`๐^{}`$ with the symmetric expression $`\frac{1}{2}(๐+๐^{})`$; furthermore we use the variables $`๐_\mu `$, $`๐๐ธ_{\mu \eta }`$, and put into evidence the ratio $`\alpha =m/M`$ between the masses, thus coming to
$`{\displaystyle \frac{4\pi m}{\mathrm{}}}{\displaystyle \underset{pp^{}}{}}{\displaystyle \underset{q}{}}{}_{}{}^{}{\displaystyle \frac{1}{q}}\left|\stackrel{~}{t}(\left[{\displaystyle \frac{๐+๐^{}}{2}}+๐\right]^2/\mathrm{\hspace{0.17em}2}M+i\epsilon ,q)\right|^2{\displaystyle \underset{\mu }{}}n_\mu \delta [(1+\alpha )q+\alpha (๐+๐^{}){\displaystyle \frac{๐}{q}}2๐_\mu {\displaystyle \frac{๐}{q}}]e^{\frac{i}{\mathrm{}}๐\widehat{๐}}|๐๐|\widehat{\varrho }|๐^{}๐^{}|e^{\frac{i}{\mathrm{}}๐\widehat{๐}}.`$
The anticommutator term can be treated in an analogous way, so that the final expression for the dissipative contributions in (1) becomes, neglecting for simplicity the slow energy dependence of the T-matrix
$`{\displaystyle \frac{4\pi m}{\mathrm{}}}{\displaystyle \underset{q}{}}{}_{}{}^{}{\displaystyle \frac{|\stackrel{~}{t}(q)|^2}{q}}{\displaystyle \underset{\mu }{}}n_\mu \{`$ $`{\displaystyle \underset{pp^{}}{}}\delta \left[(1+\alpha )q+\alpha (๐+๐^{}){\displaystyle \frac{๐}{q}}2๐_\mu {\displaystyle \frac{๐}{q}}\right]e^{\frac{i}{\mathrm{}}๐\widehat{๐}}|๐๐|\widehat{\varrho }|๐^{}๐^{}|e^{\frac{i}{\mathrm{}}๐\widehat{๐}}`$
$`{\displaystyle \frac{1}{2}}{\displaystyle \underset{p}{}}\delta [(1+\alpha )q+2\alpha ๐{\displaystyle \frac{๐}{q}}2๐_\mu {\displaystyle \frac{๐}{q}}]\{|๐๐|,\widehat{\varrho }\}\}.`$
We can now go over to the continuum limit in $`๐_\mu `$ and $`๐`$, evaluating the integral with respect to $`๐_\mu `$ in the case of a Boltzmann gas, using $`n_\mu =n\lambda _m^3\mathrm{exp}[\beta (๐_\mu ^2/2m)]`$, $`\lambda _m`$ being the thermal wavelength of the gas particles, $`n`$ their density and $`\beta =1/(k_\mathrm{B}T)`$ giving the temperature dependence. We thus obtain, in the Brownian limit $`\alpha 1`$
$`{\displaystyle \frac{4\pi ^2m^2}{\beta \mathrm{}}}n\lambda _m^3{\displaystyle d^3๐\frac{|\stackrel{~}{t}(q)|^2}{q}e^{\frac{\beta }{8m}q^2}\left[e^{\frac{i}{\mathrm{}}๐\widehat{๐}}e^{\frac{\beta }{4M}๐\widehat{๐}}\widehat{\varrho }e^{\frac{\beta }{4M}๐\widehat{๐}}e^{\frac{i}{\mathrm{}}๐\widehat{๐}}\frac{1}{2}\{e^{\frac{\beta }{2M}๐\widehat{๐}},\widehat{\varrho }\}\right]}.`$
To get to the master equation describing quantum dissipation we want to extract the temperature dependence of this expression, in the limit of small momentum transfer $`q`$. We therefore expand the exponential operators up to second order in $`๐`$, which is also equivalent to keeping terms at most bilinear in the operators $`\widehat{๐}`$ and $`\widehat{๐}`$. Due to symmetry properties of the coefficients only terms bilinear in $`๐`$ and of the form $`๐_i^2`$ ($`i`$ denoting Cartesian coordinates) need be retained, so that we have
$`{\displaystyle \frac{2\pi ^2m^2}{\beta \mathrm{}}}n\lambda _m^3{\displaystyle }d^3๐{\displaystyle \frac{|\stackrel{~}{t}(q)|^2}{q}}e^{\frac{\beta }{8m}q^2}{\displaystyle \underset{i=1}{\overset{3}{}}}๐_i^2\times `$
$`\left\{{\displaystyle \frac{1}{\mathrm{}^2}}[\widehat{๐}_i,[\widehat{๐}_i,\widehat{\varrho }]]+{\displaystyle \frac{\beta ^2}{16M^2}}[\widehat{\text{p}}_i,[\widehat{\text{p}}_i,\widehat{\varrho }]]+{\displaystyle \frac{i}{\mathrm{}}}{\displaystyle \frac{\beta }{2M}}[\widehat{๐}_i,\{\widehat{\text{p}}_i,\widehat{\varrho }\}]\right\}.`$
Let us note that in the derivation important cancellations and compensations arise between the terms coming from the anticommutator and incoherent part of (1), necessary in order to obtain the final structure, thus confirming the fact that in quantum theory we cannot have separate friction and diffusion terms . Decisive for the determination of the final structure of the equation is also the Brownian limit $`\alpha =m/M1`$. Supposing without loss of generality the medium to be isotropic, so that $`๐_i^2=\frac{1}{3}q^2`$, we can define the following coefficient
$$D_{pp}=\frac{2}{3}\frac{\pi ^2m^2}{\beta \mathrm{}}n\lambda _m^3d^3๐|\stackrel{~}{t}(q)|^2qe^{\frac{\beta }{8m}q^2}$$
(3)
depending on the collision cross-section through the T-matrix and obtain the compact expression
$`{\displaystyle \underset{i=1}{\overset{3}{}}}\left\{{\displaystyle \frac{D_{pp}}{\mathrm{}^2}}[\widehat{๐}_i,[\widehat{๐}_i,\widehat{\varrho }]]+{\displaystyle \frac{D_{qq}}{\mathrm{}^2}}[\widehat{\text{p}}_i,[\widehat{\text{p}}_i,\widehat{\varrho }]]+{\displaystyle \frac{i}{\mathrm{}}}\gamma [\widehat{๐}_i,\{\widehat{\text{p}}_i,\widehat{\varrho }\}]\right\}`$
where $`D_{qq}=(\beta \mathrm{}/4M)^2D_{pp}`$ and $`\gamma =(\beta /2M)D_{pp}`$. Exploiting (2) for the T-matrix we simply obtain for the potential term in the continuum limit $`\widehat{\text{V}}=n\frac{2\pi \mathrm{}^2}{m}d^3๐|๐๐|\text{Re}f(E_p,\theta =0)`$, $`n`$ being the density of the gas particles, so that it essentially depends on the forward scattering amplitude $`f(E_p,\theta =0)`$ as expected, and vanishes if the latter does not depend on energy. The complete master equation then becomes
$`{\displaystyle \frac{d\widehat{\varrho }}{dt}}=`$ $``$ $`{\displaystyle \frac{i}{\mathrm{}}}[\widehat{\text{H}}_0+\widehat{\text{V}},\widehat{\varrho }]{\displaystyle \frac{D_{pp}}{\mathrm{}^2}}{\displaystyle \underset{i=1}{\overset{3}{}}}[\widehat{๐}_i,[\widehat{๐}_i,\widehat{\varrho }]]`$ (4)
$``$ $`{\displaystyle \frac{D_{qq}}{\mathrm{}^2}}{\displaystyle \underset{i=1}{\overset{3}{}}}[\widehat{\text{p}}_i,[\widehat{\text{p}}_i,\widehat{\varrho }]]{\displaystyle \frac{i}{\mathrm{}}}\gamma {\displaystyle \underset{i=1}{\overset{3}{}}}[\widehat{๐}_i,\{\widehat{\text{p}}_i,\widehat{\varrho }\}].`$ (5)
This is the main result of this Letter, a CP time evolution for a quantum Brownian particle derived at a fundamental level, using a different, new approach with respect to the usual independent oscillator model. The equation obtained is translationally invariant and has the correct Lindblad form . In particular the requirement of CP amounts to check that $`D_{pp}D_{qq}\mathrm{}^2\gamma ^2/4`$, which in our case is verified with the equal sign, thus uniquely determining the different coefficients in a structure like (4) apart from an overall multiplying factor. Let us note that the requirement of a stationary thermal equilibrium solution only determines the ratio between $`D_{pp}`$ and $`\gamma `$, and also CP simply indicates that the coefficient $`D_{qq}`$ should be different from zero and within some range, without actually fixing it. This explains the wide variety of different contributions that have been added to the Caldeira equation to make it preserve positive definiteness. The fact that $`D_{pp}D_{qq}=\mathrm{}^2\gamma ^2/4`$ has as a consequence the following interesting distinctive feature: in order to write (4) in a manifest Lindblad form, only one generator for each Cartesian direction has to be introduced, instead of two. In fact using the thermal wavelength $`\lambda _M=\sqrt{\mathrm{}^2/MkT}`$ associated to the Brownian particle and defining the operators $`\widehat{\text{a}}_i=\frac{\sqrt{2}}{\lambda _M}\left(\widehat{\text{x}}_i+\frac{i}{\mathrm{}}\frac{\lambda _M^2}{4}\widehat{\text{p}}_i\right)`$, satisfying $`[\widehat{\text{a}}_i,\widehat{\text{a}}_j^{}]=\delta _{ij}`$, we can rewrite (4) in the form
$`{\displaystyle \frac{d\widehat{\varrho }}{dt}}=`$ $``$ $`{\displaystyle \frac{i}{\mathrm{}}}[\widehat{\text{H}}_0+\widehat{\text{V}},\widehat{\varrho }]{\displaystyle \frac{D_{pp}}{\mathrm{}^2}}{\displaystyle \frac{\lambda _M^2}{4}}{\displaystyle \underset{i=1}{\overset{3}{}}}{\displaystyle \frac{i}{\mathrm{}}}[\{\widehat{๐}_i,\widehat{\text{p}}_i\},\widehat{\varrho }]`$ (6)
$`+`$ $`{\displaystyle \frac{D_{pp}}{\mathrm{}^2}}\lambda _M^2{\displaystyle \underset{i=1}{\overset{3}{}}}\left[\widehat{\text{a}}_i\widehat{\varrho }\widehat{\text{a}}_i^{}{\scriptscriptstyle \frac{1}{2}}\{\widehat{\text{a}}_i^{}\widehat{\text{a}}_i,\widehat{\varrho }\}\right].`$ (7)
This makes an important qualitative difference with a more phenomenological model derived by Diรณsi , also linked to the fact that he obtains an equation with the asymmetric expression $`๐_\mu `$ instead of the momentum transfer $`๐`$. This connection between number of generators and relationships among the coefficients in a master-equation of the form (4) has not been stressed in the literature, even though it provides an important qualitative feature, helpful in providing clearcut distinctions. In this spirit our work also sheds some light on the recent phenomenological work of Gao and the subsequent following debate . Gao works from the very beginning with a single generator $`V=\mu \widehat{๐}+i\nu \widehat{\text{p}}`$ and this automatically leads him to obtain a generalized Caldeira equation with $`D_{qq}=\gamma /(8Mk_\mathrm{B}T)`$, so that $`D_{pp}D_{qq}=\mathrm{}^2\gamma ^2/4`$ is verified. This explains the difference from the diffusion coefficient $`D_{qq}=\gamma /(6Mk_\mathrm{B}T)`$ found in . Despite the fact that the coefficients in the master-equation with which Gao starts are actually completely fixed by the requirement of thermal equilibrium and his choice of a single generator, our work provides some fundamental evidence in favor of this structure, giving through (3) the quantitative estimate $`\gamma =(\beta /2M)D_{pp}`$ for the relaxation coefficient. The heavily criticized Hamiltonian term $`\frac{\gamma }{2}\{\widehat{๐},\widehat{\text{p}}\}`$ which Gao obtains, however, does not appear in (4), so that no fictitious counterterm is necessary: its appearance in rewriting (4) in the form (6) clarifies why the initial choice of Gao led to this trouble.
We have thus obtained a new fundamental derivation of CP dissipative evolution, driven by collisions with the environment, with temperature dependent friction and diffusion coefficients expressed in terms of physical quantities such as the collision cross-section. The associated master-equation has the peculiarity of being expressible in Lindblad form with only a single generator for each Cartesian direction, thus giving some evidence in favor of a recent phenomenological model , though being deprived of its unphysical features . The underlying calculations, even though recovering a single particle description by tracing over matter, are rooted in a second quantization formalism conceived for the description of a subset of reduced degrees of freedom slowly varying on a given time scale . A major extension of this model could consist in considering the effect of correlated initial conditions on the dynamics and on the property of CP, as we intend to do in the future.
I am very indebted to Prof. L. Lanz who followed the whole work and Iโd like to thank Prof. A. Barchielli and Prof. O. Melsheimer for useful discussions. This work was supported by the Alexander von Humboldt-Stiftung. |
warning/0002/hep-ph0002233.html | ar5iv | text | # Effects of a dynamical role for exchanged quarks and nuclear gluons in nuclei: multinucleon correlations in deep-inelastic lepton scattering
## Abstract
It is shown that new data from the HERMES collaboration, as well as all of the earlier improved data from experiments concerning the EMC effect and shadowing in deep-inelastic scattering of leptons from nuclei, provide strong evidence for an explicit dynamical role played by exchanged quarks and nuclear gluons in the basic, tightly-bound systems of three and four nucleons, $`{}_{}{}^{3}\mathrm{He}`$ and $`{}_{}{}^{4}\mathrm{He}`$. This opens the way for specific quark-gluon dynamics instigating multinucleon correlations in nuclei.
New experimental data have appeared , relating to the unusual behavior of nucleons and their constituents in atomic nuclei. The unusual behavior appears when nuclei are probed by deep-inelastic scattering of charged leptons. The original discovery, called the โEMC effectโ, is a depression below unity of the ratio of cross sections (per nucleon) $`\frac{\sigma _\mathrm{A}}{\sigma _\mathrm{D}}`$, for nucleus A as compared to deuterium, D. This occurs in the domain of momentum-fraction $`x`$, which is characteristic of valence quarks in a free nucleon, $`x\stackrel{>}{_{}}0.25`$. A depression in $`\frac{\sigma _\mathrm{A}}{\sigma _\mathrm{D}}`$ also occurs in the domain of small $`x\stackrel{<}{_{}}0.06`$, at low and moderate $`Q^2`$, the negative of the squared four-momentum transfer from the lepton. This effect is usually referred to as โshadowingโ. An effect of this kind is expected at small $`x`$, at and near to $`Q^2=0`$, as a geometrical effect of the nuclear surface, which diminishes the intensity of vector-meson-like components of the photon through strong interactions with nucleons near to this surface <sup>F1</sup><sup>F1</sup>F1This particular effect, accomodated by conventional concepts in nuclear physics, is not dealt with in this paper, except for an observation at the end concerning a possible difference between โshadowingโ for real photons $`(Q^2=0,\stackrel{}{Q}0)`$ and for virtual photons with $`Q^2`$ near zero ($`\stackrel{}{Q}0`$). . From the very beginning of the experimental measurements, there have been two important observations (which often have not been emphasized in subsequent work):
* The โEMC effectโ is conspicuously prominent in the basic nucleus $`{}_{}{}^{4}\mathrm{He}`$; its growth with A in going to <sup>40</sup>Ca and to <sup>119</sup>Sn is moderate. A corollary to this is the matter of possible dependence upon $`Q^2`$ of the loss of momentum fraction from valence quarks in nuclei. There have been indications in the data, that a small additional loss occurs as $`Q^2`$ is increased, that is $`\frac{\sigma _\mathrm{A}}{\sigma _\mathrm{D}}`$ is depressed further below unity. This effect, although small, is โimportant because of its direction: if the loss of valence-quark momentum fraction increases with $`Q^2`$, the underlying dynamics can hardly be solely conventional nuclear physicsโ.
* There have been indications in data of a rather sharp and strong drop of $`\frac{\sigma _\mathrm{A}}{\sigma _\mathrm{D}}`$ below unity, in the domain $`0.01<x<0.1`$ for $`0.05<Q^2<1.5(\frac{\mathrm{GeV}}{c})^2`$. This suggests the possibility that an important part of the โshadowingโ effect has an origin in specifically partonic dynamics within nuclei, an origin that should be a natural extension of the partonic dynamics which leads to the loss of momentum fraction in the $`x`$-domain of the valence quarks.
The above facts point to the possibility that quark-gluon degrees of freedom contribute explicitly to multinucleon correlations and forces in nuclei. These correlations and forces are prominently present already in the basic, tightly-bound He nuclei. When one recalls that the maximum nucleon densities encountered in <sup>3</sup>He and <sup>4</sup>He are about two times higher than those in any other nucleus, it is not surprising that if relatively short-range, three-body correlations of quark-gluonic origin exist, there are significant effects observable in these systems, in particular when probed by deep-inelastic lepton scattering. This was the basis for a detailed phenomenological model for a three-nucleon correlation (force) which involved as dynamics the exchange of one quark from each nucleon, under the influence of nuclear gluon interactions <sup>F2</sup><sup>F2</sup>F2A diagrammatic representation is in Fig. 2 of Ref. 7 and Fig. 2 of Ref. 10. Momentum fraction is lost from charged constituents to the quanta of a nuclear gluonic field.
The dynamical model offers the possibility of understanding the essential EMC effect at $`x>0.25`$, and also a part of the shadowing-like effect, in a unified picture. A quantitative description (on a $`\chi ^2`$ basis) of all of the nuclear data in the domain $`0.02x0.7`$ was achieved with a few parameters whose values are estimated a priori from physical considerations, and whose fit values agree with these estimates. The A-dependence of $`\frac{\sigma _\mathrm{A}}{\sigma _\mathrm{D}}`$ is successfully predicted, as has been shown by the detailed analyses of data by Smirnov. New predictions are made. One, which has been experimentally verified, involves the close similarity between $`\frac{\sigma _\mathrm{A}}{\sigma _\mathrm{D}}`$ for the โirregularโ nucleus <sup>6</sup>Li and that for <sup>4</sup>He. The nucleus <sup>6</sup>Li is irregular in having an r. m. s. radius of $`2.5`$ fm, about the same as that for <sup>12</sup>C. Thus, if viewed as a system of uniform nucleon density, this density is low, about one-half that of <sup>4</sup>He. From this point of view, one would expect little reduction in $`\frac{\sigma _\mathrm{A}}{\sigma _\mathrm{D}}`$ for <sup>6</sup>Li. Instead the data show that on โa bin to bin comparison the $`\frac{\mathrm{Li}}{\mathrm{He}}`$ ratio is consistent with unity over the common $`x`$ rangeโ $`(0.01\stackrel{<}{_{}}x\stackrel{<}{_{}}0.5)`$. (Note Fig. 1 of Ref. 7.) Thus, as probed by deep-inelastic lepton scattering, <sup>6</sup>Li behaves as if composed of a <sup>4</sup>He-like group (and/or <sup>3</sup>He-like groups) of nucleons.
Quantitative examples of the increased loss of momentum fraction from valence quarks (reaching a limit) with increasing $`Q^2`$, were given in Figs. 7, 8 of Ref. 10, for <sup>4</sup>He and <sup>12</sup>C in particular. The new data from HERMES give the most distinct indication of this effect, for <sup>3</sup>He and <sup>14</sup>N, over the range $`1<Q^2<20(\frac{\mathrm{GeV}}{c})^2`$, for $`x0.4`$. We describe this below. The new data from HERMES exhibit a sharp and strong drop of $`\frac{\sigma _\mathrm{A}}{\sigma _\mathrm{D}}`$ below unity for $`x<0.06`$, The data show a marked $`Q^2`$ dependence of this effect for (average) $`x`$ values in the domain $`0.01<x<0.06`$. What is strikingly unusual is that the depression of $`\frac{\sigma _\mathrm{A}}{\sigma _\mathrm{D}}`$ in this domain of $`x`$, first increases as $`Q^2`$ moves up from near zero, $`0.3<Q^2<1.5(\frac{\mathrm{GeV}}{c})^2`$. A principal purpose of this paper is to describe this effect quantitatively, and to discuss physically two possible, complementary origins of such unusual $`Q^2`$ dependence. This is done below, within the unified framework provided by the model for three-nucleon correlations which originate in explicit quark-gluon dynamics, in <sup>3</sup>He in particular. The new HERMES data show that $`\frac{\sigma _\mathrm{A}}{\sigma _D}`$ already falls significantly below unity in <sup>3</sup>He, as in <sup>4</sup>He, which was anticipated.
To start the analysis, we briefly explain the essential physical features of the formula<sup>F3</sup><sup>F3</sup>F3In the comparision with data carried out in Ref. 10, this formula is identified with the ratio of structure functions $`\frac{F_2^\mathrm{A}(x,Q^2)}{F_2^\mathrm{D}(x,Q^2)}`$. This is identical with $`\frac{\sigma _\mathrm{A}}{\sigma _\mathrm{D}}`$ if the ratio of longitudinal to transverse deep-inelastic scattering cross sections $`\stackrel{~}{R}^\mathrm{A}=\frac{\sigma _\mathrm{L}^\mathrm{A}}{\sigma _\mathrm{T}^\mathrm{A}}`$ satisfies $`\stackrel{~}{R}^\mathrm{A}(x,Q^2)=\stackrel{~}{R}^\mathrm{D}(x,Q^2)`$ (or independently of such an equality, if $`Q^20`$). Here, we use the phenomenological formula for $`\frac{\sigma _\mathrm{A}}{\sigma _\mathrm{D}}`$, and we explicitly discuss the matter (Ref. 1) of $`\stackrel{~}{R}^\mathrm{A}(x,Q^2)`$ possibly not equal to $`\stackrel{~}{R}^\mathrm{D}(x,Q^2)`$ later in this paper. for $`\frac{\sigma _\mathrm{A}}{\sigma _\mathrm{D}}`$ which provides a quantitative description of all of the nuclear data in the domain $`0.02x0.7`$, for $`1<Q^2<100(\frac{\mathrm{GeV}}{c})^2`$.
$`{\displaystyle \frac{\sigma _\mathrm{A}(x,Q^2)}{\sigma _\mathrm{D}(x,Q^2)}}=1`$ $``$ $`\delta (A)\left\{{\displaystyle \frac{3.3\sqrt{x}(1x)^3C(Q^2)f(Q^2,x^2)\sqrt{x}\mathrm{e}^{Bx^2}}{3.3\sqrt{x}(1x)^3+1.1(1x)^7}}\right\}`$ (1)
$``$ $`\delta (A)\left\{{\displaystyle \frac{1.1(1x)^7\stackrel{~}{C}x^\beta \mathrm{e}^{B^{}x^2}}{3.3\sqrt{x}(1x)^3+1.1(1x)^7}}\right\}`$
with $`B=11.3`$, $`B^{}=35`$, $`\beta =0.35`$,
$$\delta (A)=0.27\left\{1\frac{1}{A^{1/3}}\frac{1.145}{A^{2/3}}+\frac{0.93}{A}+\frac{0.88}{A^{4/3}}\frac{0.59}{A^{5/3}}\right\}$$
and
$`C(Q^2)`$ $`=`$ $`{\displaystyle \frac{3}{_0^1\frac{dx}{x}\sqrt{x}\mathrm{e}^{11.3x^2}f(Q^2,x^2)}}`$
$`f(Q^2,x^2)`$ $`=`$ $`\mathrm{e}^{48x^4\left(\frac{\mathrm{ln}(Q^2/2)}{\mathrm{ln}(Q^2/0.04)}\right)}`$
$`\stackrel{~}{C}`$ $``$ $`{\displaystyle \frac{1.1_{0.06}^1\frac{dx}{x}(1x)^7}{_{0.06}^1\frac{dx}{x}x^{0.35}\mathrm{e}^{35x^2}}}`$
In this formula for $`\frac{\sigma _\mathrm{A}}{\sigma _\mathrm{D}}`$, the second term modifies the usual valence-quark distribution of momentum fraction. While maintaining the valence-quark number through $`C(Q^2)`$, a fraction $`\delta (A)`$ of the valence quarks (per nucleon) are removed from the usual distribution for a free nucleon, and are distributed instead in a Gaussian form<sup>F4</sup><sup>F4</sup>F4The approximation of using a Gaussian form in momentum space is based upon an approximate Gaussian form for quark motion in configuration space in the three-nucleon correlation. The Gaussian form is not zero at $`x=1`$, but the contribution of this term to $`\sigma _\mathrm{A}`$ is then $`\mathrm{e}^{11.3}`$, which is negligible compared to the non-zero contribution at $`x=1`$ due to the nuclear Fermi motion. The Gaussian form does simulate the effect of Fermi motion in the ratio $`\frac{\sigma _\mathrm{A}}{\sigma _\mathrm{D}}`$, causing it to move sharply upward above unity as $`x1`$. characterized by a size parameter $`B`$; the size characterizes the spatial motion of a fraction of the quarks in a correlated, three-nucleon system. The fit parameter corresponds to a spatial dimension $`(\sqrt{3B}/m_N)`$ of about 1.2 fm, the same as that estimated a priori from geometrical considerations concerning a three-nucleon correlation, (and is well within the r. m. s. charge radius of <sup>3</sup>He, which is about 1.7 fm). The modification results in a gradual depression of $`\frac{\sigma _\mathrm{A}}{\sigma _\mathrm{D}}`$ below unity in the valence-quark domain, $`x\stackrel{>}{_{}}0.3`$. This depression increases a little with increasing $`Q^2`$, in particular in the domain $`1<Q^2<20(\frac{\mathrm{GeV}}{c})^2`$, reaching a limiting depression from above, as is illustrated in Figs. 7, 8 of Ref. 10 for <sup>4</sup>He and <sup>12</sup>C, respectively. The effect, which is due to $`f(Q^2,x^2)`$ in the second term in Eq. (1), arises physically because the correlation among valence quarks and nuclear gluons, as seen by a high-$`Q^2`$ probe, can be considered to be an additional nuclear degree of freedom, whose $`x`$-distribution โsoftensโ as $`Q^2`$ increases due to more momentum fraction being lost to nuclear gluons. In $`f(Q^2,x^2)`$, the loss is limited by the decreasing strength of the gluon coupling, $`\alpha _s(Q^2)\frac{1}{\mathrm{ln}(Q^2/0.04)}`$ (for $`\mathrm{\Lambda }_{\mathrm{QCD}}200\frac{\mathrm{MeV}}{c}`$)<sup>F5</sup><sup>F5</sup>F5 $`f(Q^2,x^2)=1`$ for $`Q^2=2(\frac{\mathrm{GeV}}{c})^2`$, and is taken as 1 for $`Q^2`$ below this value in the present analysis of the new data down to $`Q^20.3(\frac{\mathrm{GeV}}{c})^2`$.. The A dependence is given by the function $`\delta (A)`$, which uniquely arises by excluding from the three-nucleon correlation the number of nucleons $`\mathrm{A}_S`$, which reside in the (relatively diffuse) surface of a large nucleus, that is $`\delta (A)=N(1\frac{\mathrm{A}_S}{\mathrm{A}})`$. The single, overall normalization parameter, $`N=(27\pm 6)\%`$, is fit to $`\delta (\mathrm{A}=4)=6_1^{+1.5}\%`$. (Then, $`\delta (\mathrm{A}=14)12.6\%`$ and $`\delta (\mathrm{A}=3)4.2\%`$.) Let us look now at the HERMES data, shown in Figs. 1,2 for <sup>14</sup>N and <sup>3</sup>He, respectively. For $`x0.4`$, well into the domain of valence-quark momentum fraction, the data indicate that $`\frac{\sigma _{\mathrm{A}=14}}{\sigma _\mathrm{D}}`$ falls up to $`4\%`$ more below unity as $`Q^2`$ goes up from $`1`$ to $`20(\frac{\mathrm{GeV}}{c})^2`$. For <sup>3</sup>He, the additional drop may be up to $`2\%`$.<sup>F6</sup><sup>F6</sup>F6 These are small changes which might be ignored, within the experimental uncertainties. Nevertheless, the direction is in accord with earlier experimental indications, and the changes are expected within a model of explicit quark-gluon correlations in nuclei. Note also the $`Q^2`$ dependence of the <sup>119</sup>Sn data for $`0.4<x<0.5`$, in Fig. 4 of M. Arneodo et. al. (NMC), Nucl. Phys. B481, 23 (1996). These numbers appear similar to those calculated from Eq. (1) ten years ago, in the Figs. 8, 7 for <sup>12</sup>C and <sup>4</sup>He, respectively.
The physical picture of a three-nucleon correlation of quark-gluonic origin allows for the natural prediction of a shadowing-like effect, in fact a rather sharp drop of $`\frac{\sigma _\mathrm{A}}{\sigma _\mathrm{D}}`$ below unity, for $`x<0.06`$. This is contained in the last term in Eq. (1). A colored agglomeration of sea quarks (antiquarks) is exchanged from each nucleon, instead of a valence quark, under the influence of nuclear gluon interactions. The fraction of the sea (taken as $`\delta (A)`$ in Eq. (1)) involved in this dynamics is removed from the relevant momentum-fraction distribution for a free nucleon and is distributed instead in a form given by the $`x`$-dependence multiplying $`\stackrel{~}{C}`$ in Eq. (1). The crucial point is the non-zero power of $`x`$, $`\beta >0`$. Physically, this means that whereas the number of sea quarks and antiquarks in an individual nucleon is formally infinite (i. e. $`_0^1\frac{dx}{x}\mathrm{}`$), the number of these which are involved in the three-nucleon correlation is finite. In the representation for $`\frac{\sigma _\mathrm{A}}{\sigma _\mathrm{D}}`$ in Eq. (1), it is the factor $`x^\beta `$ which results in the relatively sharp drop below unity for $`x<0.06`$. <sup>F7</sup><sup>F7</sup>F7The fit parameter $`\beta `$ is restricted by $`0<\beta <\frac{1}{2}`$, because the distribution of momentum fraction in the sea is โsofterโ than that for the valence quarks ($`\sqrt{x}`$), as $`x0`$. The initial sea is taken for $`x0.06`$. The factor $`\stackrel{~}{C}(Q^2)`$ ensures that the number of sea quarks and antiquarks involved in the larger spatial domain of the three-nucleon correlation is equal to the number removed from the individual nucleon distribution. The โsizeโ parameter $`B^{}`$ in the Gaussian, characterizing the spatial motion of the sea, is expected to be larger than $`B`$; the fit number corresponds to a dimension of about $`2.2`$ fm ($`B^{}`$ may increase somewhat with $`A`$.) Eq. (1) resulted in a representation with $`\chi ^2\stackrel{<}{_{}}1`$ per degree of freedom, for all of the nuclear data taken with $`Q^2`$ between about 2 and 100 $`(\frac{\mathrm{GeV}}{c})^2`$, for $`x`$ in the domain $`0.02x0.7`$.<sup>F8</sup><sup>F8</sup>F8At smaller $`x`$, one encounters the need to deal explicitly with the geometrical shadowing effect which is present already for real photons ($`Q^2=0`$, $`\stackrel{}{Q}0`$).<sup>F1</sup>
Now the HERMES data exhibit a new $`Q^2`$ dependence for $`0.3<Q^2<1.5(\frac{\mathrm{GeV}}{c})^2`$, with $`x`$ in the domain $`0.0126x0.055`$, as shown by the first 6 graphs in Figs. 1,2 for <sup>14</sup>N and <sup>3</sup>He, respectively. When averaged over the relevant $`Q^2`$, the fall of $`\frac{\sigma _\mathrm{A}}{\sigma _\mathrm{D}}`$ below unity is even sharper and deeper than in most of the earlier data. We show here that the unusual behavior as a function of $`Q^2`$, at small $`x`$, can be well represented by a simple, overall form factor, and a normalization change, multiplying the last term in Eq. (1). The modification is
$$\delta (A)\times 16.3\delta (A)F(Q^2)$$
(2)
with
$$F(Q^2)=\frac{4Q^2}{Q^2+1}\frac{1}{Q^2+1}=4\left\{1\frac{1}{Q^2+1}\right\}\frac{1}{Q^2+1},Q^2\mathrm{in}\left(\frac{\mathrm{GeV}}{c}\right)^2$$
Keeping all parameters fixed, the curves which result from Eqs. (1,2) are superimposed on data in Figs. 1,2. Significant aspects of the data both in the valence-quark $`x`$ domain and in the shadowing domain are evident in the curves. At the lowest values of $`x`$, the curves are meant to refer specifically to the new HERMES data. The NMC and E665 data involve a much higher lepton-beam energy, whereas possible explicit dependence upon this energy is suppressed in the phenomenological Eqs. (1,2). The data at these higher energies and at low $`Q^2`$, may be expected to be closer to unity because the kinematic variable $`ฯต(x,Q^2,E)`$ which appears in the ratio of cross sections (as discussed in the next paragraph), is close to unity, thus removing sensitivity<sup>F3</sup> to $`\stackrel{~}{R}^\mathrm{A}(x,Q^2)`$. A possible physical reason for the presence of a form factor which vanishes not only as $`Q^2\mathrm{}`$, but also as $`Q^20`$ (this is the unusual aspect in Eq. (2)), lies in a high degree of coherence associated with the system of exchanged sea partons, and nuclear gluons (which are individually, colored quanta). This occurs when the system at small values of $`x`$ responds to a low-$`Q^2`$ probe. Qualitatively, exchange of colored quanta (or aggregates) initially at small $`x`$, involve spatial dimensions of the order of $`\frac{1}{m_Nx}=`$ many fermis. Thus coherence over a nucleus, and consequently a form factor falling at least like approximately $`1/Q^2`$, may be expected. However, in addition there is the confinement axiom. Starting with a system of nucleons, individual colored entities cannot be present over extended dimensions within this system. This suggests that a probe with $`Q^2=0`$ ($`\stackrel{}{Q}=0`$) โseesโ no color-induced correlation. Then, the remaining shadowing is related to colorless aggregates of quarks which constitute the low-mass vector meson components of the photon. Such behavior is incorporated phenomenologically in a simple way in Eq. (2); with $`F(Q^2)`$ normalized to unity at its maximum, which occurs at $`Q^2=1(\frac{\mathrm{GeV}}{c})^2`$ for this form. Note that averaging $`6.3\delta (A)F(Q^2)`$ over $`Q^2`$ in the interval $`0.3<Q^2<1.5(\frac{\mathrm{GeV}}{c})^2`$, gives an effective, A-dependent coefficient of the last term in Eq. (1) of $`6.1\delta (A)`$ (instead of $`\delta (A)`$); this accounts quantitatively for the increased depth of the sharp drop for $`x<0.06`$ in the HERMES data (their Fig. 1) as compared to the earlier data <sup>F9</sup><sup>F9</sup>F9 Averaging $`6.3\delta (A)F(Q^2)`$ over $`Q^2`$ in the interval $`1<Q^2<100(\frac{\mathrm{GeV}}{c})^2`$ gives an effective, A-dependent coefficient of the last term in Eq. (1) of $`\delta (A)`$, in agreement with the representation of the earlier data. It is worth noting that the fit normalization parameter $`N0.3`$ in the factor $`\delta (A)`$ which multiplies the valence-quark term in Eq. (1), contains a factor of $`\frac{1}{3}`$ corresponding to the a priori probability for involving a single valence quark from each nucleon in the correlation (at any instant). This constraint does not hold for the sea quark normalization as changed in Eq. (3); the fit normalization here is of order unity.. At $`x0.01`$ the earlier depression of $`\frac{\sigma _\mathrm{A}}{\sigma _\mathrm{D}}`$ below unity by about $`6\%`$ for <sup>12</sup>C is increased to about $`35\%`$ for <sup>14</sup>N in the HERMES data!
Within the physical picture embodied in Eq. (1), we consider the interpretation of the new data in terms<sup>F3</sup> of $`\stackrel{~}{R}^\mathrm{A}(x,Q^2)>\stackrel{~}{R}^\mathrm{D}(x,Q^2)`$, as given in the analysis by the HERMES collaboration. The ratio $`\frac{\sigma _\mathrm{A}}{\sigma _\mathrm{D}}`$ is written in terms of the ratio of structure functions $`F_2`$, times a multiplicative factor
$$\frac{\sigma _\mathrm{A}(x,Q^2)}{\sigma _\mathrm{D}(x,Q^2)}=\frac{F_2^\mathrm{A}(x,Q^2)}{F_2^\mathrm{D}(x,Q^2)}\left\{\frac{(1+ฯต\stackrel{~}{R}^\mathrm{A}(x,Q^2))(1+\stackrel{~}{R}^\mathrm{D}(x,Q^2))}{(1+ฯต\stackrel{~}{R}^\mathrm{D}(x,Q^2))(1+\stackrel{~}{R}^\mathrm{A}(x,Q^2))}\right\}$$
(3)
where $`ฯต(x,Q^2,E)`$ is the kinematic variable related to the virtual-photon polarization, defined by HERMES in their Eq. (2), with $`0ฯต1`$; ($`E`$ is the initial-lepton laboratory energy). The factor multiplying $`\frac{F_2^\mathrm{A}}{F_2^\mathrm{D}}`$ is unity independently of the value of $`(\stackrel{~}{R}^\mathrm{A}\stackrel{~}{R}^\mathrm{D})`$, when the kinematic variable $`ฯต`$ becomes unity; this is formally at $`Q^2=0`$. As $`ฯต`$ moves towards zero, which occurs for larger $`Q^2`$ in the HERMES data, the factor becomes less than unity if $`\stackrel{~}{R}^\mathrm{A}(x,Q^2)>\stackrel{~}{R}^\mathrm{D}(x,Q^2)`$. However, if $`(\stackrel{~}{R}^\mathrm{A}\stackrel{~}{R}^\mathrm{D})`$ becomes effectively zero already at moderate values of $`Q^2`$, then the multiplicative factor is again unity, independently of the value of $`ฯต`$. <sup>F10</sup><sup>F10</sup>F10This is in fact, the situation for the values of $`\frac{\stackrel{~}{R}^\mathrm{A}(x,Q^2)}{\stackrel{~}{R}^\mathrm{D}(x,Q^2)}`$ extracted by HERMES under certain assumptions concerning the behavior of $`\frac{F_2^\mathrm{A}(x,Q^2)}{F_2^\mathrm{D}(x,Q^2)}`$ and $`\frac{\stackrel{~}{R}^\mathrm{A}(x,Q^2)}{\stackrel{~}{R}^\mathrm{D}(x,Q^2)}`$. Thus, a strong variation at low $`Q^2`$ which is evident in the HERMES data in Figs. 1,2, then occurs in $`\frac{\stackrel{~}{R}^\mathrm{A}}{\stackrel{~}{R}^\mathrm{D}}`$. Assuming that an interpretation of the new data in terms of $`(\stackrel{~}{R}^\mathrm{A}\stackrel{~}{R}^\mathrm{D})>0`$ is valid, we make an approximate identification of Eq. (3) with the successful phenomenological formula in Eq. (1), modified as in Eq. (2). Then, for $`\stackrel{~}{R}^\mathrm{D}`$ and $`\stackrel{~}{R}^\mathrm{A}`$ less than unity, we obtain an approximate relation in the domain $`x<0.06`$, with $`Q^2<1.5(\frac{\mathrm{GeV}}{c})^2`$, for a function <sup>F11</sup><sup>F11</sup>F11The more complete form of the right-hand side of Eq. (4) resulting from Eqs. (1,2) has two zeros: one at $`Q^20.043(\frac{\mathrm{GeV}}{c})^2`$ which corresponds to $`ฯต1`$; the other at $`Q^223(\frac{\mathrm{GeV}}{c})^2`$ which corresponds to $`\stackrel{~}{R}^\mathrm{A}(x,Q^2)\stackrel{~}{R}^\mathrm{D}(x,Q^2)`$. $`\mathrm{\Delta }(x,Q^2)`$, at given $`E`$
$$\mathrm{\Delta }(x,Q^2)=\left\{1ฯต\right\}(\stackrel{~}{R}^\mathrm{A}\stackrel{~}{R}^\mathrm{D})6.3\delta (A)s(x)\left\{1\frac{1}{Q^2+1}\right\}\left(\frac{1}{Q^2+1}\right)$$
(4)
where
$$s(x)=\frac{1.1(1x)^7\stackrel{~}{C}x^{0.35}\mathrm{e}^{35x^2}}{3.3\sqrt{x}(1x)^3+1.1(1x)^7}$$
The quantity on the right in $`\{\mathrm{}\}`$ has correct limiting behaviors to phenomenologically, approximately represent $`\{1ฯต\}`$ i. e. $`ฯต=1`$ at $`Q^2=0`$, and $`ฯต0`$ for large $`Q^2`$ (at given $`E`$ and $`x`$). The explicit dependence upon $`E`$ and $`x`$ is suppressed in this approximation; the numerical values of $`ฯต`$ are roughly given as in the HERMES data for the lowest $`x`$-bins. At fixed small $`x`$, Eq. (4) suggests for the ratio
$$\left\{\frac{\stackrel{~}{R}^\mathrm{A}(x,Q^2)}{\stackrel{~}{R}^\mathrm{D}(x,Q^2)}1\right\}\delta (A)\left(\frac{1}{Q^2+1}\right)$$
(5)
Eq. (5) is an approximate representation<sup>F11</sup> of the growth at low $`Q^2`$ for the ratio at small $`x`$, which is extracted from the $`Q^2`$ variation of $`\frac{\sigma _\mathrm{A}}{\sigma _\mathrm{D}}`$ by HERMES using some assumptions<sup>F10</sup> (in the upper part of their Fig. 5). It is possible that $`\frac{\stackrel{~}{R}^\mathrm{A}}{\stackrel{~}{R}^\mathrm{D}}`$ falls again as $`Q^2`$ approaches zero. Furthermore it is notable that the extracted increase in the ratio in going from He to <sup>14</sup>N is approximately given by $`\frac{\delta (A=14)}{\delta (A=4)}2`$, in accord with Eq. (5). The interpretation of the data in terms of $`\stackrel{~}{R}^\mathrm{A}>\stackrel{~}{R}^\mathrm{D}`$ is strengthened by the likelihood of an enhancement in the interaction of longitudinal photons with the correlated sea among three nucleons, for the reasons given below. With respect to the $`Q^2`$ domain of the HERMES data at small $`x`$, one might assume that the multiplicative factor in Eq. (3) is unity, and attribute the behavior of their basic data for $`\frac{\sigma _\mathrm{A}(x,Q^2)}{\sigma _\mathrm{D}(x,Q^2)}`$ shown here in Figs. 1,2, to an unusual, strong variation with low $`Q^2`$, at small $`x`$, of $`\frac{F_2^\mathrm{A}(x,Q^2)}{F_2^\mathrm{D}(x,Q^2)}`$, as in a form factor with the general limiting behavior of $`F(Q^2)`$ in Eq. (2). However, then the $`Q^2`$ dependence of the data taken at higher beam energies is not represented at small $`x`$ and moderate $`Q^2`$ (by our simple $`F(Q^2)`$).
In fact, assuming the validity of the interpretation in terms of $`\stackrel{~}{R}^\mathrm{A}>\stackrel{~}{R}^\mathrm{D}`$, also allows for a definite argument that there is new physics involving multinucleon correlations from explicit quark-gluon dynamics in nuclei. Interaction of the virtual photon with a single โisolatedโ quark, largely involves no helicity-flip; then angular momentum conservation requires that the interaction be with a transverse photon (as is most easily seen in the โbrick wallโ system for the photon-quark collision). This constraint does not hold in a multiquark correlation, where the struck quark is not โfreeโ, but rather propagates, emits a gluon, and propagates to another nucleon. An exercise involving only Clebsch-Gordon coefficients indicates that if the spin-orientation of an interacting system of three quarks (nucleons) does not change, then interaction with a longitudinal photon (behaving like a spin-zero quanta) is favored over a transverse photon by an a priori factor of $`\frac{3}{2}`$. There can also be a dynamical enhancement of longitudinal photons in interaction with aggregates of sea. By comparing initial-interaction coefficients of longitudinal and transverse virtual-photon polarization, we estimate an enhancement factor of the order of
$$\left\{\frac{\left((m_Nx)^2+\stackrel{~}{m}_q^2\right)^{1/2}}{(m_Nx)}\stackrel{~}{f}(Q^2)\right\}^22,$$
in the small-$`x`$, low-$`Q^2`$ domain of HERMES, for an effective aggregate mass $`\stackrel{~}{m}_q`$ of $`20`$ MeV ($`\stackrel{~}{f}`$ is a form factor), with $`m_N`$ the nucleon mass (therefore $`m_Nx20`$ MeV for $`x0.02`$).
A final remark concerns the shadowing at $`Q^2=0`$ for real photons, when compared to that for virtual photons with $`Q^2`$ near to zero. If the explicit partonic effect first gives increasing shadowing as $`Q^2`$ moves away from zero, then this effect on top of the conventional geometric effect already present for real photons, would give somewhat more shadowing for the virtual photons. There may be some experimental evidence for this. In addition, since the geometric shadowing effect at a given very small $`x`$ is expected to become less with a marked increase in $`Q^2`$, if the partonic shadowing effect initially increases away from $`Q^2=0`$, then the overall shadowing effect will show a reduced variation with increasing $`Q^2`$.
In summary, when viewed and quantitatively correlated within a unified physical picture, the new data, and all of the earlier improved data from experiments over the last ten years concerning the EMC effect and shadowing, provide the best evidence yet for an explicit dynamical role played by exchanged quarks and nuclear gluons in atomic nuclei.
We thank the HERMES collaboration for information, in particular Prof. M. Dรผren and Dr. G. van der Steenhoven. |
warning/0002/nucl-th0002018.html | ar5iv | text | # The ๐โข๐โข๐ coupling constant
## I Introduction
Nowadays quantum chromodynamics (QCD) is widely believed to be the underlying theory of the strong interaction. Yet the non-abelian nature of the gauge group makes analytical calculation extremely difficult in the low energy sector. A typical example is the various coupling constants of meson nucleon interaction. These couplings are inputs for the one boson exchange potentials for the nuclear forces and the analysis of the important pseudoscalar and vector meson photo- and electro-production experiments currently underway in MAMI (Mainz) and Spring8 (JHF) etc. For the pion nucleon sector there is enough precise data to extract these couplings. Then they are used as inputs to make predictions and analyze other experimental data. In the kaon nucleon hyperon sector the situation is not so encouraging. But there is still some data available. The worst occurs in the $`\eta NN`$ and $`\eta ^{}NN`$ sector, where knowledge of them is rather poor. In the present paper we shall focus on the calculation of $`\eta NN`$ coupling constant.
There were some theoretical papers on this issue. But the results from various approaches differed greatly. With $`SU_f(3)`$ symmetry it was found $`\alpha _{\eta NN}=\frac{g_{\eta NN}^2}{4\pi }=3.68`$ from the analysis of the nucleon nucleon potential . Similar values was obtained in the non-relativistic model . From the analysis of forward nucleon nucleon scattering using the dispersion relation it was found that $`\alpha _{\eta NN}<1`$, consistent to be zero . In the author was able to relate the proton matrix element of flavor singlet current in the large $`N_c`$ limit to the pseudoscalar meson nucleon coupling constants, leading to $`\alpha _{\eta NN}1.3`$. Typical values of $`\alpha _{\eta NN}`$ obtained in fits with one boson exchange potentials range from 3 to 7 since the eta meson does not contribute significantly to the $`NN`$ phase shifts and nuclear binding at normal densities . However this coupling is smaller than 1 and can be neglected in the full Bonn potential . $`\alpha _{\eta NN}`$ extracted from the reaction $`\pi ^{}p\eta n`$ lies between 0.6โ1.7 . An interesting indirect constraint of $`\alpha _{\eta NN}`$ comes from the $`\pi `$-$`\eta `$ mixing amplitude generated by $`\overline{N}N`$ loops and neutron proton mass difference using hadronic models. In order to let this amplitude agree with results from chiral perturbation theory, $`\alpha _{\eta NN}`$ is required to be in the range $`0.32`$-$`0.53`$ . Eta meson photo-production did not fix $`\alpha _{\eta NN}`$ either. In $`\alpha _{\eta NN}`$ was suggested to around $`1.0`$ or $`1.4`$. Yet a recent analysis of more precise eta meson photo-production experiments in Mainz suggested smaller value of $`\alpha _{\eta NN}`$ . In other words, the eta nucleon coupling constant is still very controversial. To derive it within an independent and reliable theoretical framework shall prove valuable. We shall use the now well developed light cone QCD sum rules (LCQSR) technique to calculate $`\alpha _{\eta NN}`$ in this work. Note our approach differs from all the above ones in that it starts microscopically from the QCD Lagrangian.
QCD sum rules (QSR) are successful when applied to the low-lying hadron masses and couplings. In this approach the nonperturbative effects are introduced via various condensates in the vacuum. The light cone QCD sum rule differs from the conventional short-distance QSR in that it is based on the expansion over the twists of the operators. The main contribution comes from the lowest twist operator. Matrix elements of nonlocal operators sandwiched between a hadronic state and the vacuum defines the hadron wave functions. When the LCQSR is used to calculate the coupling constant, the double Borel transformation is always invoked so that the excited states and the continuum contribution can be subtracted quite cleanly. Moreover, the final sum rule depends only on the value of the hadron wave function at the middle point $`u_0=1/2`$ for the diagonal case, which is much better known than the whole wave function . In the present case our sum rules involve with the eta wave function (EWF) $`\phi _\eta (u_0=\frac{1}{2})`$ etc. These parameters are universal in all processes at a given scale.
We have used QCD sum rules to study the meson nucleon strong interactions. In the pion is treated as the external field to analyze the possible isospin symmetry violations of the pion nucleon coupling constant. Later the light cone QSR (LCQSR) was employed to extract the $`\pi NN(1535)`$ coupling constant, which was found to be strongly suppressed . The same formalism was extended to the case of vector meson nucleon interaction . The values of the vector and tensor coupling constants and their ratios of $`\rho NN`$ and $`\omega NN`$ interaction from the LCQSR agree well with the ones from the experimental data and the dispersion relation analysis. With the advent of the eta meson distribution amplitudes up to twist four , we are now able to calculate the $`\eta NN`$ coupling constant with a theoretically well developed formalism. Although $`\eta `$ meson is a Goldstone boson, its mass is not small in the real world and comparable with the typical hadronic scale due to the explicit breaking of $`SU_f(3)`$ flavor symmetry. We have included the eta mass correction in our calculation. Moreover the eta meson is an isoscalar, which leads to the big difference of the LCQSR for the $`\eta NN`$ coupling constant from that for the $`\pi ^0NN`$ coupling. We arrive at $`\alpha _{\eta NN}=\frac{g_{\eta NN}^2}{4\pi }=(0.3\pm 0.15)`$. The numerically small value is due to the cancellation between the leading term and mass correction terms. This point can be seen clearly in later sections.
Our paper is organized as follows: Section I is an introduction. We introduce the two point function for the $`\eta NN`$ vertex and saturate it with nucleon intermediate states in section II. The definitions of the eta wave functions (EWF) are also presented. Numerical analysis and a short summary is given in the last section.
## II The LCQSR for the $`\eta NN`$ coupling
We start with the two point function
$$\mathrm{\Pi }(p_1,p_2,q)=id^4xe^{ipx}0|๐ฏ\eta _p(x)\overline{\eta _p}(0)|\eta (q)$$
(1)
with $`p_1=p`$, $`p_2=pq`$ and the Ioffeโs nucleon interpolating field
$$\eta _p(x)=ฯต_{abc}[u^a(x)๐\gamma _\mu u^b(x)]\gamma _5\gamma ^\mu d^c(x),$$
(2)
$$\overline{\eta }_p(y)=ฯต_{abc}[\overline{u}^b(y)\gamma _\nu C\overline{u}^{aT}(y)]\overline{d}^c(y)\gamma ^\nu \gamma ^5,$$
(3)
where $`a,b,c`$ is the color indices and $`๐=i\gamma _2\gamma _0`$ is the charge conjugation matrix. For the neutron interpolating field, $`ud`$.
$`\mathrm{\Pi }(p_1,p_2,q)`$ has the general form
$$\mathrm{\Pi }(p_1,p_2,q)=F(p_1^2,p_2^2,q^2)\widehat{q}\gamma _5+F_1(p_1^2,p_2^2,q^2)\gamma _5+F_2(p_1^2,p_2^2,q^2)\widehat{p}\gamma _5+F_3(p_1^2,p_2^2,q^2)\sigma _{\mu \nu }\gamma _5p^\mu q^\nu $$
(4)
The sum rules derived from the chiral even tensor structure yield better results than those from the chiral even ones in the QSR analysis of the nucleon mass . We shall focus on the tensor structure $`\widehat{q}\gamma _5`$ and study the function $`F(p_1^2,p_2^2,q^2)`$ as in the QSR analysis of the pion nucleon coupling constant.
The eta nucleon coupling constant $`g_{\eta NN}`$ is defined by the $`\eta N`$ interaction Lagrangian:
$$_{\eta NN}=g_{\eta NN}\overline{N}i\gamma _5\eta N..$$
(5)
At the phenomenological level the eq.(1) can be expressed as:
$$\mathrm{\Pi }(p_1,p_2,q)=i\lambda _N^2m_Ng_{\eta NN}(q^2)\frac{\gamma _5\widehat{q}}{(p_1^2M_N^2)(p_2^2M_N^2)}+\mathrm{}$$
(6)
where we include only the tensor structure $`\gamma _5\widehat{q}`$ only. The ellipse denotes the continuum and the single pole excited states to nucleon transition contribution. $`\lambda _N`$ is the overlapping amplitude of the interpolating current $`\eta _N(x)`$ with the nucleon state
$$0|\eta _N(0)|N(p)=\lambda _Nu_N(p)$$
(7)
Neglecting the four particle component of the eta wave function, the expression for $`F(p_1^2,p_2^2,q^2)`$ with the tensor structure at the quark level reads,
$`i{\displaystyle d^4xe^{ipx}0|T\eta _p(x)\overline{\eta }_p(0)|\eta (q)}=`$ (8)
$`2\mathrm{i}{\displaystyle d^4xe^{ipx}ฯต^{abc}ฯต^{a^{}b^{}c^{}}Tr\{\gamma _\nu CiS_{}^{T}{}_{u}{}^{bb^{}}(x)C\gamma _\mu iS_u^{aa^{}}(x)\}\gamma _5\gamma _\mu 0|d^c(x)\overline{d}^c^{}(0)|\eta (q)\gamma _\nu \gamma _5}`$ (9)
$`+4\mathrm{i}{\displaystyle d^4xe^{ipx}ฯต^{abc}ฯต^{a^{}b^{}c^{}}Tr\{\gamma _\nu CiS_{}^{T}{}_{u}{}^{bb^{}}(x)C\gamma _\mu 0|u^a(x)\overline{u}^a^{}(0)|\eta (q)\}\gamma _5\gamma _\mu iS_d^{cc^{}}(x)\gamma _\nu \gamma _5}`$ (10)
where $`iS(x)`$ is the full light quark propagator with both perturbative term and contribution from vacuum fields .
By the operator expansion on the light-cone the matrix element of the nonlocal operators between the vacuum and eta state defines the two and three particle eta wave function. In order to simplify the notations we use $`\overline{q}\mathrm{\Gamma }_\mu q`$ to denote $`(\overline{u}\mathrm{\Gamma }_\mu u+\overline{d}\mathrm{\Gamma }_\mu d2\overline{s}\mathrm{\Gamma }_\mu s)/\sqrt{6}`$. We also introduce $`F_\eta =\frac{f_\eta }{\sqrt{6}}`$, where $`f_\eta `$ is defined as
$$<0|\overline{q}(0)\gamma _\mu \gamma _5q(0)|\eta (q)>=if_\eta q_\mu .$$
(11)
Up to twist four the Dirac components of this wave function can be written as :
$`<0|\overline{q}(0)\gamma _\mu \gamma _5q(x)|\eta (q)>=if_\eta q_\mu {\displaystyle _0^1}๐ue^{iuqx}[\phi _\eta (u)+{\displaystyle \frac{1}{16}}m_\eta ^2x^2A(u)]`$ (12)
$`+{\displaystyle \frac{i}{2}}f_\eta m_\eta ^2{\displaystyle \frac{q_\mu }{qx}}{\displaystyle _0^1}๐ue^{iuqx}B(u)+O(x^4),`$ (13)
$$<0|\overline{q}(0)i\gamma _5q(x)|0>=f_\eta \mu _\eta _0^1๐ue^{iuqx}\phi _P(u),$$
(14)
$$<0|\overline{q}(0)\sigma _{\mu \nu }\gamma _5q(x)|0>=\frac{i}{6}f_\eta \mu _\eta (q_\mu x_\nu q_\nu x_\mu )_0^1๐ue^{iuqx}\phi _\sigma (u),$$
(15)
$`<0|\overline{q}(0)\sigma _{\alpha \beta }\gamma _5g_sG_{\mu \nu }(ux)q(x)|\eta (q)>=`$ (16)
$`if_\eta \mu _\eta \eta _3[(q_\mu q_\alpha g_{\nu \beta }q_\nu q_\alpha g_{\mu \beta })(q_\mu q_\beta g_{\nu \alpha }q_\nu q_\beta g_{\mu \alpha })]{\displaystyle ๐\alpha _i\phi _{3\eta }(\alpha _i)e^{iqx(\alpha _1+v\alpha _3)}},`$ (17)
$`<0|\overline{q}(0)\gamma _\mu \gamma _5g_sG_{\alpha \beta }(vx)q(x)|\eta (q)>=`$ (18)
$`f_\eta m_\eta ^2\left[q_\beta \left(g_{\alpha \mu }{\displaystyle \frac{x_\alpha q_\mu }{qx}}\right)q_\alpha \left(g_{\beta \mu }{\displaystyle \frac{x_\beta q_\mu }{qx}}\right)\right]{\displaystyle ๐\alpha _i\phi _{}(\alpha _i)e^{iqx(\alpha _1+v\alpha _3)}}`$ (19)
$`+f_\eta m_\eta ^2{\displaystyle \frac{q_\mu }{qx}}(q_\alpha x_\beta q_\beta x_\alpha ){\displaystyle ๐\alpha _i\phi _{}(\alpha _i)e^{iqx(\alpha _1+v\alpha _3)}}`$ (20)
and
$`<0|\overline{q}(0)\gamma _\mu g_s\stackrel{~}{G}_{\alpha \beta }(vx)q(x)|\eta (q)>=`$ (21)
$`if_\eta m_\eta ^2\left[q_\beta \left(g_{\alpha \mu }{\displaystyle \frac{x_\alpha q_\mu }{qx}}\right)q_\alpha \left(g_{\beta \mu }{\displaystyle \frac{x_\beta q_\mu }{qx}}\right)\right]{\displaystyle ๐\alpha _i\stackrel{~}{\phi }_{}(\alpha _i)e^{iqx(\alpha _1+v\alpha _3)}}`$ (22)
$`if_\eta m_\eta ^2{\displaystyle \frac{q_\mu }{qx}}(q_\alpha x_\beta q_\beta x_\alpha ){\displaystyle ๐\alpha _i\stackrel{~}{\phi }_{}(\alpha _i)e^{iqx(\alpha _1+v\alpha _3)}}.`$ (23)
The operator $`\stackrel{~}{G}_{\alpha \beta }`$ is the dual of $`G_{\alpha \beta }`$: $`\stackrel{~}{G}_{\alpha \beta }=\frac{1}{2}ฯต_{\alpha \beta \delta \rho }G^{\delta \rho }`$; $`๐\alpha _i`$ is defined as $`๐\alpha _i=d\alpha _1d\alpha _2d\alpha _3\delta (1\alpha _1\alpha _2\alpha _3)`$. Due to the choice of the gauge $`x^\mu A_\mu (x)=0`$, the path-ordered gauge factor $`P\mathrm{exp}\left(ig_s_0^1๐ux^\mu A_\mu (ux)\right)`$ has been omitted.
The EWF $`\phi _\eta (u)`$ is of twist two , $`\phi _P(u)`$, $`\phi _\sigma (u)`$, and $`\phi _{3\eta }`$ are of twist three, while $`A(u)`$, part of $`B(u)`$ and all the EWFs appearing in eqs.(20), (23) are of twist four. The EWFs $`\phi (x_i,\mu )`$ ($`\mu `$ is the renormalization point) describe the distribution in longitudinal momenta inside the eta meson, the parameters $`x_i`$ ($`_ix_i=1`$) representing the fractions of the longitudinal momentum carried by the quark, the antiquark and gluon.
The normalization and definitions of the various constants can be found in . Some of them are $`_0^1๐u\phi _\eta (u)=_0^1๐u\phi _\sigma (u)=1`$, $`๐\alpha _i\phi _{}(\alpha _i)=๐\alpha _i\phi _{}(\alpha _i)=0`$ etc.
Since the steps to derive LCQSRs are very similar to those in , we present final sum rule directly. Interested readers may consult the above papers for details.
$`m_N\lambda _N^2g_{\eta NN}e^{\frac{M_N^2}{M^2}}=`$ (24)
$`e^{\frac{u_0(1u_0)m_\eta ^2}{M^2}}\{{\displaystyle \frac{F_\eta ^u}{2\pi ^2}}\phi _\eta (u_0)M^6f_2({\displaystyle \frac{s_0}{M^2}})+{\displaystyle \frac{m_\eta ^2}{4\pi ^2}}2u_0[F_\eta ^uA(u_0)+F_\eta ^d\varphi _B(u_0)]M^4f_1({\displaystyle \frac{s_0}{M^2}})`$ (25)
$`{\displaystyle \frac{F_\eta ^u}{9\pi ^2}}a\mu _\eta [\phi _\sigma (u_0)+{\displaystyle \frac{u_0}{2}}\phi _\sigma ^{}(u_0)]M^2f_0({\displaystyle \frac{s_0}{M^2}})`$ (26)
$`+{\displaystyle \frac{1}{12\pi ^2}}F_\eta ^u\mu _\eta \eta _3am_\eta ^2I_1[\phi _{3\eta }]{\displaystyle \frac{1}{6\pi ^2}}F_\eta ^u\mu _\eta \eta _3aI_2[\phi _{3\eta }]M^2f_0({\displaystyle \frac{s_0}{M^2}})`$ (27)
$`+{\displaystyle \frac{1}{2\pi ^2}}m_\eta ^2(F_\eta ^u+F_\eta ^d)\{{\displaystyle \frac{1}{4}}I_2[\phi _{}]{\displaystyle \frac{1}{4}}I_2[\phi _{}]+I_1[\phi _{}]I_4[\stackrel{~}{\phi }_{}]{\displaystyle \frac{1}{4}}I_7[\stackrel{~}{\phi }_{}]\}M^4f_1({\displaystyle \frac{s_0}{M^2}})`$ (28)
$`+{\displaystyle \frac{1}{2\pi ^2}}m_\eta ^4(F_\eta ^u+F_\eta ^d)\{I_3[\phi _{}]I_3[\phi _{}]I_5[\stackrel{~}{\phi }_{}]I_6[\stackrel{~}{\phi }_{}]+I_5[\stackrel{~}{\phi }_{}]+I_6[\stackrel{~}{\phi }_{}]\}M^2f_0({\displaystyle \frac{s_0}{M^2}})\}`$ $`,`$ (29)
where $`f_n(x)=1e^x\underset{k=0}{\overset{n}{}}\frac{x^k}{k!}`$ is the factor used to subtract the continuum, $`s_0`$ is the continuum threshold. $`u_0=\frac{M_1^2}{M_1^2+M_2^2}`$, $`M^2\frac{M_1^2M_2^2}{M_1^2+M_2^2}`$, $`M_1^2`$, $`M_2^2`$ are the Borel parameters, and $`\phi _\sigma ^{}(u_0)=\frac{d\phi _\sigma (u)}{du}|_{u=u_0}`$. In order to make comparison with the sum rule for $`\pi ^0NN`$ coupling constant $`g_{\pi NN}`$, we have labeled the eta meson decay constant $`F_\eta `$ with the flavor index.
The functions $`I_i[\phi _{3\eta }]`$ etc are defined as:
$$\varphi _B(u_0)=_0^{u_0}๐uB(u),$$
(30)
$$I_1[F]=2u_0_0^{u_0}๐\alpha _1_0^{1u_0}๐\alpha _2\frac{F(\alpha _1,\alpha _2,1\alpha _1\alpha _2)}{(1\alpha _1\alpha _2)^2}(12u_0+\alpha _1\alpha _2),$$
(31)
$`I_2[F]={\displaystyle _0^{u_0}}๐\alpha _1{\displaystyle \frac{F(\alpha _1,1u_0,u_0\alpha _1)}{u_0\alpha _1}}+{\displaystyle _0^{1u_0}}๐\alpha _2{\displaystyle \frac{F(u_0,\alpha _2,1u_0\alpha _2)}{1u_0\alpha _2}}`$ (32)
$`2{\displaystyle _0^{u_0}}๐\alpha _1{\displaystyle _0^{1u_0}}๐\alpha _2{\displaystyle \frac{F(\alpha _1,\alpha _2,1\alpha _1\alpha _2)}{(1\alpha _1\alpha _2)^2}},`$ (33)
$$I_3[F]=2u_0_0^{u_0}๐\alpha _1_0^{1u_0}๐\alpha _2F(\alpha _1,\alpha _2,1\alpha _1\alpha _2)\frac{(u_0\alpha _1)(1u_0\alpha _2)}{(1\alpha _1\alpha _2)^2},$$
(34)
$$I_4[F]=2u_0_0^{u_0}๐\alpha _1_0^{1u_0}๐\alpha _2\frac{F(\alpha _1,\alpha _2,1\alpha _1\alpha _2)}{1\alpha _1\alpha _2},$$
(35)
$$I_5[F]=2u_0_0^{u_0}๐\alpha _1_0^{1u_0}๐\alpha _2F(\alpha _1,\alpha _2,1\alpha _1\alpha _2)\frac{u_0\alpha _1}{1\alpha _1\alpha _2},$$
(36)
$$I_6[F]=2u_0_0^{u_0}๐\alpha _1_0^{u_0\alpha _1}๐\alpha _3F(\alpha _1,1\alpha _1\alpha _3,\alpha _3),$$
(37)
$$I_7[F]=2u_0\{_0^{u_0}๐\alpha _1\frac{F(\alpha _1,1u_0,u_0\alpha _1)}{u_0\alpha _1}_0^{1u_0}๐\alpha _2\frac{F(u_0,\alpha _2,1u_0\alpha _2)}{1u_0\alpha _2}\},$$
(38)
where $`F=\phi _{3\eta },\phi _{},\phi _{},\stackrel{~}{\phi }_{},\stackrel{~}{\phi }_{}`$.
## III Discussion
Since eta meson is an isoscalar, we have $`F_\eta ^u=F_\eta ^d=F_\eta =\frac{f_\eta }{\sqrt{6}}`$. Replacing the $`\eta `$ index with $`\pi `$ and $`F_\eta `$ by $`f_\pi `$ in (24), we recover the sum rule for $`g_{\pi NN}`$ . Note $`f_\pi ^u=f_\pi ^d=f_\pi `$. In other words, the twist four terms involved with three particle pion wave functions vanish due to isospin symmetry. The first term in (24) is the leading twist two term. The third term is of twist three and of the same sign as the leading term. The second term comes from two particle EWF and the remaining terms all come from three particle EWFs. Although they are of twist four except the fourth term, their contribution is greatly enhanced by the factor $`m_\eta ^2`$ in contrast with $`m_\pi ^2`$ in the $`\pi NN`$ coupling case. Moreover they are of the opposite sign as the leading twist two and three terms, which leads to strong cancellation. In other words, large mass and isoscalar structure of eta meson causes $`g_{\eta NN}`$ to be much smaller than $`g_{\pi NN}`$.
The sum rule (24) is symmetric and diagonal, which requires the Borel parameters $`M_1^2=M_2^2`$, i.e, $`u_0=\frac{1}{2}`$. The working interval for analyzing the QCD sum rule (24) is $`0.9\text{GeV}^2M_B^21.8\text{GeV}^2`$, a standard choice for analyzing the various QCD sum rules associated with the nucleon. In order to diminish the uncertainty due to $`\lambda _N`$, we shall divide (24) by the Ioffeโs mass sum rule for the nucleon:
$$32\eta ^4\lambda _N^2e^{\frac{M_N^2}{M^2}}=M^6f_2(\frac{s_0}{M^2})+\frac{b}{4}M^2f_0(\frac{s_0}{M^2})+\frac{4}{3}a^2\frac{a^2m_0^2}{3M^2}.$$
(39)
The various parameters which we adopt are $`f_\eta =(0.133\pm 0.01)`$ GeV , $`\eta _3=0.013`$, $`a=4\pi ^2<0|\overline{q}q|0>=0.67\text{GeV}^3`$, $`\mu _\eta =2.13`$GeV at the scale $`\mu =1`$GeV, $`s_0=2.25`$GeV<sup>2</sup>, $`m_N=0.938`$GeV, $`\lambda _N=0.026`$GeV<sup>3</sup> .
At $`u_0=\frac{1}{2}`$ the values of various eta meson wave functions are: $`\phi _\eta (u_0)=1.05`$, $`A(u_0)=4.14`$, $`\varphi _B(u_0)=0`$, $`\phi _\sigma (u_0)=1.44`$, $`\phi _\sigma ^{}(u_0)=0`$, $`I_1[\phi _{3\eta }]=0`$, $`I_1[\phi _{}]=0.026`$, $`I_2[\phi _{3\eta }]=0.9375`$, $`I_2[\phi _{}]=0`$, $`I_2[\phi _{}]=0`$, $`I_3[\phi _{}]=0`$, $`I_3[\phi _{}]=0`$, $`I_4[\stackrel{~}{\phi }_{}]=0.313`$, $`I_5[\stackrel{~}{\phi }_{}]=0.032`$, $`I_5[\stackrel{~}{\phi }_{}]=0.0396`$, $`I_6[\stackrel{~}{\phi }_{}]=0.052`$, $`I_6[\stackrel{~}{\phi }_{}]=0.044`$, $`I_7[\stackrel{~}{\phi }_{}]=0`$ at $`u_0=\frac{1}{2}`$ and $`\mu =1`$GeV.
The dependence on the Borel parameter $`M^2`$ of $`g_{\eta NN}`$ are shown in FIG 1 with $`s_0=2.35,2.25,2.15`$ GeV<sup>2</sup>. The final sum rule is stable in the working region of the Borel parameter $`M^2`$. We obtain:
$$g_{\eta NN}=(1.7\pm 0.3).$$
(40)
In the above numerical analysis we have used relatively large quark condensate value $`<\overline{q}q>=(240\pm 10)^3`$MeV<sup>3</sup>, which corresponds to $`a=0.67`$GeV<sup>3</sup>. In the literatures another value $`<\overline{q}q>=(225\pm 10)^3`$MeV<sup>3</sup> and $`a=0.55`$GeV<sup>3</sup> is also used. Since we are not able to know very precisely the quark condensate value, we also present the variation of $`g_{\eta NN}`$ with $`M^2,s_0`$ with $`a=0.55`$GeV<sup>3</sup> in FIG 2. In this case we have,
$$g_{\eta NN}=(2.1\pm 0.3).$$
(41)
We have included the uncertainty due to the variation of the continuum threshold and the Borel parameter $`M^2`$ in (40) and (41). In other words, only the errors arising from numerical analysis of the sum rule (24) are considered. Other sources of uncertainty include: (1) the truncation of OPE on the light cone at the twist four operators. For example the four particle component of EWF is discarded explicitly; (2) the EWFs are estimated with QCD sum rule, which also induces some errors; (3) the continuum model used in the subtraction of contribution from the higher resonances and continuum spectrum; (4) errors in $`f_\eta `$ etc.
With all these uncertainties we arrive at
$$\alpha _{\eta NN}=(0.3\pm 0.15).$$
(42)
For the $`\eta ,\eta ^{}`$ sector instanton effects might be important. Itโs well known that a large part of $`\eta ^{}`$ mass comes from the $`U_A(1)`$ anomaly. Through $`\eta \eta ^{}`$ mixing instantons also affect eta meson mass and decay constant $`f_\eta `$. Fortunately we know from phenomenological analysis that the $`\eta \eta ^{}`$ mixing angle is about $`20`$ degrees . So for eta meson such effects may be not so large as in the $`\eta ^{}`$ channel. Direct instantons favor strongly the scalar and pseudoscalar channel and might affect the mass sum rules for the mesons in these channels. In our QCD sum rule analysis of eta NN coupling constant we have chosen the tensor structure $`\widehat{q}\gamma _5`$. Moreover we have used the experimental values for $`m_\eta ,f_\eta `$ as inputs instead of invoking the eta meson mass sum rules to extract them. Hence the possible correction from instantons is expected to be relatively small.
In short summary we have calculated the eta nucleon coupling constant with the light cone QCD sum rules. The continuum and the excited states contribution is subtracted rather cleanly through the double Borel transformation. Our approach differs from all the available methods in the extraction of $`g_{\eta NN}`$ and starts from the quark gluon level. So it is independent and more reliable to some extent. Our result of $`\alpha _{\eta NN}`$ favors the small value. Except the nonrelativistic quark model and fits with one boson exchange potentials, other approaches tend to yield small values for $`\alpha _{\eta NN}`$. However in such potentials the eta meson was treated as some effective degree of freedom to model other multi-meson correlations. Hence the eta meson in these potentials can not be related to the real eta meson seen in the photo- or electro-production experiments in a simple way. In other words, the $`\alpha _{\eta NN}`$ in these potentials may be not the same quantity as the coupling we have calculated. We hope our extraction of $`\alpha _{\eta NN}`$ can be used to analyze future eta meson photo- and electro-production experiments.
Figure Captions
FIG 1. The sum rule for $`g_{\eta NN}`$ as a function of the Borel parameter $`M^2`$ with $`a=0.67`$GeV<sup>3</sup> and the continuum threshold $`s_0=2.35,2.25,2.15`$GeV<sup>2</sup>.
FIG 2. The same notations as in FIG 1 except $`a=0.55`$GeV<sup>3</sup>. |
warning/0002/hep-th0002014.html | ar5iv | text | # References
SKYRMIONS FROM SU(3) HARMONIC MAPS AND THEIR QUANTIZATION
V.B. Kopeliovich<sup>โ โ</sup>, B.E. Stern and W.J.Zakrzewski,
Institute for Nuclear Research of the Russian Academy of Sciences,
Moscow 117312, Russia
Department of Mathematical Sciences, University of Durham,
Durham DH1 3LE, UK
Static properties of $`SU(3)`$ multiskyrmions with baryon numbers up to $`6`$ are estimated. The calculations are based on the recently suggested generalization of the $`SU(2)`$ rational map ansรคtze applied to the $`SU(3)`$ model. Both $`SU(2)`$ embedded skyrmions and genuine $`SU(3)`$ solutions are considered and it is shown that although, at the classical level, the energy of the embeddings is lower, the quantum corrections can alter these conclusions. This correction to the energy of the lowest state, depending on the Wess-Zumino term, is presented in the most general case. 1. Topological soliton models, and the Skyrme model among them , have recently generated a fair amount of interest because they may be able to describe various properties of low energy baryons. So far, most of such studies have involved the $`SU(2)`$ Skyrme model, whose solutions were then embedded into the $`SU(3)`$ model. This is justified as the solutions of the $`SU(2)`$ model are also solutions of $`SU(N)`$ models. However, there exist solutions of the $`SU(N)`$ models which are not embeddings of $`SU(2)`$ fields and it is important to assess their contribution.
As with solutions of the $`SU(2)`$ Skyrme models, the solutions of the $`SU(3)`$ model, with very few exceptions, can only be determined numerically. Like for the $`SU(2)`$ case one starts with a harmonic map ansatz which gives fields โcloseโ to the genuine solutions . Then one can use these fields as starting configurations of various numerical minimization schemes. For the $`SU(2)`$ fields such an approach was carried out in where it reproduced the results of the earlier numerical approaches and so showed that the harmonic map approximations are very close to the final fields. The real solutions have energies only few hundreds of $`MeV`$ lower than the harmonic approximations and the baryonic charge and energy distributions do not look very different etc. This has justified the use of harmonic field approximants in calculating various quantities when estimating quantum corrections to the classical results.
In this paper we take a look at the nonembedding solutions of the $`SU(3)`$ model. First we recall the results of and use them in a $`SU(3)`$ numerical minimization program to estimate the values of energies of true solutions. The $`SU(3)`$ variational minimization program has been rearranged for this purpose to allow the consideration of quite general field ansรคtze. This is discussed in the next section.
In the following section we discuss various quantum corrections and compare our results with the similar results for the embeddings. The paper ends with a short section presenting our conclusions.
2. The Lagrangian of the $`SU(3)`$ Skyrme model, in its well known form, depends on the parameters $`F_\pi `$ and $`e`$ and is given by :
$$=\frac{F_\pi ^2}{16}Trl_\mu l^\mu +\frac{1}{32e^2}Tr[l_\mu ,l_\nu ]^2+\frac{F_\pi ^2m_\pi ^2}{16}Tr\left(U+U^{}2\right)$$
$`\left(1\right)`$
Here $`USU(3)`$ is a unitary matrix describing the chiral (meson) fields, and $`l_\mu =U^{}_\mu U`$. In the model $`F_\pi `$ should be fixed at the physical value: $`F_\pi `$ = $`186`$ Mev . The flavour symmetry breaking $`(FSB)`$ in the Lagrangian will be considered below.
The Wess-Zumino term, which should to be added to the action, is given by
$$S^{WZ}=\frac{iN_c}{240\pi ^2}_\mathrm{\Omega }d^5xฯต^{\mu \nu \lambda \rho \sigma }Tr\left(\stackrel{~}{l}_\mu \stackrel{~}{l}_\nu \stackrel{~}{l}_\lambda \stackrel{~}{l}_\rho \stackrel{~}{l}_\sigma \right),$$
$`\left(2\right)`$
where $`\mathrm{\Omega }`$ is a 5-dimensional region with the 4-dimensional space-time as its boundary and where $`\stackrel{~}{l}_\mu `$ is a 5-dimensional analogue of $`l_\mu =U^{}_\mu U`$. As is well known, this extra term does not contribute to the static masses of classical configurations, but it defines important topological properties of skyrmions and plays an important role in their quantization .
In an ansatz was presented which allows us to find approximate solutions of the $`SU(3)`$ model. This ansatz involves parametrising the static field as
$$U\left(\stackrel{}{x}\right)=\mathrm{exp}\left\{f\left(r\right)\left(P(\theta ,\phi )\frac{1}{3}\right)\right\},$$
$`\left(3\right)`$
where $`r,\theta `$ and $`\phi `$ are polar coordinates, $`f(r)`$ is a radial profile function which has to be determined numerically and $`P(\theta ,\phi )`$ is a projector involved in the harmonic map ansatz. As shown in this projector is given by
$$P=\frac{FF^{}}{\left|F\right|^2},$$
$`\left(4\right)`$
where $`F`$ is a 3-component vector, whose entries are polynomials in $`z=tan(\frac{\theta }{2})e^{i\phi }.`$ The largest degree of the polynomial gives the baryon number $`B`$ of the final $`SU(3)`$ Skyrme field configuration. For an $`SU(2)`$ embedding the approach is similar except that this time we put
$$U\left(\stackrel{}{x}\right)=\left(\begin{array}{cc}U_2& \stackrel{}{0}\\ \stackrel{}{0}& 1\end{array}\right),$$
$`\left(5\right)`$
where $`U_2`$ is an $`U(2)`$ matrix determined in an analogous way as $`(3)`$ except that the vector $`F`$ has only two components.
The fields $`F`$ in Eq. $`(4)`$ are so chosen that the configuration $`(3)`$ has the smallest energy; this provides us with the harmonic map approximation to the real minimal energy static solution of the Skyrme model. We have then taken the expressions for $`F`$ determined in and used the corresponding $`U`$ in $`(3)`$ as an initial field in our minimization program.
For $`B=2`$ $`F=(1,\sqrt{2}z,z^2)^T`$, for $`B=3`$ $`F=(1/\sqrt{2},1.576z,z^3)^T`$, for $`B=4`$ $`F=(1,2.72z^2,z^4)^T`$,
for $`B=5`$ $`F=(4.5z^2,2z^4+1,z^52.7z)^T`$ and for $`B=6`$ $`F=(kz^3,13z^5,z^6+3z)^T`$ with $`k=7.06`$ .
We have performed many minimizations of the field configurations involving baryon numbers $`3,4,5`$ and $`6`$. In all cases the initial $`SU(3)`$ harmonic map fields had energies somewhat higher than the corresponding harmonic map embeddings.
Our $`SU(3)`$ minimization program, which uses the parametrization of the $`SU(3)`$ field in terms of two, mutually orthogonal, complex unit vectors<sup>1</sup><sup>1</sup>1The authors thank W.K. Baskerville for making them aware of this parametrisation was then used to minimize the energy further. The constraints of the orthogonality and of the vectors being of unit length were replaced by the introduction of extra positive contributions to the energy (with large coefficients) \- the so called โpenalty termsโ.
The program itself minimized the energy using a mixture of a finite element and finite step methods, and we varied the coefficients of the penalty terms so as not to be trapped in a local minimum. In each case the final energy of the field configuration was lower by a few hundreds of $`MeV`$ from the energy of the harmonic approximant. However, due to the smallness of the lattice the baryon number was also somewhat lower; so our values in the Table were obtained by a linear extrapolation: $`M=M+\mathrm{\Delta }M`$, where $`\mathrm{\Delta }M=\alpha (BB_{obs})`$ and where $`B`$ is the baryon number and $`B_{obs}`$ is its value on the lattice. $`\alpha `$ was determined by looking at various minimizations (with different values of the coefficients of the penalty terms).
To check the stability of the program we also performed a series of minimizations when the initial field was a mixture of the $`SU(3)`$ harmonic map field and of the embedding. When the mixture was close to the embedding or to the $`SU(3)`$ harmonic map the minimization program took it down to these fields. When the mixture was far from either of these special fields it evolved to a new configuration of higher energy than either of the special ones. Although some of these new configurations can be numerical artifacts it is clear that the spectrum of static solutions of the $`SU(3)`$ model is very complicated, with most states having energies larger than the energies of the embeddings or the fields derived by the harmonic map ansatz. Hence our estimates in the Table have some physical justification.
In the Table we present some of our results. It is difficult to assess their accuracy; it is probably within $`1\%`$ for the masses and several $`\%`$ for other quantities. In all numerical minimizations we have worked on a small tetrahedral lattice and hence we had a small โleakageโ of both the energy and of the baryon number.
| $`B`$ | $`M_{cl}`$ | $`M_{cl}^{SU_2}`$ | $`\mathrm{\Delta }E_{M.t.}`$ | $`C_S`$ | $`\mathrm{\Theta }_{1,2}`$ | $`\mathrm{\Theta }_3`$ | $`\mathrm{\Theta }_4`$ | $`\mathrm{\Theta }_5`$ | $`\mathrm{\Theta }_{6,7}`$ | $`\mathrm{\Theta }_8`$ | $`\mathrm{\Theta }_{38}`$ | $`WZ_3`$ | $`WZ_8`$ | $`\mathrm{\Delta }E`$ |
| --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- |
| 1 | $`1.70`$ | $`1.70`$ | $`46`$ | | $`5.56`$ | $`5.56`$ | $`2.04`$ | $`2.04`$ | $`2.04`$ | $``$ | $``$ | - | - | $`368`$ |
| 2 | $`3.38`$ | $`3.26`$ | $`89`$ | $`0.33`$ | $`5.70`$ | $`6.40`$ | $`7.10`$ | $`7.10`$ | $`5.70`$ | $`5.0`$ | $`1.2`$ | $`0.00`$ | $`0.00`$ | $`0.0`$ |
| 3 | $`4.85`$ | $`4.80`$ | $`110`$ | $`0.33`$ | $`7.15`$ | $`7.90`$ | $`8.60`$ | $`8.60`$ | $`9.20`$ | $`5.1`$ | $`1.9`$ | $`0.7`$ | $`0.15`$ | $`35`$ |
| 4 | $`6.48`$ | $`6.20`$ | $`174`$ | $`0.31`$ | $`12.2`$ | $`9.80`$ | $`11.6`$ | $`11.6`$ | $`12.1`$ | $`6.2`$ | $`3.0`$ | $`1.0`$ | $`0.6`$ | $`80`$ |
| 5 | $`7.90`$ | $`7.78`$ | $`211`$ | $`0.36`$ | $`11.5`$ | $`11.6`$ | $`15.4`$ | $`12.1`$ | $`12.7`$ | $`15.2`$ | $`0.1`$ | $`0.2`$ | $`0.8`$ | $`23`$ |
| 6 | $`9.38`$ | $`9.24`$ | $`251`$ | $`0.33`$ | $`15.4`$ | $`14.8`$ | $`15.5`$ | $`15.5`$ | $`14.2`$ | $`16.3`$ | $`1.2`$ | $`0.01`$ | $`0.03`$ | $`0.0`$ |
Table. The values of the masses $`M_{cl}`$ in $`GeV`$, the mass term $`\mathrm{\Delta }E_{M.t.}`$ (in $`MeV`$), the strangeness content $`C_S`$ and the moments of inertia (in $`GeV^1`$) for the $`SU(3)`$ projector configurations. The quantities for $`B=1`$ hedgehog and the masses of $`SU(2)`$ embeddings $`M_{cl}^{SU_2}`$, in $`GeV`$ are given for comparison. The components of the Wess-Zumino term $`WZ_3`$ and $`WZ_8`$ which are different from zero for the $`SU(3)`$ projector ansatz are also shown. $`\mathrm{\Delta }E`$ in the last column is the quantum correction (in $`MeV`$) due to zero modes for the lowest energy state. The parameters of the model are $`F_\pi =186MeV,e=4.12`$.
As can be seen from this Table, the moments of inertia are pairwise equal, except for the case $`B=5`$ when $`\mathrm{\Theta }_4\mathrm{\Theta }_5`$. These equalities are a consequence of the symmetry properties of our multiskyrmion configurations.
It should be noted that the energy of the $`B=4`$ configuration is very close to that of the $`B=4`$ toroidal soliton obtained in $`1988`$ within the axially symmetrical generalization of the ansatz by Balachandran et al. . Within the accuracy of our calculations we cannot be certain which configuration has the lowest energy in the chirally symmetrical limit. For $`B=6`$ the torus-like $`SO(3)`$ configuration has energy considerably higher than the energy of the $`SU(3)`$ projector ansatz presented in the Table.
Let us add here that the configuration of a double-torus form was considered also within the $`SU(2)`$-version of the model several years ago . It has the energy somewhat higher than the energy of the known minimal energy configuration.
3. Following we consider the contribution of the Wess-Zumino $`(WZ)`$ term which defines the quantum numbers of the system in the quantization procedure. Its expression is given by $`(2)`$ above. As usually, we introduce the time-dependent collective coordinates for the quantization of zero modes according to the relation: $`U(\stackrel{}{r},t)=A(t)U_0(\stackrel{}{r})A^{}(t)`$. Next we perform some integration by parts and rewrite the expression for the WZ-term as:
$$L^{WZ}=\frac{iN_c}{48\pi ^2}ฯต_{\alpha \beta \gamma }TrA^{}\dot{A}\left(R_\alpha R_\beta R_\gamma +L_\alpha L_\beta L_\gamma \right)d^3x,$$
$`\left(6\right)`$
where $`L_\alpha =U_0^{}_\alpha U_0=iL_{k,\alpha }\lambda _k`$ and $`R_\alpha =_\alpha U_0U_0^{}=U_0L_\alpha U_0^{}`$, or
$$L^{WZ}=\frac{N_c}{24\pi ^2}\underset{k=1}{\overset{k=8}{}}\omega _kWZ_kd^3x=\underset{k=1}{\overset{k=8}{}}\omega _kL_k^{WZ},$$
$`\left(7\right)`$
with the angular velocities of rotation in the configuration space defined in the usual way:
$`A^{}\dot{A}=\frac{i}{2}\omega _k\lambda _k`$. Summation over repeated indices is assumed here and below. The functions $`WZ_k`$ can be expressed through the chiral derivatives $`\stackrel{}{L}_k`$:
$$WZ_i=WZ_i^R+WZ_i^L=\left(R_{ik}\left(U_0\right)+\delta _{ik}\right)WZ_k^L,$$
$`\left(8a\right)`$
$`i,k=1,\mathrm{}8`$, and are given by
$$WZ_1^L=(L_1,L_4L_5+L_6L_7)\left(L_2L_3L_8\right)/\sqrt{3}2(L_8,L_4L_7L_5L_6)/\sqrt{3}$$
$$WZ_2^L=(L_2,L_4L_5+L_6L_7)\left(L_3L_1L_8\right)/\sqrt{3}2(L_8,L_4L_6+L_5L_7)/\sqrt{3}$$
$$WZ_3^L=(L_3,L_4L_5+L_6L_7)\left(L_1L_2L_8\right)/\sqrt{3}2(L_8,L_4L_5L_6L_7)/\sqrt{3}$$
$$WZ_4^L=(L_4,L_1L_2L_6L_7)\left(L_3L_5L_8\right)/\sqrt{3}+2(\stackrel{~}{L}_8,L_1L_7+L_2L_6)/\sqrt{3}$$
$$WZ_5^L=(L_5,L_1L_2L_6L_7)+\left(L_3L_4L_8\right)/\sqrt{3}2(\stackrel{~}{L}_8,L_1L_6L_2L_7)/\sqrt{3}$$
$$WZ_6^L=(L_6,L_1L_2+L_4L_5)+\left(L_3L_7L_8\right)/\sqrt{3}2(\stackrel{~}{\stackrel{~}{L}}_8,L_1L_5L_2L_4)/\sqrt{3}$$
$$WZ_7^L=(L_7,L_1L_2+L_4L_5)\left(L_3L_6L_8\right)/\sqrt{3}+2(\stackrel{~}{\stackrel{~}{L}}_8,L_1L_4+L_2L_5)/\sqrt{3}$$
$$WZ_8^L=\sqrt{3}\left(L_1L_2L_3\right)+\left(L_8L_4L_5\right)+\left(L_8L_6L_7\right),$$
$`\left(9\right)`$
where $`(L_1,L_2L_3)`$ denotes the mixed product of vectors $`\stackrel{}{L}_1`$, $`\stackrel{}{L}_2`$, $`\stackrel{}{L}_3`$, ie $`(L_1,L_2L_3)=(\stackrel{}{L}_1\stackrel{}{L}_2\stackrel{}{L}_3)`$ and $`\stackrel{~}{L}_3=(L_3+\sqrt{3}L_8)/2`$, $`\stackrel{~}{L}_8=(\sqrt{3}L_3L_8)/2`$, $`\stackrel{~}{\stackrel{~}{L}}_3=(L_3+\sqrt{3}L_8)/2`$, $`\stackrel{~}{\stackrel{~}{L}}_8=(\sqrt{3}L_3+L_8)/2`$ are the third and eighth components of the chiral derivatives in the $`(u,s)`$ and $`(d,s)`$ $`SU(2)`$-subgroups. $`R_{ik}(U_0)=\frac{1}{2}Tr\lambda _iU_0\lambda _kU_0^{}`$ is a real orthogonal matrix, and $`WZ_i^R`$ are defined by the expressions $`(9)`$ with the substitution $`\stackrel{}{L}_k\stackrel{}{R}_k`$. Relations similar to $`(9)`$ can be obtained for $`\stackrel{~}{WZ}_3`$ and $`\stackrel{~}{WZ}_8`$; they are analogs of $`WZ_3`$ and $`WZ_8`$ for the $`(u,s)`$ or $`(d,s)`$ $`SU(2)`$ subgroups, thus clarifying the symmetry of the $`WZ`$-term in the different $`SU(2)`$ subgroups of $`SU(3)`$.
The baryon number of the $`SU(3)`$ skyrmions can be written also in terms of $`\stackrel{}{L_i}`$ in a form where its symmetry in the different $`SU(2)`$ subgroups of $`SU(3)`$ is more explicit:
$$B=\frac{1}{2\pi ^2}\left((L_1,L_2L_3)+(L_4,L_5\stackrel{~}{L_3})+(L_6,L_7\stackrel{~}{\stackrel{~}{L_3}})+\frac{1}{2}\left[(L_1,L_4L_7L_5L_6)+(L_2,L_4L_6+L_5L_7)\right]\right)d^3r.$$
$`\left(10\right)`$
The contributions of the three $`SU(2)`$ subgroups enter the baryon number on an equal footing. In addition, mixed terms corresponding to the contribution of the chiral fields from different subgroups are also present.
The results of calculating the $`WZ`$-term according to $`(9)`$ depend on the orientation of the soliton in the $`SU(3)`$ configuration space. The Guadagniniโs quantization condition was generalized in for configurations of the โmolecularโ type to
$$Y_R^{min}=\frac{2}{\sqrt{3}}\frac{L^{WZ}}{\omega _8}\frac{N_cB\left(13C_S\right)}{3},$$
$`\left(11\right)`$
where $`Y_R`$ is the so-called right hypercharge characterizing the $`SU(3)`$ irrep under consideration, and the scalar strangeness content $`C_S`$ is defined in terms of the real parts of the diagonal matrix elements of the matrix $`U`$:
$$C_S=\frac{<1ReU_{33}>}{<3Re\left(U_{11}+U_{22}+U_{33}\right)>},$$
$`\left(12\right)`$
where $`<>`$ denotes the averaging or integration over the whole 3-dimensional space. When solitons are located in the $`(u,d)`$ $`SU(2)`$ subgroup of $`SU(3)`$ only $`\stackrel{}{L}_1`$, $`\stackrel{}{L}_2`$ and $`\stackrel{}{L}_3`$ are different from zero, $`C_S=0`$, $`WZ_8^R`$ and $`WZ_8^L`$ are both proportional to the $`B`$-number density, and the well known quantization condition takes place
$$Y_R=N_cB/3.$$
$`\left(13\right)`$
The interpolation $`(11)`$ does not work so well for configurations we consider here.
The expression for the rotation energy density of the system depending on the angular velocities of rotations in the $`SU(3)`$ collective coordinate space defined in Section 2 can be written in the following compact form :
$$L_{rot}=\frac{F_\pi ^2}{32}(\stackrel{~}{\omega }_1^2+\stackrel{~}{\omega }_2^2+..+\stackrel{~}{\omega }_8^2)+$$
$$+\frac{1}{16e^2}\{(\stackrel{}{s}_{12}+\stackrel{}{s}_{45})^2+(\stackrel{}{s}_{45}+\stackrel{}{s}_{67})^2+(\stackrel{}{s}_{67}\stackrel{}{s}_{12})^2+\frac{1}{2}((2\stackrel{}{s}_{13}\stackrel{}{s}_{46}\stackrel{}{s}_{57})^2+(2\stackrel{}{s}_{23}+\stackrel{}{s}_{47}\stackrel{}{s}_{56})^2+$$
$$+(2\stackrel{~}{\stackrel{}{s}_{34}}+\stackrel{}{s}_{16}\stackrel{}{s}_{27})^2+(2\stackrel{~}{\stackrel{}{s}_{35}}+\stackrel{}{s}_{17}+\stackrel{}{s}_{26})^2+(2\stackrel{~}{\stackrel{}{s}_{36}}+\stackrel{}{s}_{14}+\stackrel{}{s}_{25})^2+(2\stackrel{~}{\stackrel{}{s}_{37}}+\stackrel{}{s}_{15}\stackrel{}{s}_{24})^2)\},$$
$`\left(14a\right)`$
or
$$L_{rot}=V_{ik}\stackrel{~}{\omega }_i\stackrel{~}{\omega }_k/2=V_{ik}g_{li}g_{mk}\omega _l\omega _m/2$$
$`\left(14b\right)`$
Here $`\stackrel{}{s}_{ik}=\stackrel{~}{\omega }_i\stackrel{}{L}_k\stackrel{~}{\omega }_k\stackrel{}{L}_i`$, $`i,k=1,2\mathrm{}8`$ are the $`SU(3)`$ indices, and $`\stackrel{~}{\stackrel{}{s}_{34}}=(\stackrel{}{s}_{34}+\sqrt{3}\stackrel{}{s}_{84})/2`$, $`\stackrel{~}{\stackrel{}{s}_{35}}=(\stackrel{}{s}_{35}+\sqrt{3}\stackrel{}{s}_{85})/2`$, $`\stackrel{~}{\stackrel{}{s}_{36}}=(\stackrel{}{s}_{36}+\sqrt{3}\stackrel{}{s}_{86})/2`$, $`\stackrel{~}{\stackrel{}{s}_{37}}=(\stackrel{}{s}_{37}+\sqrt{3}\stackrel{}{s}_{87})/2`$, ie similarly to the definitions of $`\stackrel{~}{L}_3`$ and $`\stackrel{~}{L}_8`$. $`g_{li}`$ are given in $`(15)`$. To get $`(14)`$ we have used the identity: $`\stackrel{}{s}_{ab}\stackrel{}{s}_{cd}\stackrel{}{s}_{ad}\stackrel{}{s}_{cb}=\stackrel{}{s}_{ac}\stackrel{}{s}_{bd}`$. The formula $`(14)`$ possesses remarkable symmetry relative to the different $`SU(2)`$ subgroups of $`SU(3)`$. The functions $`L_8`$ or $`\stackrel{~}{L}_8`$ do not enter $`(14)`$ nor the expression $`(10)`$ for the baryon number density. The functions $`\stackrel{~}{\omega }_i`$ are connected to the body-fixed angular velocities of $`SU(3)`$ rotations by the transformations $`\widehat{\stackrel{~}{\omega }}=U_0^{}\widehat{\omega }U_0\widehat{\omega },`$ or, equivalently
$$\stackrel{~}{\omega }_i=\left(R_{ik}\left(U_0^{}\right)\delta _{ik}\right)\omega _k=g_{ki}\omega _k.$$
$`\left(15\right)`$
Here $`R_{ik}(V^{})=R_{ki}(V)`$ is a real orthogonal matrix, $`i,k=1,\mathrm{}8`$, and $`\stackrel{~}{\omega }_i^2=2(\omega _i^2R_{kl}(U_0)\omega _k\omega _l)`$.
The expression for the static energy can be obtained from $`(14)`$ by means of the substitution $`\stackrel{~}{\omega }_i2L_i`$ and $`\stackrel{}{s}_{ik}2\stackrel{}{n_{ik}}`$, with $`\stackrel{}{n}_{ik}`$ being the cross product of $`\stackrel{}{L}_i`$ an $`\stackrel{}{L}_k`$ ie $`\stackrel{}{n}_{ik}=\stackrel{}{L}_i\stackrel{}{L}_k`$. From $`(14)`$ we have then the inequality
$$E_{stat}M.t.3\pi ^2B(F_\pi /e),$$
$`\left(16\right)`$
which was obtained first by Skyrme for the $`SU(2)`$ model.
Eight diagonal moments of inertia and $`28`$ off-diagonal ones define the rotation energy, a quadratic form in $`\omega _i\omega _k`$ as follows from $`(14)`$ and $`(15)`$. The analytical expressions for the moments of inertia are too lengthy to be reproduced here. Fortunately, it is possible to perform calculations without explicit analytical formulas, by substituting $`(15)`$ into $`(14)`$.
For configurations generated by $`SU(3)`$ projectors the Lagrangian of the system can be written in terms of angular velocities of rotation and moments of inertia in the form (in the body-fixed system):
$$L_{rot}=\frac{\mathrm{\Theta }_1}{2}\left(\omega _1^2+\omega _2^2\right)+\frac{\mathrm{\Theta }_3}{2}\omega _3^2+\frac{\mathrm{\Theta }_4}{2}\omega _4^2+\frac{\mathrm{\Theta }_5}{2}\omega _5^2+\frac{\mathrm{\Theta }_6}{2}\left(\omega _6^2+\omega _7^2\right)+\frac{\mathrm{\Theta }_8}{2}\omega _8^2+\mathrm{\Theta }_{38}\omega _3\omega _8+WZ_3\omega _3+WZ_8\omega _8.$$
$`\left(17\right)`$
After the standard quantization procedure the Hamiltonian of the system,
$`H=\omega _iL/\omega _iL`$, is a bilinear function of the generators $`R_i`$ of $`SU(3)`$ rotations:
$$H=\frac{R_1^2+R_2^2}{2\mathrm{\Theta }_1}+\mathrm{\Theta }_8\frac{\left(R_3WZ_3\right)^2}{2D_{38}}+\frac{R_4^2}{2\mathrm{\Theta }_4}+\frac{R_5^2}{2\mathrm{\Theta }_5}+\frac{R_6^2+R_7^2}{2\mathrm{\Theta }_6}+$$
$$+\mathrm{\Theta }_3\frac{\left(R_8WZ_8\right)^2}{2D_{38}}\frac{\mathrm{\Theta }_{38}}{D_{38}}\left(R_3WZ_3\right)\left(R_8WZ_8\right),$$
$`\left(18\right)`$
where $`D_{38}=\mathrm{\Theta }_3\mathrm{\Theta }_8\mathrm{\Theta }_{38}^2`$. For the states belonging to a definite $`SU(3)`$ irrep the rotation energy can be written in terms of the second order Casimir operators of the $`SU(2)`$ and $`SU(3)`$ groups:
$$E_{rot}=\frac{N\left(N+1\right)R_3^2}{2\mathrm{\Theta }_1}+\frac{U\left(U+1\right)R_{3,us}^2}{2\overline{\mathrm{\Theta }}_4}+\frac{V\left(V+1\right)R_{3,ds}^2}{2\mathrm{\Theta }_6}+$$
$$+\frac{\mathrm{\Theta }_8\left(R_3WZ_3\right)^2}{2D_{38}}+\frac{\mathrm{\Theta }_3\left(R_8WZ_8\right)^2}{2D_{38}}\frac{\mathrm{\Theta }_{38}}{D_{38}}\left(R_3WZ_3\right)\left(R_8WZ_8\right)+\frac{\left(R_4^2R_5^2\right)\left(\mathrm{\Theta }_5\mathrm{\Theta }_4\right)}{4\mathrm{\Theta }_4\mathrm{\Theta }_5}$$
$`\left(19\right)`$
with $`\overline{\mathrm{\Theta }_4}=2\mathrm{\Theta }_4\mathrm{\Theta }_5/(\mathrm{\Theta }_4+\mathrm{\Theta }_5)`$.
$$R_{3,us}=R_3/2+\sqrt{3}R_8/2,R_{8,us}=\sqrt{3}R_3/2R_8/2,$$
$$R_{3,ds}=R_3/2+\sqrt{3}R_8/2,R_{8,ds}=\sqrt{3}R_3/2+R_8/2,$$
$`\left(20\right)`$
with $`R_3=R_{3,ud}`$, $`R_8=R_{8,ud}`$; $`R_{3,ab},R_{8,ab}`$ denote the $`3^{rd}`$ and $`8^{th}`$ generators of the $`(a,b)`$ $`SU(2)`$ subgroup of $`SU(3)`$.
$`N,U`$ and $`V`$ are the values of the right isospin in the $`(u,d)SU(2)`$ subgroup, and so called $`U`$-spin and $`V`$-spin. They are connected with the second order Casimir operator of the $`SU(3)`$ group $`C_2(SU_3)=\frac{1}{3}(p^2+q^2+pq)+p+q`$; $`p,q`$ being the numbers of the upper and low indices in the tensor describing the $`SU(3)`$ irrep $`(p,q)`$ via the symmetric relation $`C_2(SU_3)=N(N+1)+U(U+1)+V(V+1)(R_3^2+R_8^2)/2`$. The hypercharge $`Y=2R_8/\sqrt{3}`$. The terms linear in the angular velocities present in the Lagrangian due to the Wess-Zumino term cancel in the Hamiltonian.
For a state with the lowest energy which belongs to the $`SU(3)`$ singlet with $`(p,q)=0`$ the quantum correction simplifies to:
$$\mathrm{\Delta }E=\frac{1}{2D_{38}}\left(\mathrm{\Theta }_8WZ_3^2+\mathrm{\Theta }_3WZ_8^22\mathrm{\Theta }_{38}WZ_3WZ_8\right).$$
$`\left(21a\right)`$
The quantities $`D_{38}`$, $`\mathrm{\Theta }_3+\mathrm{\Theta }_8`$, $`WZ_3^2+WZ_8^2`$ do not depend on the choice of the particular subgroup of $`SU(3)`$ group, therefore $`Eq.(21)`$ also is invariant, as one expects from general arguments, since
$$\mathrm{\Theta }_8WZ_3^2+\mathrm{\Theta }_3WZ_8^22\mathrm{\Theta }_{38}WZ_3WZ_8=\left(\mathrm{\Theta }_8+\mathrm{\Theta }_3\right)\left(WZ_3^2+WZ_8^2\right)\mathrm{\Theta }_3WZ_3^2\mathrm{\Theta }_8WZ_8^22\mathrm{\Theta }_{38}WZ_3WZ_8.$$
The generalization of $`(21a)`$ to an arbitrary case is
$$\mathrm{\Delta }E^{min}=\mathrm{\Theta }_{ij}^1WZ_iWZ_j/2,$$
$`\left(21b\right)`$
where $`\mathrm{\Theta }_{ij}^1`$ is the matrix inverse to the tensor of inertia $`\mathrm{\Theta }_{ij}`$. Eq. $`(21b)`$ is valid for all cases except the degenerate case when $`Det\mathrm{\Theta }_{ij}=0`$. Note, that the case of the $`SU(2)`$ embedding just corresponds to the degenerate case, with $`\mathrm{\Theta }_8=0`$, and this leads to the rigorous quantization condition .
As can be seen from the Table, for the configurations we consider here and for our choice of the subgroups $`\mathrm{\Theta }_{38}^2\mathrm{\Theta }_3\mathrm{\Theta }_8`$, so a small correction due to $`\mathrm{\Theta }_{38}`$ can be neglected, and $`\mathrm{\Delta }EWZ_3^2/(2\mathrm{\Theta }_3)+WZ_8^2/(2\mathrm{\Theta }_8)`$. Since the moments of inertia are of the order of magnitude $`10GeV^1`$ and both $`WZ_3`$ and $`WZ_8`$ are not greater $`1`$, the quantum correction for the lowest states does not exceed several tens of $`MeV`$: see the Table.
The quantum correction from the $`SU(3)`$ zero modes which should be added to the classical energy of the soliton located originally in the $`(u,d)`$ $`SU(2)`$ subgroup equals to $`\mathrm{\Delta }E=3B/4\mathrm{\Theta }_F365MeV`$ .
Note, that for the $`SO(3)`$ solitons considered in the $`WZ`$-term is equal to zero, and for this reason the lowest quantized state was an $`SU(3)`$ singlet without any quantum correction to its mass.
4. The $`FSB`$ part of the mass terms in the Lagrangian density which defines, in particular, the mass splittings inside $`SU(3)`$ multiplets is defined, as usual, by
$$L_{FSB}=\frac{F_D^2m_D^2F_\pi ^2m_\pi ^2}{24}Tr\left(1\sqrt{3}\lambda _8\right)\left(U+U^{}2\right)+\frac{F_D^2F_\pi ^2}{48}Tr\left(1\sqrt{3}\lambda _8\right)\left(Ul_\mu l^\mu +l_\mu l^\mu U^{}\right)$$
$`\left(22\right)`$
Here $`m_D`$ is the mass of $`K,D`$ or $`B`$ meson. The ratios $`F_D/F_\pi `$ are known to be $`1.22`$ and $`1.7\pm 0.2`$ for, respectively, kaons and $`D`$-mesons. The $`L_{FSB}`$ given by $`(22)`$ is sufficient to describe the mass splittings of the octet and decuplet of baryons within the collective coordinate quantization approach .
The contribution to the classical mass of the solitons from the $`SU(3)`$ projector ansatz which comes from $`FSB`$ part of the Lagrangian $`(22)`$, considered as a perturbation, is not small even for strangeness:
$$\mathrm{\Delta }M=2C_S\mathrm{\Delta }E_{M.t}\left(\frac{F_K^2m_K^2}{F_\pi ^2m_\pi ^2}1\right)$$
$`\left(23\right)`$
where $`\mathrm{\Delta }E_{M.t}`$ is the $`FS`$ mass term contribution shown in the Table. Numerically, it equals $`1.13GeV`$ for $`B=2`$ and $`3.2GeV`$ for $`B=6`$. Thus it makes sense to include the $`FSB`$ part of the mass term into the minimization procedure.
We do not consider here the quantum corrections due to the rotations in the ordinary space because they are explicitily zero for the lowest states with even $`B`$ ($`J=0`$), or small .
Our investigations have shown that the space of local minima for the $`SU(3)`$ skyrmions has a rather complicated form. In addition to the known local minima - the $`SU(2)`$ embeddings, $`SO(3)`$ solitons and skyrmion molecules - there are also other local minima best approximated by the $`SU(3)`$ projector ansaetze .
It is possible to quantize these configuratons and to estimate the spectrum of states by means of the collective coordinate quantization procedure. The total sum of the classical mass and of the zero modes quantum corrections can be the smallest for some baryon numbers. However, to draw the final conclusions one needs to calculate the Casimir energies of solitons which are essentially different for solitons of different form and size. This is a very complicated problem solved approximately - only for the $`B=1`$ hedgehog configurations. For the case of $`SU(2)`$ embeddings it was, however, possible to draw conclusions of physical relevance under the natural assumption that the unknown Casimir energy, or loop corrections, cancel in the difference of energies of states with different flavours, and that the states with $`u,d`$ flavours can be identified with ordinary nuclei . But this assumption may be incorrect for the states generated by the $`SU(3)`$ projector ansatz.
We are thankful to W.K. Baskerville, T. Ioannidou, B. Piette and P.M. Sutcliffe for their interest and helpful discussions.
This work has been supported by the UK PPARC grant: PPA/V/S/1999/00004.
References |
warning/0002/quant-ph0002028.html | ar5iv | text | # Activating bound entanglement in multiโparticle systems
## I Introduction
The existence of bound entanglement (BE) has been one of the most intriguing surprises in quantum information during the last few years. A state is BE if (despite of being entangled) one cannot distill maximally entangled states (MES) out of it by using local operations and classical communication (LOCC). From this definition it follows that with bound entangled states (BES) one cannot perform reliable teleportation , quantum communication , etc, i.e. they seem not to be useful for quantum information purposes. However, it has been shown that under certain conditions BE can be โactivatedโ . That is, with the help of some other entangled states they enable to carry out certain tasks which could not be performed by using these other entangled states and LOCC alone.
The first examples of BE arose in the context of two systems. In particular, the Horodecki showed that a necessary condition for a state of two systems to be distillable is that the corresponding density operator had nonpositive partial transposition . Thus, all entangled states with a positive partial transposition cannot be distilled, and therefore they are examples of BES . Very recently, we showed that a different kind of BE can also appear in multiparticle systems . In particular, we considered states of three systems $`A`$, $`B`$, and $`C`$, spatially separated, that cannot be prepared locally and that have the following properties: (i) if we would be allowed to join systems $`A`$ and $`C`$ in one place, the state could be prepared locally; (ii) if we would be allowed to join systems $`B`$ and $`C`$ in one place, the state could be prepared locally; (iii) even if we would be allowed to join systems $`A`$ and $`B`$ in one place, the state could not be prepared locally. The property (i) immediately implies that one cannot distill a MES between the systems $`A`$ and $`B`$ (or $`B`$ and $`C`$) by using LOCC, since even if we would allow nonlocal operations between $`A`$ and $`C`$ we would not be able to do it. Analogously, the property (ii) implies that one cannot distill a MES between the systems $`A`$ and $`C`$. Thus, no MES can be distilled between any of the systems. The property (iii), however, indicates that the state is entangled, since even if we would allow nonlocal operations between $`A`$ and $`B`$ we could not prepare it locally. Thus, despite the state is entangled, one cannot distill any MES out of it, and therefore it is a BES. Nevertheless, the state presented in can be activated. This follows from the fact that if we allow $`A`$ and $`B`$ to share some maximally entangled states (which is equivalent to allow $`A`$ and $`B`$ to join), then one can indeed distill a MES shared between $`A`$, $`B`$ and $`C`$. Another example of activable BE has been given by Smolin in the context of unlockable entanglement for the case of four systems . In that case, the state has the extra property that the entanglement between $`A`$ and $`B`$ can be activated using a single copy of the state only and by letting $`C`$ and $`D`$ share only one MES. These examples show that a new kind of BE can arise when we split some free entanglement among more than two parties. They also show that the possibility of activating BE is sometimes counterintuitive and may have some applications related to the process of secret sharing .
In this paper we construct multiparticle states which display novel properties related to BE. They illustrate the richness of multiparticle as compared to twoโparticle entangled states regarding the possibility of activating BE. This paper is organized as follows. In Sec. II, we state the problem of activation of BE in multiparticle systems and present various intriguing examples of activable BE. The rest of the paper provides the tools needed to construct states having the properties given in all those examples, as well as the guidelines to construct other states with different properties. In Sec. III, we review both the concept of bipartite splittings and the family of $`N`$ qubit states introduced in Ref. . These states will play a central role throughout the remaining of the paper; as we will show, they give rise to a vast variety of activable BES. In Sec. IV, we consider the distillability properties of a multiparticle state when we allow the particles to join into two groups (i.e. according to a bipartite splitting of the particles). We will show that it is possible to find BES that can be activated iff the particles are joint according to certain bipartite splittings. Moreover, we will show that one can choose the bipartite splittings for which the BE can be activated without any restriction, and that our family of states covers all possible examples of this kind. In Sec. V, we consider a more general framework where the particles are allowed to join into more than two groups. We derive necessary conditions for distillability and activation based on the distillability properties of bipartite splittings. We show that these conditions are also sufficient for our family of states, which implies that they are an important ingredient to construct generic examples of activable BES. In particular, they allow us to construct BES fulfilling the properties introduced in the examples of Sec. II. We finally summarize our results in Sec. VI.
## II Activation of bound entanglement
Let us consider $`N`$ parties, $`A_1,\mathrm{},A_N`$, at different locations, each of them possessing several qubits. We will assume that the state of the qubits is described by a density operator of the form $`\rho ^M`$, where $`\rho `$ denotes an entangled state of exactly $`N`$ qubits, each of them belonging to a different party. Thus, the parties have $`M`$ copies of the same state, where $`M`$ can be as large as we wish. This ensures that the parties can use distillation protocoles in order to obtain MES between some of them. In that case we will say that the state $`\rho `$ is distillable (with respect to the specific parties that can obtain a MES). A state $`\rho `$ is BE if it is not distillable when the parties remain separated from each other. We say that a BES can be activated if it becomes distillable once some of the parties join and form groups to act together. Note that instead of allowing some parties to join we could have allowed them to share some MES. In that case we would have the same situation given the fact that separated parties sharing MES can perform any arbitrary joint operation by simply teleporting back and forth the states of their particles.
### A Examples
In this subsection we introduce some relevant examples of BES, by specifying their properties with respect to activation. In the following sections we will explicitely construct density operators that fulfill the properties given here for each of the examples. The goal of this subsection is to display the underlaying properties of activable BES.
Example I: The state $`\rho _I`$ becomes distillable iff the parties form two groups with exactly $`j`$ and $`Nj`$ members, respectively. Furthermore, it does not matter which of the parties join in each group, but only the number of members. For example, if $`N=8`$ and $`j=3`$ (see Fig. 1a), we have that $`\rho _I`$ is distillable if exactly $`3`$ and $`5`$ parties join, but remains undistillable when the parties form two groups with 1-7, 2-6, 4-4 members, or if they form more than two groups. In particular, $`\rho _I`$ is not distillable if the parties remain separated from each other, which corresponds to having $`8`$ groups.
Example II: The state $`\rho _{II}`$ is distillable iff the parties form two groups, each of them containing say between $`40\%`$ and $`60\%`$ of the parties (and where each party is contained in one of the two groups). In particular, the state $`\rho _{II}`$ remains undistillable if the parties form more than two groups or if they form two groups, with one of them containing less than $`40\%`$ of the parties. Again, it does not matter whether certain parties belong to the same group; only the total number of particles within each group is important. We can view $`\rho _{II}`$ as a state having โmacroscopicโ entanglement, but no (distillable) โmicroscopicโ entanglement (see Fig. 1b). Note that one can also have the opposite effect, i.e a state $`\rho _{II}^{}`$ which is distillable iff the parties form two groups, and one of the groups contains less than say $`20\%`$ of the parties.
Example III: The state $`\rho _{III}`$ becomes distillable iff the parties form two groups, where the first group contains a specific set of $`M`$ parties $`A=\{A_{k_1},\mathrm{}A_{k_M}\}`$, and the second group contains the remaining parties. For all other configurations in groups $`\rho _{III}`$ remains undistillable. For example, we have for $`N=5`$ and $`A=\{A_1,A_3,A_5\}`$ that $`\rho _{III}`$ is distillable iff the the parties form two groups, $`(A_1A_3A_5)(A_2A_4)`$, and not distillable otherwise.
Example IV: The state $`\rho _{IV}`$ has the following properties: Given any two groups, each of them with a specific number $`j`$ (or more) parties, they can distill (with the help of the other parties) a MES between them. In particular, if there are several groups of $`j`$ or more parties, they can distill a GHZโlike state between all the groups. When a group with less than $`j`$ members is formed, it cannot distill a MES with any other group (see Fig. 2b). We call this effect clustering, since the parties have to form clusters with at least $`j`$ members in order to be able to create MES. In a similar way, one can also choose the state such that not only the number but also the specific parties which have to form groups in order to create a MES is given.
In the preceeding examples we have that some number of parties have to join into groups in order to distill some MES with some other parties. In the following, we will give examples in which the parties that join enable the remaining ones to distill entanglement.
Example V: The state $`\rho _V`$ is such that once any $`(N2)`$ parties form a group, the remaining two parties can distill a MES, but $`\rho _V`$ remains undistillable if less than $`(N2)`$ parties join (see Fig. 2b).
Example VI: $`\rho _{VI}`$ is a state of $`N=4`$ parties for which once the parties $`(A_3A_4)`$ form a group a GHZโlike state can be distilled among $`A_1`$, $`A_2`$, and the group $`(A_3A_4)`$, whereas it is undistillable whenever any other parties but $`(A_3A_4)`$ are joint (see Fig. 3a). In contrast to the previous example, in this one it is not only possible to create a MES between $`A_1`$ and $`A_2`$ by joining $`(A_3A_4)`$, but also this last group can distill a MES with the remaining parties.
Example VII: $`\rho _{VII}`$ is a state of $`N=5`$ parties such that a MES between parties $`A_1`$ and $`A_2`$ can be created iff either the parties $`(A_3A_4)`$ or $`(A_3A_5)`$ join (see Fig. 3b), but no entanglement can be distilled if the parties $`(A_4A_5)`$ join.
In the following, we will explain how all those examples can be constructed and understood by giving necessary conditions for distillation and activation in multiparticle systems. These conditions also provide us with the tools to construct other examples of activable BES. We will also introduce a family of states which includes all the examples I-VII as well as all those examples which can be constructed using the rules following from the necessary conditions for distillation and activation, which are also sufficient for this family.
## III Definitions and notation
In this section we first review the concept of bipartite splittings, and introduce some notation that will be used in the following sections. Then we review the properties of the family of $`N`$โqubit states $`\rho _N`$ introduced in Ref. . As mentioned above, these states give rise to a wide range of activable BES. In particular, among them one can find states corresponding to the examples of activable BES introduced in the previous section.
### A Bipartite splittings
Let us denote by $`๐ซ`$ the set of all possible (bipartite) splittings of $`N`$ parties into two groups. For example, for $`3`$ parties $`๐ซ`$ contains the splittings $`(A_1A_3)`$$`(A_2)`$, $`(A_2A_3)`$$`(A_1)`$, and $`(A_3)`$$`(A_1A_2)`$. We will denote these bipartite splittings by $`P_k`$, where $`k=k_1k_2\mathrm{}k_{N1}`$ is a chain of $`N1`$ bits, such that $`k_n=0,1`$ if the $`n`$โth party belongs to the same group as the last party or not. For example, for $`3`$ parties the splittings $`(A_1A_3)`$$`(A_2)`$, $`(A_2A_3)`$$`(A_1)`$, and $`(A_3)`$$`(A_1A_2)`$ will be denoted by $`P_{01}`$, $`P_{10}`$, and $`P_{11}`$, respectively. We will denote by $`A`$ the side of the splitting to which the party $`N`$ belongs and by $`B`$ the other side. In general, there exist $`s=2^{N1}1`$ of such splittings. In the following, when we consider bipartite splittings the parties in each of the groups will be allowed to act together (i.e. to perform joint operations).
### B Family of states $`\rho _N`$
Let us consider $`\rho _N`$, the family of $`N`$โqubit states introduced in . We have that $`\rho \rho _N`$ if it can be written as
$`\rho `$ $`=`$ $`{\displaystyle \underset{\sigma =\pm }{}}\lambda _0^\sigma |\mathrm{\Psi }_0^\sigma \mathrm{\Psi }_0^\sigma |`$ (2)
$`+{\displaystyle \underset{k0}{}}\lambda _k(|\mathrm{\Psi }_k^+\mathrm{\Psi }_k^+|+|\mathrm{\Psi }_k^{}\mathrm{\Psi }_k^{}|),`$
where
$$|\mathrm{\Psi }_k^\pm \frac{1}{\sqrt{2}}(|k_1k_2\mathrm{}k_{N1}0\pm |\overline{k}_1\overline{k}_2\mathrm{}\overline{k}_{N1}1),$$
(3)
are GHZโlike states with $`k=k_1k_2\mathrm{}k_{N1}`$ being a chain of $`N1`$ bits, and $`\overline{k}_i=0,1`$ if $`k_i=1,0`$, respectively. We have that $`\rho _N`$ is parametrized by $`2^{N1}`$ independent real numbers. The labeling is chosen such that $`\mathrm{\Delta }\lambda _0^+\lambda _0^{}0`$. As we will see below, both the separability and distillability properties of the states belonging to this family are completely determined by the coefficients
$$s_k\{\begin{array}{c}\text{1 if }\lambda _k<\mathrm{\Delta }/2\hfill \\ \text{0 if }\lambda _k\mathrm{\Delta }/2\text{.}\hfill \end{array}$$
(4)
Let us emphasize that the notation used for the states of this family parallels the one used to denote the partitions $`P_k`$. In fact, as shown in the separability properties of $`\rho _N`$ for a given partition $`P_k`$ are directly related to the coefficient $`s_k`$:
Lemma 0: For any bipartite splitting $`P_k๐ซ`$, and $`\rho \rho _N`$ we have $`\rho ^{T_A}0s_k=0\rho `$ is separable with respect to this splitting<sup>*</sup><sup>*</sup>*$`\rho ^{T_A}`$ denotes the partial transposition with respect to the parties $`A`$. For the definition of partial transpostion in mulitparticle systems see .
Thus the coefficient $`s_k`$ determines whether $`\rho `$ is separable or not with respect the bipartite splitting $`P_k`$. Note that there are no restrictions to the values of these coefficients; that is, for any choice of $`\{s_k\}`$ there always exists a state $`\rho \rho _N`$ with these values. This automatically implies that the family $`\rho _N`$ provides us with all possible examples of states in which the separability properties of the bipartite splittings are completely specified. As we will see in the next section, the same is true regarding distillability with respect to bipartite splittings.
## IV Activation of BE for bipartite splittings: examples I-III
With the states introduced in the previous section, we are at the position of examining the examples IโIII. If we take a closer look at them, we find that they all have a common feature: the states $`\rho _I`$, $`\rho _{II}`$ and $`\rho _{III}`$ are distillable with respect to certain bipartite splittings, and not distillable with respect to other bipartite splittings. In this Section we will show that given a subset of all possible bipartite splittings we can always find states in $`\rho _N`$ such that they are distillable iff the parties join according to a bipartite splitting contained in this set. This general result clearly allows to find states within our family which correspond to examples IโIII.
Let us be more specific. We consider the set of all possible bipartite splittings $`๐ซ`$. Let us specify for each splitting $`P_k`$ whether we want that one can distill a MES or not. To do that, we will assign to each splitting a $`0`$ if we do not want it to be possible and a $`1`$ otherwise. That is, each possible specification corresponds to a function $`f:๐ซ\{0,1\}`$. There are $`2^s`$ of such functions, which we will call specifications. Now, given a specification $`f`$ we define the set of splittings $`๐ฎ_f=\{P๐ซ\text{ such that }f(P)=1\}`$. Thus, the problem reduces to finding a state $`\rho `$ such that only for the splitings contained in $`๐ฎ_f`$ one can distill a maximally entangled state. Note that examples IโIII are just different specifications $`f`$. What we will show here is that there exist states $`\rho \rho _N`$ fulfilling any given specification.
### A Family of states $`\rho _N`$
We will show here that for the family of states $`\rho _N`$, bipartite inseparability is equivalent to distillability. This equivalence is expressed in the following Lemma:
Lemma 1: Given the bipartite splitting $`P_k๐ซ`$ and $`\rho \rho _N`$, we have that $`\rho `$ is distillable with respect to $`P_k`$ $`s_k=1`$.
Proof: Using Lemma 0 we have that if $`s_k=0`$ then $`\rho ^{T_A}0`$, and therefore the state is not distiallable, since nonโpositive partial transposition is a necessary condition for distillation . Thus, we just have to show that if $`s_k=1`$ then the state is distillable. We denote by $`|0_{A,B}`$ and $`|1_{A,B}`$ the states in which all qubits in side $`A`$ or $`B`$ are in state $`0`$ and $`1`$, respectively. We have that $`|\mathrm{\Psi }_0^\pm =\frac{1}{\sqrt{2}}(|0_A|0_B\pm |1_A|1_B)`$ and $`|\mathrm{\Psi }_k^\pm =\frac{1}{\sqrt{2}}(|0_A|1_B\pm |1_A|0_B)`$. By measuring the projectors $`|0_{A,B}0|+|1_{A,B}1|`$ in $`A`$ and $`B`$ respectively, we only get contributions from the states $`|\mathrm{\Psi }_0^\pm `$ and $`|\mathrm{\Psi }_k^\pm `$ and one obtains that the state after a successful measurement is
$`\rho `$ $``$ $`\lambda _0^+|\mathrm{\Psi }_0^+\mathrm{\Psi }_0^+|+\lambda _0^{}|\mathrm{\Psi }_0^{}\mathrm{\Psi }_0^{}|`$ (6)
$`+\lambda _k(|0,10,1|+|1,01,0|.`$
This state is known to be distillable for $`s_k=1`$ . $`\mathrm{}`$
This Lemma tells us that for a given specification we just have to find a state in $`\rho _N`$ such that $`s_k=f(P_k)`$ for all $`k`$. Since the values $`s_k`$ are not restricted in any way we have that for any specification we can find states $`\rho \rho _N`$ fulfilling it, and thus our family provides all different kinds of bipartite distillable and notโdistillable states.
Let us now come back to the examples IโIII. We can take as state $`\rho _I`$ \[example I\] one from the family $`\rho _N`$ which has $`s_k=f(P_k)=1`$ iff the number of ones in $`k`$ is $`j`$ or $`(N1j)`$ and $`s_k=f(P_k)=0`$ otherwise (this means that all bipartite splittings which contain exactly $`j`$ members in one group are distillable, and all others are separable). In the example II, we choose $`\rho _{II}\rho _N`$ such that $`s_k=f(P_k)=1`$ iff $`P_k`$ has between $`40\%`$ and $`60\%`$ of the parties in $`B`$. In the example III we take $`\rho _{III}\rho _N`$ such that $`s_k=f(P_k)=1`$ only for one specific $`P_k`$. We also emphasize that in a similar way one may construct many other interesting examples.
## V Activation of BE by external action: Examples IV-VII
While the separability and distillability properties with respect to bipartite splittings of a state $`\rho `$ were sufficient to completely understand and construct examples I-III, in the remaining examples we are faced with a slightly more complicated situation. Now we have that the parties can join into more than two groups, and it may even happen that not all of the parties are needed in order to activate the BE. Thus, the bipartite splittings $`P_k`$ do no longer provide a complete description of the problem. However, we can still use them to derive necessary conditions regading distillability and even activation of BE between any number of groups. And what is more important, we will show that these necessary conditions turn out to be also sufficient for the family of states $`\rho _N`$. On the one hand, this allows us to construct various different kinds of activable BES, such as those corresponding to examples IV-VII. On the other hand, it ensures that the family $`\rho _N`$ provides examples for all possible kinds of BES which can be constructed by using the necessary conditions for distillation based on the bipartite properties of a state $`\rho `$. We wish to emphasize that due to the fact that the conditions we obtain are only necessary and not sufficient in general, there may exist other kind of activable BES that cannot be obtained with the methods developed here. Nevertheless, these methods also give indications about some other kinds of BES.
### A Necessary conditions for distillation
The separability properties of the bipartite splittings $`P_k`$ of a multiparticle state $`\rho `$ provide necessary conditions for the distillability properties of $`\rho `$. This is expressed in the following result:
Theorem 1: Let $`C=\{A_{i_1},\mathrm{},A_{i_M}\}`$ and $`D=\{A_{j_1},\mathrm{},A_{j_L}\}`$ be two disjoint groups of $`M`$ and $`L`$ parties respectively, whereas the rest of the parties are separated. If a MES between $`C`$ and $`D`$ can be distilled then $`\rho `$ has to be nonโseparable with respect to all those bipartite splittings $`P_k`$ in which the groups $`C`$ and $`D`$ are located on different sides.
Proof: Let us assume that $`\rho `$ is separable with respect to one of those bipartite splittings $`P_k`$, so that $`CA`$ and $`DB`$. This means that even if we allow the groups $`C`$ and $`D`$ to join some other parties (belonging to $`A`$ and $`B`$, respectively), they will not be able to distill a MES. This is due to the fact that nonseparability is a necessary condition for distillability. $`\mathrm{}`$
Theorem 1 relates the distillability properties of $`\rho `$ to the classification with respect to the separability properties of a multiparticle state given in . We also note here that the $`k`$โseparability properties with respect to $`k`$โpartite splittings ($`k>2`$) provide no additional information about the distillability properties of a state $`\rho `$. Theorem 1 also determines the necessary conditions for the creation of GHZโlike states, since the possibility of creating maximally entangled pairs between any two out of $`l`$ parties is a necessary and sufficient condition for the creation of a GHZโlike state among those $`l`$ parties. We can also change โnonโseparableโ in the theorem to โdistillableโ, which provides an even stronger condition. It is not clear whether this condition is then also sufficient. One may think of the existence of bound entangled states which are distillable with respect to all possible bipartite splittings, but which are not distillable when considering the parties independently. Recently, it has been reported that such states in fact exist .
We also see from Theorem 1 that in order for a MES between two specific, separated parties, to be distillable, it is necessary that the corresponding state is inseparable with respect to at least $`2^{N2}`$ different bipartite splittings. Thus there exist many states which are inseparable, but do not fulfill the necessary condition for distillability between any two parties. All those states are obviously BE. For example, any state which has less than $`2^{N2}`$ (and more than one) inseparable bipartite splittings is clearly bound entangled. Naturally, the question arises under which conditions this BE can be activated.
### B Necessary conditions for activation of BES
Clearly, the situation changes when the parties are allowed to form groups and act together. In this case, the parties may be able to change the separability properties of certain bipartite splittings $`P_k`$. For example, for $`N=3`$ parties $`A,B`$, and $`C`$, if the parties $`A`$ and $`B`$ form a group they may be able to change the separability properties of the bipartite splittings (i) $`(B)`$$`(AC)`$ and (ii) $`(A)`$$`(BC)`$. This is due to the fact that โjoiningโ is equivalent to having some extra entanglement available. In this case, this extra entanglement between $`A`$ and $`B`$ can be used to change the separability properties of the bipartite splittings in question. Note, however, that this extra entanglement does not help to change the separability properties of the bipartite splitting (iii) $`(C)`$$`(AB)`$, since $`(AB)`$ where already joint in this bipartite splitting. This also allows us to understand the example of an activable BE three party state given in the introduction , where we had that (i) and (ii) are separable, while (iii) is inseparable and the state is thus BE according to Theorem 1. By joining the parties $`A`$ and $`B`$, one may however change the bipartite splittings (i) and (ii) from separable to inseparable, so the necessary conditions for distillation may now be fulfilled, since all bipartite splittings can now be inseparable in principle. As shown in , this activation can in fact be achieved, i.e. the change of the separability properties of the splittings (i) and (ii) as well as the distillation are possible. In a similar way, one can explain the example given in for $`N=4`$. The effect of activation of BE for a state $`\rho `$ can thus be viewed as a consequence of the following theorem:
Theorem 2: Consider a state $`\rho `$ which is separable with respect to a given bipartite splitting $`(A)`$$`(B)`$. When joining $`M`$ parties $`C=\{A_{i_1},\mathrm{},A_{i_M}\}`$, a necessary condition that we can make $`\rho `$ distillable with respect to the splitting $`(A)`$$`(B)`$ is that: (i) $`CA,B`$; (ii) by using an operation acting on $`C`$ one can transform the state such that it is now nonseparable with respect to the bipartite splitting $`(A)`$$`(B)`$.
Proof: (ii) follows trivially from Theorem 1, whereas (i) follows from (ii). $`\mathrm{}`$
According to this theorem, when joining some parties into a group $`C`$, they may change the separability and distillability properties of certain bipartite splittings, namely all those splittings $`(A)`$$`(B)`$ for which $`CA,B`$. So, it may happen that the conditions for distillability, which were not fulfilled before joining the parties, are now fulfilled, and a MES shared between some parties can now be created in principle. Theorems 1 and 2 together provide necessary conditions for the activation of multiparticle BES and provide thus the framework for the construction of generic examples of different kinds of activable bound entangled states. To achieve this, one chooses at the beginning the separability and distillability of the bipartite splittings $`P_k`$ of a state $`\rho `$ such that the state is not distillable if the parties remain separated from each other, but that the distillability conditions may be fulfilled when some of the parties are allowed to form groups. We also note here that due to the fact that both theorems only provide necessary conditions in general, we do not have that the distillabity and activation properties of a state $`\rho `$ are completely determined by the separability properties of its bipartite splittings. On the one hand, there might exist states which cannot be distilled or activated eventhough they fulfill the necessary conditions for distillation/activation, i.e. those states are further protected against activation. On the other hand, we already see that there are various kinds of nonโactivable bound entangled states. For example, all states which are biseparable with respect to all bipartite splittings but are inseparable with respect to any $`k`$โpartite splitting ($`k>2`$) (see also classification proposed in ) are clearly bound entangled and can be neither distilled nor activated. For $`N=3`$, such an example is known .
Using Theorems 1 and 2, it is now straightforward to check that in example I the state $`\rho _I`$ reamains undistillable whenever the parties form more than two groups, since the necessary condition for distillation of a MES between any two groups can not be fulfilled. On the other hand, it is already clear that if the parties form two groups with a different number of members than $`j`$ and $`Nj`$, we have by construction that the state $`\rho _I\rho _N`$ is separable with respect to this bipartite splitting and thus not distillable. In a similar way, one can check that in examples II and III the states $`\rho _{II}`$ and $`\rho _{III}`$ are only distillable iff the parties join in two groups of required size (example II) or required members (example III), respectively.
### C Family of states $`\rho _N`$
In this section, we show that the necessary condition for distillation and activation expressed in Theorems 1 and 2 are also sufficient for the family of states $`\rho _N`$. We first show that if the necessary conditions for the distillation of MES between any two groups of parties, both not including a certain party $`A_l`$, are fulfilled, one can disentangle party $`A_l`$ from the remaining system while keeping the necessary conditions for distillation between the two groups in question. In order to achieve this, the party $`A_l`$ has to cooperate, i.e. its help is requried. This is expressed in the following lemma:
Lemma 2: One can convert any $`N`$โqubit state $`\rho \rho _N`$ to a $`(N1)`$โqubit state $`\stackrel{~}{\rho }\rho _{N1}`$ by measuring a certain party $`A_l`$, such that for all bipartie splittings $`P_k`$ the following property is fulfilled: If $`\rho `$ is inseparable with respect to the bipartite splittings $`(A_lC)`$โ(rest) and $`(C)`$$`(A_l`$ rest) $`\stackrel{~}{\rho }`$ is inseparable with respect to the bipartite splitting $`(C)`$โ(rest).
Proof: We assume without loss of generality that $`A_l=A_1`$. If we measure in $`A_1`$ the Projector $`P_+=|+_{A_1}+|`$ where $`|+=(|0+|1)/\sqrt{2}`$, we find that the remaining $`(N1)`$ particles are in a state of the form $`\rho _{N1}`$ with new coefficients $`\stackrel{~}{\lambda }_{j_2j_3\mathrm{}j_{N1}}=\lambda _{0j_2j_3\mathrm{}j_{N1}}+\lambda _{1j_2j_3\mathrm{}j_{N1}}`$ and similar for $`\lambda _0^\pm `$, so that $`\stackrel{~}{\mathrm{\Delta }}=\mathrm{\Delta }`$. The property we want to show is simply that iff both
$$\lambda _{0j_2j_3\mathrm{}j_{N1}}<\mathrm{\Delta }/2\text{ and }\lambda _{1j_2j_3\mathrm{}j_{N1}}<\mathrm{\Delta }/2,$$
(7)
it follows that
$$\stackrel{~}{\lambda }_{j_2j_3\mathrm{}j_{N1}}<\stackrel{~}{\mathrm{\Delta }}/2.$$
(8)
It may happen that although (7) is fulfilled, (8) is not. In this case, we apply the first step of the distillation procedure proposed in first, where one measures certain POVM elements on $`M`$ copies of the initial state and is left with a new (unnormalized) state of the form $`\rho _N`$ with new coefficients $`\lambda _k^{}=\lambda _k^M`$ and $`(\mathrm{\Delta }^{}/2)=(\mathrm{\Delta }/2)^M`$. For $`M`$ sufficiently large, we can now allways have that after applying $`P_+`$, the remaining state is such that (8) is fulfilled, since the new $`\mathrm{\Delta }/2`$ is exponentially amplified compared to any $`\lambda _{k_{1,2}}<\mathrm{\Delta }/2`$ and thus $`(\mathrm{\Delta }/2)^M>\lambda _{k_1}^M+\lambda _{k_2}^M`$ for $`M`$ sufficiently large. $`\mathrm{}`$
Theorem 3: For the family of states $`\rho _N`$, we have that the necessary condition for distillability given in Theorem 1 is also sufficient.
Proof: To show this, one just has to sequentially apply Lemma 2 to all particles $`A_i\{C,D\}`$, which leaves us with a $`(M+L)`$ qubitโstate $`\rho \rho _{M+L}`$ which has $`\rho ^{T_C}0`$ (which, according to Lemma 0 means that the corresponding $`s_k=1`$) and thus Lemma 1 can be applied, ensuring that any state $`\rho \rho _N`$ fullfilling the necessary conditions for any kind of distillation is in fact distillable. $`\mathrm{}`$
Note that the help of all parties is required, independent of whether they finally share a MES or not. This follows from the fact that Lemma 2 is based on the cooperation of party $`A_l`$, which is separated from the remaining system after the procedure described above. If one would just trace out party $`A_l`$ (i.e. the party does not cooperate with the remaining parties), one finds that the remaining parties cannot create any entanglement at all.
Theorem 4: For $`\rho \rho _N`$, and the situation of Theorem 2, we can in fact change all those bipartite splittings $`(A)`$-$`(B)`$ for which $`CA,B`$ from separable to inseparable (which is equivalent to distillable in this case) without changing the separability properties of the remaining bipartite splittings.
Proof: We assume without loss of generality that we join the first $`M`$ parties $`C=\{A_1,\mathrm{},A_M\}`$. Let $`j=j_1j_2\mathrm{}j_M`$ $`(l=l_1\mathrm{}l_{NM1})`$ be $`M`$ \[$`(NM1)`$\] digit binary numbers. We have to show that we can change all those bipartite splittings which do not contain all parties $`C`$ on one side from separable to inseparable. This is equivalent to showing that for all $`\lambda _{jl}`$ with $`j\{0,2^M1\}`$, we can have $`\mathrm{\Delta }/2>\lambda _{jl}`$. When applying the POVM element $`\stackrel{~}{P}=_{j=0}^{2^M1}\sqrt{y_j}|j_Cj|`$ , we find that we obtain again an (unnormalized) state of the form $`\rho _N`$ with new coefficients $`\stackrel{~}{\lambda }_{jl}=y_j\lambda _{jl}`$. Choosing $`y_0=y_{2^M1}=1`$ and all other $`y_j`$ sufficiently small, we have that $`\stackrel{~}{\mathrm{\Delta }}=\mathrm{\Delta }`$ and for $`j\{0,2^M1\}`$ we can obtain that $`\stackrel{~}{\mathrm{\Delta }}/2>\stackrel{~}{\lambda }_{jl}`$ as required. Furthermore, all other relations $`\mathrm{\Delta }/2\times \lambda _{jl}`$, with $`\times \{>,\}`$ and thus the separability properties of all those bipartite splittings for which $`CA,B`$ remain unchanged. $`\mathrm{}`$
Theorems 3 and 4 together ensure that any BES within the family $`\rho _N`$ which is activable in principle can in fact be activated, i.e. the necessary conditions for the activation of bound entangled states given in Theorem 1 and 2 are also sufficient for the family $`\rho _N`$.
### D Examples IV-VII
We are now at the position to construct and explain examples IVโVII, as well as to provide many other interesting examples of activable BES.
Let us start with example IV: In this case, we choose the state $`\rho _{IV}\rho _N`$ such that it is separable with respect to all bipartite splittings $`P_k`$ where either the group $`A`$ or $`B`$ has less than $`j`$ members. All other bipartite splittings are chosen to be distillable. It is clear that if a group with less than $`j`$ members is formed, they cannot distill a MES with any other group (since the corresponding bipartite splitting is separable). However once the parties form two groups, each having more than $`j`$ members, it is straightforward to check (using Theorems 4 and 3) that these two groups can in fact create a MES.
In example V, we choose the state $`\rho _V\rho _N`$ such that it is distillable with respect to all bipartite splittings $`P_k`$ where either the group $`A`$ or $`B`$ contains exactly one particle (so that the number of ones in $`k`$ is either 1 or $`N1`$), and separable with respect to all other bipartite splittings $`P_k`$. It is clear that it is sufficient to join any $`N2`$ particles, since this allows - according to Theorem 4 - to make inseparable (and thus distillable) all those bipartite splittings where the remaining particles are located in different groups, so that distillation of a MES between the remaining two particles becomes possible, according to Theorem 3. On the other hand, if less than $`N2`$ parties join, one can easily verify that for any two of the remaining parties, there always remains at least one bipartite splitting separable which has to be inseparable in order that distillation can be possible. Thus the remaining parties cannot share entanglement if less than $`N2`$ parties join.
In example VI we choose $`\rho _{VI}\rho _4`$ such that it is inseparable with respect to the bipartie splittings $`(A_1A_2)`$$`(A_3A_4)`$, $`(A_1)`$$`(A_2A_3A_4)`$ and $`(A_2)`$$`(A_1A_3A_4)`$ and separable with respect to all other bipartite splittings. Clearly, this state is BE since at least $`2^{N2}=2^2=4`$ bipartite splittings have to be inseparable so that a pair between any two separated parties can be distilled. Furthermore, one can check (using Theorems 3 and 4) that $`\rho _{VI}`$ remains undistillable when joining any two parties but $`(A_3A_4)`$, but one can create a GHZ state shared among $`A_1`$$`A_2`$$`(A_3A_4)`$ once one joins the parties $`(A_3A_4)`$.
Finally, in example VII we have $`N=5`$ and choose $`\rho _{VII}\rho _5`$ such that the state is inseparable with respect to all bipartite splittings which contain $`A_1`$ and $`A_2`$ in different groups, except the splitting $`(A_1A_3)`$$`(A_2A_4A_5)`$ which is chosen to be separable as well as all other bipartite splittings. One can readily check that $`\rho _{VII}`$ has the desired properties, i.e. it is BE and can be activated when joining the parties $`(A_3A_4)`$ or $`(A_3A_5)`$.
If we demand however that entanglement between $`A_1`$ and $`A_2`$ should be created once any two of the remaining parties join, one finds that such a a state cannot be constructed using our method. In this case, we have to demand that the initial state $`\rho _{VII}\rho _N`$ has to be separable with respect to at least one of the bipartite splittings where $`A_1`$ and $`A_2`$ belong to different groups. Otherwise, a MES between $`A_1`$ and $`A_2`$ could be distilled from the beginning and the state is thus not BE. Let us assume without loss of generality that the separable splitting is either (i) $`(A_1A_3)`$$`(A_2A_4A_5)`$ or (ii) $`(A_1)`$$`(A_2A_3A_4A_5)`$. In both cases, the separability properties of this splitting cannot be changed when joining the parties $`(A_4A_5)`$ and so no entanglement can be created between $`A_1`$ and $`A_2`$ although two of the remaining parties join (note that if we had taken any other bipartite splitting with $`A_1`$ and $`A_2`$ belonging to different groups to be separable, we would have found some other two of the remaining parties which can join but not change the properties of this splitting).
Due to the fact that the conditions for distillation and activation we give here are only necessary in general, they do not allow us to rule out the possibility that a state having the desired properties can exist. In this case, the activation would not be based on the change of the separability properties of the bipartite splittings, but on some other mechanism.
## VI Summary
In summary, we have given rules to construct activable bound entangled states using the separability and distillablity properties of a density operator with respect to bipartite splittings. This method allows us to construct examples for all possible kinds of activable BES where the parties join into exactly two groups. In particular, the family of states introduced in Ref. contains examples of all these kinds of activable BES. We have also given some relevant examples of activation of BE in which the parties join into more than two groups, and where the role of some of the groups is just to help the others to create a MES.
We thank G. Vidal and J. Smolin for discussions. This work was supported by the Austrian Science Foundation under the SFB โcontrol and measurement of coherent quantum systems (Project 11), the European Community under the TMR network ERBโFMRXโCT96โ0087, the European Science Foundation and the Institute for Quantum Information GmbH. |
warning/0002/quant-ph0002053.html | ar5iv | text | # Cloning and quantum computation
## 1 Introduction and overview
Since it became clear that it is impossible to make perfect copies of an unknown quantum state , much effort has been put into developing optimal cloning processes. As cloning represents a distribution of quantum information over a larger system, it can be seen as a type of quantum information processing tool. In this article we discuss the usefulness of quantum cloning to enhance the performance of some quantum computation tasks.
The COPY operation in classical computing is very useful, as it allows one to make multiple copies of the output of some computation, that can be fed as the input to further multiple processes. In quantum computing, however, the copying (quantum cloning) is imperfect, introducing some noise in the second round of computation. This situation is pictured in Fig. 1(a).
We can represent the first part of the quantum computation as an unitary $`U_0`$ applied to the initial state $`|0`$, resulting in an output state that we clone. We then feed the clones to two different computational branches, represented by unitaries $`U_1`$ and $`U_2`$. At the end of the process we make a measurement on the two final states, obtaining some information about the two computational branches we want to perform ($`U_1U_0|0`$ and $`U_2U_0|0`$). The problem with this quantum scenario is that the copies are imperfect, resulting in lower chances of getting the correct results at the end. Nevertheless, we will show that at least for some tasks, the use of cloning improves our chances of correctly computing both branches, if there are constraints on the number of times we can run the first part $`U_0`$. In the two examples we discuss below we will be comparing approaches in which there is some distribution of quantum information (done by the cloning process) with any approach that does not resort to this (see Fig 1(b)).
We may discriminate between two main approaches to quantum cloning. The first relies on adding some ancillary quantum system in a known state and unitarily evolving the resulting combined system, deterministically obtaining a pure state with partial mixed density matrices $`\rho _c`$ (the clones) that are as close as possible to the original state $`|\psi `$, as measured by the fidelity $`F=\psi \left|\rho _c\right|\psi `$ . The clones are of the following form:
$$\rho _c=\frac{๐}{d}(1\eta )+\eta |\psi \psi |$$
Optimal universal cloning machines are the unitaries that result in the largest state-independent clone fidelities $`F`$. The efficiency of these machines have been shown to be characterized by :
$`\eta `$ $`={\displaystyle \frac{N}{M}}{\displaystyle \frac{\left(M+d\right)}{\left(N+d\right)}};`$ (1)
$`F`$ $`={\displaystyle \frac{(1\eta )}{d}}+\eta ={\displaystyle \frac{MN+N(M+d)}{M(N+d)}}`$ (2)
where $`d`$ is the dimensionality of $`|\psi `$, $`M`$ is the number of clones and $`N`$ is the number of original copies of $`|\psi `$.
The second kind of cloning procedure is non-deterministic, consisting in adding an ancilla, performing unitary operations and measurements, with a postselection of the measurement results. Duan and Guo have shown that linearly independent pure states can be probabilistically cloned that way, and proved some theorems that allow one to calculate the optimal efficiencies. The resulting clones are perfect, but the procedure only succeeds with a certain probability $`p<1`$, which depends on the particular set of states which we are trying to clone.
Cloning machines can be viewed as a way of encoding quantum information contained in the input state $`|\psi `$ into a number of clones. In a sense, it accomplishes this more successfully than any procedure that relies on obtaining information about $`|\psi `$ through measurement. In order to see this, consider the universal cloning machines described above, operating on $`N`$ copies of $`|\psi `$, producing $`M`$ identical clones described by reduced density matrices $`\rho _c`$. Each clone has fidelity given by eq. 2; notice that the fidelity of the clones is a decreasing function of $`M`$. The best โclassical clonesโ that we can produce through measurement on $`|\psi `$ followed by state preparation have a lower fidelity, also given by 2, but with $`M\mathrm{}`$ . It is in this sense that we can say that the cloning process distributes quantum information about $`|\psi `$ in a way that direct measurement on $`|\psi `$ cannot.
With these considerations in mind, it is natural to wonder if, and how, cloning can be used to improve the performance of quantum information processing tasks. One might think that cloning could be helpful in state estimation, but it has been shown that this task is equivalent to cloning, when the number of copies $`M\mathrm{}`$ . As a result, in order to obtain some improvement our strategy needs to rely on using the quantum information present in the clones for further coherent quantum information processing, in the same spirit as the circuit in Fig. 1a. In what follows we give two examples of tasks which can be better performed if we use quantum cloning. In the first example we apply optimal universal cloning machines, whereas in the second we rely on the probabilistic cloning discussed by Duan and Guo .
## 2 Examples
In this section we present two examples of quantum computational tasks whose performance is enhanced if we distribute quantum information using quantum cloning. The first task makes use of state-independent universal quantum cloning, whereas the second task relies on state-dependent probabilistic quantum cloning.
### 2.1 First example
The first example we present is based on the scenario introduced in Fig. 1a. It models the general situation in which we want to perform $`M`$ different quantum computations, all of them with some first computational steps $`U_0`$ in common. Suppose that we are constrained to run $`U_0`$ only once. This may happen if $`U_0`$ is a complex, lengthy computation. In this case, we will be forced to find a scheme that obtains the $`M`$ computation results with the largest probability, using $`U_0`$ only once. Possible schemes may or may not resort to cloning to distribute quantum information; the example below is one in which cloning enables us to improve our performance, in relation to any scheme in which there is no information distribution using cloning.
In order to specify our task, suppose that we are given $`(M+1)`$ quantum blackboxes. What blackbox $`j`$ does is to accept one $`d`$-level quantum system as an input and apply a unitary operator $`U_j`$ to it, producing the evolved state as the output. We may think of the blackboxes as quantum oracles, or quantum sub-computations. The $`U_j`$ are chosen randomly from all possible $`U(d)`$ unitaries, using the unique uniform distribution invariant under action of $`U(d)`$ (see ). Our task will be to build quantum circuits that use each $`U_j`$ at most once to create $`M`$ mixed quantum states $`\rho _j`$, each as close as possible to $`|\varphi _j=U_jU_0|0,`$ $`(j=1,2,\mathrm{}M)`$, where $`|0`$ is an arbitrary reference state. Our score will be given by the average fidelity of our guesses:
$$\overline{F}=\frac{1}{M}\underset{j=1}{\overset{M}{}}\varphi _j\left|\rho _j\right|\varphi _j.$$
If we are not allowed to clone the state, there are two possible strategies. The first no-cloning strategy is to start by finding one of the $`|\varphi _j`$, say $`|\varphi _1=U_1U_0|0`$, with fidelity one. Now that we have used $`U_0`$ and $`U_1`$ once already, we must make guesses about the other $`(M1)`$ states $`|\varphi _j`$ by using only the remaining $`(M1)`$ blackboxes. As they were drawn from an uniformly random distribution, the best we can do is to make random guesses (each, on average, with $`F=1/d`$), obtaining, on average, a score
$$\overline{F}_1=\frac{1}{M}\left(1+\frac{(M1)}{d}\right).$$
The second no-cloning strategy starts by running $`U_0`$ , followed by measurements that accomplish an optimal estimation of the resulting state $`U_0|0`$. After this, we can use the information gathered to build the $`M`$ imperfect copies necessary to proceed to the second part of the computation with the $`U_j`$ $`(j=1,2,\mathrm{}M)`$. As we have mentioned, this second approach yields clones with fidelity given by eq. 2 with $`N=1,M\mathrm{}`$ (see ):
$$\overline{F}_2=\frac{2}{(d+1)}.$$
(3)
We obtain our guesses for states $`|\varphi _j`$ $`(j=1,2,\mathrm{}M)`$ by applying each $`U_j`$ $`(j=1,2,\mathrm{}M)`$ to a clone, resulting in a score also given by eq. 3. The best no-cloning strategy will be either of the two presented above, depending on the parameters $`M`$ and $`d`$.
Now let us see how cloning allows us to obtain a higher score $`\overline{F}`$ . We accomplish this by using a quantum circuit that first applies $`U_0`$ to the initial state $`|0`$, followed by an optimal universal cloning machine to obtain $`M`$ imperfect copies $`\rho _c`$ of state $`U_0|0`$. We then apply each $`U_j`$ $`(j=1,2,\mathrm{}M)`$ to a clone, obtaining reduced density matrices
$$\rho _j=\frac{๐}{d}(1\eta )+\eta |\varphi _j\varphi _j|$$
with $`\eta `$ given by eq. 1, with $`N=1`$. Using the resulting $`\rho _j`$โs as our guesses for states $`|\varphi _j`$ $`(j=1,2,\mathrm{}M)`$, we obtain an overall score
$$\overline{F}_{cloning}=\frac{2M+d1}{M(d+1)}$$
(4)
which is always higher than $`\overline{F}_1`$ and $`\overline{F}_2`$. In fact, eq. 4 represents the optimal score obtainable for this task, at least in the case $`M=2`$. In order to see this, we first note that asymmetric cloning (arising when the factors $`\eta `$ are in general different for each copy) is of no help in raising the score. This can be deduced from and , where the authors consider asymmetric cloning with $`M=2`$ and show that the sum of the fidelities of the copies is maximized by symmetric cloning. Furthermore, the optimality of the universal cloning procedure we have used entails optimality for the fidelity of each of the $`\rho _j`$, and therefore a maximal value of the score $`\overline{F}`$. This shows that this task is optimally performed (with optimal score given by eq. 4) if and only if we are allowed to use cloning. It is straightforward to generalize the result to the case where we are allowed to run $`U_0`$ $`N`$ times $`(N<M)`$, instead of just once, and quantum cloning still offers an advantage.
The scenario described above models the situation in which we have a series of quantum computations with some computational steps $`U_0`$ in common. We must note that we have assumed complete lack of knowledge about the intermediate state $`U_0|0`$ and about the final target states $`U_jU_0|0`$ $`(j=1,2,\mathrm{}M)`$. In the general case this will not be a good assumption, as many quantum computations will output states picked from a limited set of states. This can be taken into account with state-dependent quantum cloning and a different choice of scoring functions. In the next section we give an example of this.
### 2.2 Second example
In our second example we take the blackboxes of the previous example to consist of arbitrary quantum circuits that query a given function only once. The query of function $`f_i`$ is the unitary that performs
$$|x|y|x|yf_i(x),$$
where we have used the symbol $``$ to represent the bitwise $`XOR`$ operation. For ease of analysis, we restrict ourselves to the case $`M=2`$ and also restrict the set of possible functions $`f_0`$, $`f_1`$ and $`f_2`$. Our task will involve determining two functionals, one which depends only on $`f_0`$ and $`f_1`$, and the other on $`f_0`$ and $`f_2`$. As in the previous example, we will compare the performances of cloning and no-cloning strategies.
In order to precisely state our task, let us start by considering all functions $`h_i`$ which take two bits to one bit. We may represent each such function with four bits $`a,b,c`$ and $`d`$, writing $`h_{a,b,c,d}`$ to represent the function $`h`$ such that $`h(00)=a,h(01)=b,h(10)=c,`$ and $`h(11)=d`$. Let us now define some sets of functions that will be helpful in stating our task:
$`S_{f0}=\{h_{0010},h_{0101},h_{1001}\},`$
$`S_1=\{h_{0001},h_{0010},h_{0100},h_{1000}\},S_2=\{h_{0000},h_{0011},h_{0101},h_{1001}\}`$
$`S_{f12}=S_1S_2,`$
$`S_{0000}=\{h_{0000},h_{1111}\},S_{0011}=\{h_{0011},h_{1100}\},`$ (5)
$`S_{0101}=\{h_{0101},h_{1010}\},S_{1001}=\{h_{1001},h_{0110}\},`$ (6)
$`S_f=S_{0000}S_{0011}S_{0101}S_{1001}.`$
Now we randomly pick a function $`f_0`$ $`S_{f0}`$, after which two other functions $`f_1`$ and $`f_2`$ are picked from the set $`S_{f12}`$, also in a random fashion but obeying the constraints:
$$f_0f_1,f_0f_2S_f.$$
(7)
Here we use the symbol $``$ to represent addition modulo 2, which is equivalent to the bitwise $`XOR`$ operation. Our task will be to find in which of the four sets $`S_{0000},S_{0011},S_{0101}`$ and $`S_{1001}`$ lie each of the functions $`f_0f_1`$ and $`f_0f_2`$, using quantum circuits that query $`f_0,f_1`$ and $`f_2`$ at most once each. Our score will be given by the average probability of successfully guessing both correctly.
The best no-cloning strategy we have found goes as follows. Firstly, note that if $`f_0=h_{0010}`$ then both $`f_1`$ and $`f_2`$ must be in set $`S_1`$, because of the constraints given by eq. 7; similarly, if $`f_0`$ is either $`h_{0101}`$ or $`h_{1001}`$, then $`f_1`$ and $`f_2`$ must be in set $`S_2`$. Since we have drawn the function $`f_0`$ randomly, we will have both functions $`f_1`$ and $`f_2`$ in set $`S_2`$ with probability $`p=2/3`$. We will assume that this is the case; then we can discriminate between the two possibilities for $`f_0`$ with a single, classical function call. Furthermore, by using the quantum circuit in Fig. 2 twice (once with each of $`f_1`$ and $`f_2`$) we can distinguish the four possibilities for functions $`f_1`$ and $`f_2`$.
This happens because this quantum circuit results in four orthogonal states $`|\varphi _i=_{x=00}^{11}(1)^{f_i(x)}|x`$, depending on which function in set $`S_2`$ was queried. This allows us to determine functions $`f_0,f_1`$ and $`f_2`$ correctly with probability $`p=2/3`$, in which case we can determine which sets contain $`f_0f_1`$ and $`f_0f_2`$ and accomplish our task. Even in the case where our initial assumption about $`f_0`$ was wrong, we may still have guessed the right sets by chance; a simple analysis shows that our chances of getting both right this way are only $`1/16`$. On average, then, by using this no-cloning strategy we obtain a score:
$$p_1=\frac{2}{3}+\frac{1}{3}\frac{1}{16}=0.6875.$$
This is the best no-cloning score we could find for this task.
We can do better than that with quantum cloning. The idea now is to devise a quantum circuit that queries function $`f_0`$ only once, makes two clones of the resulting state and then queries functions $`f_1`$ and $`f_2`$, one in each branch of the computation. Since we have some information about the state produced by one query of $`f_0`$, the best cloning strategy will no longer be the universal, deterministic cloning derived in ; the probabilistic cloning machines discussed by Duan and Guo , will suit this task better.
The quantum circuit that we apply to solve this problem is given in Fig. 3.
Immediately after querying function $`f_0`$, we have one of three possible linearly independent states (each corresponding to one of the possible $`f_0`$โs):
$`|h_{0010}`$ $`{\displaystyle \frac{1}{2}}\left[|00+|01|10+|11\right],`$ (8)
$`|h_{0101}`$ $`{\displaystyle \frac{1}{2}}\left[|00|01+|10|11\right],`$ (9)
$`|h_{1001}`$ $`{\displaystyle \frac{1}{2}}\left[|00+|01+|10|11\right].`$ (10)
We can build probabilistic cloning machines with different cloning efficiencies (defined as the probability of cloning successfully) for each of the states 8-10. Theorem 2 of provides us with inequalities that allow us to derive achievable efficiencies for the probabilistic cloning process. We did a numerical search that yielded the following achievable efficiencies for probabilistically cloning the states in eqs. 8-10:
$`\gamma _1`$ $`\gamma (|h_{0010})=0.14165,`$ (11)
$`\gamma _2`$ $`\gamma (|h_{0101})=\gamma (|h_{1001})=0.57122.`$ (12)
After the cloning process we can measure a โflagโ subsystem and know whether the cloning was successful or not. For this particular cloning process, the probability of success is, on average, $`p_{\text{success}}=(\gamma _1+2\gamma _2)/30.4280`$. Let us suppose that it was successful. Then each of the cloning branches goes through the second part of the circuit in Fig. 3, to yield one of the four orthogonal states:
$`|h_{0000}`$ $`{\displaystyle \frac{1}{2}}\left[|00+|01+|10+|11\right],`$ (13)
$`|h_{0011}`$ $`{\displaystyle \frac{1}{2}}\left[|00+|01|10|11\right],`$ (14)
$`|h_{0101}`$ $`{\displaystyle \frac{1}{2}}\left[|00|01+|10|11\right],`$ (15)
$`|h_{1001}`$ $`{\displaystyle \frac{1}{2}}\left[|00+|01+|10|11\right],`$ (16)
which can be discriminated unambiguously. We obtain state $`|h_{0000}`$ if and only if the combined function $`f_0f_i`$ is one of the two in set $`S_{0000}`$, as can be checked by calculating the effect of the circuit in Fig. 3 for all possible $`f_0,f_1`$ and $`f_2`$. The situation is similar for the other three states; the detection of each of them signals precisely which one of the four sets $`\{S_{0000},S_{0011},S_{0101},S_{1001}\}`$ contains $`f_0f_i`$ . As a result, if the cloning process is successful, we manage to accomplish our task.
However, the cloning process will fail with probability $`(1p_{\text{success}})`$. If this happens, a simple evaluation of the posterior probabilities for function $`f_0`$ shows that it is more likely to be $`h_{0010}`$ than the other two, thanks to the relatively low cloning efficiency for the state in eq. 8, in relation to the states in eqs. 9 and 10 (see eqs. 11-12). If we then guess that $`f_0=h_{0010}`$, we will be right with probability
$$p_{0010}=\frac{(1\gamma _1)}{(1\gamma _1)+2(1\gamma _2)}0.5002.$$
What is more, we are still free to design quantum circuits to obtain information about $`f_1`$ and $`f_2`$, since at this stage we still have not queried them. Given our guess that $`f_0=h_{0010}`$, only the four functions in $`S_1`$ can be candidates for $`f_1`$ and $`f_2`$, because of the constraints given by eq. 7. These four possibilities can be discriminated unambiguously by running a circuit like that of Fig. 2 twice, once with $`f_1`$ and once with $`f_2`$. The circuit produces one of four orthogonal states, each corresponding to one of the four possibilities for $`f_i`$. Therefore, if our guess that $`f_0=h_{0010}`$ was correct, we are able to find the correct $`f_1`$ and $`f_2`$ and therefore accomplish our task. In the case that $`f_0h_{0010}`$ after all, we may still have guessed the right sets by chance; a simple analysis shows that this will happen with probability $`1/16`$.
The above considerations lead to an overall probability of success given by
$$p_2=p_{\text{success}}+(1p_{\text{success}})\left[p_{0010}+(1p_{0010})\frac{1}{16}\right]0.7320>p_1=0.6875,$$
thus showing that our cloning approach is more efficient than the previous one, which does not use cloning. We have not proven that the first approach is the most efficient among those that do not resort to cloning, but we conjecture that it is.
Besides this larger probability of obtaining the correct result, our cloning approach offers another advantage: the measurement of the โflagโ state allows us to be confident about having the correct result in a larger fraction of our attempts. For the probabilistic cloning machines described above this fraction was $`0.428`$, but this can be improved by choosing a different cloning machine, characterized by $`\gamma _1=0.3485,\gamma _2=0.5258`$. This latter machine signals a guaranteed correct result in a fraction $`(\gamma _1+2\gamma _2)/30.467`$ of the runs. The best no-cloning approach for obtaining these guaranteed correct results would involve unambiguous discrimination of the function $`f_0`$, followed by the distinction among the four possibilities for functions $`f_1`$ and $`f_2`$ (this second step is simple if we know $`f_0`$ for certain). Theorem 4 of provides us with a tool to numerically determine the best efficiency for unambiguous discrimination of $`f_0`$. A numerical search indicates that this can be done only with efficiency $`1/3`$, and therefore this is the limit for the fraction of runs for which we can obtain a guaranteed correct result for the task at hand, if we do not resort to cloning.
## 3 Conclusion
We have given two examples of tasks whose performance is enhanced by the use of quantum cloning. As we have discussed, cloning may offer advantages for a whole class of quantum computational tasks. Cloning need not be made only once during the course of a computation; nor does it necessarily need be one of the two kinds discussed above. For example, asymmetric cloning may also be useful, depending on the nature of the task at hand.
We must note that general quantum algorithms already manipulate quantum information, distributing it among different parts of the quantum register during a computation. What we have shown here is that quantum cloning can be taken as a natural quantum information processing tool to do this quantum information distribution, in order to optimize our use of computational resources. It would be interesting to find other tasks that could profit from cloning, perhaps by combining already known quantum algorithms with some intermediate cloning steps.
In this paper we have not discussed how the cloning circuit complexity scales with input size. Some authors have developed quantum circuits for deterministically cloning single qubits , and networks for state-dependent cloning . Further work on circuits for deterministically cloning $`d`$-dimensional systems ($`d>2`$) is still required.
## 4 Acknowledgments
We acknowledge support from the Royal Society, the ORS Award Scheme and the Brazilian agency Coordenaรงรฃo de Aperfeiรงoamento de Pessoal de Nรญvel Superior (CAPES). |
warning/0002/hep-th0002226.html | ar5iv | text | # 1 Introduction
## 1 Introduction
In these notes we study correlation functions in certain four-dimensional field theories by examining fluctuations about kink solutions of gauged 5-dimensional $`๐ฉ=8`$ supergravity . In light of the AdS/CFT correspondence , these gravitational backgrounds have been interpreted as the duals of $`๐ฉ=4`$ Super Yang-Mills theory either perturbed by various relevant operators , or given expectation values for some scalar fields . Two-, three- and four-point correlation functions in pure $`๐ฉ=4`$ SYM are by now well-understood , while in kink backgrounds a few two-point functions have been calculated for operators whose dual scalars are inert . Here we focus on the more complex case of field theory operators whose dual supergravity scalars are varying radially in the kink background, although we obtain some new results for inert scalars also.
Implicit in our use of the AdS/CFT correspondence is the assumption that 5D maximally supersymmetric gauged supergravity arises as a consistent truncation of ten-dimensional Type IIB supergravity on $`AdS_5\times S^5`$. In the case where the dual field theory is on the Coulomb branch, the โliftsโ of certain five-dimensional solutions to ten-dimensional type IIB supergravity are known and correspond to continuous distributions of D3-branes in the compact extra dimensions . Other ten-dimensional configurations of non-coincident D-branes have also been studied . However, we shall largely take a five-dimensional viewpoint.
Five-dimensional $`๐ฉ=8`$ gauged supergravity possesses 42 scalar fields $`\phi ^I`$, and a potential $`V(\phi )`$ which depends on 40 of them. The scalars fall into various representations of the $`SO(6)`$ gauge group: there is a $`20^{}`$ representation dual to the dimension 2 primary operators $`\mathrm{Tr}X^2`$, a $`10\overline{10}`$ representation corresponding to their dimension 3 descendents, and finally the singlet axion-dilaton, which does not enter into $`V(\phi )`$ and is dual to the $`๐ฉ=4`$ Lagrangian $`=\mathrm{Tr}F^2+\mathrm{}`$, and thus corresponds to the exactly marginal coupling $`\theta +4\pi i/g^2`$.
The part of the supergravity action essential for our concerns is
$`S={\displaystyle d^5x\sqrt{g}\left[\frac{1}{4}R+\frac{1}{2}g^{\mu \nu }_\mu \phi ^I_\nu \phi ^IV(\phi )\right]},`$ (1)
where the signature of the bulk metric is $`(+)`$. In these conventions the scalar fields are dimensionless. Tractable kink solutions typically have 1 or 2 โactiveโ scalars $`\phi ^I(r)`$, which depend on the radial coordinate $`r`$, while the remaining โinertโ scalars $`\phi ^\alpha `$ are constant (and typically vanish). The bulk geometry has an $`r`$-dependent scale factor, and the kink background takes the form
$`ds^2`$ $`=`$ $`e^{2A(r)}(\eta _{ij}dx^idx^j)dr^2,`$ (2)
$`\phi ^I`$ $`=`$ $`\phi ^I(r),`$
which is invariant under 4-dimensional Poincarรฉ transformations of the $`x^i,i=0,1,2,3`$. The geometries are asymptotically like $`AdS_5`$, i.e. $`A(r)r/L`$ as $`r\mathrm{}`$. In most cases there is a curvature singularity in the interior, i.e. $`A(r)\mathrm{}`$ as $`rr_s`$. The significance of the singularity will be briefly discussed below.
We are primarily interested in fluctuations and correlation functions of SYM operators $`๐ช_I(x)`$ dual to the bulk scalars $`\phi ^I(x)`$. Due to the reduced symmetry in the kink geometries, only two-point correlators are currently tractable to analysis. For inert scalars $`\phi ^\alpha `$, the procedure is straightforward in principle. One must obtain solutions of the equations of motion of the fluctuations $`\stackrel{~}{\phi }^\alpha (r,x)`$, governed by the action (1) expanded to second order in $`\stackrel{~}{\phi }^\alpha `$. Notice that the effective mass $`U(\phi ^I)`$ depends on the active scalars:
$`S_\alpha [\phi _\alpha ]={\displaystyle \frac{1}{2}}{\displaystyle }d^5x\sqrt{g}(g^{\mu \nu }_\mu \stackrel{~}{\phi }^\alpha _\nu \stackrel{~}{\phi }^\alpha U(\phi ^I(r))\stackrel{~}{\phi }^\alpha {}_{}{}^{2}).`$ (3)
The solution $`\stackrel{~}{\phi }^\alpha `$ must be chosen to vanish at the singularity, $`r=r_s`$, and obey Dirichlet boundary conditions $`\stackrel{~}{\phi }^\alpha (R,x)=\mathrm{\Phi }_\alpha (x)`$ at some large cutoff radius $`R`$. The action (3) is then interpreted as the generating functional for the field theory. Evaluated on the solutions to the equations of motion it reduces to a boundary term, and the correlation function is obtained from the $`R\mathrm{}`$ limit,
$`\stackrel{~}{S}_\alpha [\phi _\alpha ]={\displaystyle \frac{1}{2}}{\displaystyle d^4x\sqrt{g}g^{rr}\stackrel{~}{\phi }^\alpha (R,x)_r\stackrel{~}{\phi }^\alpha (R,x)},`$ (4)
as prescribed for two-point functions in pure $`AdS`$ by .
Since $`U(\phi ^I(r))`$ vanishes for the dilaton, the dilaton two-point function is usually the simplest to compute. Analyticity properties (in the momentum $`p^i`$ conjugate to $`x^i`$) reveal the spectrum of the boundary theory in the large $`N`$, large $`\lambda =g^2N`$ limit. For the various Coulomb branch flows of , one finds either a mass gap and continuous spectrum or an infinite discrete spectrum of poles. Since the supergravity description formally breaks down at the singularity, it is arguable whether these are actual features of the field theory dynamics or artifacts of a poor approximation. We have little to add to these arguments except to point out that a similar discrete spectrum has also been found for a distribution of D3-branes with non-singular metric . Some new examples of correlators of the dilaton and other inert scalars will be given below.
For active scalars the situation is more complex because the fluctuations $`\stackrel{~}{\phi }^I(r,x)`$ are coupled by the equations of motion to metric perturbations $`h_{ij}(r,x)`$, defined by
$`ds^2`$ $`=`$ $`e^{2A(r)}\left(\eta _{ij}+h_{ij}(r,x)\right)dx^idx^jdr^2,`$ (5)
$`\phi ^I`$ $`=`$ $`\overline{\phi }^I(r)+\stackrel{~}{\phi }^I(r,x).`$
It is especially interesting that the $`\stackrel{~}{\phi }_I`$ couple to trace components of $`h_{ij}`$, i.e.
$`h_{ij}={\displaystyle \frac{1}{4}}h(r,x)\eta _{ij}{\displaystyle \frac{}{x^i}}{\displaystyle \frac{}{x^j}}H(r,x).`$ (6)
This is expected if the flow $`\overline{g}_{\mu \nu }(r),\overline{\phi }^I(r)`$ describes a relevant deformation of $`๐ฉ=4`$ SYM. In this case the stress tensor satisfies the operator relation
$`T_i^i(x)={\displaystyle \underset{I}{}}\beta ^I๐ช_I(x),`$ (7)
where $`๐ช_I(x)`$ is the relevant operator dual to $`\phi ^I`$, and $`\beta ^I`$ is the beta-function for its coupling. Hence it is natural that the trace of $`h_{ij}`$, which is dual to $`T_i^i`$, should not be independent of the $`\phi ^I`$. One might hope that the present line of investigation interfaces with more general studies of holographic RG flows and perhaps with ideas concerning the c-theorem, either holographic or field-theoretic .
On the other hand for a flow describing the Coulomb branch of $`๐ฉ=4`$ SYM, one expects that $`T_i^i(x)=0`$, since the trace vanishes as an operator if the Lagrangian contains only marginal couplings. One then hopes that despite the fact that $`\stackrel{~}{\phi }^I`$ and $`h_i^i`$ are coupled, $`h_i^i`$ will not excite $`T_i^i`$. Since it is well-known that scalars dual to vevs fall off more rapidly on the boundary than those dual to operator perturbations, it is reasonable to expect that the excitation of $`h_i^i`$ caused by a fluctuation $`\stackrel{~}{\phi }^I`$ scales away too quickly to excite an operator in the field theory.
In the next section we derive the coupled linear fluctuation equations of $`\stackrel{~}{\phi }^I`$, $`h`$, and $`H`$. Other components of $`h_{ij}`$ decouple consistently; the transverse traceless modes are known each to obey the same wave equation as the dilaton, and the remaining components can be set to zero by a choice of coordinates. We discuss a pure diffeomorphism solution to the equations, the universal solution. A more detailed analysis of the coupled fluctuation equations is then undertaken. In the case where the background flow involves a single active scalar $`\overline{\phi }(r)`$, we derive an uncoupled third order differential equation for its fluctuation $`\stackrel{~}{\phi }(r,x)`$. The universal solution allows us to apply reduction of order methods to study this equation. It is also instructive to convert the fluctuation equations to equivalent Schrรถdinger form where positivity properties of supersymmetric quantum mechanics can be applied. Indeed it turns out that the role of SUSY QM is ubiquitous.
In the following two sections we solve the fluctuation equations in detail for two examples of RG flows in which the background is explicitly known. In section 3 we treat the flow obtained in , which was interpreted as a relevant perturbation of $`๐ฉ=4`$ SYM leading to pure $`๐ฉ=1`$ SYM theory in the infrared. In section 4, we discuss the Coulomb branch flow of corresponding to a distribution of D3-branes in a 2-dimensional disc. In both cases we discuss the two-point correlators of the dilaton<sup>1</sup><sup>1</sup>1The dilaton correlator for the Coulomb branch flow was calculated in , while for the N=1 flow it was obtained in the very recent which appeared during the preparation of this manuscript. and of a second inert scalar, and obtain a consistent picture of the spectrum of excitations in the boundary field theory. The fluctuation equation for active scalars is reduced to a hypergeometric equation and solved in both cases, and the metric component $`h(r,x)`$ is also determined.
At this point we attempt to determine the two-point correlation function of the active scalars from (4) by the standard procedure of identifying the most singular term as $`R\mathrm{}`$ which is non-analytic in the momentum $`p_i`$ conjugate to $`x^i`$, but encounter difficulties. For the $`๐ฉ=1`$ flow, the standard procedure does not have the expected structure of leading non-analytic term plus more singular polynomials in $`p^2`$. This remains true even if an integration constant initially chosen so that the fluctuation vanishes at the curvature singularity is allowed to vary. For the Coulomb branch flow, it appears that if the integration constant is chosen so the active scalar fluctuation is regular, the metric perturbation $`h_i^i`$ does not vanish on the boundary as the argument above for vanishing $`T_i^i`$ suggests. If the integration constant is chosen to make $`h_i^i`$ behave as expected on the boundary, then we can extract a correlator which is physically reasonable except for $`p^2=0`$ poles. However, regularity of the fluctuations in the interior then fails.
It is conceivable that the interior curvature singularity is the source of the difficulties above, but it is also possible that the standard procedure used to calculate correlators must be modified in the coupled active scalar/graviton sector. We discuss our attempts to do this in Sec 5. This includes an evaluation of the on-shell action (1), which reduces to boundary terms linear and quadratic in fluctuations. Expected supplementary boundary terms in the gravitational action are added. An issue of diffeomorphism invariance is also discussed. None of this leads to a resolution of the problem, which is left for future work, perhaps by people who can approach it with fresh energy and ideas.
It is worthwhile to point out two byproducts of our analysis of the fluctuation equations. First, the boundary values of the metric and active scalar fluctuations are not independent, as shown in equation (29). It is not clear that this constraint has been incorporated in a recent Hamilton-Jacobi analysis of holographic RG flows . Second, the SUSY QM analysis of the final uncoupled active scalar equations appears to be applicable to the stability of domain walls in the brane-world scenario .
## 2 The coupled graviton/scalar system
We first review the general structure of kink solutions of 5D $`๐ฉ=8`$ gauged supergravity which preserve partial supersymmetry. We then discuss the equations obeyed by linear fluctuations of the kink backgrounds with particular attention to the equations which couple fluctuations $`\stackrel{~}{\phi }^I(r,x)`$ of the active scalars $`\phi ^I(r)`$. Fluctuations of inert scalars and transverse traceless metric fluctuations satisfy uncoupled equations of a simpler structure.
### 2.1 Background Flows
The scalar field manifold of the supergravity theory can be described as the coset $`E_{6(6)}/USp(8)`$. The potential $`V(\phi )`$ is invariant under $`SO(6)\times SL(2,R)E_{6(6)}`$, where $`SO(6)SU(4)`$ is the gauge group and the $`SL(2,R)`$ factor corresponds to the axion-dilaton. The potential is a very complicated function of 40 scalars, and progress has come from restrictions to smaller sets of fields which preserve some subgroup $`HSO(6)`$. Symmetry properties insure that such restrictions are consistent with the full dynamics . The full theory also has a complicated coset metric $`G_{IJ}(\phi )`$ which simplifies to $`\delta _{IJ}`$ on the restricted scalar subspaces. This metric was therefore omitted in (1).
The critical points of $`V(\phi )`$ are significant for the problem of RG flows, and those extrema which preserve at least an $`SU(2)`$ subgroup of $`SO(6)`$ have been classified . The potential takes negative values at these extrema, so the solutions of the equations include exact anti-de Sitter geometries with scalars fixed at their critical values. In particular there is a critical point with full $`SO(6)`$ symmetry at the origin of field space, and the corresponding $`AdS_5`$ solution with $`SO(4,2)`$ isometry group is the well known holographic dual of the conformal phase of $`๐ฉ=4`$, d=4 SYM theory. Kink solutions of the form (2) are topologically like $`AdS_5`$ but have a smaller isometry group. As $`r`$ approaches the boundary at $`r=\mathrm{}`$, the kink metric approaches that of the maximally symmetric $`AdS_5`$ solution, and scalars vanish at rates determined by the scale dimensions of their operator duals in the SYM theory. As $`r`$ decreases toward the interior of the geometry certain known kinks flow toward a second critical point of the potential with more negative cosmological constant. The field theory interpretation is that of $`๐ฉ=4`$ SYM perturbed with a relevant operator such that the perturbed theory has a non-trivial IR fixed point. It would be very desirable to study the fluctuations about such flows, but this cannot be done at present because the explicit form of these kinks is not yet known<sup>2</sup><sup>2</sup>2Very recently the 10D lift of the fixed point itself has been constructed by .. All explicitly known flows have a curvature singularity at some finite $`r=r_s`$. This is associated with a breakdown of the supergravity description. The standard prescription is to impose the boundary condition that fluctuations vanish at the singularity and to proceed to extract correlation functions from the on-shell action (4). This procedure leads to problematic results as we will see, but we do not yet have a better alternative to propose.
All explicitly known flow solutions are supersymmetric in the sense that there are associated Killing spinors in the supergravity theory. It is known that the potential $`V(\phi )`$ can be derived from a superpotential $`W(\phi )`$,
$`V(\phi )={\displaystyle \frac{g^2}{8}}{\displaystyle \underset{I}{}}\left({\displaystyle \frac{W(\phi )}{\phi ^I}}{\displaystyle \frac{W(\phi )}{\phi ^I}}\right){\displaystyle \frac{g^2}{3}}W(\phi )^2.`$ (8)
on restricted subsets of the field space which include the active scalars $`\phi ^I(r)`$. Here $`g`$ is the $`SU(4)`$ gauge coupling of the 5D supergravity theory. It has dimensions of 1/length, and it is related to the length scale $`L`$ of the boundary $`AdS_5`$ space and the cosmological constant $`\mathrm{\Lambda }`$ by
$`g=2/L,\mathrm{\Lambda }=12/L^2=4V(\phi =0).`$ (9)
The main simplification of supersymmetric flows is that any solution of the first order Killing spinor conditions
$`{\displaystyle \frac{dA(r)}{dr}}={\displaystyle \frac{g}{3}}W(\phi ),{\displaystyle \frac{\phi ^I(r)}{r}}={\displaystyle \frac{g}{2}}{\displaystyle \frac{W(\phi )}{\phi ^I}},`$ (10)
automatically gives a Poincarรฉ-invariant kink solution of the more complicated second order field equations of the action (1).
It is worth noting that non-supersymmetric solutions of the form (2) can also be obtained by solving (10), with a choice of $`W(\phi )`$ other than the true superpotential of the theory. It was conjectured in that all solutions of the field equations of the form (2) can be obtained this way. Recently this was explained in terms of Hamilton-Jacobi theory . However, our examples will be flows where $`W(\phi )`$ is in fact the true superpotential.
### 2.2 Fluctuation equations
We now consider fluctuations of the scalars and the metric around such a background geometry. We use the freedom to choose coordinates to fix axial gauge,
$`h_{\mu 5}=0,`$ (11)
where $`\mu =0,1,2,3,5`$. With this choice the metric, including fluctuations, has the form
$`ds^2`$ $`=`$ $`e^{2A(r)}\left((\eta _{ij}+h_{ij}(x,r))dx^idx^j\right)dr^2.`$ (12)
There is still gauge freedom remaining. We find two classes of residual diffeomorphisms $`\delta g_{\mu \nu }=_\mu ฯต_\nu +_\nu ฯต_\mu `$ preserving axial gauge (11). There are four-dimensional diffeomorphisms
$`ฯต_i`$ $`=`$ $`e^{2A(r)}\omega _i(x),ฯต_5=0,`$ (13)
$`\delta h_{ij}`$ $`=`$ $`_i\omega _j(x)+_j\omega _i(x).`$
Additionally, there is the transformation
$`ฯต_5`$ $`=`$ $`ฯต_5(x),ฯต_i=(_iฯต_5(x))\left(e^{2A(r)}{\displaystyle ^r}๐r^{}e^{2A(r^{})}\right),`$ (14)
$`\delta h_{ij}`$ $`=`$ $`2A^{}(r)\eta _{ij}ฯต_5(x)2(_i_jฯต_5(x))\left({\displaystyle ^r}๐r^{}e^{2A(r)}\right).`$
Here we are viewing the result of an infinitesimal coordinate transformation on the background as a metric fluctuation.
Both active and inert scalars can fluctuate around their backgrounds: $`\phi ^I(x,r)=\overline{\phi }^I(r)+\stackrel{~}{\phi }^I(x,r)`$. For ease of notation, we shall in the future drop the bar and use $`\phi (r)`$ to denote the classical background. When we discuss inert scalars, whose background values are zero, we will drop the tilde on the fluctuation.
Einsteinโs equations for linear fluctuations were derived in as equations relating the first order Ricci tensor to its scalar source. These equations are
$`\mathrm{`}\mathrm{`}R_{ij}^{(1)}\text{}=e^{2A}\left({\displaystyle \frac{1}{2}}_r^2+2A^{}_r\right)h_{ij}+{\displaystyle \frac{1}{2}}\eta _{ij}e^{2A}A^{}_r(\eta ^{kl}h_{kl}){\displaystyle \frac{1}{2}}\text{ }\text{ }\text{ }\text{ }h_{ij}`$ (15)
$`{\displaystyle \frac{1}{2}}\eta ^{kl}\left(_i_jh_{kl}_i_kh_{jl}_j_kh_{il}\right)={\displaystyle \frac{4}{3}}e^{2A}{\displaystyle \frac{V(\phi )}{\phi ^I}}\stackrel{~}{\phi ^I}\eta _{ij},`$
$`R_{55}^{(1)}={\displaystyle \frac{1}{2}}(_r^2+2A^{}_r)\eta ^{kl}h_{kl}=4\phi ^I{}_{}{}^{}\stackrel{~}{\phi }_{}^{I}{}_{}{}^{}+{\displaystyle \frac{4}{3}}{\displaystyle \frac{V(\phi )}{\phi ^I}}\stackrel{~}{\phi ^I},`$ (16)
$`R_{j5}^{(1)}={\displaystyle \frac{1}{2}}\eta ^{kl}_r(_kh_{jl}_jh_{kl})=2\phi ^I_j\stackrel{~}{\phi }^I.`$ (17)
The notation $`\mathrm{`}\mathrm{`}R_{ij}^{(1)}`$โ indicates that a simplification of the actual $`R_{ij}^{(1)}`$ equation has been made in (15) as in . The scalar fluctuation equation is
$`e^{2A}\text{ }\text{ }\text{ }\text{ }\stackrel{~}{\phi }^I\stackrel{~}{\phi }^I{}_{}{}^{\prime \prime }4A^{}\stackrel{~}{\phi }^I{}_{}{}^{}+{\displaystyle \frac{^2V(\phi )}{\phi ^I\phi ^J}}\stackrel{~}{\phi }^J={\displaystyle \frac{1}{2}}\phi ^I{}_{}{}^{}\eta _{}^{ij}h_{ij}^{}.`$ (18)
A prime above denotes $`_r`$. In general primes on fields $`h_{ij},\stackrel{~}{\phi }`$ and backgrounds $`A,\phi `$ will denote a derivative with respect to the radial coordinate, while primes on $`V(\phi )`$ or $`W(\phi )`$ refer to $`\phi `$-derivatives.
For inert scalars $`\stackrel{~}{\phi }^\alpha `$, the right-hand side of (18) vanishes, and consequently they do not couple to the graviton. Any coupling to the active scalars due to the potential term must vanish, otherwise the inert scalars could not have been zero in the background:
$`{\displaystyle \frac{^2V(\phi )}{\phi ^\alpha \phi ^I}}=0.`$ (19)
Thus they satisfy decoupled second-order equations, several of which we consider explicitly in sections 3 and 4.
It is clear from the equations above that active scalar fluctuations couple to the graviton. To simplify the coupled system we consider an arbitrary graviton fluctuation, decomposed in a complete momentum-space basis. We use the vectors
$`p^i=(p,0,0,0),\epsilon ^0=(0,0,0,1),\epsilon ^\pm =(0,1,\pm i,0)/\sqrt{2},`$ (20)
and the decomposition
$`h_{ij}(r,p)`$ $`=`$ $`\epsilon _i^+\epsilon _j^+h^{++}(r,p)+\epsilon _i^{}\epsilon _j^{}h^{}(r,p)+(\epsilon _i^+\epsilon _j^0+\epsilon _i^0\epsilon _j^+)h^{+0}(r,p)+`$
$`(\epsilon _i^{}\epsilon _j^0+\epsilon _i^0\epsilon _j^{})h^0(r,p)+(\epsilon _i^+\epsilon _j^{}+\epsilon _i^{}\epsilon _j^+2\epsilon _i^0\epsilon _j^0)h^{00}(r,p)+`$
$`(p_i\epsilon _j^++\epsilon _i^+p_j)a^+(r,p)+(p_i\epsilon _j^{}+\epsilon _i^{}p_j)a^{}(r,p)+(p_i\epsilon _j^0+\epsilon _i^0p_j)a^0(r,p)+`$
$`\eta _{ij}h(r,p)/4+p_ip_jH(r,p).`$
In this decomposition, the five $`h^{xx}`$ are transverse traceless, the $`a^x`$ are traceless and longitudinal, and $`h`$ and $`H`$ are the trace components.
The $`h^{xx}`$ contribute only in (15), and it can be shown that each $`h^{xx}`$ satisfies the same uncoupled equation as a free massless scalar in the background (2), as is well-known . The $`h^{xx}`$ modes are the expected five physical components of a graviton in five bulk dimensions, and they have the same fluctuation spectrum as the inert dilaton. This can be obtained from an equivalent Schrรถdinger equation with supersymmetric potential . These modes do not couple to the active scalars.
Further, we can examine equation (17) and see that the $`a^x`$ must be $`r`$-independent. One may then gauge them to zero using the residual coordinate invariance (13). One is still free to use the transformation (14).
We are left with the components $`h`$ and $`H`$, which do couple to active scalar fluctuations. Before we present the equations of this system, consider the action of the remaining gauge freedom (14). A transformation parameterized by $`ฯต_5(x)=e^{ipx}E(p)`$ modifies the fields by
$`\delta h=8A^{}(r)E(p),\delta H^{}=2e^{2A(r)}E(p),\delta \stackrel{~}{\phi }=\phi ^{}(r)E(p).`$ (22)
One may check that (22) solves the equations of motion (1518). We will refer to this pure gauge solution as the universal solution. It will have a vital technical role in our study of the coupled system.
The fields that remain in our reduced system are
$`h_{ij}(r,x)=e^{ipx}\left({\displaystyle \frac{1}{4}}h(r,p)\eta _{ij}+p_ip_jH(r,p)\right),\stackrel{~}{\phi }^I(r,x)=e^{ipx}\stackrel{~}{\phi }^I(r,p).`$ (23)
and we now begin the process of simplifying the coupled equations which they satisfy.
Substituting the ansatz (23) into the $`R_{j5}`$ equation (17), $`H`$ drops out and we are left with
$`h^{}(r)={\displaystyle \frac{16}{3}}\phi ^I{}_{}{}^{}\stackrel{~}{\phi }_{}^{I},`$ (24)
where we have canceled a uniform factor of $`p_j`$.
We now consider linear combinations of the equations (15) traced with the tensors $`\eta ^{ij}`$ and $`p^ip^j`$. It turns out that scalars decouple in the difference of the $`\eta ^{ij}`$ trace and the $`p^ip^j`$ trace. Dividing by $`p^2`$ leaves a simple equation relating $`H`$ and $`h`$:
$`H^{\prime \prime }+4A^{}H^{}={\displaystyle \frac{1}{2}}e^{2A}h.`$ (25)
One can also show that an independent linear combination of the two traces is trivial given the $`R_{j5}`$ equation (24), so the only new relation obtained is (25).
Substituting the ansatz (23) into the $`R_{55}`$ equation (16), we obtain
$`{\displaystyle \frac{1}{2}}(h^{\prime \prime }+2A^{}h^{}+p^2H^{\prime \prime }+2p^2A^{}H^{})=4\phi ^I{}_{}{}^{}\stackrel{~}{\phi }_{}^{I}{}_{}{}^{}+{\displaystyle \frac{4}{3}}_IV(\phi )\stackrel{~}{\phi }^I.`$ (26)
Combining (25) and (26), we can eliminate $`H^{\prime \prime }`$ and obtain an algebraic relation for $`H^{}`$. Further removing derivatives of $`h`$ in favor of $`\stackrel{~}{\phi }^I`$ using (24) and substituting $`\phi `$-derivatives of $`W`$ for $`\phi ^{}`$ and $`V^{}(\phi )`$, we obtain:
$`2A^{}H^{}={\displaystyle \frac{1}{2}}e^{2A}h+{\displaystyle \frac{2g_IW(\phi )}{3p^2}}\left[2\stackrel{~}{\phi }^I{}_{}{}^{}g_I_JW(\phi )\stackrel{~}{\phi }^J\right].`$ (27)
The equations (24), (25), (27), together with the Klein-Gordon equation (18), constitute $`n+3`$ equations for the $`n+2`$ fields $`h,H`$ and $`\stackrel{~}{\phi }^I`$. However, one can show that the gravitational equations are related by the Bianchi identity: the expected combination of derivatives of (24), (25), (27) vanishes identically if the Klein-Gordon equation is satisfied.
It also turns out that we can construct an algebraic equation for $`h`$ which will be of later use. The Klein-Gordon equation (18) becomes
$`p^2e^{2A}\stackrel{~}{\phi }^I\stackrel{~}{\phi }^I{}_{}{}^{\prime \prime }4A^{}\stackrel{~}{\phi }^I{}_{}{}^{}+{\displaystyle \frac{^2V(\phi )}{\phi ^I\phi ^J}}\stackrel{~}{\phi }^J={\displaystyle \frac{1}{2}}\phi ^I{}_{}{}^{}(h^{}+p^2H^{}).`$ (28)
The left-hand side is just derivatives of $`\stackrel{~}{\phi }^I`$. On the right, $`h^{}`$ can be written in terms of $`\stackrel{~}{\phi }^I`$ using (24), while we can eliminate $`H^{}`$ in favor of $`\stackrel{~}{\phi }`$ and a factor of $`h`$ using (27). Thus $`h`$ is determined solely by $`\stackrel{~}{\phi }^I`$ and derivatives. To avoid tedious notation, we write only the case of a single active scalar, although the general case is straightforward:
$`3p^2e^{2A}h={\displaystyle \frac{16W}{W^{}}}\stackrel{~}{\phi }^{\prime \prime }+g\left({\displaystyle \frac{64W^2}{3W^{}}}+8W^{}\right)\stackrel{~}{\phi }^{}+`$ (29)
$`4g^2\left({\displaystyle \frac{WW^{\prime \prime 2}}{W^{}}}+WW^{\prime \prime \prime }(\stackrel{~}{\phi }){\displaystyle \frac{8W^2W^{\prime \prime }}{3W^{}}}W^{}W^{\prime \prime }{\displaystyle \frac{4p^2W}{g^2W^{}}}e^{2A}\right)\stackrel{~}{\phi }.`$
This equation shows that the boundary values of $`h(r,p)`$ and $`\stackrel{~}{\phi }(r,p)`$ are not independent, a fact which may have implications for the formulation of Hamilton-Jacobi dynamics proposed in .
Unfortunately our equations are still coupled. The next step is to derive a third order equation involving only the $`\stackrel{~}{\phi }^I`$. In the case of flows with only a single active scalar, this becomes an uncoupled equation which is the key to the present exploration and can be solved for some specific flows.
The Klein-Gordon equation (18) has the conventional form of a scalar field in the curved background, with an additional source term involving the graviton fluctuation. We can express the source in terms of scalars by noting that the $`R_{55}`$ equation (16) can be written
$`{\displaystyle \frac{1}{2}}e^{2A}_r(e^{2A}\eta ^{ij}h_{ij}^{})=4\phi ^I{}_{}{}^{}\stackrel{~}{\phi }_{}^{I}{}_{}{}^{}+{\displaystyle \frac{4}{3}}{\displaystyle \frac{V(\phi )}{\phi ^I}}\stackrel{~}{\phi }^I,`$ (30)
which can be integrated to obtain
$`\eta ^{ij}h_{ij}^{}=8e^{2A}{\displaystyle ^r}dr^{}e^{2A(r^{})}\left(\phi ^I{}_{}{}^{}\stackrel{~}{\phi }_{}^{I}{}_{}{}^{}+{\displaystyle \frac{1}{3}}_IV(\phi )\stackrel{~}{\phi }^I\right).`$ (31)
Substituting into (18) and taking an $`r`$-derivative, we find the third-order equation
$`_r\left\{{\displaystyle \frac{\phi ^I^{}}{|\phi ^K{}_{}{}^{}|_{}^{2}}}e^{2A}\left(\left(_r^2+4A^{}_re^{2A}\text{ }\text{ }\text{ }\text{ }\right)\delta _{IJ}_I_JV(\phi )\right)\stackrel{~}{\phi }^J\right\}=`$ (32)
$`4e^{2A}\left(\phi ^J{}_{}{}^{}\stackrel{~}{\phi }_{}^{J}{}_{}{}^{}+{\displaystyle \frac{1}{3}}_JV(\phi )\stackrel{~}{\phi }^J\right).`$
Although a simplification, (32) still couples all active scalar fluctuations. Since several known flows involve only one active scalar, we specialize to the case of a single active scalar $`\phi `$. Then (32) reduces to an uncoupled equation,
$`_r\left\{{\displaystyle \frac{1}{\phi ^{}}}e^{2A}\left(_r^2+4A^{}_rV^{\prime \prime }(\phi )e^{2A}\text{ }\text{ }\text{ }\text{ }\right)\stackrel{~}{\phi }\right\}=4e^{2A}\left(\phi ^{}\stackrel{~}{\phi }^{}+{\displaystyle \frac{1}{3}}V^{}(\phi )\stackrel{~}{\phi }\right).`$ (33)
In the case of a single active scalar, it is also possible to combine (24, 25, 27) to obtain an uncoupled third order equation for $`h`$, but we have chosen to work with (33) in applications to specific flows.
One may easily verify that the universal solution, $`\stackrel{~}{\phi }W^{}(\phi )e^{ipx}`$, satisfies (33) by using (8) and (10) to relate $`V(\phi )`$, $`\phi ^{}(r)`$ and $`A^{}(r)`$ to the superpotential $`W(\phi )`$ and its derivatives. This allows us to use the method of reduction of order to obtain a second-order equation for the remaining solutions. Writing
$`\stackrel{~}{\phi }=W^{}(\phi )e^{ipx}{\displaystyle ^r}๐r^{}R(r^{},p),`$ (34)
for some unknown $`R(r,p)`$, and using the properties of the flow (8) and (10), we find
$`R^{\prime \prime }(r,p)+g(W^{\prime \prime }2W)R^{}(r,p)`$ $`+`$ (35)
$`g^2\left({\displaystyle \frac{2}{3}}WW^{\prime \prime }+{\displaystyle \frac{8}{9}}W^2+{\displaystyle \frac{1}{2}}W^{}W^{\prime \prime \prime }W^2\right)R(r,p)`$ $`+`$ $`p^2e^{2A}R(r,p)=0.`$
We remind the reader that primes on $`R(r)`$ are $`r`$-derivatives while those on $`W`$ are with respect to $`\phi `$. To determine the boundary scaling rates of the three independent solutions of (33), we may approximate $`W(\phi )=3/2+\rho \phi ^2/2`$ where $`\rho `$ is related to the operator scale dimension $`\mathrm{\Delta }`$ by $`\rho =\mathrm{\Delta }4`$ for operator deformations and $`\rho =\mathrm{\Delta }`$ for vacuum expectation values. Standard Frobenius analysis then gives the exponential rates
$`\stackrel{~}{\phi }(r)\mathrm{exp}(\rho r/L),\mathrm{exp}((\rho 2)r/L),\mathrm{exp}((\rho 4)r/L),`$ (36)
for the universal solution, and the two independent solutions of (35), respectively.
By solving (35), we can obtain the solution to $`\stackrel{~}{\phi }`$ from (34). The integration constant is just the freedom to add a multiple of the universal solution to our result for $`\stackrel{~}{\phi }`$. This will be important later on for assuring regularity of the solutions at the singularity.
### 2.3 Supersymmetric Quantum Mechanics
We now discuss the transformation of the fluctuation equations into Schrรถdinger form, which gives considerable intuition into the nature of the fluctuation spectrum. For the dilaton and for the transverse traceless metric fluctuations the Schrรถdinger potential is that of a supersymmetric quantum mechanics, and this gives a simple proof that the spectrum of normalizable fluctuations contains no tachyons. We will see that SUSY-QM provides a useful framework in which to consider all fluctuations.
The Schrรถdinger form obtains after a change of radial coordinate $`r`$ to the horospherical coordinate $`z`$, defined by the line element
$`ds^2=e^{2A(z)}\left[(\eta _{ij}+h_{ij}(x,z))dx^idx^jdz^2\right],`$ (37)
and thus the Jacobian must be
$`{\displaystyle \frac{dz}{dr}}=\pm e^A.`$ (38)
We choose the $``$ sign for the present application, so that the boundary region $`r\mathrm{}`$ is mapped to $`z=0`$ and the deep interior to $`z\mathrm{}`$.
Consider first the fluctuation equation for an inert scalar $`\stackrel{~}{\phi }(r,p)`$ obtained from (3),
$`\left(\text{ }\text{ }\text{ }\text{ }+U(\phi )\right)\stackrel{~}{\phi }(r,p)=\stackrel{~}{\phi }^{\prime \prime }(r,p)+4A^{}(r)\stackrel{~}{\phi }^{}(r,p)U(\phi (r))\stackrel{~}{\phi }(r,p)+e^{2A}p^2\stackrel{~}{\phi }(r,p)=0.`$ (39)
Performing the field transformation $`\stackrel{~}{\phi }(r,p)=\mathrm{exp}(3A/2)\psi (z,p)`$ as well as the change of coordinate $`rz`$, one finds an equation in Schrรถdinger form
$`\psi ^{\prime \prime }(z,p)+๐ฑ(z)\psi (z,p)=p^2\psi (z,p),`$ (40)
with potential
$`๐ฑ(z)=\left({\displaystyle \frac{3}{2}}A^{}(z)\right)^2+{\displaystyle \frac{3}{2}}A^{\prime \prime }(z)+e^{2A(z)}U(\phi (z)),`$ (41)
in which $`A^{}(z)=dA/dz=\pm \mathrm{exp}(A)A^{}(r)`$. (Note that either sign choice in (38) leads to the same form for (40) and (41).) The first two terms in (41) are in the form of a supersymmetric quantum mechanics potential $`๐ฑ(z)=๐ฐ(z)^2+๐ฐ^{}(z)`$ derived from a prepotential $`๐ฐ(z)=(3/2)A^{}(z)`$, while $`U(\phi (z))`$ is determined by the coupling of inert and active scalars in the supergravity potential. If $`U(\phi (z))`$ is positive (or absent as it is for the dilaton fluctuation), then one has an immediate argument that normalizable fluctuations occur only for $`p^20`$. In general, the behavior of the potential (41) near the limits $`z0`$ and $`z\mathrm{}`$ is usually enough information to ascertain whether the fluctuation spectrum is discrete or continuous, with or without gap .
In our examples in sections 3 and 4, the two inert fluctuations with nonzero $`U(\phi (z))`$ turn out to have $`U(\phi (z))<0`$, so that positivity of $`p^2`$ is not obvious with the potential in the form (41). However in both cases we are able to rewrite the potentials in an exact SUSY QM form with a modified $`๐ฐ`$. The two examples work differently, and we leave the question of the existence of an exact SUSY QM potential for general coupled inert scalars for the future.
The norm is a potentially delicate issue. The Hamiltonian in (40) is self-adjoint with respect to the โSchrรถdinger normโ $`๐z\psi (z)^2`$, but not with respect to the transformed covariant norm $`๐z\mathrm{exp}(2A)\psi (z)^2`$ which is correct in the present setting. (This results because a factor of $`\mathrm{exp}(2A)`$ was dropped in passing from (39) to (40).) Because of this delicacy, solutions which have infinite Schrรถdinger norm but finite covariant norm would not necessarily have $`p^2>0`$.
The treatment of active scalar fluctuations is more complex. The goal is to transform the second order equation (35) to Schrรถdinger form. The change from $`r`$ to $`z`$ produces
$`R^{\prime \prime }(z)+ge^A\left(W^{\prime \prime }{\displaystyle \frac{5}{3}}W\right)R^{}(z)+g^2e^{2A}\left(W^2+{\displaystyle \frac{2}{3}}WW^{\prime \prime }{\displaystyle \frac{8}{9}}W^2{\displaystyle \frac{1}{2}}W^{}W^{\prime \prime \prime }\right)R=p^2R,`$ (42)
in which $`R^{}`$ and $`R^{\prime \prime }`$ are derivatives with respect to $`z`$, but $`W^{}`$, $`W^{\prime \prime }`$ etc. indicate $`\phi `$-derivatives, as has been our practice. The further transformation $`R(z)=\mathrm{exp}(S(z))\psi (z)`$, with
$`{\displaystyle \frac{dS}{dz}}={\displaystyle \frac{1}{2}}ge^{A(z)}\left(W^{\prime \prime }(\phi (z)){\displaystyle \frac{5}{3}}W(\phi (z))\right),`$ (43)
produces the Schrรถdinger form (40) with the (ugly) potential
$`๐ฑ(z)={\displaystyle \frac{1}{4}}g^2e^{2A(z)}\left[{\displaystyle \frac{1}{3}}W^2+{\displaystyle \frac{7}{3}}W^2+(W^{\prime \prime })^2{\displaystyle \frac{4}{3}}WW^{\prime \prime }W^{}W^{\prime \prime \prime }\right].`$ (44)
Unpromising as it seems, one may try to express (44) in SUSY-QM form perhaps with a simple remainder. To do this we use the ansatz
$`๐ฐ=e^A(aW+bW^{}+cW^{\prime \prime }),`$ (45)
where $`a,b,c`$ are free parameters. We find that an exact SUSY form cannot be achieved, but analogously to ($`\text{41})`$ one may write
$`๐ฑ(z)=๐ฐ(z)^2+๐ฐ^{}(z)+{\displaystyle \frac{1}{3}}g^2e^{2A(z)}(W^{})^2,`$ (46)
with prepotential
$`๐ฐ(z)={\displaystyle \frac{1}{2}}ge^{A(z)}(W^{\prime \prime }W).`$ (47)
The discussion below (41) applies here as well, but here it is manifest that the additional term beyond the SUSY-QM structure $`\mathrm{\Delta }๐ฑ=(1/3)e^{2A}(W^{})^2`$ is positive-definite. This indicates that there are no tachyons in the normalized fluctuation spectrum of the active scalar. In particular examples one must be careful to use the correct norm which is the transformation of $`๐re^{4A(r)}\stackrel{~}{\phi }^2`$ with $`\stackrel{~}{\phi }(r,p)`$ given by (34).
Two further comments can be made. Although in general $`๐ฑ(z)`$ only can be cast in the form (46), in specific examples it may well have an exact SUSY form with a prepotential different from (47). An example of this appears in section 4.3. Additionally, it is also clear that the analysis above has applications to the stability of domain walls in the brane world scenario which we plan to pursue.
## 3 The $`๐ฉ=1`$ Super Yang-Mills flow
Several kink solutions involving a single real flowing scalar have been considered in the literature. There are only two inequivalent single real scalars transforming in the $`\mathrm{๐๐}\overline{\mathrm{๐๐}}`$ of $`SU(4)`$, as there are two Weyl orbits in the $`\mathrm{๐๐}`$. Representative scalars in the longer and shorter orbits, respectively, have been called $`\sigma `$ and $`m`$, and we will use this convention. There are three inequivalent choices in the $`\mathrm{๐๐}^{}`$, one of which we will discuss in the next section.
Interesting renormalization group flows have been considered in in which the scalar manifold of bulk $`๐ฉ=8`$ gauged supergravity was restricted to a subspace of two particular inequivalent scalars in the $`\mathrm{๐๐}\overline{\mathrm{๐๐}}`$ of $`SU(4)`$. A kink solution with only $`m(r)`$ active was obtained and interpreted as the holographic dual of a perturbation of $`๐ฉ=4`$ SYM theory by a dimension 3 operator which leads to a field theoretic RG flow to pure $`๐ฉ=1`$ SYM theory at long distances. The 10-dimensional lift of this solution is not known, and the 5-dimensional geometry has an interior curvature singularity. The significance of this is not clear, but a scenario for its stringy resolution has been outlined .
We will use the kink background of as a testing ground for computations of fluctuations and correlation functions as discussed above. For this purpose we review the background solution in the next subsection. We then go on to study fluctuations and correlators of two inert scalars, and finally focus on perturbations $`\stackrel{~}{m}(r,p)`$ of the active scalar.
### 3.1 The background flow
In the consistent sub-sector of $`๐ฉ=8`$ supergravity considered in , the two scalar fields $`m`$ and $`\sigma `$ have canonical kinetic terms and superpotential and potential given by:
$`W(m,\sigma )`$ $`=`$ $`{\displaystyle \frac{3}{4}}\left[\mathrm{cosh}\left({\displaystyle \frac{2m}{\sqrt{3}}}\right)+\mathrm{cosh}\left(2\sigma \right)\right],`$ (48)
$`V(m,\sigma )`$ $`=`$ $`{\displaystyle \frac{3g^2}{8}}\left[{\displaystyle \frac{1}{4}}\mathrm{cosh}^2\left({\displaystyle \frac{2m}{\sqrt{3}}}\right)+\mathrm{cosh}\left({\displaystyle \frac{2m}{\sqrt{3}}}\right)\mathrm{cosh}(2\sigma ){\displaystyle \frac{1}{4}}\mathrm{cosh}^2(2\sigma )+1\right].`$ (49)
Since the $`\sigma `$ field appears quadratically in $`V(m,\sigma )`$ there is a supersymmetric flow in which it vanishes. This flow is a solution of (10) with superpotential
$`W(m)={\displaystyle \frac{3}{4}}\left[1+\mathrm{cosh}\left({\displaystyle \frac{2m}{\sqrt{3}}}\right)\right].`$ (50)
In the following, we set $`g=2/L`$, and $`W`$ always refers to $`W(m)`$. The kink solution is
$`m(r)={\displaystyle \frac{\sqrt{3}}{2}}\mathrm{log}{\displaystyle \frac{1+e^{r/L}}{1e^{r/L}}},`$ (51)
$`A(r)={\displaystyle \frac{1}{2}}\left({\displaystyle \frac{r}{L}}+\mathrm{log}2\mathrm{sinh}{\displaystyle \frac{r}{L}}\right).`$ (52)
There is a singularity at finite proper distance, which by choice of an additive integration constant we have located at $`r_s=0`$.
The horospheric coordinate $`z`$ is obtained from $`dr/dz=e^A`$, which can be solved to give
$`e^{2A(z)}=\mathrm{cot}^2\left({\displaystyle \frac{z}{L}}\right)=e^{\frac{2r}{L}}1.`$ (53)
Neither $`r`$ nor $`z`$ proves useful for solving the various equations. Instead we are able to make progress by using a variable $`u`$ in which the boundary is mapped to $`u_b=1`$ and the singularity to $`u_s=0`$. This is achieved by
$`u\mathrm{cos}^2\left({\displaystyle \frac{z}{L}}\right).`$ (54)
The equations for fluctuations become hypergeometric in $`u`$ as we will see, and the quantities which enter these equations can be expressed as:
$`W(u)`$ $`=`$ $`{\displaystyle \frac{3}{2u}},W^{}(u)=\sqrt{3}{\displaystyle \frac{\sqrt{1u}}{u}},`$
$`W^{\prime \prime }(u)`$ $`=`$ $`{\displaystyle \frac{u2}{u}},W^{\prime \prime \prime }(u)={\displaystyle \frac{4}{\sqrt{3}}}{\displaystyle \frac{\sqrt{1u}}{u}},`$ (55)
$`e^{2A(u)}`$ $`=`$ $`{\displaystyle \frac{u}{1u}},{\displaystyle \frac{du}{dr}}={\displaystyle \frac{2}{L}}(1u).`$
### 3.2 Correlators of inert scalars
Let us begin our study of fluctuations in this geometry by calculating the two-point function for the dilaton $`\varphi `$. The dilaton couples to $`๐ช_\varphi =\mathrm{Tr}F^2+\mathrm{}`$ and hence provides information about the glueball spectrum. The Klein-Gordon equation for this field is
$`{\displaystyle \frac{1}{\sqrt{g}}}_\mu (\sqrt{g}g^{\mu \nu }_\nu )\varphi =\left({\displaystyle \frac{}{r}}+4A^{}(r)\right){\displaystyle \frac{\varphi }{r}}+p^2e^{2A}\varphi =0.`$ (56)
In terms of $`u`$ this becomes
$`\varphi ^{\prime \prime }(u)+{\displaystyle \frac{1}{1u}}\left({\displaystyle \frac{2}{u}}1\right)\varphi ^{}(u)+{\displaystyle \frac{p^2L^2}{4}}{\displaystyle \frac{1}{u(1u)}}\varphi (u)=0.`$ (57)
This equation is hypergeometric; the solution regular in the deep interior ($`u=0`$) is
$`\varphi (u)=(u1)^2F(2{\displaystyle \frac{pL}{2}},2+{\displaystyle \frac{pL}{2}};2;u).`$ (58)
This is a simpler, but presumably equivalent, form of the solution very recently obtained in , where a different independent variable was used.
We can calculate the two-point function in the established way . Specifically, the $`p`$-space correlator is obtained by imposing a cutoff at $`z=ฯต`$ in the horospheric coordinate. One then finds
$`๐ช(p)๐ช(p)=\underset{ฯต0}{lim}\left({\displaystyle \frac{1}{ฯต^{2(\mathrm{\Delta }4)}}}\right)\left[{\displaystyle \frac{1}{z^3}}{\displaystyle \frac{d}{dz}}\mathrm{ln}(\varphi (z,p))\right]_{z=ฯต},`$ (59)
where an energy-momentum conserving $`\delta `$-function has been dropped, and $`\mathrm{\Delta }`$ is the UV dimension of the operator $`๐ช(x)`$. The only term which is kept in the limit is the most singular term in $`ฯต`$ which is non-analytic in $`p^2`$. Other more singular terms are multiplied by polynomials in $`p^2`$, and are dropped because they correspond to counter terms. This procedure is equivalent to calculating the second variation of the on-shell action (4). Field theory considerations imply that correlators behave as $`(p)^{2\nu }\mathrm{ln}(p^2)`$ for large spacelike momentum, where $`\nu =\mathrm{\Delta }2`$. Our calculations lead directly to correlation functions which are normalized so that the leading term has the same coefficient as in the $`AdS_5`$ geometry. The analysis of the Appendix of , redone for Bessel functions with integer $`\nu >2`$, gives
$`๐ช(p)๐ช(p)=\left[{\displaystyle \frac{2\nu }{\mathrm{\Gamma }(\nu )\mathrm{\Gamma }(\nu +1)4^\nu }}\right](p)^{2\nu }\mathrm{ln}(p^2)`$ (60)
For the dilaton this procedure gives (with $`s=p^2`$)
$`๐ช_\varphi (p)๐ช_\varphi (p)={\displaystyle \frac{1}{8}}s\left(s{\displaystyle \frac{4}{L^2}}\right)\left[\psi \left(2+{\displaystyle \frac{L\sqrt{s}}{2}}\right)+\psi \left(2{\displaystyle \frac{L\sqrt{s}}{2}}\right)\right],`$ (61)
where $`\psi (z)d\mathrm{ln}\mathrm{\Gamma }(z)/dz`$. This result agrees with . The correlator has a discrete spectrum of poles at $`s=4(n+2)^2/L^2`$ as is consistent with the expected glueball spectrum in a confining theory.
We can also consider the two-point function of the inert scalar $`\sigma `$. One must now take into account the coupling $`U(m)`$ to the active scalar that results from the potential (49). The wave equation is (39) with
$`U(m){\displaystyle \frac{^2V(m,\sigma )}{\sigma ^2}}|_{\sigma =0}={\displaystyle \frac{3g^2}{4}}\left[2\mathrm{cosh}\left({\displaystyle \frac{2m}{\sqrt{3}}}\right)1\right],`$ (62)
which leads to the hypergeometric wave equation
$`\sigma ^{\prime \prime }(u)+{\displaystyle \frac{1}{1u}}\left({\displaystyle \frac{2}{u}}1\right)\sigma ^{}(u)+{\displaystyle \frac{1}{(1u)^2}}\left({\displaystyle \frac{p^2L^2}{4}}{\displaystyle \frac{1u}{u}}+{\displaystyle \frac{3}{4}}{\displaystyle \frac{43u}{u}}\right)\sigma (u)=0.`$ (63)
The solution regular at the horizon is
$`\sigma (u)=(1u)^{3/2}F({\displaystyle \frac{3}{2}}+{\displaystyle \frac{1}{2}}\sqrt{9+L^2p^2},{\displaystyle \frac{3}{2}}{\displaystyle \frac{1}{2}}\sqrt{9+L^2p^2};2;u).`$ (64)
The corresponding two-point function is
$`๐ช_\sigma (p)๐ช_\sigma (p)={\displaystyle \frac{s+8/L^2}{2}}\left[\psi \left({\displaystyle \frac{3}{2}}+{\displaystyle \frac{1}{2}}\sqrt{9+L^2p^2}\right)+\psi \left({\displaystyle \frac{3}{2}}{\displaystyle \frac{1}{2}}\sqrt{9+L^2p^2}\right)\right].`$ (65)
There is a discrete spectrum of poles at $`s=4n(n+3)/L^2,n=0,1,2\mathrm{}`$. We note the presence a massless state, which we shall not attempt to reconcile with the interpretation that the flow describes a confining field theory.
These results may be compared with expectations based on the form of the Schrรถdinger potentials (41) discussed in section 2. For both fields, the one discovers that $`๐ฑ(z)+\mathrm{}`$ in both the boundary and deep interior limits, which is indicative of the discrete spectrum found above. In section 2 we showed that the Schrรถdinger potential of an inert dilaton always possesses a SUSY QM form $`๐ฑ(z)=๐ฐ(x)^2+๐ฐ^{}(z)`$, where
$`๐ฐ_\varphi (z)={\displaystyle \frac{1}{2}}ge^AW,`$ (66)
independent of the flow. Thus we can argue that the spectrum is positive, and we are in the situation of broken supersymmetry since the candidate 0-mode (which is just $`\varphi =const`$) is not normalizable. The $`\sigma `$ field has an additional term proportional to $`U(m(z))`$, which turns out to be negative; the $`m`$-flow sits atop a ridge in the $`\sigma `$-direction of field space. Thus SUSY quantum mechanics does not manifestly apply with the potential written as in (41). However, one may show that $`๐ฑ_\sigma (z)`$ can be rewritten in an exact SUSY QM form with prepotential
$`๐ฐ_\sigma (z)={\displaystyle \frac{1}{2}}ge^A\left(W+3\right),`$ (67)
and in this case the zero-mode can be shown to be normalizable, consistent with the spectrum of the correlator (65).
### 3.3 The graviton/active scalar system
It is straightforward to show that (35) becomes
$`R^{\prime \prime }(u)+{\displaystyle \frac{1}{u(1u)}}R^{}(u)+{\displaystyle \frac{1}{u(1u)}}\left({\displaystyle \frac{2u1}{u(1u)}}+p^2L^2/4\right)R(u)=0,`$ (68)
which has the solution
$`R(u)=u(1u)F({\displaystyle \frac{3}{2}}+{\displaystyle \frac{1}{2}}\zeta ,{\displaystyle \frac{3}{2}}{\displaystyle \frac{1}{2}}\zeta ;3;u),`$ (69)
where $`\zeta =\sqrt{1+p^2L^2}`$. We reconstruct $`\stackrel{~}{m}`$ following (34). Using the relations (3.1), we integrate to find
$`\stackrel{~}{m}={\displaystyle \frac{\sqrt{1u}}{u}}\left[4F_1(p;u)f_0+p^2L^2uF_2(p;u)\right],`$ (70)
where $`F_n(p;u)`$ is the hypergeometric function
$`F_n(p;u)F(n{\displaystyle \frac{3}{2}}+{\displaystyle \frac{1}{2}}\zeta ,n{\displaystyle \frac{3}{2}}{\displaystyle \frac{1}{2}}\zeta ;n;u).`$ (71)
Here $`f_0`$ is an integration constant, corresponding to the addition of a multiple of the universal solution. The choice $`f_0=4`$ ensures regularity at the singularity $`u=0`$, but we will keep $`f_0`$ as a parameter in our further discussion. The second solution to (68) is singular at $`u=0`$ and cannot be made regular. Near the boundary $`u=1`$, we observe $`\stackrel{~}{m}\sqrt{1u}z`$ which is the proper scaling for an operator with dimension $`\mathrm{\Delta }=3`$.
Let us check to see whether the form found for the fluctuation agrees with the intuition based on Schrรถdinger form (42) of (68), with potential $`๐ฑ(z)`$ given by (44). Near $`z=0`$, the leading term of the potential is $`๐ฑ(z)=1/4z^2`$, which is the limiting strength of an allowed attractive $`1/z^2`$ potential. Approaching the singularity at $`z\pi /2`$, there is a repulsive $`1/(\pi /2z)^2`$ behavior. Thus one expects a positive discrete spectrum, and this is reflected in the values of $`\zeta `$ which give a terminating hypergeometric series in (69), namely $`\zeta =2n+3`$ or $`s=p^2=4(n+1)(n+2)`$ for $`n=0,1,2\mathrm{}`$. Examining (70) at these values of $`p^2`$, we find the behavior near the boundary
$`\stackrel{~}{m}(u,p)|_{\zeta =2n+3}\sqrt{1u}\left[f_0+c_n(1u)+๐ช(1u)^2\right],`$ (72)
where $`c_n`$ is a constant. These functions are normalizable (in the covariant norm $`๐re^{4A}(\stackrel{~}{m})^2`$) only if $`f_0=0`$.
We now attempt to obtain a correlation function by applying the previously used standard procedure. Hypergeometric analytic continuation formulae give the boundary behavior
$`\stackrel{~}{m}(u,p)`$ $``$ $`{\displaystyle \frac{\sqrt{1u}}{u}}{\displaystyle \frac{1}{\mathrm{\Gamma }\left(\frac{3+\zeta }{2}\right)\mathrm{\Gamma }\left(\frac{3\zeta }{2}\right)}}[4+p^2+{\displaystyle \frac{f_0\pi p^2}{4\mathrm{cos}(\pi \zeta /2)}}`$
$`+`$ $`{\displaystyle \frac{p^4(1u)}{4}}(h_0^{\prime \prime }\mathrm{ln}(1u))+๐ช(1u)^2],`$
where $`h_0^{\prime \prime }=\psi (1)+\psi (2)\psi ((3+\zeta )/2)\psi ((3\zeta )/2)`$. Applying (59) we see that the โwould-be correlatorโ should be the most singular term (in $`(1u)`$) that is non-analytic in $`p^2`$ in
$`\mathrm{`}\mathrm{`}๐ช(p)๐ช(p)\text{}={\displaystyle \frac{s^2}{4s+16/L^2+\frac{f_0\pi s}{\mathrm{cos}(\pi \zeta /2)}}}\left[\psi \left({\displaystyle \frac{3}{2}}+{\displaystyle \frac{1}{2}}\zeta \right)+\psi \left({\displaystyle \frac{3}{2}}{\displaystyle \frac{1}{2}}\zeta \right)\mathrm{ln}\left((1u)e^{2\gamma }\right)\right].`$ (74)
At this point we hit the first major barrier of our program; we see no way to extract the leading non-analytic term. For example the coefficient of $`\mathrm{ln}(1u)`$ should be a polynomial in $`s`$ in order to be interpreted as a contact term. However here we have (for general $`f_0`$) an entire function of $`s`$ with essential singularity at infinity, while at $`f_0=0`$ there is a pole at the space-like value $`s=4/L^2`$. One might nevertheless try to identify the correlator by simply dropping the log term. However for general $`f_0`$ the expected poles of the discrete spectrum, at $`\zeta =2s+3`$, cancel between the numerator and denominator, leaving a function with no singularities in the finite $`s`$-plane. For $`f_0=0`$ one regains the expected poles, but there are also unphysical tachyon poles.
We are rather sure that our treatment of the fluctuation equations leading to (34), (35) is correct. The solution $`\stackrel{~}{m}(u,p)`$, especially at $`f_0=4`$, has the properties expected in the $`AdS`$/CFT correspondence, both at the boundary and the interior singularity. It is reassuring that (35) admits solutions with unphysical behavior in both limits, yet the actual solution $`\stackrel{~}{m}(r,p)`$ behaves correctly at both ends. It thus appears that it is the procedure to obtain the correlation function from the fluctuation which must be modified.
We can additionally calculate the associated graviton fluctuations using equations (27) and (29). We find
$`h`$ $`=`$ $`{\displaystyle \frac{1}{3\sqrt{3}u}}(24f_096F_124L^2p^2uF_2+24L^2p^2u(1u)F_3`$
$`+`$ $`L^2p^2u^2(1u)(8L^2p^2)F_4),`$
$`H^{}`$ $`=`$ $`{\displaystyle \frac{L(1u)}{12\sqrt{3}u}}(24f_096F_124L^2p^2uF_2+12L^2p^2u(1u)F_3`$
$`+`$ $`L^2p^2u^2(1u)(8L^2p^2)F_4).`$
Only $`H^{}(r)`$ is determined by the equations. Note that if $`f_0=4`$, $`\stackrel{~}{m}`$, $`h`$ and $`H^{}`$ are all regular at the singularity $`u=0`$.
In principle, one would like to use (3.3) and (3.3) to calculate correlation functions of the trace of the energy-momentum tensor, such as $`T_i^i(p)T_j^j(p)`$, which is related to renormalization group flows and proposals for c-theorems . The fact that the solutions for $`\stackrel{~}{m}`$, $`h`$ and $`H`$ are all determined by a single set of boundary data is consistent with the field theory reality that the operator $`T_i^i`$ is not independent of the operator $`๐ช_m`$, as in (7). Unfortunately, correlators constructed from the solutions (3.3), (3.3) seem to suffer from the same pathologies as (74).
The reader may justifiably complain that two-point functions associated to tensor fluctuations should require a more complex treatment than the inert scalar case, since the actions are different. This is further evidence that the procedure (4) must be generalized in the case of the graviton/active scalar system. In section 5, we evaluate the full gravity + scalar action to second order in fluctuations, in the hope of finding such a generalization.
## 4 The Coulomb branch
Flows involving a single scalar from the $`\mathrm{๐๐}^{}`$ of $`SO(6)`$, which preserve 16 of the 32 bulk supercharges, were discussed in . Unlike the flow involving $`m`$, which is dual to an operator deformation of the Lagrangian of $`๐ฉ=4`$ SYM, these flows change the vacuum of the field theory, moving it out onto the Coulomb branch. The distinction can be perceived by examining the scaling of the scalar profile near the boundary. The $`\mathrm{๐๐}^{}`$ scalars are dual to operators with $`\mathrm{\Delta }=2`$, and the Coulomb branch profiles scale as $`z^2`$ rather than the $`z^2\mathrm{ln}z`$ associated with an operator flow.
The five flows discussed in all have known lifts to 10D configurations, namely the geometries produced by discs of D3-branes of various dimensionalities. There are three inequivalent single scalars in the $`\mathrm{๐๐}^{}`$, and two of these have distinct flows in the positive and negative directions in field space, for five in all. We consider the flow called $`n=2`$ in , and examine fluctuations of two inert scalars and the graviton/active scalar system.
### 4.1 The background flow
The scalar field $`\phi `$ involved in the $`n=2`$ Coulomb branch flow also played a role in a two-scalar flow to an $`๐ฉ=1`$ superconformal fixed point , where it was denoted $`\phi _3`$. The other scalar, there called $`\phi _1`$, is a member of the same $`\mathrm{๐๐}\overline{\mathrm{๐๐}}`$ Weyl orbit as $`\sigma `$. The potential and superpotential of this two-scalar subspace is
$`W(\phi ,\sigma )={\displaystyle \frac{1}{4}}e^{2\phi /\sqrt{6}}\left[e^{\sqrt{6}\phi }(3\mathrm{cosh}(2\sigma ))+2(\mathrm{cosh}(2\sigma )+1)\right],`$ (77)
$`V(\phi ,\sigma )={\displaystyle \frac{g^2}{4}}e^{2\phi /\sqrt{6}}[e^{\sqrt{6}\phi }({\displaystyle \frac{3}{4}}+{\displaystyle \frac{1}{2}}\mathrm{cosh}(2\sigma ){\displaystyle \frac{1}{4}}\mathrm{cosh}^2(2\sigma ))+(1+\mathrm{cosh}(2\sigma ))+`$ (78)
$`{\displaystyle \frac{1}{16}}e^{\sqrt{6}\phi }(1\mathrm{cosh}^2(2\sigma ))].`$
In the notation of , $`\phi =\mu `$ is the active scalar in the $`n=2`$ Coulomb branch flow, while $`\sigma `$ is inert. The superpotential for $`\phi `$ alone is
$`W(\phi )=e^{2\phi /\sqrt{6}}{\displaystyle \frac{1}{2}}e^{4\phi /\sqrt{6}}.`$ (79)
The solution to the equations (10) involving (79) can be obtained, but they are not of direct use since it is more convenient to use a radial coordinate that is a function of the scalar itself:
$`ve^{\sqrt{6}\phi }.`$ (80)
The boundary is at $`v=1`$. For this flow $`\phi \mathrm{}`$, so $`v[0,1]`$ with a curvature singularity at $`v_s=0`$. We have the relations
$`W`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \frac{v+2}{v^{1/3}}},W^{}={\displaystyle \frac{2}{\sqrt{6}}}{\displaystyle \frac{1v}{v^{1/3}}},`$
$`W^{\prime \prime }`$ $`=`$ $`{\displaystyle \frac{2}{3}}{\displaystyle \frac{1+2v}{v^{1/3}}},W^{\prime \prime \prime }={\displaystyle \frac{4}{3\sqrt{6}}}{\displaystyle \frac{14v}{v^{1/3}}},`$ (81)
$`e^{2A}`$ $`=`$ $`{\displaystyle \frac{\mathrm{}^2}{L^2}}{\displaystyle \frac{v^{2/3}}{1v}},{\displaystyle \frac{v}{r}}={\displaystyle \frac{2}{L}}v^{2/3}(1v).`$
One can calculate the horospheric variable $`z`$ in terms of $`v`$ and the relationship is
$`v=\text{sech}^2\left({\displaystyle \frac{z\mathrm{}}{L^2}}\right).`$ (82)
Following we have introduced the length $`\mathrm{}`$, the radius of the disc of D3-branes in ten dimensions. Taking $`\mathrm{}/L0`$ with $`z/L`$ fixed removes the flow and restores pure anti-de Sitter space. This definition is reminiscent of that for the variable $`u`$ in the $`๐ฉ=1`$ flow (54); one difference is that for (82) the singularity is at $`z\mathrm{}`$, and is thus at infinite proper distance.
### 4.2 Correlators of inert scalars
The two-point function of the dilaton $`\varphi `$ has been calculated previously in this background . The Klein-Gordon equation becomes
$`\varphi ^{\prime \prime }(v)+{\displaystyle \frac{2v}{v(1v)}}\varphi ^{}(v)+{\displaystyle \frac{p^2L^4}{4\mathrm{}^2}}{\displaystyle \frac{1}{v^2(1v)}}\varphi (v)=0,`$ (83)
which has the solution
$`\varphi (v)=v^aF(a,a;2+2a;v),`$ (84)
where
$`a{\displaystyle \frac{1}{2}}+{\displaystyle \frac{1}{2}}\sqrt{1{\displaystyle \frac{L^4p^2}{\mathrm{}^2}}}.`$ (85)
This solution is regular at the singularity for spacelike $`p^2`$; the second solution of (83) has a leading $`v^{a1}`$ and is not regular. The result for the 2-point function calculated from (59) is
$`๐ช_\varphi (p)๐ช_\varphi (p)={\displaystyle \frac{1}{4}}p^4\psi \left({\displaystyle \frac{1}{2}}+{\displaystyle \frac{1}{2}}\sqrt{1{\displaystyle \frac{L^4p^2}{\mathrm{}^2}}}\right),`$ (86)
which as remarked in has a branch cut along the positive real axis, indicating a continuous spectrum with gap or threshold at $`m_{gap}^2=\mathrm{}^2/L^4`$.
For the $`\sigma `$ field, there is a $`U(\phi )`$ term in the equation of motion (39) owing to the potential (78),
$`U(\phi ){\displaystyle \frac{^2V(\phi ,\sigma )}{\sigma ^2}}|_{\sigma =0}={\displaystyle \frac{g^2}{4}}e^{2\phi /\sqrt{6}}\left[4e^{\sqrt{6}\phi }\right],`$ (87)
and the wave equation becomes
$`\sigma ^{\prime \prime }(v)+{\displaystyle \frac{2v}{v(1v)}}\sigma ^{}(v)+{\displaystyle \frac{1}{v(1v)}}\left({\displaystyle \frac{p^2L^4}{4\mathrm{}^2}}{\displaystyle \frac{1}{v}}+{\displaystyle \frac{(4v)}{4(1v)}}\right)\sigma (v)=0,`$ (88)
which has the regular solution
$`\sigma (v)=v^a\sqrt{1v}F(a,a+1;2+2a;v),`$ (89)
with $`a`$ as (85). The correlation function is
$`๐ช_\sigma (p)๐ช_\sigma (p)=p^2\psi \left({\displaystyle \frac{1}{2}}+{\displaystyle \frac{1}{2}}\sqrt{1{\displaystyle \frac{L^4p^2}{\mathrm{}^2}}}\right),`$ (90)
This correlator also has a continuous spectrum, with the same gap as (86). The momentum-dependence is appropriate for an operator of $`\mathrm{\Delta }=3`$.
We may check that this behavior is consistent with the SUSY QM viewpoint. Calculating the Schrรถdinger potentials $`๐ฑ_\varphi (z)`$ and $`๐ฑ_\sigma (z)`$ (41), we find that $`๐ฑ\mathrm{}`$ at the boundary for both, while they asymptote to $`\mathrm{}^2/L^4`$ in the deep interior. This is precisely the expected behavior for a potential with a continuous spectrum and a gap at $`m_{gap}=\mathrm{}/L^2`$. (For the dilaton, this was already noticed in .)
As always, the dilaton wave equation can be rewritten in Schrรถdinger form with an exact SUSY QM potential given by the prepotential (66). Normalizability fails near the boundary for the dilaton zero-mode. As in the $`๐ฉ=1`$ case, we find that the Coulomb branch flow is along a ridge of the potential (78) in the $`\sigma `$-direction, and thus $`U(\phi )`$ contributes negatively to the $`\sigma `$ Schrรถdinger potential (41). However, again we find that the potential can be cast into exact SUSY QM form, this time with the modified prepotential
$`๐ฐ_\sigma (z)={\displaystyle \frac{1}{2}}ge^AW.`$ (91)
Unlike the $`๐ฉ=1`$ case (67), this doesnโt have a new term beyond that of (66), but the coefficient is modified instead. In this case normalizability of the zero-mode fails at the singularity.
### 4.3 The graviton/active scalar system
One can show that (35) becomes in this case
$`R^{\prime \prime }(v)+{\displaystyle \frac{2}{v}}R^{}(v)+{\displaystyle \frac{p^2L^2}{\mathrm{}^2}}{\displaystyle \frac{1}{v^2(1v)}}R(v)=0.`$ (92)
This particularly simple form results because the potential-type term $`(2/3)WW^{\prime \prime }+(8/9)W^2+(1/2)W^{}W^{\prime \prime \prime }W^2`$ evaluates to zero; we have no ready explanation for this cancellation. The relevant solution to (92) is
$`R(v)=v^a(1v)F(1+a,2+a;2+2a;v).`$ (93)
The second solution also has $`v^{a1}`$ behavior and cannot be made regular in the interior. One can integrate (93) to obtain the solution for $`\stackrel{~}{\phi }`$:
$`\stackrel{~}{\phi }(v)=v^a(1v)_3F_2(1+a,2+a,{\displaystyle \frac{1}{3}}+a;2+2a,{\displaystyle \frac{4}{3}}+a;v).`$ (94)
We could add a multiple of the universal solution to (94), but this would destroy regularity at the interior singularity at $`v=0`$.
One can see that (92) has a constant zero-mode solution, and that (93) reduces to a constant as $`p^20`$ which is the same as $`a0`$. Although this constant mode is not normalizable, one might suspect that supersymmetric quantum mechanics is again at work. Indeed one can again look back to (35) in the present case where the potential-type term vanishes. The equivalent Schrรถdinger equation is then supersymmetric with prepotential $`U(z)=g/2e^A(W^{\prime \prime }5W/3)`$.
We may use (27),(29) to determine the graviton modes associated with (94):
$`h`$ $`=`$ $`{\displaystyle \frac{4\sqrt{2}v^a}{3\sqrt{3}L^4p^2(a+1/3)}}[4\mathrm{}^2(1+3a)(2+4v+v^2+a(2+v+v^2))\times `$
$`F(1+a,2+a;2+2a;v)2\mathrm{}^2v(1v)(2+v)(2+7a+3a^2)F(2+a,3+a,3+2a,v)`$
$``$ $`3(2+v)L^4p_3^2F_2(1+a,2+a,{\displaystyle \frac{1}{3}}+a;2+2a,{\displaystyle \frac{4}{3}}+a;v)],`$
$`H^{}`$ $`=`$ $`{\displaystyle \frac{\sqrt{2}v^{a1/3}(1v)}{3\sqrt{3}L\mathrm{}^2p^2(a+1/3)}}[4\mathrm{}^2(1+3a)(2+5v+3a(v1))\times `$
$`F(1+a,2+a;2+2a;v)+6\mathrm{}^2v(1v)(2+7a+3a^2)F(2+a,3+a,3+2a,v)`$
$`+`$ $`3L^4p_3^2F_2(1+a,2+a,{\displaystyle \frac{1}{3}}+a;2+2a,{\displaystyle \frac{4}{3}}+a;v)].`$
Our analysis of the fluctuations (94), (4.3), (4.3), all of which contain the generalized hypergeometric function $`{}_{3}{}^{}F_{2}^{}`$, is hampered because the relevant analytic continuation formulae are not in the literature. In particular we need the expansion of (94) near the boundary, i.e. a series in $`(1v)`$. We proceed by expanding (93) in series first, and then integrate to obtain a series for $`\stackrel{~}{\phi }`$. Doing this we see the leading logarithmic singularity:
$`\stackrel{~}{\phi }={\displaystyle \frac{(a+1/3)\mathrm{\Gamma }(2+2a)}{\mathrm{\Gamma }(1+a)\mathrm{\Gamma }(2+a)}}(1v)\left(\mathrm{ln}(1v)+\stackrel{~}{\phi }_0(a)+๐ช(1v)\right),`$ (97)
and while all $`๐ช(1v)`$ terms could be calculated from higher order terms in the expansion of (93), the integration constant $`\stackrel{~}{\phi }_0`$ remains undetermined. In principle it is fixed by our choice not to add a multiple of the universal solution to (94), but it is nontrivial to calculate it.
Let us consider $`h`$. We expect a graviton mode not to be singular on the boundary, and indeed $`1/(1v)`$ and $`\mathrm{ln}(1v)`$ terms cancel between the various terms in (4.3). The series expansion of $`h(v)`$ is then
$`h(v)={\displaystyle \frac{16\sqrt{2}\mathrm{}^2}{\sqrt{3}L^4p^2}}{\displaystyle \frac{\mathrm{\Gamma }(2+2a)}{\mathrm{\Gamma }(1+a)\mathrm{\Gamma }(a)}}\left(3\stackrel{~}{\phi }_0(a){\displaystyle \frac{2}{a(1+a)}}+6\gamma +6\psi (1+a)\right)+๐ช(1v)^2.`$ (98)
What form do we expect for $`h(v)`$? In the field theory, conformal invariance is only spontaneously broken, and consequently $`T_i^i=0`$ continues to hold as an operator equation on the Coulomb branch. As a result, one might expect $`h(v)`$ to fall off more rapidly on the boundary than in the case of an operator deformation, so that it does not excite a dual field theory operator. A hint of this behavior is present in the fact that no $`๐ช(1v)`$ term appears in (98). We are therefore led to postulate that the constant term vanishes as well, and
$`\stackrel{~}{\phi }_0(a)={\displaystyle \frac{2}{3a(1+a)}}2\gamma 2\psi (1+a).`$ (99)
This is a tempting assumption, since (97) then produces the two-point function<sup>3</sup><sup>3</sup>3Formulas (59), (60) are modified for a dimension $`\mathrm{\Delta }=2`$ operator. See (25) of .
$`๐ช_\phi (p)๐ช_\phi (p)=\psi \left({\displaystyle \frac{1}{2}}+{\displaystyle \frac{1}{2}}\sqrt{1{\displaystyle \frac{L^4p^2}{\mathrm{}^2}}}\right)+{\displaystyle \frac{4L^4}{3p^2\mathrm{}^2}}`$ (100)
In addition to the expected branch point there is a pole at $`p^2=0`$. This is a problem for the formalism since the associated constant zero-mode is non-normalizable, so it should not show up in the correlator. It is also a problem for the field theory interpretation. To see this recall that an extra factor of $`N^2`$ must be inserted to agree with the short distance form $`๐ช_\phi (x)๐ช_\phi (0)N^2/x^4`$ in field theory. Although the field theory contains massless Goldstone bosons from the breaking of $`SO(6)`$ flavor symmetry to $`SO(4)\times SO(2)`$ by the disc of D3-branes, it appears very unlikely that these states should couple to $`๐ช_\phi =\mathrm{Tr}X^2`$ with strength $`N`$<sup>4</sup><sup>4</sup>4We thank S. Gubser and A. Hanany for discussions of this issue..
This motivates us to test the assumption (99) numerically. Specifically we took the defining power series $`{}_{3}{}^{}F_{2}^{}=_nc_nv^n`$ about $`v=0`$ which is logarithmically divergent at $`v=1`$ and subtracted the explicitly summable series of the large $`n`$ limit of the $`c_n`$ (obtained using Stirlingโs formula). The result is a convergent series whose value at $`v=1`$ gives the unknown constant $`\stackrel{~}{\phi }_0(a)`$. The numerical series agrees remarkably well with (98) for parameter values $`a1`$, but agreement fails for small $`a`$ since the poles of the power series coefficients for $`a0`$ do not coincide with those of (98). We must therefore conclude that the physically motivated assumption above does not agree with the properties of $`{}_{3}{}^{}F_{2}^{}`$. It should also be observed that the same correlator (100) can be obtained by subtracting a multiple of the universal solution from $`h(v)`$ and $`\stackrel{~}{\phi }(v)`$, so as to impose the condition $`h(1)=0`$. This diffeomorphism has exactly the effect of inserting (99) in (97). However the diffeomorphism also makes $`h(v)`$ and $`\stackrel{~}{\phi }(v)`$ singular at the origin, and we cannot presently justify it.
The situation may be summarized as follows. We imposed the condition that $`h(v)`$ vanish on the boundary because it seems to be physically required for Coulomb branch flows. This leads to a correlation function with apparently unphysical zero-mode poles. Further, the condition $`h(1)=0`$ does not seem to be compatible with the properties of $`{}_{3}{}^{}F_{2}^{}`$ obtained by numerical study. Perhaps the analytic form of the constant $`\stackrel{~}{\phi }(a)`$ could illuminate the situation.
## 5 Calculation of correlation functions
Our main purpose here is to discuss some of our attempts to determine the correct prescription for the calculation of active scalar and graviton correlators. First we will outline a calculation of the on-shell supergravity action (1) through second order in the fields $`\stackrel{~}{\phi }`$, $`h`$, and $`H`$ which are coupled by the fluctuation equations of Sec 2. Because of the coupling, a single set of boundary data determines the boundary behavior of all the fields. Accordingly, factors of $`h`$ and $`H`$ in the action will contribute to the correlation functions of the active scalars. Hence to calculate two-point functions we should keep terms quadratic in any of $`\stackrel{~}{\phi }`$, $`h`$ and $`H`$. One then varies the on-shell action with respect to the boundary data to obtain the correlation function.
The bulk action is given in (1). In addition, it is well-known that this should be supplemented with certain boundary terms . First among these is the Gibbons-Hawking term, which is included to generate a well-posed Hamiltonian formalism :
$`S_{GH}`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle _{}}d^4x\sqrt{g_4}๐ฆ={\displaystyle \frac{1}{2}}{\displaystyle _{}}d^4x\sqrt{g_4}_\mu n^\mu ={\displaystyle \frac{1}{2}}{\displaystyle _{}}_r\sqrt{g_4}`$
$`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle _{}}e^{4A}[{\displaystyle \frac{1}{2}}h_i^i{}_{}{}^{}+{\displaystyle \frac{1}{4}}h_i^ih_j^j{}_{}{}^{}{\displaystyle \frac{1}{2}}h^{ij}h_{ij}^{}+4A^{}(1+{\displaystyle \frac{1}{2}}h_i^i+{\displaystyle \frac{1}{8}}(h_i^i)^2{\displaystyle \frac{1}{4}}h^{ij}h_{ij})+\mathrm{}],`$
where $`๐ฆ=_\mu n^\mu `$ is the trace of the second fundamental form of the boundary, with $`n^\mu `$ a normal to the boundary; in the second line we expand this term to second order in the fluctuations. In addition, a boundary cosmological term has been recommended for canceling the leading volume divergence:
$`S_{vol}`$ $`=`$ $`\beta {\displaystyle _{}}\sqrt{g_4}A^{}`$ (102)
$`=`$ $`\beta {\displaystyle _{}}e^{4A}A^{}\left(1+{\displaystyle \frac{1}{2}}h_i^i+{\displaystyle \frac{1}{8}}(h_i^i)^2{\displaystyle \frac{1}{4}}h^{ij}h_{ij}+\mathrm{}\right),`$ (103)
for some $`\beta `$ which we will fix momentarily. Note that $`A^{}=1/L`$ is a constant on boundary, which we include explicitly to simplify the formulas that follow.
We proceed to expand the bulk action to second order in fluctuations, making use of the equations of motion as necessary. After a remarkably tedious calculation, we find that as in simpler cases, the bulk action can be reduced entirely to a set of boundary terms:
$`S_{bulk}`$ $`=`$ $`S_0+S_1+S_2,`$
$`S_0`$ $`=`$ $`{\displaystyle d^4xe^{4A}\left(\frac{1}{2}A^{}\right)},`$ (104)
$`S_1`$ $`=`$ $`{\displaystyle }d^4xe^{4A}({\displaystyle \frac{1}{4}}h_i^i{}_{}{}^{}{\displaystyle \frac{1}{4}}A^{}h_i^i\phi ^{}\stackrel{~}{\phi }),`$
$`S_2`$ $`=`$ $`{\displaystyle d^4xe^{4A}\left(\frac{1}{2}\stackrel{~}{\phi }\stackrel{~}{\phi }^{}\frac{1}{4}\phi ^{}\stackrel{~}{\phi }h_i^i+\frac{3}{16}h^{ij}h_{ij}^{}\frac{1}{16}h_i^i(h_j^j)^{}+\frac{1}{8}A^{}h^{ij}h_{ij}\frac{1}{16}A^{}h_i^ih_j^j\right)},`$
evaluated on the boundary $`r=R`$. Here we have organized the action by the order of the fluctuations. The zeroth-order term is the volume divergence, which also receives contributions from the boundary terms (5) and (102); it is canceled by the choice $`\beta =1/2`$, which in fact removes all terms proportional to $`A^{}`$. The total action then reduces to
$`S_{tot}={\displaystyle }d^4xe^{4A}(\phi ^{}\stackrel{~}{\phi }{\displaystyle \frac{1}{2}}\stackrel{~}{\phi }\stackrel{~}{\phi }^{}{\displaystyle \frac{1}{4}}\varphi ^{}\stackrel{~}{\phi }h_i^i+{\displaystyle \frac{1}{16}}h_i^ih_j^j{}_{}{}^{}{\displaystyle \frac{1}{16}}h^{ij}h_{ij}^{}).`$ (105)
We can express this in momentum space in terms of $`h`$ and $`H`$,
$`S_{tot}`$ $`=`$ $`\phi ^{}(R)\stackrel{~}{\phi }(R,p=0)+{\displaystyle }d^4pe^{4A}({\displaystyle \frac{1}{2}}\stackrel{~}{\phi }(R,p)\stackrel{~}{\phi }^{}(R,p)`$
$`+`$ $`{\displaystyle \frac{3}{32}}h(R,p)h^{}(R,p)+{\displaystyle \frac{3}{32}}p^2H(R,p)h^{}(R,p)+{\displaystyle \frac{3}{64}}p^2h(R,p)H^{}(R,p)).`$
The first term in (5) is linear in the fluctuation $`\stackrel{~}{\phi }`$, and is thus suggestive of a one-point function. It does not clearly discriminate between the $`๐ฉ=1`$ flow, where no one-point function is expected, and the Coulomb branch flow, where one is. We find this puzzling.
On the other hand, the quadratic terms suggest a modified calculation of the scalar correlator in which $`h`$ and $`H`$ are related to the boundary data for $`\stackrel{~}{\phi }`$. This is straightforward but complicated. We have performed such a calculation for the case of the $`๐ฉ=1`$ active scalar, but ultimately encountered the same difficulties as in section 3. The resolution of the problem presumably involves understanding (5) better, but other pieces of the puzzle may still be missing.
A further uncertainty is the issue of the diffeomorphism invariance. We refer particularly to (14) which has been interpreted as describing bending of the cutoff surface and horizon. It is unsettling that the calculation of the active correlator in the Coulomb branch flow was so markedly changed by such a diffeomorphism. The on-shell action must be diffeomorphism invariant, and it is not clear to us that this is manifest in (5). In particular, a term $`H`$ with no derivatives appears; this quantity is absent from the equation of motion and its constant term is thus not determined, but has the form of a pure diffeomorphism (13). A better understanding will have to involve coming to terms with these issues.
## Acknowledgements
We are grateful for discussions with D. Anselmi, S. Gubser, A. Hanany, A. Karch, J. Minahan, G. Moore, J. Polchinski, L. Rastelli, and N. Warner. The research of O.D. was supported by the U.S. Department of Energy under contract #DE-FC02-94ER40818. The research of D.Z.F. was supported in part by the NSF under grant number PHY-97-22072. |
warning/0002/math0002031.html | ar5iv | text | # 1 Essential backgrounds and notations for physicists.
## 1 Essential backgrounds and notations for physicists.
In this section, we collect some basic facts and notations that will be needed in the discussion. The part that is related to equivariant vector bundles over toric manifolds is singled out in Sec. 2. Readers are referred to the listed literatures for more details.
$``$ Toric geometry. (\[A-G-M\], \[C-K\], \[Da\], \[Ew\], \[Fu\], \[Gre\], \[G-K-Z\], \[Ke\], and \[Od1,Od2\].) Physicists are referred particularly to \[A-G-M\] or \[Gre\] for a nice expository of toric geometry. Let us fix the terminology and notations here and refer the details to \[Fu\].
Notation :
> $`N\text{}^n`$: a lattice;
>
> $`M=\text{Hom}(N,\text{})`$: the dual lattice of $`N`$;
>
> $`T_N=\text{Hom}(M,\text{}^{})`$: the (complex) $`n`$-torus;
>
> $`\mathrm{\Sigma }`$: a fan in $`N_{\text{}}`$;
>
> $`X_\mathrm{\Sigma }`$: the toric variety associeted to $`\mathrm{\Sigma }`$;
>
> $`\mathrm{\Sigma }(i)`$: the $`i`$-skeleton of $`\mathrm{\Sigma }`$;
>
> $`U_\sigma `$: the local affine chart of $`X_\mathrm{\Sigma }`$ associated to $`\sigma `$ in $`\mathrm{\Sigma }`$;
>
> $`x_\sigma U_\sigma `$: the distinguished points associated to $`\sigma `$;
>
> $`O_\sigma `$: the $`T_N`$-orbit of $`x_\sigma `$ under the $`T_N`$-action on $`X_\mathrm{\Sigma }`$;
>
> $`V(\sigma )`$: the orbit closure of $`O_\sigma `$;
>
> $`M(\sigma )=\sigma ^{}M`$.
Recall that points $`v`$ in the interior of $`\sigma N`$ represent one-parameter subgroups $`\lambda _v`$ in $`\text{๐}_N`$ such that $`lim_{z0}\lambda _v(z)=x_\tau `$. Recall also that the normal cones associated to a polyhedron $`\mathrm{\Delta }`$ with vertices in $`M`$ form a fan in $`N_{\text{}}`$, called the normal fan of $`\mathrm{\Delta }`$. This determines a projective toric variety.
$``$ Toric surface. (\[Od2\].) Any complete nonsingular toric surface is obtained from consecutive equivariant blowups of either $`\text{}\mathrm{P}^2`$ or one of the Hurzebruch surfaces $`\text{๐ฝ}_a`$ at $`T_N`$ fixed points. Indeed, one has a complete classification of them as follows:
Fact 1.1 \[toric surface\]. (\[Od2\].) The set of isomorphism classes of complete nonsingular toric surfaces $`X_\mathrm{\Sigma }`$ is in one-to-one corresponce with the set of equivalent classes of weighted circular graphs $`w=(w_1,\mathrm{},w_s)`$ (under rotation and reflection) of the following form: (Figure 1-1.)
> (1) The circular graph having $`3`$ vertices with weights $`1,\mathrm{\hspace{0.17em}1},\mathrm{\hspace{0.17em}1}`$.
>
> (2) The circular graph having $`4`$ vertices with weights in circular order $`0,a,\mathrm{\hspace{0.17em}0},a`$.
>
> (3) The weighted circular graphs with $`s5`$ vertices that is obtained from one with $`(s1)`$ vertices by adding a vertex of weight $`1`$ and reducing the weight of each of its two adjacent vertices by $`1`$.
Let $`\mathrm{\Sigma }(1)=(v_1,\mathrm{},v_s)`$ in, say, counterclockwise order in $`N_{\text{}}`$, then for each $`i`$, there exists a unique integer $`a_i`$ such that $`v_{i1}+v_{i+1}+a_iv_i=0`$ (here, $`s+11`$). The correspondence is then given by $`X_\mathrm{\Sigma }(a_1,\mathrm{},a_s)`$.
$``$ Line bundles with positive/negative $`c_1`$. (\[C-K\], \[Hi\], and \[Re\].) Recall that the Kรคhler cone in $`H^2(X_\mathrm{\Sigma },\text{})`$ consists of all the Kรคhler classes of $`X_\mathrm{\Sigma }`$ and the Mori cone of $`X_\mathrm{\Sigma }`$ consists of all the classes in $`H_2(X_\mathrm{\Sigma },\text{})`$ representable by effective $`2`$-cycles. From \[Re\], the Mori cone of $`X_\mathrm{\Sigma }`$ is generated by $`V(\tau )`$, $`\tau \mathrm{\Sigma }(n1)`$. We call a class $`\omega H^2(X_\mathrm{\Sigma },\text{})`$ positive (resp. negative), in notation $`\omega >0`$ (resp. $`\omega <0`$), if $`\omega `$ (resp. $`\omega `$) lies in the Kรคhler cone of $`X_\mathrm{\Sigma }`$. The fact that the Kรคhler cone and the Mori cone of a complete toric manifold are dual to each other gives us a criterion for a line bundle over $`X_\mathrm{\Sigma }`$ to have positive $`c_1`$:
Fact 1.2 \[positive/negative line bundle\]. Given an $`n`$-dimensional toric manifold $`X_\mathrm{\Sigma }`$. Let $`L\text{Pic}(X_\mathrm{\Sigma })`$ be a line bundle over $`X_\mathrm{\Sigma }`$ and $`D`$ be the associated divisor class as an element in $`H_{2n2}(X_\mathrm{\Sigma },\text{})`$. Then $`c_1(L)>0`$ (resp. $`<0`$) if and only if $`DV(\tau )>0`$ (resp. $`<0`$) for all $`\tau \mathrm{\Sigma }(n1)`$.
$``$ The augmented intersection matrix. (\[Fu\].) Given an $`n`$-dimensional complete nonsingular fan $`\mathrm{\Sigma }`$. Let $`\mathrm{\Sigma }(1)=\{v_1,\mathrm{},v_J\}`$ and $`\mathrm{\Sigma }(n1)=\{\tau _1,\mathrm{},\tau _I\}`$. Let $`A_1`$ and $`A_{n1}`$ be respectively the first and $`(n1)`$-th Chow group of $`X_\mathrm{\Sigma }`$. Then $`A_1`$ is generated by $`V(\tau _i)`$, $`i=1,\mathrm{},I`$, and $`A_{n1}`$ is generated by $`D(v_j)`$, $`j=1,\mathrm{},J`$. There is a nondegenerate pairing $`A_1\times A_{n1}\text{}`$ by taking the intersection number. Let $`Q`$ be the $`I\times J`$ matrix whose $`(i,j)`$-entry is the intersection number $`V(\tau _i)D(v_j)`$. Since the generators for $`A_1`$ and $`A_{n1}`$ used here may not be linearly independent, we shall call $`Q`$ the augmented intersection matrix (with respect to the generators). Explicitly, $`Q`$ can be determined as follows.
Let $`\tau _i=[v_{j_1},\mathrm{},v_{j_{n1}}]\mathrm{\Sigma }(n1)`$. Then $`\tau _i`$ is the intersection of two $`n`$-cones
$$\sigma _1=[v_{j_1},\mathrm{},v_{j_{n1}},v_{j_n}]\text{and}\sigma _2=[v_{j_1},\mathrm{},v_{j_{n1}},v_{j_n^{}}]$$
in $`\mathrm{\Sigma }`$. These vertices in $`\sigma _1\sigma _2`$ satisfy a linear equation of the form
$$v_{j_n}+v_{j_n^{}}+a_1v_{j_1}+\mathrm{}+a_{n1}v_{j_{n1}}=\mathrm{\hspace{0.33em}0},$$
for some unique integers $`a_1,\mathrm{},a_n`$ determined by $`\sigma _1\sigma _2`$ . In terms of this, the $`i`$-th row of $`Q`$ is simply the coefficient (row) vector of the above equation. I.e. $`V(\tau _i)D(v_{j_k})=a_k`$ for $`k=1,\mathrm{},n1`$; $`V(\tau _i)D(v_{j_n})=V(\tau _i)D(v_{j_n^{}})=1`$; and $`V(\tau _i)D(v_j)=0`$ for all other $`j`$.
$``$ Cox homogeneous coordinate ring of a toric manifold. (\[Co\] and \[C-K\], also \[Au\] and \[Do\].) Let $`\mathrm{\Sigma }`$ be a fan in $`\text{}^n`$ with $`\mathrm{\Sigma }(1)`$ generated by $`\{v_1,\mathrm{},v_a\}`$ and $`A_{n1}(X_\mathrm{\Sigma })`$ be the Chow group of $`X_\mathrm{\Sigma }`$. Let $`(z_1,\mathrm{},z_a)`$ be the coordinates of $`\text{}^a`$. For $`\sigma =[v_{j_1},\mathrm{},v_{j_k}]\mathrm{\Sigma }`$, denote by $`z^{\widehat{\sigma }}`$ the monomial from $`(z_1\mathrm{}z_a)/(z_{j_1}\mathrm{}z_{j_k})`$ after cancellation. Then $`X_\mathrm{\Sigma }`$ can be realized as the geometric quotient
$$X_\mathrm{\Sigma }=(\text{}^{\mathrm{\Sigma }(1)}Z(\mathrm{\Sigma }))/G,$$
where $`Z(\mathrm{\Sigma })`$ is the exceptional subset $`\{(z_1,\mathrm{},z_a)|z^{\widehat{\sigma }}=0\text{for all }\sigma \text{ in }\mathrm{\Sigma }\}`$ in $`\text{}^a`$ and $`G`$ is the group $`\text{Hom}_{\text{}}(A_{n1}(X_\mathrm{\Sigma }),\text{}^{})`$ that acts on $`\text{}^a`$ via the embedding in $`(\text{}^{})^a`$, obtained by taking $`\text{Hom}(,\text{}^{})`$ of the following exact sequence
$$\begin{array}{cccccccccc}0& & M& & \text{}^a& & A_{n1}(X_\mathrm{\Sigma })& & 0& \\ & & m& & (m(v_1),\mathrm{},m(v_a))& & & & & .\hfill \end{array}$$
More facts will be recalled along the way when we need them. Their details can be found in Sec. 1-3 in \[Co\].
## 2 Basics of equivariant vector bundles over toric manifolds.
Equivariant vector bundles over a toric manifold have been classified by Kaneyama and Klyachko independently, using differently sets of data (\[Ka1\] and \[Kl\]). In this article, we use Kaneyamaโs data as the starting point to compute splitting types. Some necessary facts from \[Ka1\] and \[Kl\] are summarized below with possibly slight modification/rephrasing to make the geometric picture more transparent.
Equivariant vector bundles over a toric manifold.
A vector bundle $``$ over a toric manifold $`X_\mathrm{\Sigma }`$ is equivariant if $`g^{}=`$ for all $`g\text{๐}_N`$. An equivariant bundle is linearizable if the action on the base can be lifted to a fiberwise linear action on the total space of the bundle.
Fact 2.1 \[linearizability\]. Every equivaraint vector bundle $``$ over a toric manifold $`X_\mathrm{\Sigma }`$ is linearizable.
In general, a bundle can be described by its local trivializations and the pasting maps. For an equivariant vector bundle $``$ over a toric manifold $`X_\mathrm{\Sigma }`$, these data can be integrated with the linearization of the toric action.
> (1) Local trivializations : Over each invariant affine chart $`U_\sigma `$ for $`\sigma \mathrm{\Sigma }`$, the bundle splits:
>
> $$|_{U_\sigma }=_\chi U_\sigma \times E_\sigma ^\chi ,$$
>
> where $`E_\sigma ^\chi `$ is the representation of $`T_N`$ associated to the weight $`\chi M`$.
>
> (2) The pasting maps : Over each orbit $`O_\tau U_{\sigma _1}U_{\sigma _2}`$, the pasting map $`\phi _{\sigma _2\sigma _1}:U_{\sigma _1}U_{\sigma _2}\text{GL}(r,\text{})`$ is determined by its restriction at a point, say, $`x_\tau `$, in $`O_\tau `$, due to the equivariant requirement. The pasting maps over different orbits in $`U_{\sigma _1}U_{\sigma _2}`$ are related to each other by the holomorphicity requirement that $`\phi _{\sigma _2\sigma _1}:U_{\sigma _1}U_{\sigma _2}\text{GL}(r,\text{})`$ must be holomorphic for every pair of $`\sigma _1`$, $`\sigma _2`$ in $`\mathrm{\Sigma }`$. This implies that indeed $`\phi _{\sigma _2\sigma _1}`$ is completely determined by its restriction to a point, say, $`x_0`$, in the dense open orbit $`U_0U_{\sigma _1}U_{\sigma _2}`$. Together with the local splitting propertities in (1) above, in fact $`\phi _{\sigma _2\sigma _1}`$ is a regular matrix-valued function on $`U_{\sigma _1}U_{\sigma _2}`$. Pasting maps also have to satisfy the cocycle condition : $`\phi _{\sigma _3\sigma _1}=\phi _{\sigma _3\sigma _2}\phi _{\sigma _2\sigma _1}`$ over $`U_{\sigma _1}U_{\sigma _2}U_{\sigma _3}`$ for every triple of $`\sigma _1`$, $`\sigma _2`$, $`\sigma _3`$ in $`\mathrm{\Sigma }`$. This implies that the full set of pasting maps between affine charts of $`X_\mathrm{\Sigma }`$ is determined by the set of pasting maps $`\phi _{\sigma _2\sigma _1}`$ with $`\sigma _1,\sigma _2\mathrm{\Sigma }(n)`$ that satisfy the cocycle condition.
These observations lead to Kaneyamaโs data for equivariant vector bundles over $`X_\mathrm{\Sigma }`$.
The bundle data and the classification after Kaneyama.
(a) The data of local trivialization : a collection of weight systems.
For $`\sigma \mathrm{\Sigma }(n)`$, $`x_\sigma `$ is a fixed point of the toric action; thus $`_{x_\sigma }`$ is an invariant fiber of the lifted toric action. Associated to the representation of $`\text{๐}^n`$ on $`_{x_\sigma }`$ is the weight system $`๐ฒ_\sigma M`$. $`๐ฒ_\sigma `$ determones the local trivialization of $`|_{U_\sigma }`$: $`|_{U_\sigma }=_{\chi ๐ฒ_\sigma }U_\sigma \times E_\sigma ^\chi `$.
(b) The data of pasting : net of weight systems and pasting maps.
These weight systems must satisfy a compatibility condition as follows. Let $`\tau =\sigma _1\sigma _2\mathrm{\Sigma }(n1)`$ be the common codimension-$`1`$ wall of two maximal cones $`\sigma _1`$ and $`\sigma _2`$, then the stabilizer $`\text{Stab}(x_\tau )`$ of $`x_\tau `$ is an $`(n1)`$-subtorus in $`\text{๐}_N`$ associated to the sublattice in $`N`$ spanned by $`\tau `$. Associated to the representation of $`\text{Stab}(x_\tau )`$ on the fiber $`_{x_\tau }`$ is a weight system $`๐ฒ_\tau M/M(\tau )`$, where $`M(\tau )=\tau ^{}M`$. The projection map $`MM/M(\tau )`$ induces the maps
$$๐ฒ_{\sigma _1}\stackrel{\pi _{\tau \sigma _1}}{}๐ฒ_\tau \stackrel{\pi _{\tau \sigma _2}}{}๐ฒ_{\sigma _2}$$
between weight systems. Since they correspond to the refinement of the $`\text{Stab}(\tau )`$-weight spaces to the $`\text{๐}_N`$-weight spaces, the holomorphicity requiremnent of equivariant pasting maps implied that both of these maps are surjective. Thus, $`\{๐ฒ_\sigma |\sigma \mathrm{\Sigma }(n)\}`$ form a net of weight systems. Figure 2-1 indicates how the net of weight systems may look like for $`\mathrm{\Sigma }`$ that comes from the normal cone of a strong convex polyhedron $`\mathrm{\Delta }`$ in $`M`$.
The equivariant pasting maps are given by a map
$$P:\mathrm{\Sigma }(n)\times \mathrm{\Sigma }(n)\text{GL}(r,\text{})$$
that satisfies
$$P(\sigma _3,\sigma _2)P(\sigma _2,\sigma _1)=P(\sigma _3,\sigma _1).$$
$`P`$ gives a set of compatible pasting maps for the fiber $`_{x_0}`$ in different $`|_{U_\sigma }`$ with respect to the bases given by the weight space decomposition. Holomorphicity condition for its equivariant extension to over $`U_{\sigma _1}U_{\sigma _2}`$ requires that:
> For any $`\tau =\sigma _1\sigma _2\mathrm{\Sigma }(n1)`$, let $`๐ฒ_{\sigma _i}=(\chi _{\sigma _i1},\mathrm{},\chi _{\sigma _ir})`$, written with multiplicity given by the dimension of the corresponding weight space. Then $`P(\sigma _2,\sigma _1)_{ij}=0`$ if $`\chi _{\sigma _2i}\chi _{\sigma _1j}M\tau ^{}`$.
(c) Equivalence of the bundle data.
Given $`\mathrm{\Sigma }`$, two bundle data $`(๐ฒ,P)`$ and $`(๐ฒ^{},P^{})`$ are said to be equivariant if $`๐ฒ=๐ฒ^{}`$ and there is a map $`\rho :\mathrm{\Sigma }(n)\text{GL}(r,\text{})`$ such that $`P^{}(\sigma _2,\sigma _1)=\rho (\sigma _2)P(\sigma _2,\sigma _1)\rho (\sigma _1)^1`$. Equivariant data determine isomorphic linearized equivariant vector bundles over $`X_\mathrm{\Sigma }`$.
## 3 The splitting type of an equivariant vector bundle.
Recall first a theorem of Grothendieck (\[Gro\] Theorem 2.1), which says that any holomorphic vector bundle over $`\text{}\mathrm{P}^1`$ splits into the direct sum of a unique set of line bundles. The following definition follows from \[L-L-Y2\]:
Definition 3.1 \[splitting type\]. Let $``$ be an equivariant vector bundle of rank $`r`$ over a toric manifold $`X_\mathrm{\Sigma }`$. Suppose that there exist nontrivial equivariant line bundles $`L_1,\mathrm{},L_r`$ over $`X_\mathrm{\Sigma }`$ such that each $`c_1(L_i)`$ is either $`0`$ or $`<0`$ and that the restriction $`|_{V(\tau )}`$ is isomorphic to the direct sum $`(_{i=1}^rL_i)|_{V(\tau )}`$ for any $`\tau \mathrm{\Sigma }(n1)`$. Then $`\{L_1,\mathrm{},L_r\}`$ is called a splitting type of $``$.
Definition 3.2 \[system of splitting numbers\]. For each $`\tau \mathrm{\Sigma }(n1)`$, suppose that
$$|_{V(\tau )}=_{i=1}^r๐ช(d_i^\tau )\text{with}d_1^\tau d_2^\tau \mathrm{}d_r^\tau .$$
From Grothendieckโs theorem, $`(d_1^\tau ,\mathrm{},d_r^\tau )`$ is uniquely determined by $``$. We shall call the set
$$\Xi ()=\{(d_1^\tau ,\mathrm{},d_r^\tau )|\tau \mathrm{\Sigma }(n1)\}$$
the system of splitting numbers associated to $``$.
To compute the splitting types of $``$, we first extract the bundle data of $`|_{V(\tau )}`$ from that of $``$ and then compute $`\Xi ()`$ from the bundle data of $`|_{V(\tau )}`$ by weight bootstrapping. Using these numbers, one can then determine all the splitting types of $``$ by the augmented intersection matrix $`Q`$ associated to $`\mathrm{\Sigma }`$. Let us now turn to the details.
We shall assume that the rank $`r`$ of $``$ $`2`$.
The bundle data of $`|_{V(\tau )}`$.
Let $`\tau =\sigma _1\sigma _2\mathrm{\Sigma }(n1)`$, $`E=|_{V(\tau )}`$, $`U_1=U_{\sigma _1}V(\tau )`$, and $`U_2=U_{\sigma _1}V(\tau )`$. Let $`v_\tau `$ be a lattice point in the interior of $`\tau `$, then the pasting map $`\phi _{21}(x_\tau )`$ for $`E`$ from $`E|_{U_1}`$ to $`E|_{U_2}`$ over $`x_\tau `$ is given by $`lim_{z0}(\lambda _{v_\tau }(z)P(\sigma _2,\sigma _1)\lambda _{v_\tau }(z)^1)`$, where $`z\text{}^{}`$. Since $`\text{Stab}(V(\tau ))`$ acts on $`E`$ via the $`T_N`$-action on $`E`$ and the pasting map for $`E`$ is $`T_N`$-equivariant and, hence, commutes with the $`\text{Stab}(V(\tau ))`$-action, $`E`$ can be decomposed into a direct sum of $`\text{Stab}(V(\tau ))`$-weight subbundles:
$$E=_{\chi ๐ฒ_\tau }E^\chi .$$
Indeed, for $`\chi ๐ฒ_\tau `$, $`E^\chi |_{U_1}=_{\chi ^{}\pi _{\tau \sigma _1}^1(\chi )}U_1\times E_{\sigma _1}^\chi ^{}`$; and similarly for $`E^\chi |_{U_2}`$. Thus, up to a permutation of elements in the basis, we may assume that $`\phi _{21}(x_\tau )`$ is in a block diagonal form with each block labelled by a distinct $`\chi ๐ฒ_\tau `$. Let us now turn to the weight system for $`E`$ at $`x_{\sigma _1}`$ and $`x_{\sigma _2}`$.
Let $`\tau _{\sigma _1}^{}`$ be the primitive lattice point of $`\sigma _1^{}\tau ^{}`$ in $`M`$ and $`v_{\sigma _1}`$ be a lattice point in $`N`$ such that $`\tau _{\sigma _1}^{},v_{\sigma _1}=1`$. Let $`\lambda _{v_{\sigma _1}}`$ be the corresponding one-parameter subgroup in $`\text{๐}_N`$. Then $`\lambda _{v_{\sigma _1}}`$ acts on $`O_\tau `$ freely and transitively with $`lim_{z0}\lambda _{v_{\sigma _1}}(z)x_\tau =x_{\sigma _1}`$. Let
$$๐ฒ_{\sigma _1}=\{\chi _{\sigma _11},\chi _{\sigma _12},\mathrm{}\}\text{and}๐ฒ_{\sigma _2}=\{\chi _{\sigma _21},\chi _{\sigma _22},\mathrm{}\}$$
be the set of $`\text{๐}_N`$-weights at $`x_{\sigma _1}`$ and $`x_{\sigma _2}`$ respectively. Then, as a $`\lambda _{v_{\sigma _1}}`$-equivariant bundle with the induced linearization from the linearization of $``$, the corresponding $`\lambda _{\sigma _1}`$-weights of $`E_{x_{\sigma _1}}`$ and $`E_{x_{\sigma _2}}`$ are given respectively by
$$๐ฒ_{\sigma _1}^{\text{๐}^1}=\{\chi _{\sigma _11},v_{\sigma _1},\chi _{\sigma _12},v_{\sigma _1},\mathrm{}\}\text{and}๐ฒ_{\sigma _2}^{\text{๐}^1}=\{\chi _{\sigma _21},v_{\sigma _1},\chi _{\sigma _22},v_{\sigma _1},\mathrm{}\}.$$
Notice that both sets depend on the choice of $`v_{\sigma _1}`$; however, different choices of $`v_{\sigma _1}`$ will lead only to an overall shift of $`๐ฒ_{\sigma _1}^{\text{๐}^1}๐ฒ_{\sigma _2}^{\text{๐}^1}`$ by an integer.
From bundle data to splitting numbers : weight bootstrapping.
Recall first the following fact by Grothendieck \[Gro\]:
Fact 3.3 \[Grothendieck\]. Given a holomorphic vector bundle $`E`$ of rank $`r`$ over $`\text{}\mathrm{P}^1`$. Let $`E_0=\{0\}E_1\mathrm{}E_r=E`$ be a filtration of $`E`$ such that $`E_i/E_{i1}`$ is a line bundle for $`1ir`$ and that the degree $`d_i`$ of $`E_i/E_{i1}`$ form a non-increasing sequence. Then $`E`$ is isomorphic to the direct sum $`_{i=1}^r(E_i/E_{i1})`$.
Following previous discussions and notations, we only need to work out the splitting numbers for each $`E^\chi `$. Thus, without loss of generality, we may assume that $`๐ฒ_\tau =\chi `$ in the following discussion.
Fix a $`v_{\sigma _1}`$. Note that in terms of the one-parameter subgroup $`\lambda _{v_{\sigma _1}}`$ acting on $`V(\tau )`$, $`x_{\sigma _1}`$ has coordinate $`0`$ while $`x_{\sigma _2}`$ has coordinate $`\mathrm{}`$. Let
$$E|_{U_1}=_{i=1}^aU_1\times E_1^{\chi _{1i}}\text{and}E|_{U_2}=_{j=1}^bU_2\times E_2^{\chi _{2j}}$$
be the induced $`\text{๐}^1`$-weight space decomposition of $`E|_{U_1}`$ and $`E|_{U_2}`$ respectively, and $`\phi _{12}(x_\tau )`$ be the pasting map at $`x_\tau `$ from $`(E|_{U_1})|_{x_\tau }`$ to $`(E|_{U_2})|_{x_\tau }`$. We assume that $`\chi _{11}>\mathrm{}>\chi _{1a}`$ and $`\chi _{21}<\mathrm{}<\chi _{2b}`$. Our goal is now to work out a filtration of $`E`$, using the given bundle data, that satisfies the property in the above fact.
Let $`v`$ be a non-zero vector in the fiber $`E_{x_\tau }`$ over $`x_\tau `$. Then associated to the $`\text{๐}^1`$-orbit $`\text{๐}^1v`$ of $`v`$ is a line bundle $`_v`$ that contains $`\text{๐}^1v`$ as a meromorphic section $`s_v`$. Let $`v_1^{\chi _{1i^{}}}`$ be the lowest $`\text{๐}^1`$-weight component of $`v`$ in $`E|_{U_1}`$ and $`v_2^{\chi _{2j^{}}}`$) be the highest $`\text{๐}^1`$-weight component of $`v`$ in $`E|_{U_2}`$. Then $`_v|_{x_{\sigma _1}}`$ (resp. $`_v|_{x_{\sigma _2}}`$) lies in $`E_1^{\chi _1i^{}}`$ (resp. $`E_2^{\chi _{2j^{}}}`$) and the meromorphic section $`s_v`$ is holomorphic over $`O_\tau `$ with a zero at $`x_{\sigma _1}`$ (resp. $`x_{\sigma _2}`$) of order $`\chi _{1i^{}}`$ (resp. $`\chi _{2j^{}}`$). (Here, a zero of order $`k`$ means the same as a pole of order $`k`$.) This shows that
$$_v=๐ช(\chi _{1i^{}}\chi _{2j^{}}).$$
Notice that, from the previous discussion, this degree is independent of the choices of $`v_{\sigma _1}`$. We shall now proceed to construct a $`\text{๐}^1`$-equivariant line subbundle in $`E`$ that achieves the maximal degree.
Fix a basis for the $`\text{๐}^1`$-weight spaces in the local trivialization of $`E`$, then $`\phi _{21}(x_\tau )`$ is expressed by a matrix $`A`$ which admits a weight-block decomposition $`A=[A_{\chi _{2j},\chi _{1i}}]_{ji}`$. Consider the chain of submatrices $`B_{kl}`$ in $`A`$ that consists of weight blocks $`A_{\chi _{2j},\chi _{1i}}`$ with $`j=k,\mathrm{},b`$ and $`i=1,\mathrm{},l`$, as indicated in Figure 3-1.
Each $`B_{kl}`$ gives the linear map from $`_{i=1}^lE_1^{\chi _{1i}}`$ to $`_{j=k}^bE_2^{\chi _{2j}}`$ induced by $`A`$. Let $`N_{kl}`$ be the kernel of $`B_{kl}`$, as a subspace in $`(E|_{U_1})_{x_\tau }`$. Define
$$_i=E_1^{\chi _{11}}\mathrm{}E_1^{\chi _{1i}}\text{at }x_\tau ,i=1,\mathrm{},a.$$
Then one has the following sequence of filtrations:
$$\begin{array}{ccccccccc}\{0\}& & _1& & \mathrm{}& & _a& =(E|_{U_1})_{x_\tau }\hfill & \\ & & & & \mathrm{}& & & & \\ & & N_{b1}& & \mathrm{}& & N_{ba}& & \\ & & & & \mathrm{}& & & & \\ & & \mathrm{}& & \mathrm{}& & \mathrm{}& & \\ & & & & \mathrm{}& & & & \\ & & N_{11}& & \mathrm{}& & N_{1a}& & .\hfill \end{array}$$
Let
$$_{ij}=_i_{i1}N_{ji},\text{for}i=1,\mathrm{},a,j=1,\mathrm{},b;$$
then
$$(E|_{U_1})_{x_\tau }=_{i,j}_{ij}.$$
By construction, if $`_{ij}`$ is non-empty, then, for any $`v`$ in $`_{ij}`$, $`\text{deg}(_v)=\chi _{1i}\chi _{2j}`$. Thus we may define the characteristic number $`\mu (_{ij})`$ for $`_{ij}`$ non-empty to be $`\chi _{1i}+\chi _{2j}`$. Now let
$$d_1=\text{max}\{\mu (_{ij})|_{ij}\text{ non-empty, }i=1,\mathrm{},a,j=1,\mathrm{},b\}$$
and $`v_{ij}`$ for some $`(ij)`$ that realizes $`d_1`$; then by construction $`E_1=_v`$ is a $`\text{๐}^1`$-equivariant line subbundle of $`E`$ that achieves the maximal possible degree.
Suppose the basis for the weight spaces in the local trivialization of $`E`$ are given by $`(e_{11},\mathrm{},e_{1r})`$ and $`(e_{21},\mathrm{},e_{2r})`$ respectively. Let $`v=_{i=1}^kc_ie_{1i}`$ with $`c_k0`$ in $`E|_{U_1}`$ and $`v=_{i=1}^k^{}c_i^{}e_{2i}`$ with $`c_k^{}^{}0`$ in $`E|_{U_2}`$. Then, by replacing $`e_{1k}`$ and $`e_{2k^{}}`$ by $`v`$ and putting it as the first element in the bases, one shows
Lemma 3.4. There exist a new weight space decomposition of $`E|_{U_1}`$ and $`E|_{U_2}`$ respectively and a choice of the basis for the new weight spaces such that $`E_1|_{x_\tau }`$ is spanned by the first element in the basis.
This renders the pasting map into the form:
$$A=\left[\begin{array}{cc}1& \\ 0& A_1\end{array}\right],$$
where $`0`$ is the $`(r1)`$-dimensional zero vector and $`A_1`$ is a nondegenerate $`(r1)\times (r1)`$ matrix. With respect to the new trivialization, the $`\text{๐}^1`$-weight spaces and their basis descends then to the quotient $`E/E_1`$ with pasting map given by $`A_1`$.
Repeating the discussion $`r`$ times, one obtains $`\text{๐}^1`$-equivariant line subbundles
$$E_1E,E_2/E_1E/E_1,\mathrm{},E_{r1}/E_{r2}E/E_{r2},E_r/E_{r1}.$$
of maximal degree in each pair and the associated filtration
$$E_0=\{0\}E_1\mathrm{}E_r=E.$$
Lemma 3.5. The filtration of $`E`$ obtained above satisfies the condition of Grothendieck in Fact 3.3.
Proof. Since $`E_i/E_{i1}`$ is a $`\text{๐}^1`$-equivariant line subbundle of maximal degree in $`E/E_{i1}`$, it must be so also in $`E_{i+1}/E_{i1}`$. Together with the fact that $`E_{i+1}/E_i=(E_{i+1}/E_{i1})/(E_i/E_{i1})`$, we only need to justify the claim for the rank $`2`$ case. Thus, assume that $`E`$ is of rank $`2`$ and the filtration is given by $`\{0\}E_1E`$. By Lemma 3.4, we may choose a basis compatibe with the $`\text{๐}^1`$-weight space decomposition and the filtration such that the pasting map $`\phi _\mathrm{}0`$ and its inverse are given respectively by the following matrices with the $`\text{๐}^1`$-weight indicated:
$$\begin{array}{c}\chi _{11}\chi _{12}\\ A=\left[\begin{array}{cc}1& \\ 0& 1\end{array}\right]\begin{array}{c}\chi _{21}\hfill \\ \chi _{22}\hfill \end{array}\end{array}\text{and}\begin{array}{c}\chi _{21}\chi _{22}\\ A^1=\left[\begin{array}{cc}1& \hfill \\ 0& \hfill 1\end{array}\right]\begin{array}{c}\chi _{11}\hfill \\ \chi _{12}\hfill \end{array}\end{array}\text{.}$$
This implies that $`\text{deg}E_1=\chi _{11}\chi _{21}`$ and $`\text{deg}(E/E_1)=\chi _{12}\chi _{22}`$. The linear independency of the line bundles associated to the column vectors of $`A`$ at $`z=0`$ requires that $`\chi _{22}\chi _{21}`$. Similarly, the linear independency of the line bundles associated to the column vectors of $`A^1`$ at $`z=\mathrm{}`$ requires that $`\chi _{12}\chi _{11}`$. Consequently, $`\text{deg}E_1\text{deg}(E/E_1)`$. This concludes the proof.
$`\mathrm{}`$
Consequently, by Fact 3.3 one obtains the decomposition of $`E`$ as the direct sum of line bundles and the splitting numbers.
Remark 3.6 \[ weight matching \]. For $`\tau =\sigma _1\sigma _2\mathrm{\Sigma }(n1)`$, if all the $`\text{Stab}(x_\tau )`$-weight spaces are $`1`$-dimensional, then up to a permutation of the elements in bases, the pasting map $`\phi _{21}(x_\tau )`$ becomes diagonal and the correpondences $`๐ฒ_{\sigma _1}๐ฒ_\tau ๐ฒ_{\sigma _2}`$ are bijective. Let $`\chi _{\sigma _1i}\chi _{\tau i}\chi _{\sigma _2i}`$, $`i=1,\mathrm{},r`$, be the correpondences of weights. Then, up to a permutation, the splitting number of $``$ over $`V(\tau )`$ is given by
$$(\chi _{\sigma _11}\chi _{\sigma _21},v_{\sigma _1},\mathrm{},\chi _{\sigma _1r}\chi _{\sigma _2r},v_{\sigma _1}),$$
where recall that $`\tau _{\sigma _1}^{}=\tau ^{}\sigma _1^{}`$ and $`\tau _{\sigma _1}^{},v_{\sigma _1}=1`$.
From the system of splitting numbers to the splitting types of $``$ .
Following the notation in Sec. 1, let $`\mathrm{\Sigma }(n1)=\{\tau _1,\mathrm{},\tau _I\}`$ and
$$\Xi ()=\{(d_1^{\tau _i},\mathrm{},d_r^{\tau _i})|i=1,\mathrm{},r\}.$$
be the system of splitting numbers associated to $``$. Let $`R`$ be the $`I\times r`$ matrix whose $`i`$-th row is $`(d_1^{\tau _i},\mathrm{},d_r^{\tau _i})`$. Recall the augmented intersection matrix $`Q`$ from Sec. 1. Then the problem of finding splitting types of $``$ is equivalent to finding out matrices $`R^{}`$ obtained by row-wise permutations of $`R`$ such that:
> (1) Each column of $`R^{}`$ has only all positive, all zero, or all negative entries.
>
> (2) The following matrix linear equation has an integral solution:
>
> $$QX=R^{},$$
>
> where $`X`$ is an $`J\times r`$ matrix.
Let $`X_{kl}`$ be the $`(k,l)`$-entry of $`X`$. Then associated to the $`r`$-many column vectors of the solution matrix $`X`$ are the line bundles $`L_l`$ represented by $`_{k=1}^JX_{kl}D(v_k)`$, for $`l=1,\mathrm{},r`$. From Fact 1.2 in Sec.1, Condition (1) above for $`R^{}`$ means that $`c_1(L_l)`$ is either $`0`$ or $`<0`$. Such set of line bundles gives then a splitting type of $``$ by construction. If there exist no such $`(R^{},X)`$, then $``$ does not admit a splitting type. Finding all such $`R^{}`$ and solving the matrix $`X`$ can be achieved by using a computer.
## 4 The splitting type of some examples.
In this section, we compute the splitting types of some equivaraint vector bundles over toric manifolds to illustrate the ideas in previous sections and also for future use. The details of the toric manifolds used here can be found in \[Fu\] and \[Od2\].
Example 4.1 \[equivariant vector bundles of rank 2 over $`\text{}\mathrm{P}^2`$\]. Recall first (\[Fu\]) the toric data for $`\text{}\mathrm{P}^2`$, as illustrated in Figure 4-1(a). Let $``$ be an indecomposable equivariant vector bundle of rank 2 over $`\text{}\mathrm{P}^2=\text{Proj}(\text{}[u_0,u_1,u_2])`$. From \[Ka1\], $``$ is isomorphic to $`_{a,b,c,n}=(a,b,c)๐ช(n)`$ or its dual bundle for some positive integers $`a,b,c`$ and integer $`n`$, where $`(a,b,c)`$ is the rank 2 bundle defined by the exact sequence
$$\begin{array}{cccccccccc}0& & ๐ช_{\text{}\mathrm{P}^2}& & ๐ช(a)๐ช(b)๐ช(c)& & (a,b,c)& & 0& \\ & & 1& & (u_0^a,u_1^b,u_2^c)& & & & & .\hfill \end{array}$$
From the bundle data of $`(a,b,c)`$ as worked out in \[Ka1\], the weight systems for $`(a,b,c)`$ at the distinguished points $`x_{\sigma _1},x_{\sigma _1}`$ and $`x_{\sigma _3}`$ are given respectively by (cf. Figure 4-1(b))
$$W_1=\{(a,0),(0,b)\},W_2=\{(b,b),(c,0)\},W_3=\{(a,a),(0,c)\}.$$
Comparing with the toric data for $`\text{}\mathrm{P}^2`$ and following the discussions in Sec. 3, in particular Remark 3.6, one has
$$\begin{array}{c}(a,b,c)|_{\overline{x_{\sigma _1}x_{\sigma _2}}}=๐ช(a+c)๐ช(b),(a,b,c)|_{\overline{x_{\sigma _2}x_{\sigma _3}}}=๐ช(a+b)๐ช(c),\text{and}\hfill \\ (a,b,c)|_{\overline{x_{\sigma _1}x_{\sigma _3}}}=๐ช(b+c)๐ช(a).\hfill \end{array}$$
Consequently,
$$\begin{array}{c}|_{\overline{x_{\sigma _1}x_{\sigma _2}}}=๐ช(a+c+n)๐ช(b+n),|_{\overline{x_{\sigma _2}x_{\sigma _3}}}=๐ช(a+b+n)๐ช(c+n),\text{and}\hfill \\ |_{\overline{x_{\sigma _1}x_{\sigma _3}}}=๐ช(b+c+n)๐ช(a+n).\hfill \end{array}$$
Thus, up to permutations, the system of splitting numbers associated to $``$ is
$$\Xi ()=\{(a+c+n,b+n),(a+b+n,c+n),(b+c+n,a+n)\}$$
and
$$R=\left[\begin{array}{ccc}a+c+n& b+n& \\ a+b+n& c+n& \\ b+c+n& a+n& \end{array}\right].$$
From Figure 4-1(a), the augmented intersection matrix $`Q`$ for $`\text{}\mathrm{P}^2`$ is given by
$$Q=\left[\begin{array}{ccc}\hfill 1& \hfill 1& \hfill 1\\ \hfill 1& \hfill 1& \hfill 1\\ \hfill 1& \hfill 1& \hfill 1\end{array}\right].$$
Performing row-wise permutations to $`R`$, one obtains four possible $`R^{}`$, up to overall permutations of column vectors. Solving the matrix equation $`QX=R^{}`$, one concludes that
Corollary. For the indecomposable equivaraint rank $`2`$ bundle $`_{a,b,c,n}`$ over $`\text{}\mathrm{P}^2`$ to admit a splitting type, one must have $`a=b=c`$. In this case, the splitting type of $`_{a,a,a,n}`$ is unique and is given by $`(๐ช(2a+n),๐ช(a+n))`$.
This concludes the example.
$`\mathrm{}`$
Example 4.2 \[(co)tangent bundle of toric manifolds\]. Notice that, since $`T_{}X`$ and $`T^{}X`$ are dual to each other, their splitting types are negative to each other. Thus, we only need to consider $`T_{}X`$. Let us compute first the splitting numbers of $`T_{}X_\mathrm{\Sigma }`$. Let $`\tau =\sigma _1\sigma _2\mathrm{\Sigma }(n1)`$ with
$$\sigma _1=[v_1,\mathrm{},v_{n1},v_n]\text{and}\sigma _2=[v_1,\mathrm{},v_{n1},v_n^{}].$$
These vertices in $`\sigma _1\sigma _2`$ satisfy a linear equation of the form
$$v_{j_n}+v_{j_n^{}}+a_1v_{j_1}+\mathrm{}+a_{n1}v_{j_{n1}}=\mathrm{\hspace{0.33em}0},$$
for some unique integers $`a_1,\mathrm{},a_n`$ determined by $`\sigma _1\sigma _2`$. Let $`(e^1,\mathrm{},e^n)`$ be the dual basis in $`M`$ with respect to $`(v_1,\mathrm{},v_n)`$. Then
$$\sigma _1^{}=[e^1,\mathrm{},e^{n1},e^n]\text{and}\sigma _2^{}=[e^1a_1e^n,\mathrm{},e^{n1}a_{n1}e^n,e^n].$$
Consequently, $`\chi _{\sigma _1i}=e^i`$ for $`i=1,\mathrm{},n`$, $`\chi _{\sigma _2i}=e^ia_ie^n`$ for $`i=1,\mathrm{},n1`$; and $`\chi _{\sigma _2n}=e^n`$. Choosing $`v_{\sigma _1}=v_n`$ and by the discussion in Sec. 3, one concludes that the splitting number of $`T_{}X_\mathrm{\Sigma }`$ over $`V(\tau )`$ is given by
$$(a_1,\mathrm{},a_{n1},\mathrm{\hspace{0.17em}2}),$$
up to a permutation. Thus, the system $`\Xi (T_{}X_\mathrm{\Sigma })`$ of spliting numbers associated to $`T_{}X_\mathrm{\Sigma }`$ is already coded in $`\mathrm{\Sigma }`$, as it should be.
In the dual picture, if $`X_\mathrm{\Sigma }`$ is projective and, hence, $`\mathrm{\Sigma }`$ is realized as the normal fan of a strongly convex polyhedron $`\mathrm{\Delta }`$ in $`M`$. Let $`m_\sigma `$ be the vertex of $`\mathrm{\Delta }`$ associated to $`\sigma \mathrm{\Sigma }(n)`$. Then $`๐ฒ_\sigma `$ is the set of primitive vectors that generate the tangent cone of $`\mathrm{\Delta }`$ at $`m_\sigma `$. If $`\tau =\sigma _1\sigma _2\mathrm{\Sigma }(n1)`$, then $`m_{\sigma _1}`$ and $`m_{\sigma _2}`$ are connected by an edge $`\overline{m_{\sigma _1}m_{\sigma _2}}`$ of $`\mathrm{\Delta }`$ that is parallel to $`\tau ^{}`$. Thus the projection $`MM/M(\tau )`$ is given by the projection along the $`\overline{m_{\sigma _1}m_{\sigma _2}}`$-direction. The strong convexity of $`\mathrm{\Delta }`$ implies that $`๐ฒ_{\sigma _1}`$ and $`๐ฒ_{\sigma _2}`$ match up bijectively under this projection. Thus $`\Xi (X_\mathrm{\Sigma })`$ can be also read off directly from $`\mathrm{\Delta }`$.
We can now compute the splitting type of some concrete examples. The result shows that : Not every tangent bundle of a toric manifold admits a splitting type.
(a) The projective space $`\text{}\mathrm{P}^n`$. Let $`(v_1,\mathrm{},v_n)`$ be a basis of $`N`$ and $`v_{n+1}=(v_1\mathrm{}v_n)`$ and $`\mathrm{\Sigma }`$ be the fan whose maximal cones are generated by every independent $`n`$ elements in $`\{v_1,\mathrm{}v_{n+1}\}`$. Then $`\text{}\mathrm{P}^n=X_\mathrm{\Sigma }`$. By construction, $`\mathrm{\Sigma }(1)`$ is given by $`\{v_1,\mathrm{},v_{n+1}\}`$ with $`v_1+\mathrm{}v_{n+1}=0`$. Thus the splitting number for $`T_{}\text{}\mathrm{P}^n`$ over any invariant $`\text{}\mathrm{P}^1`$ is given by
$$(\mathrm{\hspace{0.17em}2},\underset{n1}{\underset{}{1,\mathrm{},\mathrm{\hspace{0.17em}1}}}).$$
The augmented intersection matrix $`Q`$ has all of its entries equal to $`1`$. From this, one concludes that the splitting type of $`\text{}\mathrm{P}^n`$ is unique and is given by
$$(๐ช(2),\underset{n1}{\underset{}{๐ช(1),\mathrm{},๐ช(1)}}).$$
(b) The Hirzebruch surface $`\text{๐ฝ}_a`$. The toric data for the Hurzebruch $`\text{๐ฝ}_a`$ and its weighted circular graph \[Od2\] is given in Figure 4-2.
Consequently,
$$\Xi (T_{}\text{๐ฝ}_a)=\{(2,0),(2,a),(2,0),(2,a)\}$$
and its augmented intersection matrix is given by
$$Q=\left[\begin{array}{cccc}\hfill 0& \hfill 1& \hfill 0& \hfill 1\\ \hfill 1& \hfill a& \hfill 1& \hfill 0\\ \hfill 0& \hfill 1& \hfill 0& \hfill 1\\ \hfill 1& \hfill 0& \hfill 1& \hfill a\end{array}\right].$$
Thus, the only Hirzebruch surface whose tangent bundle can admit a splitting type is when $`a=0`$ since the line bundle in a splitting type must be either positive or negative. For $`a=0`$, $`\text{๐ฝ}_0=\text{}\mathrm{P}^1\times \text{}\mathrm{P}^1`$. Let $`H^2(\text{}\mathrm{P}^1\times \text{}\mathrm{P}^1,\text{})=\text{}\text{}`$ from the product structure. Then direct computations as in Example 4.1 concludes that $`T_{}(\text{๐ฝ}_0)`$ admits a unique splitting type
$$(๐ช(2,2),๐ช(\text{๐ฝ}_0)),$$
where $`๐ช(2,2)`$ is the line bundle associated to $`(2,2)`$ in $`H^2(\text{๐ฝ}_0,\text{})`$ and $`๐ช(\text{๐ฝ}_0)`$ is the trivial line bundle.
(c) The blowups of $`\text{}\mathrm{P}^2`$ or $`\text{๐ฝ}_a`$. Recall from \[Fu\] and \[Od2\] that every complete nonsingular toric surface $`X`$ is obtained from $`\text{}\mathrm{P}^2`$ or $`\text{๐ฝ}_a`$, $`a>0`$, by a succession of blowups at the $`T_N`$-fixed points. Let $`(a_1,\mathrm{},a_s)`$ be the sequence of weights that appear in the weighted circular graph for $`X`$. Then
$$\Xi (T_{}X)=\{(2,a_1),\mathrm{},(2,a_s)\}.$$
Consequently, a necessary condition for $`T_{}X`$ to admit a splitting type is that the weights that appear in the graph must be all positive, all zero, or all negative. The augmented intersection matrix is given by
$$Q=\left[\begin{array}{cccccc}a_1& 1& & & & 1\\ 1& a_2& 1& & & \\ & 1& a_3& 1& & \\ & & \mathrm{}& \mathrm{}& \mathrm{}& \\ & & & 1& a_{n1}& 1\\ 1& & & & 1& a_n\end{array}\right],$$
where all the missing entries are $`0`$. From these data, the splitting type of $`T_{}X`$, if exists, can be worked out.
To provide more examples and also for the interest of string theory, equivariant blowups of $`\text{}\mathrm{P}^2`$ up to $`9`$ points that admits a splitting type are searched out by computer. These include del Pezzo type and $`\frac{1}{2}`$K3 type surfaces. It turns out that there are only $`8`$ of them, $`4`$ of del Pezzo type and $`4`$ of $`\frac{1}{2}`$K3 type. Their splitting types are listed in Table 4-1.
| topology of $`X`$ | $`k`$ | $`w=(a_1,\mathrm{},a_s)`$ | $`H_2(X;\text{})`$ | splitting type of $`T_{}X`$ | Remark |
| --- | --- | --- | --- | --- | --- |
| $`\text{}\mathrm{P}^2\mathrm{}\mathrm{\hspace{0.17em}3}\overline{\text{}\mathrm{P}^2}`$ | $`3`$ | $`(1,1,1,1,1,1)`$ | $`\text{}^4`$ | $`\begin{array}{c}((2,4,4,2),(1,2,2,1))\hfill \end{array}`$ | del Pezzo type |
| $`\text{}\mathrm{P}^2\mathrm{}\mathrm{\hspace{0.17em}5}\overline{\text{}\mathrm{P}^2}`$ | $`5`$ | $`(1,2,1,2,1,2,1,2)`$ | $`\text{}^6`$ | $`\begin{array}{c}((2,4,8,6,6,2),\hfill \\ (2,3,6,4,4,1))\hfill \end{array}`$ | del Pezzo type |
| $`\text{}\mathrm{P}^2\mathrm{}\mathrm{\hspace{0.17em}6}\overline{\text{}\mathrm{P}^2}`$ | $`6`$ | $`(1,2,2,1,2,2,1,2,2)`$ | $`\text{}^7`$ | $`\begin{array}{c}((2,4,8,14,8,4,2),\hfill \\ (2,3,6,11,6,3,2))\hfill \end{array}`$ | del Pezzo type |
| $`\text{}\mathrm{P}^2\mathrm{}\mathrm{\hspace{0.17em}7}\overline{\text{}\mathrm{P}^2}`$ | $`7`$ | $`(1,2,2,1,3,1,2,2,1,3)`$ | $`\text{}^8`$ | $`\begin{array}{c}((2,4,8,14,8,12,6,2),\hfill \\ (3,4,7,12,6,9,4,1))\hfill \end{array}`$ | del Pezzo type |
| $`\text{}\mathrm{P}^2\mathrm{}\mathrm{\hspace{0.17em}9}\overline{\text{}\mathrm{P}^2}`$ | $`9`$ | $`(1,2,2,2,1,4,1,2,2,2,1,4)`$ | $`\text{}^{10}`$ | $`\begin{array}{c}((2,4,8,14,22,10,20,12,6,2),\hfill \\ (4,5,8,13,20,8,16,9,4,1))\hfill \end{array}`$ | $`\frac{1}{2}`$K3 type |
| $`\text{}\mathrm{P}^2\mathrm{}\mathrm{\hspace{0.17em}9}\overline{\text{}\mathrm{P}^2}`$ | $`9`$ | $`(1,2,2,3,1,2,2,3,1,2,2,3)`$ | $`\text{}^{10}`$ | $`\begin{array}{c}((2,4,8,14,36,24,14,6,6,2),\hfill \\ (3,4,7,12,32,21,12,5,6,2))\hfill \end{array}`$ | $`\frac{1}{2}`$K3 type |
| $`\text{}\mathrm{P}^2\mathrm{}\mathrm{\hspace{0.17em}9}\overline{\text{}\mathrm{P}^2}`$ | $`9`$ | $`(1,2,3,1,2,3,1,2,3,1,2,3)`$ | $`\text{}^{10}`$ | $`\begin{array}{c}((2,4,8,22,16,12,22,12,4,2),\hfill \\ (3,4,7,20,14,10,19,10,3,2))\hfill \end{array}`$ | $`\frac{1}{2}`$K3 type |
| $`\text{}\mathrm{P}^2\mathrm{}\mathrm{\hspace{0.17em}9}\overline{\text{}\mathrm{P}^2}`$ | $`9`$ | $`(1,3,1,3,1,3,1,3,1,3,1,3)`$ | $`\text{}^{10}`$ | $`\begin{array}{c}((2,4,12,10,20,12,18,8,8,2),\hfill \\ (3,4,12,9,18,10,15,6,6,1))\hfill \end{array}`$ | $`\frac{1}{2}`$K3 type |
| (1) $`k`$ is the number of points blown up from $`\text{}\mathrm{P}^2`$. | | | | | |
| --- | --- | --- | --- | --- | --- |
| (2) $`H_2(X;\text{})`$ is generated by the first $`(s2)`$ divisors in the list $`w`$. | | | | | |
| (3) The line bundles in the splitting type are represented by divisors of $`X`$ as elements in $`H_2(X;\text{})`$. | | | | | |
Table 4-1. Complete list of $`T_{}S`$ and its splitting type for toric surfaces obtained from $`\text{}\mathrm{P}^2`$ via equivariant blowups up to $`9`$ points.
The fan for these toric surfaces are indicated in Figure 4-3.
This concludes the example.
$`\mathrm{}`$
## 5 Remarks and issues for further study.
In the previous section, we have illustrated how the data of a linearized equivariant vector bundle $``$ over a toric manifold $`X_\mathrm{\Sigma }`$ is used to determine its splitting type if it exists. It turns out that these examples are related to the following kind of exact sequence<sup>1</sup><sup>1</sup>1We thank Bong H. Lian for drawing our attention to this and the reference \[Ja\]:
$$0๐ช_{X_\mathrm{\Sigma }}\stackrel{\eta }{}_{i=1}^{r+1}๐ช_{X_\mathrm{\Sigma }}(D_i)\mathrm{\hspace{0.33em}0},$$
where $`D_i`$ are Cartier T-Weil divisors of $`X_\mathrm{\Sigma }`$ and $`\eta `$ is a holomorphic bundle inclusion. For such $``$, the system of splitting numbers $`\Xi ()`$ may be obtained directly from this exact sequence.
Example 5.1 \[equivariant vector bundles of rank $`n`$ over $`\text{}\mathrm{P}^n`$\]. Let $`[z_0:\mathrm{},z_n]`$ be the homogeneous coordinates of the projective space $`\text{}\mathrm{P}^n`$. Recall from \[Ka2\] that an indecomposable equivariant vector bundle $``$ of rank $`n`$ over $`\text{}\mathrm{P}^n`$ is isomorphic to either $`E๐ช_{\text{}\mathrm{P}^n}(d)`$ or $`E^{}๐ช_{\text{}\mathrm{P}^n}(d)`$ for some integer $`d`$ where $`E`$ is the equivariant vector bundle defined by the exact sequence
$$0๐ช_{\text{}\mathrm{P}^n}\stackrel{\eta }{}_{i=0}^n๐ช_{\text{}\mathrm{P}^n}(m_i)E\mathrm{\hspace{0.33em}0},$$
where $`m_i`$ are positive integers and $`\eta `$ sends $`1`$ to $`(z_0^{m_0},\mathrm{},z_n^{m_n})`$. Since the $`(n+1)n/2`$ many invariant $`\text{}\mathrm{P}^1`$ in $`\text{}\mathrm{P}^n`$ are given by
$$V_{ij}=\{[\mathrm{\hspace{0.17em}0},\mathrm{},\mathrm{\hspace{0.17em}0},z_i,\mathrm{\hspace{0.17em}0},\mathrm{},\mathrm{\hspace{0.17em}0},z_j,\mathrm{\hspace{0.17em}0},\mathrm{},\mathrm{\hspace{0.17em}0}]|(z_i,z_j)\text{}^2\{(0,0)\}\},$$
$`0i<jn`$, the above exact sequence, when restricted to $`V_{ij}`$, reduces to
$$\begin{array}{cccccccccc}0& & ๐ช_{V_{ij}}& & _{k=1}^{r+1}๐ช_{V_{ij}}(m_k)& & E|_{V_{ij}}& & 0& \\ & & 1& & (\mathrm{\hspace{0.17em}0},\mathrm{},\mathrm{\hspace{0.17em}0},z_i^{m_i},\mathrm{\hspace{0.17em}0},\mathrm{},\mathrm{\hspace{0.17em}0},z_j^{m_j},\mathrm{\hspace{0.17em}0},\mathrm{},\mathrm{\hspace{0.17em}0})& & & & & .\hfill \end{array}$$
It follows from the multiplicativity of total Chern class that
$$E|_{V_{ij}}๐ช_{\text{}\mathrm{P}^1}(m_i+m_j)๐ช_{\text{}\mathrm{P}^1}(m_0)\mathrm{}\widehat{๐ช_{\text{}\mathrm{P}^1}(m_i)}\mathrm{}\widehat{๐ช_{\text{}\mathrm{P}^1}(m_j)}\mathrm{}๐ช_{\text{}\mathrm{P}^1}(m_n)$$
and, hence, the system of splitting numbers of $`E`$ is
$$\Xi (E)=\{(m_i+m_j,m_0,\mathrm{},\widehat{m_i},\mathrm{},\widehat{m_j},\mathrm{},m_n)|\mathrm{\hspace{0.17em}0}i<jn\},$$
where terms with $`\widehat{}`$ are deleted. The augmented matrix in this case is
$$Q=\left[\begin{array}{ccc}1& \mathrm{}& 1\\ \mathrm{}& \mathrm{}& \mathrm{}\\ 1& \mathrm{}& 1\end{array}\right].$$
Without loss of generality, one may assume that $`0<m_0\mathrm{}m_n`$, then $`m_i<m_n+m_{n1}`$ for all $`i`$. Consequently, with the notation from Sec. 3, for $`QX=R^{}`$ to have a solution, one must have $`m_i+m_j=m_n+m_{n1}`$ for all $`i<j`$, which implies that $`m_0=\mathrm{}=m_n`$. One concludes therefore
Corollary. For the indecomposable equivaraint rank $`n`$ bundle $`=E๐ช_{\text{}\mathrm{P}^n}(d)`$ over $`\text{}\mathrm{P}^n`$ to admit a splitting type, one must have $`m_0=\mathrm{}=m_n`$. In this case, the splitting type of $``$ is unique and is given by $`(๐ช_{\text{}\mathrm{P}^n}(2m_0+d),\mathrm{},๐ช_{\text{}\mathrm{P}^n}(m_0+d))`$.
This generalizes Example 4.1. Note also that the case $`m_0=\mathrm{}=m_n=1`$ with $`d=0`$ corresponds to $`T_{}\text{}\mathrm{P}^n`$ and the above discussion double-checks part of Example 4.2.
$`\mathrm{}`$
One can generalize this example slightly to toric manifolds as follows. First, let us state a lemma, whose proof is straightforward.
Lemma 5.2 Given an exact sequence
$$0๐ช_{\text{}\mathrm{P}^1}\stackrel{\eta }{}๐ช_{\text{}\mathrm{P}^1}\left(_{i=1}^r๐ช_{\text{}\mathrm{P}^1}(m_i)\right)E\mathrm{\hspace{0.33em}0},$$
where $`\eta (1)=(s_0,s_1,\mathrm{},s_r)`$, such that $`s_0`$ is non-zero. Then $`E_{i=1}^r๐ช_{\text{}\mathrm{P}^1}(m_i)`$.
Given a toric $`n`$-fold $`X_\mathrm{\Sigma }`$, consider now the exact sequence
$$0๐ช_{X_\mathrm{\Sigma }}\stackrel{\eta }{}_{i=1}^{r+1}๐ช_{X_\mathrm{\Sigma }}(D_i)\mathrm{\hspace{0.33em}0}.$$
The restriction of the sequence to $`V(\tau )`$, $`\tau \mathrm{\Sigma }(n1)`$, is given by
$$0๐ช_{\text{}\mathrm{P}^1}\stackrel{\eta }{}_{i=1}^{r+1}๐ช_{\text{}\mathrm{P}^1}(D_iV(\tau ))E\mathrm{\hspace{0.33em}0}.$$
If this exact sequence is of the kind in Example 5.1 or Lemma 5.2 for all $`\tau \mathrm{\Sigma }(n1)`$, then $`\Xi ()`$ can be readily obtained and the splitting type of such $``$, if exists, can then be determined. Inspired from Example 5.1, to realize this, recall the Cox homogeneous coordinates of $`X_\mathrm{\Sigma }`$ from \[Co\] (cf. Sec. 1): let $`a=|\mathrm{\Sigma }(1)|`$, then $`X_\mathrm{\Sigma }`$ can be realized as a quotient $`X_\mathrm{\Sigma }=(\text{}^{\mathrm{\Sigma }(1)}Z(\mathrm{\Sigma }))/G`$. Let $`(z_1,\mathrm{},z_a)`$ be the standard coordinates of $`\text{}^a`$ and $`\tau =[v_{j_1},\mathrm{},v_{j_{n1}}]\mathrm{\Sigma }(n1)`$, then $`V(\tau )`$ can be realized as the quotient of the coordinate subspace: $`V(\tau )=\{z_{j_1}=\mathrm{}=z_{j_{n1}}\}/G`$. Furthermore, if $`\tau =\sigma _1\sigma _2`$, where
$$\sigma _1=[v_{j_1},\mathrm{},v_{j_{n1}},v_{j_n}]\text{and}\sigma _2=[v_{j_1},\mathrm{},v_{j_{n1}},v_{j_n^{}}],$$
then $`[z_{j_n}:z_{j_n^{}}]`$ serves as a homogeneous coordinates for $`V(\tau )\text{}\mathrm{P}^1`$. For all other $`i`$, $`\{z_i=0\}\{z_{j_1}=\mathrm{}=z_{j_{n1}}\}`$ lies in the exceptional subset $`Z(\mathrm{\Sigma })`$ and, hence, $`z_i`$ as an element in the homogeneous coordinate ring $`\text{}[z_1,\mathrm{},z_a]`$, graded by the Chow group $`A_{n1}(X_\mathrm{\Sigma })`$, descends to a non-zero section in $`๐ช_{X_\mathrm{\Sigma }}(D(v_i))|_{V(\tau )}๐ช_{\text{}\mathrm{P}^1}`$. In general, since $`๐ช_{X_\mathrm{\Sigma }}(D_1)๐ช_{X_\mathrm{\Sigma }}(D_2)=๐ช_{X_\mathrm{\Sigma }}(D_1+D_2)`$ for any Cartier T-Weil divisor $`D_1`$, $`D_2`$, for any monomial $`_kz_{j_k}^{\alpha _k}`$ with $`j_k\{j_1,\mathrm{},j_{n1},j_n,j_n^{}\}`$ and $`\alpha _k`$ positive integers, $`_kz_{j_k}^{\alpha _k}`$ descends to a non-zero section in $`๐ช_{X_\mathrm{\Sigma }}(_k\alpha _kD(v_{j_k}))|_{V(\tau )}๐ช_{\text{}\mathrm{P}^1}`$. This fact provides us with a guideline for defining $`\eta `$ so that exact sequences as in Lemma 5.2 can appear when restricted to invariant $`\text{}\mathrm{P}^1`$โs in $`X_\mathrm{\Sigma }`$. Such examples can be constructed plenty. Let us give an example below to illustrate the idea.
Example 5.3 \[simple rank 3 bundle over Hirzebruch surface\]. Let $`X_\mathrm{\Sigma }=\text{๐ฝ}_a`$ be a Hirzebruch surface (cf. Example 4.2 (b)). Consider the rank $`3`$ bundle $`(m_1,m_2,m_3,m_4)`$ over $`\text{๐ฝ}_a`$ defined by the exact sequence
$$\begin{array}{cccccccccc}0& & ๐ช_{\text{๐ฝ}_a}& & _{k=1}^4๐ช_{\text{๐ฝ}_a}(m_kD_{v_k})& & (m_1,m_2,m_3,m_4)& & 0& \\ & & 1& & (z_1^{m_1},z_2^{m_2},z_3^{m_3},z_4^{m_4})& & & & & ,\hfill \end{array}$$
where $`m_i`$ are positive integers. From Lemma 5.2 and the discussions above, one concludes that the system of splitting numbers of $`(m_1,m_2,m_3,m_4)`$ is given by
$$\Xi ()=\{(m_1,m_2,m_4),(m_1,m_2,m_3),(m_2,m_3,m_4),(m_1,m_3,m_4)\}.$$
Recall the augmented intersection matrix $`Q`$ for $`\text{๐ฝ}_a`$ from Example 4.2 (b). Through a tedious but straightforward algebra, one can show that the only case when $`(m_1,m_2,m_3,m_4)`$ admits a splitting type is when $`a=0`$ (i.e. $`X=\text{}\mathrm{P}^1\times \text{}\mathrm{P}^1`$) with $`m_1=m_3`$ and $`m_2=m_4`$. In this case, the splitting type is unique and is given by
$$๐ช_{\text{}\mathrm{P}^1\times \text{}\mathrm{P}^1}(m_2,m_1)๐ช_{\text{}\mathrm{P}^1\times \text{}\mathrm{P}^1}(m_1,m_2)๐ช_{\text{}\mathrm{P}^1\times \text{}\mathrm{P}^1}(m_1,m_2),$$
where we idetify $`\text{Pic}(X)`$ with $`H_2(X,\text{})\text{}\text{}`$, with generators
$$๐ช_{\text{}\mathrm{P}^1\times \text{}\mathrm{P}^1}(1,0)[\text{}\mathrm{P}^1\times ],๐ช_{\text{}\mathrm{P}^1\times \text{}\mathrm{P}^1}(0,1)[\times \text{}\mathrm{P}^1].$$
$`\mathrm{}`$
For more general $`\eta `$, the restriction of the exact sequence over $`X_\mathrm{\Sigma }`$ to each invariant $`\text{}\mathrm{P}^1`$ in $`X_\mathrm{\Sigma }`$ leads to an exact sequence of the form
$$0๐ช_{\text{}\mathrm{P}^1}\stackrel{\eta }{}_{i=1}^{r+1}๐ช_{\text{}\mathrm{P}^1}(m_i)_{\text{}\mathrm{P}^1}\mathrm{\hspace{0.33em}0}.$$
A complete study of how $`\eta `$ determines the splitting of $`_{\text{}\mathrm{P}^1}`$ as a direct sum of line bundles requires more work<sup>2</sup><sup>2</sup>2 We thank Jason Starr for discussions on this and the references \[Bri\], \[E-VV1\], and \[E-VV2\]..
We conclude the discussion of splitting types here and leave its further study and applications for another work.
Appendix. The computer code.
The computer code in Mathematica that carries out the computation in Example 4.2 (c) is attached below for reference,
```
(* This is a code in Mathematica. *)
(* The purpose of this code is to sort out and compute the splitting type of the tangent bundle
of toric surfaces. The result of computation is written to the file โma-result.txtโ. *)
(* Subroutines enclosed: BlowUp, BlowUpN, GenerateMatrix, (MAIN) SplittingType *)
(* Definition of the function โBlowUpโ. *)
(* โBlowUp[weightlist]โ generates the list of weights on the circular weighted graph obtained by
equivariant blowup at a $T_N$-fixed point of a toric surface represented by โweightlistโ.
Date of completion: 10/15/1999. Test: Tested correct. Date of last revision: 10/16/1999.
*)
BlowUp[ weightlist_ ] :=
Module[ { a1, a2, b, b1, b2, list, list1, list2, list3, m1, m2, newlist },
m1=Length[weightlist];
list1[i_] := ReplacePart[ weightlist,
{ weightlist[[i]]-1, -1, weightlist[[i+1]]-1 }, i+1 ];
list2[i_] := Delete[ list1[i], i];
list3[i_] := Flatten[ list2[i] ] ;
list=ReplacePart[ weightlist, {-1, weightlist[[1]]-1 }, 1 ];
list=ReplacePart[ list, list[[m1]]-1, m1];
list={ Flatten[list] };
newlist=Join[ list, Table[ list3[i] , {i, 1, m1-1}] ];
newlist=Union[newlist];
m2=Length[newlist];
Do[
a1=newlist[[i]];
a2=Reverse[a1];
b1=Table[ RotateRight[a1, i], {i, 1, m1+1} ];
b2=Table[ RotateRight[a2, i], {i, 1, m1+1} ];
b=Union[b1, b2];
newlist=Union[{a1}, Complement[newlist, b] ];
If[ m2>Length[newlist],
Return[newlist]
],
{i, 1, m2}
];
Return[newlist]
]
(* Definition of the function โBlowUpNโ. *)
(* โBlowUpN[weightlist, n]โ generates the list of weights on the circular weighted graph obtained
by consecutive equivariant blowup at a $T_N$-fixed points of a toric surface โnโ times, starting
from the one represented by โweightlistโ.
Date of completion: 10/15/1999. Test: Tested correct. Date of last revision: 10/16/1999.
*)
BlowUpN[ weight_, n_ ] :=
Module[ { m, newlist, oldlist, totallist},
totallist={weight};
oldlist={weight};
newlist={};
Do[
m=Length[oldlist];
Do[
newlist=Union[ newlist, BlowUp[ oldlist[[j]] ] ],
{j, 1, m}
];
oldlist=newlist;
totallist=Join[ totallist, newlist],
{i, 1, n}
];
totallist=Union[totallist];
Return[totallist];
]
(* Definition of the function โGenerateMatrixโ. *)
(* โGenerateMatrix[weightlist]โ generates a matrix following the rule discussed in the paper on
splitting types of equivariant vector bundle on toric manifolds .
Date of completion: 10/15/1999. Test: Tested correct. Date of last revision: 10/15/1999.
*)
GenerateMatrix[ weightlist_ ] :=
Module[ { listfirst, listlast, list1, list2, m, newlist, v },
m=Length[weightlist];
v=Table[ 0, {i, 1, m-2} ];
list1[i_]:=ReplacePart[ v, { 1, weightlist[[i]], 1 }, i-1 ];
list2[i_]:=Flatten[ list1[i] ];
listfirst=Flatten[ ReplacePart[ v, {1, weightlist[[1]], 1 }, 1 ] ];
listfirst={ RotateLeft[listfirst, 1] };
listlast=Flatten[ ReplacePart[ v, {1, weightlist[[m]], 1 }, 1 ] ];
listlast={ RotateLeft[listlast, 2] };
newlist=Join[ listfirst, Table[ list2[i], {i, 2, m-1} ], listlast ];
Return[newlist]
]
(* MAIN ROUTINE *)
(* Definition of the function โSplittingTypeโ. *)
(* โSplittingType[weight, n]โ sorts out from all the toric surfaces that arise from equivariant
blowups up to โnโ times of the toric surface whose associated weighted circular graph is given
by โweightโ those that admit a splitting type and computes their splitting types.
Date of completion: 10/15/1999. Test: Tested correct. Date of last revision: 10/17/1999.
*)
SplittingType[weight_, n_] :=
Module[ { b, m, matrix, t, totallist, t1, x1, x2 },
totallist=BlowUpN[weight, n];
m=Length[totallist];
Do[
t=totallist[[i]];
t1=Union[t];
If[ Complement[t1,{0}]===t1,
matrix=GenerateMatrix[t];
b=Table[2, {j, 1, Length[t]}];
x1=LinearSolve[matrix, b1];
If[ Length[x1]>=3,
x2=LinearSolve[matrix, t];
PutAppend[i, "ma-result.txt"];
PutAppend[t, "ma-result.txt"];
PutAppend[x1, "ma-result.txt"];
PutAppend[x2, "ma-result.txt"]
]
],
{i, 1, m}
];
]
(* Case of study *)
DeleteFile["ma-result.txt"];
SplittingType[{1, 1, 1}, 9];
``` |
warning/0002/quant-ph0002062.html | ar5iv | text | # Untitled Document
ON THE WEAK-COUPLING LIMIT
AND COMPLETE POSITIVITY
F. Benatti
Dipartimento di Fisica Teorica, Universitร di Trieste
Strada Costiera 11, 34014 Trieste, Italy
and
Istituto Nazionale di Fisica Nucleare, Sezione di Trieste
R. Floreanini
Istituto Nazionale di Fisica Nucleare, Sezione di Trieste
Dipartimento di Fisica Teorica, Universitร di Trieste
Strada Costiera 11, 34014 Trieste, Italy
Abstract
We consider two non-interacting systems embedded in a heat bath. If they remain dynamically independent, physical inconsistencies are avoided only if the single-system reduced dynamics is completely positive also beyond the weak-coupling limit.
1. Introduction
In a variety of different contexts, ranging from quantum optics to the foundations of quantum mechanics, the dissipative and irreversible time-evolutions of open quantum systems in weak interaction with suitable, large environments are commonly described by the so-called quantum dynamical semigroups \[1-3\]. These consist of linear maps $`\mathrm{\Lambda }_t`$, $`t0`$, that act on the physical states (density matrices) $`\rho ^S`$ of the open systems, $`S`$, and satisfy an evolution equation of Kossakowski-Lindlab form \[4-10\].
The maps $`\mathrm{\Lambda }_t`$ are linear and completely positive, that is, if $`S`$ is coupled to an arbitrary $`N`$-level system $`S_N`$, the maps $`\mathrm{\Lambda }_t\mathrm{id}_N`$ preserve the positivity of the states $`\rho ^{S+S_N}`$ of the compound system $`S+S_N`$ for all $`N`$.
Complete positivity is an algebraic property whose physical implications are better understood in the negative : if $`\mathrm{\Lambda }_t`$ is not completely positive, then for some $`N`$ an entangled initial state of $`S+S_N`$ surely exists which develops negative eigenvalues under the action of $`\mathrm{\Lambda }_t\mathrm{id}_N`$. On the other hand, if the initial state of $`S+S_N`$ is not entangled, i.e. $`\rho ^{S+S_N}=\rho ^S\rho ^{S_N}`$, it never develops negative eigenvalues under the action of $`\mathrm{\Lambda }_t\mathrm{id}_N`$.
Physical consistency demands that the eigenvalues of $`\mathrm{\Lambda }_t[\rho ^S]`$ must be positive in order to be interpretable as probabilities. It is logically necessary, but physically less compulsory, that the same should be true of all $`\mathrm{\Lambda }_t\mathrm{id}_N[\rho ^{S+S_N}]`$; this latter request would only be guaranteed by the complete positivity of $`\mathrm{\Lambda }_t`$. Yet, the system $`S_N`$ is totally arbitrary and unchanging, only statistical correlations with $`S`$ being allowed. Although such an occurrence is always possible, it is not always accepted as a justification of complete positivity as a necessary property of reduced dynamics.
It is sometimes argued \[12-14\] that complete positivity is the consequence of two auxiliary technical simplifications that are essential in the standard \[1-3\] derivation of quantum dynamical semigroups from the closed dynamics of the system $`S`$ plus its environment. It is in fact assumed $`a)`$ that the initial state of $`S`$ be uncorrelated to that of the environment, and $`b)`$ that a Markov approximation is possible on rescaled times $`\tau =\lambda ^2t`$, where $`\lambda <<1`$ is the strength of the system-environment interaction \[2-9\] (the weak-coupling limit). However, requests $`a)`$ ad $`b)`$ are not always physically plausible \[12-21\]; in particular, it might be necessary to examine the subsystem dynamics on times of the order of $`\lambda ^4t`$, hence beyond the weak-coupling limit \[15, 17-20\] and, in such instances, dynamics not of completely positive type may appear.
The standard derivation of the Redfield-Bloch equations , commonly used to describe the reduced dynamics of open two-level systems in chemical physics, fails to produce even positive dynamical maps $`\mathrm{\Lambda }_t`$ . While the danger is fully acknowledged , it is not accepted that one should end up with a completely positive time-evolution . Rather, it is argued that only those states whose positivity is preserved should be physically admissible or that a slippage in the initial conditions is needed in order to avoid inconsistencies \[13-17\]. However, from the previous considerations, it is clear that accommodating the problem of positivity does not properly address the issue of complete positivity which is strictly related to quantum entanglement.
The relevance and role of complete positivity is most clearly seen in the phenomenology of neutral K-mesons as open quantum systems in interaction with a gravitational background. Geometrical fluctuations at Planckโs scale act as a source of dissipation and decoherence \[23-29\]. Each single K-meson is thus assumed to evolve according to a semigroup of positivity-preserving, entropy increasing phenomenological linear maps $`\mathrm{\Lambda }_t`$. It turns out \[25-29\] that these maps must also be complete positive. Otherwise, physical inconsistencies would plague the resulting phenomenology of couples of K-mesons evolving in time according to the factorized dynamical maps $`\mathrm{\Lambda }_t\mathrm{\Lambda }_t`$. More important, these dissipative phenomenological models can actually be put to test in experiments performed at the so-called $`\varphi `$-factories .
In the following, we show that if the environment is such that there is no induced interaction between two otherwise non-interacting systems embedded in it, then the two-system reduced dynamics is in factorized form. Moreover this will be true at the $`2`$-nd and $`4`$-th order in the system-environment coupling constant $`\lambda `$. It then follows that if, beyond the weak-coupling limit, the reduced dynamics is not completely positive, then either the environment establishes a dynamical dependence between the two subsystems or the approximations leading to the reduced dynamics are not physically consistent.
2. Complete Positivity.
Let $`S`$ be a physical system whose time-evolution is given by a (semi)group of linear maps $`\mathrm{\Lambda }_t`$ acting on the states of $`S`$ represented by density matrices $`\rho ^S`$. Usually, $`\mathrm{\Lambda }_t`$ is positive and thus preserves the positivity of the eigenvalues of $`\rho ^S`$. ยฟFrom an abstract point of view , $`\mathrm{\Lambda }_t`$ is completely positive if and only if the map $`\mathrm{\Lambda }_t\mathrm{id}_N`$ acting on the states $`\rho ^{S+S_N}`$ of the system $`S`$ coupled with an arbitrary $`N`$-level system $`S_N`$, is also positive for all possible $`N`$.
If $`\mathrm{\Lambda }_t`$ is only positive, troubles are expected when $`S`$ is coupled to a generic $`N`$-level system and the joint state $`\rho ^{S+S_N}`$ carries correlations between $`S`$ and $`S_N`$. In fact, if $`S`$ and $`S_N`$ are not entangled, then $`\rho ^{S+S_N}=\rho ^S\rho ^{S_N}`$ (or a convex combinations of factorized states), so that
$$\mathrm{\Lambda }_t\mathrm{id}_N[\rho ^S\rho ^{S_N}]=\mathrm{\Lambda }_t[\rho ^S]\rho ^{S_N},$$
$`(2.1)`$
and the positivity of $`\mathrm{\Lambda }_t\mathrm{id}_N`$ automathically follows from the positivity of $`\mathrm{\Lambda }_t`$.
However, let $`S`$ be a two-level system and $`S_N=S_2=S`$. As a common state $`\rho ^{S+S}`$ we consider the projection
$$\begin{array}{ccc}\hfill \rho _A=\frac{1}{2}[& \left(\begin{array}{cc}1& 0\\ 0& 0\end{array}\right)\left(\begin{array}{cc}0& 0\\ 0& 1\end{array}\right)+\left(\begin{array}{cc}0& 0\\ 0& 1\end{array}\right)\left(\begin{array}{cc}1& 0\\ 0& 0\end{array}\right)\hfill & (2.2\mathrm{a})\hfill \\ \hfill & \left(\begin{array}{cc}0& 1\\ 0& 0\end{array}\right)\left(\begin{array}{cc}0& 0\\ 1& 0\end{array}\right)\left(\begin{array}{cc}0& 0\\ 1& 0\end{array}\right)\left(\begin{array}{cc}0& 1\\ 0& 0\end{array}\right)]\hfill & (2.2\mathrm{b})\hfill \end{array}$$
onto a singlet-like state of $`S+S`$ with eigenvalues $`0`$ and $`1`$. As a linear map on $`S`$, let us consider the transposition operation $`T:\left(\begin{array}{cc}a& b\\ c& d\end{array}\right)\left(\begin{array}{cc}a& c\\ b& d\end{array}\right)`$. The map $`T`$ is positive, but
$$\begin{array}{ccc}\hfill T\mathrm{id}_2[\rho _A]=\frac{1}{2}[& \left(\begin{array}{cc}1& 0\\ 0& 0\end{array}\right)\left(\begin{array}{cc}0& 0\\ 0& 1\end{array}\right)+\left(\begin{array}{cc}0& 0\\ 0& 1\end{array}\right)\left(\begin{array}{cc}1& 0\\ 0& 0\end{array}\right)\hfill & (2.3\mathrm{a})\hfill \\ \hfill & \left(\begin{array}{cc}0& 0\\ 1& 0\end{array}\right)\left(\begin{array}{cc}0& 1\\ 0& 0\end{array}\right)\left(\begin{array}{cc}0& 1\\ 0& 0\end{array}\right)\left(\begin{array}{cc}0& 0\\ 1& 0\end{array}\right)]\hfill & (2.3\mathrm{b})\hfill \end{array}$$
has eigenvalues $`\pm 1/2`$. Therefore the transposition $`T`$ is not completely positive, already the coupling to a $`2`$-level system failing to be positive . Clearly, the origin of troubles is the term (2.2b) which encodes the entanglement between the two systems and is changed by $`T\mathrm{id}_2`$ into (2.3b). Instead, the term $`(2.2a)`$, which represents an uncorrelated density matrix, is left unchanged.
Stinespringโs decomposition theorem ensures that the most general completely positive linear map $`\mathrm{\Lambda }`$ on the states of $`S`$ is of the form
$$\mathrm{\Lambda }[\rho ^S]=\underset{\mathrm{}}{}V_{\mathrm{}}^{}\rho ^SV_{\mathrm{}},$$
$`(2.4)`$
where $`V_{\mathrm{}}`$ are suitable bounded operators. Evidently, $`\mathrm{\Lambda }\mathrm{id}_N[\rho ^{S+S_N}]`$ is positive for any positive $`\rho ^{S+S_N}`$. The quantum mechanical time-evolution
$$\mathrm{\Lambda }_t[\rho ^S]=\mathrm{e}^{itH}\rho ^S\mathrm{e}^{itH^{}}$$
$`(2.5)`$
is of the form (2.4) and maps pure states into pure states thus preserving coherence. However, if one considers open quantum systems, coherence is usually lost. Hence, $`\mathrm{\Lambda }_t`$ cannot be of the form (2.5); whether it is of the form (2.4) must be decided on physical grounds.
Remark 2.1 In the approach to $`K`$-mesons as open quantum systems \[23-26\], the transpositio n map $`T`$ is replaced by a phenomenological dynamical map $`\mathrm{\Lambda }_t`$ and, in , $`T\mathrm{id}_2`$ by $`\mathrm{\Lambda }_t\mathrm{\Lambda }_t`$. It can be shown that $`\mathrm{\Lambda }_t`$ must be completely positive. Otherwise, physically realizable initial correlated states of two K-mesons as in (2.2) would develop negative eigenvalues.
3. Reduced Dynamics
The quantum open system of interest $`S`$ is assumed to be weakly interacting with a large (infinite) environment $`R`$, the dynamics of $`S+R`$ being governed by the Hamiltonian
$$H=H_S+H_R+\lambda H_{SR},$$
$`(3.1)`$
where $`H_S`$ and $`H_R`$ are the Hamiltonians of the system $`S`$, respectively environment $`R`$ and $`H_{SR}`$ is an interaction term with coupling strength $`\lambda `$. The total system $`S+R`$ is closed and its states, represented by density matrices $`\rho `$, evolve reversibly according to
$$\frac{\rho _t}{t}=L_H[\rho _t]:=i[H,\rho _t].$$
$`(3.2)`$
The environment is supposed to be in equilibrium with respect to $`H_R`$, thus $`[\rho ^R,H_R]=0`$. Let $`R`$ be a heat bath in equilibrium at temperature $`\beta ^1`$, namely we take
$$\rho ^R=\frac{\mathrm{exp}(\beta H_R)}{\mathrm{Tr}\mathrm{exp}(\beta H_R)}.$$
$`(3.3)`$
The interaction term is chosen to be of the form
$$H_{SR}=\underset{a}{}A_a^SA_a^R,$$
$`(3.4)`$
where the self-adjoint operators $`A_a^S`$ and $`A_a^R`$ refer to the system and environment, respectively. It is no restrictive to assume $`\phi _R(A_a^R)=0`$ for all $`a`$.
Remarks 3.1
i) The environment must eventually be considered infinite dimensional in order to allow for continuous spectra of $`H_R`$ and avoid recurrences. The states $`\rho ^R`$ are thus not confined to density matrices, nor are the expectation values
$$\phi _R(A^R)=\mathrm{Tr}_R(\rho ^RA^R)$$
$`(3.5)`$
always expressible via a trace operation. However, we will stick to the density matrix notation, the genuine case of infinitely many degrees of freedom being easily retrieved \[5-8\].
ii) The dissipative effects perturb the time-evolution of $`S`$ given by its own Hamiltonian $`H_S`$ and are at least of second order in the coupling $`\lambda `$. The weak-coupling limit consists in going from the fast-time variable $`t`$ to the slow-time variable $`\tau =\lambda ^2t`$, with $`\lambda 0`$. The technical procedure is physically justified when the ratio $`\tau _R/\tau _S`$ between the characteristic time $`\tau _R`$ of the environment and the characteristic time of the dissipative effects on $`S`$, $`\tau _S`$, is small . However, there might be heat bath temperatures for which one has to retain higher powers in $`\tau _R/\tau _S`$, being forced to go beyond the weak-coupling limit .
In the interaction representation, $`\stackrel{~}{\rho }_t:=\mathrm{exp}(tL_0)[\rho _t]`$ with $`L_0[\rho ]:=i[H_S+H_R,\rho ]`$, the evolution equation reads
$$\frac{\stackrel{~}{\rho }_t}{t}=\lambda \mathrm{e}^{tL_0}L_{H_{SR}}\mathrm{e}^{tL_0}[\stackrel{~}{\rho }_t].$$
$`(3.6)`$
For the sake of simplicity, we assume the spectrum of $`H_S`$ to be discrete and non-degenerate with eigenvalues $`\epsilon _r`$ and eigenvectors $`|r`$. We enumerate the operators $`|rs|`$ by denoting them as $`V_j^S`$, so that $`H_{SR}=_jV_j^SV_j^R`$ with $`V_j^R=_a\mathrm{Tr}_S(V_j^SA_a^S)A_a^R`$. Further, setting $`\omega _j:=\epsilon _r\epsilon _s`$, we get
$$H_{SR}(t)=\mathrm{e}^{tL_0}[H_{SR}]=\underset{j}{}\mathrm{e}^{i\omega _jt}V_j^S\mathrm{e}^{itH_R}V_j^R\mathrm{e}^{itH_R}.$$
$`(3.7)`$
The formal solution of (3.6) is $`\stackrel{~}{\rho }_t=\left(1+_{n=1}^{\mathrm{}}\lambda ^nU_t^{(n)}\right)[\stackrel{~}{\rho }_0]`$, with
$$U_t^{(n)}=_0^tdt_1_0^{t_1}dt_2\mathrm{}_0^{t_{n1}}dt_nL_{H_{SR}(t_1)}L_{H_{SR}(t_2)}\mathrm{}L_{H_{SR}(t_n)},$$
$`(3.8)`$
where $`L_A[]:=i[A,]`$.
In order to extract the system $`S`$ reduced dynamics, we operate on the states $`\rho `$ of $`S+R`$ with the projector $`P[\rho ]=\mathrm{Tr}_\mathrm{R}(\rho )\rho _R`$ which decouples the environment degrees of freedom. Further, we take as initial state of $`S+R`$ the state $`\rho _0=\stackrel{~}{\rho }_0=\rho ^S\rho ^R`$ with no correlation between $`S`$ and $`R`$. It follows that $`P[\rho _0]=\rho _0`$.
Remark 3.2 Despite the fact that they are the most used \[1-3\], the choice of $`P`$ made above and the assumption on the initial global state cannot be generically upheld. In particular, one cannot always benefit from a factorized initial state. However, this is in many instances plausible, as in the case of neutral $`K`$-mesons in a gravitational background. Indeed, $`K`$-mesons produced in strong $`\varphi `$-meson decays are arguably not influenced by geometrical fluctuations of gravitational origin . In general, one may be forced to adopt different projectors suited to initial states where system $`S`$ and environment $`R`$ result correlated by interactions prior $`t=0`$ \[13-17, 21\].
We now elaborate more in detail on the approach of . We assume the environment to be a Bose thermal bath described by the equilibrium state (3.3). The projector $`P`$ involves bath expectations with respect to (3.3), then only even correlation functions survive. Keeping terms up to $`\lambda ^4`$, one eventually finds
$$\frac{P[\stackrel{~}{\rho }_t]}{t}=\lambda ^2P\dot{U}_t^{(2)}P[\stackrel{~}{\rho }_t]+\lambda ^4\left[P\dot{U}_t^{(4)}P\dot{U}_t^{(2)}PU_t^{(2)}\right]P[\stackrel{~}{\rho }_t],$$
$`(3.9)`$
where $`\dot{U}_t^{(n)}`$ is the time-derivative of $`U^{(n)}`$. In particular, the $`2`$-nd and $`4`$-th order contributions, $`P\dot{U}_t^{(2)}`$ and $`P\dot{U}_t^{(2)}PU_t^{(2)}`$ read
$$\begin{array}{ccc}& P\dot{U}_t^{(2)}=_0^tdt_1PL_{H_{SR}(t)}L_{H_{SR}(t_1)},\hfill & (3.10\mathrm{a})\hfill \\ & P\dot{U}_t^{(2)}PU_t^{(2)}=_0^tdt_1_0^tdt_2_0^{t_2}dt_3PL_{H_{SR}(t)}L_{H_{SR}(t_1)}PL_{H_{SR}(t_2)}L_{H_{SR}(t_3)}.\hfill & (3.10\mathrm{b})\hfill \end{array}$$
Remark 3.3 After standard rearrangement of the integrals in (3.10b), the whole $`4`$-th order contribution in (3.9) assumes a typical cumulant expression . Thermal correlation functions are expected to factorize for large times; in such a case the cumulants vanish and allow one to operate a Markov approximation also at $`4`$-th order in $`\lambda `$.
Since $`\rho ^R`$ commutes with $`H_R`$, setting $`\rho _t^S:=\mathrm{Tr}_R(\rho _t)`$, the time-evolution equation obeyed by the open system $`S`$ has the form
$$\frac{\rho _t^S}{t}=L_{H_S}[\rho _t^S]+\lambda ^2K_t^{(2)}[\rho _t^S]+\lambda ^4K_t^{(4)}[\rho _t^S].$$
$`(3.11)`$
Let $`\mathrm{\Omega }_{jk}^\pm (t):=\phi _R\left(V_j^R(\pm t)V_k^R\right)`$ denote the environment two-point correlation functions. Then, the $`2`$-nd order dissipative contribution in (3.9) explicitly reads
$$K_t^{(2)}[\rho ^S]=\underset{j,k}{}_0^t\mathrm{d}t_1\mathrm{e}^{it_1\omega _j}(\mathrm{\Omega }_{kj}^+(t_1)[V_j^S\rho ^S,V_k^S]+\mathrm{\Omega }_{jk}^{}(t_1)[V_k^S,\rho ^SV_j^S]).$$
$`(3.12)`$
From (3.10b) it follows that, because of the chosen projector $`P`$, the $`4`$-th order dissipative operator $`K_t^{(4)}`$ involves four-point thermal correlation functions, which in turn are linear combinations of two-point ones. After a lengthy calculation one arrives at
$$K_t^{(4)}[\rho ^S]=\underset{j,k,\mathrm{},m}{}\underset{p=1}{\overset{10}{}}_0^t_0^{tt_1}_0^{tt_1t_2}d\stackrel{}{t}\mathrm{e}^{i\stackrel{}{\mathrm{\Delta }}_{k\mathrm{}m}\stackrel{}{t}}\mathrm{\Omega }_{jk\mathrm{}m}^{(p)}(\stackrel{}{t})D_{jk\mathrm{}m}^{(p)}[\rho ^S],$$
$`(3.13)`$
where $`\stackrel{}{t}=(t_1,t_2,t_3)`$ and $`\stackrel{}{\mathrm{\Delta }}_{k\mathrm{}m}\stackrel{}{t}=(\omega _k+\omega _{\mathrm{}}+\omega _m)t_1+(\omega _{\mathrm{}}+\omega _m)t_2+\omega _mt_3`$.
The quantities $`\mathrm{\Omega }^{(p)}(\stackrel{}{t})`$ are products of two point-correlation functions, while the operators $`D_{jk\mathrm{}m}^{(p)}`$ are essentially double commutators of observables of $`S`$. For instance, in the case $`p=1`$, one has
$$\mathrm{\Omega }_{jk\mathrm{}m}^{(1)}(\stackrel{}{t}):=\mathrm{\Omega }_j\mathrm{}^+(t_1+t_2)\mathrm{\Omega }_{km}^+(t_2+t_3),D_{jk\mathrm{}m}^{(1)}[\rho ^S]:=[V_j^S,[V_k^S,V_{\mathrm{}}^S]V_m^S\rho ^S].$$
$`(3.14)`$
Because of the explicit dependence of both $`K_t^{(2)}`$ and $`K_t^{(4)}`$ on time, the right hand side of (3.11) retains memory effects and does not generate a semigroup. However, a Markov approximation can be performed based on the following argument (see Remark 3.3). The two-point correlation functions $`\phi _R(V_j^R(t)V_k^R)`$ are expected to factorize for $`t`$ larger than the correlation-time $`\tau ^R`$ of the environment which is much shorter than the typical time for the dissipative effects being felt by the subsystem $`S`$. Since we assumed $`\phi _R(A_a^R)=0`$, it follows that $`\phi _R(V_j^R(t)V_k^R)=0`$ for $`t>>\tau ^R`$. Therefore, for times $`t>>\tau ^R`$, the time-dependent dissipative operators $`K_t^{(2)}`$ and $`K_t^{(4)}`$ can be replaced by time-independent dissipative operators $`K^{(2)}`$ and $`K^{(4)}`$, by extending to infinity each time-integration in (3.12) and (3.13). For a rigorous approach to this kind of Markov approximation the reader is referred to . Explicitly,
$$\begin{array}{ccc}& K^{(2)}[\rho ^S]:=\underset{j,k}{}\{\widehat{\mathrm{\Omega }}_{kj}^+(\omega _j)[V_k^S\rho ^S,V_j^S]+\widehat{\mathrm{\Omega }}_{jk}^{}(\omega _j)[V_k^S,\rho ^SV_j^S]\}\hfill & (3.15\mathrm{a})\hfill \\ & K^{(4)}[\rho ^S]:=\underset{j,k,\mathrm{},m}{}\underset{p=1}{\overset{10}{}},\widehat{\mathrm{\Omega }}_{jk\mathrm{}m}^{(p)}(\stackrel{}{\mathrm{\Delta }}_{k\mathrm{}m})D_{jk\mathrm{}m}^{(p)}[\rho ^S],\hfill & (3.15\mathrm{b})\hfill \end{array}$$
where $`\widehat{\mathrm{\Omega }}_{jk}^\pm (\omega ):=_0^{\mathrm{}}dt\mathrm{e}^{i\omega t}\mathrm{\Omega }_{jk}^\pm (t)`$ and
$$\widehat{\mathrm{\Omega }}_{jk\mathrm{}m}^{(p)}(\stackrel{}{\mathrm{\Delta }}_{k\mathrm{}m}):=_0^{\mathrm{}}dt_1_0^{\mathrm{}}dt_2_0^{\mathrm{}}dt_3\mathrm{e}^{i\stackrel{}{\mathrm{\Delta }}_{k\mathrm{}m}\stackrel{}{t}}\mathrm{\Omega }_{jk\mathrm{}m}^{(p)}(\stackrel{}{t}).$$
$`(3.16)`$
4. Non-interacting Open Quantum Systems
Let the open system $`S`$ consist of two non-interacting systems $`S_1`$ and $`S_2`$ whose dynamics, disregarding for the moment the presence of the environment $`R`$, is governed by the Hamiltonian operators $`H_{S_1}`$ and $`H_{S_2}`$. Again, we assume the energies $`\epsilon _{ar}`$ of $`H_{S_a}`$, $`a=1,2`$, to be discrete and non-degenerate and set $`\omega _{aj}:=\epsilon _{ar}\epsilon _{as}`$.
In absence of $`R`$, the system $`S=S_1+S_2`$ would evolve in time according to the Hamiltonian $`H_S=H_{S_1}\mathrm{id}_2+\mathrm{id}_1H_{S_2}`$. Instead, we suppose $`S_{1,2}`$ to interact weakly and independently with a thermal bath in the state (3.4). We chose an interaction term of the form
$$H_{SR}=\underset{j}{}\left((V_{1j}^S\mathrm{id}_2)V_{1j}^R+(\mathrm{id}_1V_{2j}^S)V_{2j}^R\right).$$
$`(4.1)`$
with $`\phi _R(V_{aj}^R)=0`$ for all $`j`$, $`a=1,2`$, and same coupling constant $`\lambda `$.
The analysis of the previous section can be repeated in this new context, the major difference being that, inserting (4.1) in (3.10) extra-indices appear identifying the system $`S_a`$, $`a=1,2`$, the various operators refer to. As a consequence, the dissipative operator in $`(3.12)`$ becomes
$$K_t^{(2)}[\rho ^S]=\underset{a,b=1}{\overset{2}{}}\underset{j,k}{}_0^t\mathrm{d}t_1\mathrm{e}^{i\omega _{aj}t_1}\{\mathrm{\Omega }_{bk;aj}^+(t_1)[V_{aj}^S\rho ^S,V_{bk}^S]+\mathrm{\Omega }_{aj;bk}^{}(t_1)[V_{bk}^S,\rho ^SV_{aj}^S]\}$$
$`(4.2)`$
where $`\mathrm{\Omega }_{aj;bk}^\pm (t)=\phi _R\left(V_{aj}^R(\pm t)V_{bk}^R\right)`$ and, for sake of simplicity, $`V_{aj}^S`$ denotes either $`V_{1j}^S\mathrm{id}_2`$ or $`\mathrm{id}_1V_{2j}^S`$. In turn, the $`4`$-th order operator reads
$$K_t^{(4)}=\underset{\genfrac{}{}{0pt}{}{a,b,c,d}{j,k,l,m}}{}\underset{p=1}{\overset{10}{}}_0^t_0^{t_1}_0^{tt_1t_2}d\stackrel{}{t}\mathrm{e}^{i\stackrel{}{\mathrm{\Delta }}_{bk;c\mathrm{};dm}\stackrel{}{t}}\mathrm{\Omega }_{aj;bk;c\mathrm{};dm}^{(p)}(\stackrel{}{t})D_{aj;bk;c\mathrm{};dm}^{(p)},$$
$`(4.3)`$
where $`\stackrel{}{\mathrm{\Delta }}_{bk;c\mathrm{};dm}\stackrel{}{t}=(\omega _{bk}+\omega _c\mathrm{}+\omega _{dm})t_1+(\omega _c\mathrm{}+\omega _{dm})t_2+\omega _{dm}t_3`$ and
$$\begin{array}{ccc}& \mathrm{\Omega }_{aj;bk;c\mathrm{};dm}^{(1)}(\stackrel{}{t}):=\mathrm{\Omega }_{aj;c\mathrm{}}^+(t_1+t_2)\mathrm{\Omega }_{bk;dm}^+(t_2+t_3)\hfill & (4.4\mathrm{a})\hfill \\ & D_{aj;bk;c\mathrm{};dm}^{(1)}[\rho ^S]:=[V_{aj}^S,[V_{bk}^S,V_c\mathrm{}^S]V_{dm}^S\rho ^S].\hfill & (4.4\mathrm{b})\hfill \end{array}$$
The two-point correlation functions $`\mathrm{\Omega }_{aj;bk}^\pm (t)`$ involve either bath operators interacting with the same system, $`a=c`$, or with different systems, $`ab`$. We can perform the Markov approximation exactly as in the case of just one system $`S`$ so that
$$\begin{array}{ccc}& K^{(2)}[\rho ^S]:=\underset{\genfrac{}{}{0pt}{}{a,b}{j,k}}{}\{\widehat{\mathrm{\Omega }}_{bk;aj}^+(\omega _{aj})[V_{aj}^S\rho ^S,V_{bk}^S]+\widehat{\mathrm{\Omega }}_{aj;bk}^{}(\omega _{aj})[V_{bk}^S,\rho ^SV_{aj}^S]\}\hfill & (4.5\mathrm{a})\hfill \\ & K^{(4)}[\rho ^S]:=\underset{\genfrac{}{}{0pt}{}{a,b,c,d}{j,k,\mathrm{},m}}{}\underset{p=1}{\overset{10}{}}\widehat{\mathrm{\Omega }}_{aj;bk;c\mathrm{};dm}^{(p)}(\stackrel{}{\mathrm{\Delta }}_{bk;c\mathrm{};dm})D_{aj;bk;c\mathrm{};dm}^{(p)}[\rho ^S],\hfill & (4.5\mathrm{b})\hfill \end{array}$$
where $`\widehat{\mathrm{\Omega }}_{aj;bk}^\pm (\omega ):=_0^{\mathrm{}}dt\mathrm{e}^{i\omega t}\mathrm{\Omega }_{aj;bk}^\pm (t)`$ and
$$\widehat{\mathrm{\Omega }}_{aj;bk;c\mathrm{};dm}^{(p)}(\stackrel{}{\mathrm{\Delta }}_{bk;c\mathrm{};dm}):=_0^{\mathrm{}}dt_1_0^{\mathrm{}}dt_2_0^{\mathrm{}}dt_3\mathrm{e}^{i\stackrel{}{\mathrm{\Delta }}_{bk;c\mathrm{};dm}\stackrel{}{t}}\mathrm{\Omega }_{aj;bk;c\mathrm{};dm}^{(p)}(\stackrel{}{t}).$$
$`(4.6)`$
Using (4.1) and (3.7), the $`2`$-nd order dissipative operator can be written as follows:
$$\begin{array}{ccc}\hfill K^{(2)}[\rho ^S]& =K_1^{(2)}\mathrm{id}_2[\rho ^S]+\mathrm{id}_1K_2^{(2)}[\rho ^S]\hfill & (4.7\mathrm{a})\hfill \\ & +\underset{\genfrac{}{}{0pt}{}{ab}{j,k}}{}\{\widehat{\mathrm{\Omega }}_{bk;aj}^+(\omega _{aj})[V_{aj}^S\rho ^S,V_{bk}^S]+\widehat{\mathrm{\Omega }}_{aj;bk}^{}(\omega _{aj})[V_{bk}^S,\rho ^SV_{aj}^S]\}.\hfill & (4.7\mathrm{b})\hfill \end{array}$$
The terms $`K_a^{(2)}`$, $`a=1,2`$ in (4.7a) are dissipative operators of the form (3.15a), involving only observables referring to the system $`S_a`$. If the contribution (4.7b) were absent, the right hand side of (4.7a) would generate a factorized time-evolution $`\rho ^S\rho _t^S=\mathrm{\Lambda }_t^1\mathrm{\Lambda }_t^2[\rho ^S]`$, with $`\mathrm{\Lambda }_t^1`$ satisfying
$$\frac{}{t}\left(\mathrm{\Lambda }_t^1\mathrm{id}_2[\rho ^S]\right)=K_1^{(2)}\mathrm{id}_2\left[\mathrm{\Lambda }_t^1\mathrm{id}_2[\rho ^S]\right]$$
$`(4.8)`$
and analogously for $`\mathrm{\Lambda }_t^2`$.
Clearly, the term (4.7b) dynamically couples the two systems $`S_1`$ and $`S_2`$ through their interaction with the same environment. However, this coupling depends on the strength of the thermal correlations between bath operators describing the interaction with different subsystems. If there are no correlatoins between them, that is if
$$\mathrm{\Omega }_{aj;c\mathrm{}}^\pm (t)=\phi _R(V_{aj}^R(\pm t))\phi _R(V_c\mathrm{}^R),$$
$`(4.9)`$
whenever $`ac`$ and $`t>0`$, then the term (4.7b) does not contribute since, with no restriction, we can assume the one-point bath correlation functions to vanish.
The same result follows from (4.9) when we include $`4`$-th order dissipative effects. This can be seen as follows. In the expression (4.4b), the right hand side vanishes if $`bc`$, since then $`V_{bk}^R`$ and $`V_c\mathrm{}^R`$ belong to different subsystems and thus commute. With $`b=c`$, using (4.9) one sees that the right hand side of (4.4a) vanishes unless $`a=c=b=d`$. Thus, if (4.9) holds, only $`a=b=c=d=1`$ and $`a=b=c=d=2`$ contribute to the sum in (4.5b). Then the $`4`$-th order dissipative operator splits as the $`2`$-nd order one,
$$K^{(4)}[\rho ^S]=K_1^{(4)}\mathrm{id}_2[\rho ^S]+\mathrm{id}_1K_2^{(4)}[\rho ^S],$$
$`(4.10)`$
where
$$K_a^{(4)}=\underset{j,k,\mathrm{},m}{}\underset{p=1}{\overset{10}{}}\widehat{\mathrm{\Omega }}_{aj;ak;a\mathrm{};am}^{(p)}(\stackrel{}{\mathrm{\Delta }}_{ak;a\mathrm{};am})D_{aj;ak;a\mathrm{};am}^{(p)},a=1,2.$$
$`(4.11)`$
As a consequence of (4.9), the evolution equation (3.11), after the Markov approximation, becomes
$$_t\rho _t^S=\left(\left(L_{H_{S_1}}+K_1^{(2)}+K_1^{(4)}\right)\mathrm{id}_2+\mathrm{id}_1\left(L_{H_{S_2}}+K_2^{(2)}+K_2^{(4)}\right)\right)[\rho _t^S]$$
$`(4.12)`$
and the dynamical maps generated by it factorize into $`\mathrm{\Lambda }_t^1\mathrm{\Lambda }_t^2`$, where $`\mathrm{\Lambda }_t^1`$ and $`\mathrm{\Lambda }_t^2`$ are generated as in (4.8) with $`K_a^{(2)}`$, $`a=1,2`$, replaced by $`K_a^{(2)}+K_a^{(4)}`$.
5. Conclusions
The request that the reduced dynamics of open quantum systems in interaction with a reservoir be completely positive is often rejected as not physically necessary \[12-14,16-21\]; it lacks physical appeal and involves trivial and uncontrollable couplings of the systems with generic $`N`$-level systems.
The issue of complete positivity assumes its full physical significance when dealing with the dynamics of correlated systems interacting with an environment. All depends on whether the reduced dynamics of the subsystems factorizes as $`\mathrm{\Lambda }_t\mathrm{\Lambda }_t`$, thus indicating that they evolve independently from each other, remaining dynamically uncorrelated.
If this is the case, then the complete positivity of the single system reduced dynamics $`\mathrm{\Lambda }_t`$ is unescapable, otherwise the joint reduced dynamics $`\mathrm{\Lambda }_t\mathrm{\Lambda }_t`$ generates unacceptable negative probabilities. ยฟFrom this point of view, the necessity of complete positivity appears to be a dynamical aspect of quantum entanglement.
On the other hand, there are situations for which the factorization of the reduced dynamics is not the case due to the subsystem-environment interaction and to the physics of the environment itself, namely the behaviour of its correlation functions. If the two-system reduced dynamics does not factorize, it is not compelling that the single-system reduced dynamics be completely positive.
In the literature there is evidence of non-completely positive reduced dynamics beyond the weak-coupling limit \[17-20\], namely taking into account contributions of order $`\lambda ^4`$ in the coupling between subsystem and environment. We have shown that if the reduced dynamics of two non-interacting subsystems factorizes at $`2`$-nd order in $`\lambda `$, it factorizes also at $`4`$-th order. Therefore, absence of complete positivity in the single-system reduced dynamics at $`4`$-th order would jeopardize the description of the physical behaviour of two of these subsystems at the same order of approximation.
Finally, we would like to stress that there are physical instances where an experimental check of complete positivity seems achievable, as in the case of $`K`$-mesons \[25-29\]; there, the conditions for a factorized reduced dynamics are plausibly fulfilled because of the weakness of the effects of the gravitational background. The reduced dynamics of $`K`$-mesons as open quantum systems must then be completely positive and experiments at $`\varphi `$-factories can explicitly clarify this fundamental request.
REFERENCES
1. E.B. Davies, Quantum Theory of Open Systems, (Academic Press, New York, 1976)
2. H. Spohn, Rev. Mod. Phys. 53 (1980) 569
3. R. Alicki and K. Lendi, Quantum Dynamical Semigroups and Applications, Lect. Notes Phys. 286, (Springer-Verlag, Berlin, 1987)
4. R. Dumcke and H. Spohn, Z. Physik B34 (1979) 419
5. V. Gorini, A. Kossakowski and E.C.G. Surdarshan, J. Math. Phys. 17 (1976) 821
6. V. Gorini, A. Frigerio, M. Verri, A. Kossakowski and E.C.G. Surdarshan, Rep. Math. Phys. 13 (1978) 149
7. V. Gorini and A. Kossakowski, J. Math. Phys. 17 (1975) 1298
8. A. Frigerio and V. Gorini, J. Math. Phys. 17 (1976) 2123
9. G. Lindblad, Commun. Math. Phys. 48 (1976) 119
10. M. Takesaki, Theory of Operator Algebras I (Springer, Berlin, 1979)
11. K. Kraus, States, Effects and Operations, Lecture Notes in Physics 190 (Springer, Belrin, 1983)
12. P. Pechukas, Phys. Rev. Lett. 73 (1994) 1060
13. V. Romero-Rochin and I. Oppenheim, Physica A 155 (1988) 52
14. V. Romero-Rochin, A. Orsky and I. Oppenheim, Physica A 156 (1989) 244
15. V. Gorini, M. Verri and A. Frigerio, Physica A 161 (1989) 357
16. A. Suarez, R. Silbey and I. Oppenheim, J. Chem. Phys. 97 (1992 ) 5101
17. T.-M. Chang and J.L. Skinner, Physica A 193 (1993) 483
18. J. Budimir and J.L. Skinner, J. Stat. Phys. 49 (1987) 1029
19. B.B. Laird, J. Budimir and J.L. Skinner, J. Chem. Phys. 94 (1991) 4391
20. B.B. Laird and J.L. Skinner, J. Chem. Phys. 94 (1991) 4405
21. A. Royer, Phys. Rev. Lett. 77 (1996) 3272
22. C.P. Slichter, Principles of Magnetic Resonance, with Examples from Solid State Physics (Springer, Berlin, New York, 1990)
23. J. Ellis, J.S. Hagelin, D.V. Nanopoulos and M. Srednicki, Nucl. Phys. B241 (1984) 381;
24. P. Huet and M.E. Peskin, Nucl. Phys. B434 (1995) 3
25. F. Benatti and R. Floreanini, Nucl. Phys. B488 (1997) 335
26. F. Benatti and R. Floreanini, Nucl. Phys. B511 (1998) 550
27. F. Benatti and R. Floreanini, Mod. Phys. Lett. A12 (1997) 1465
28. F. Benatti and R. Floreanini, Banach Centre Publ. 43 (1998) 71
29. F. Benatti and R. Floreanini, Phys. Lett. B 468 (1999) 287 |
warning/0002/hep-ex0002030.html | ar5iv | text | # 1 Introduction
## 1 Introduction
The Frascati $`\varphi `$ factory, DA$`\mathrm{\Phi }`$NE, is an $`e^+e^{}`$ collider optimized to operate at center of mass energy of M<sub>ฯ</sub> with very high luminosity (L$`{}_{peak}{}^{}=5\times 10^{32}\mathrm{cm}^2\mathrm{s}^1`$).
Since the vector meson $`\varphi `$ decays into $`K^+K^{}`$ 49 % , $`K_S`$$`K_L`$ 34 % of the times, DA$`\mathrm{\Phi }`$NE can be considered a factory of neutral and charged kaons, produced in collinear pairs, with momenta of $``$ 110 MeV/c and in a pure quantum state ($`J^{PC}=1^{})`$. Interference effects will then appear in the decays of the $`K_S`$$`K_L`$$`(`$$`K\overline{K}`$$`)`$ pairs, allowing the measurement of all (except one) CP and CPT violation parameters.
The observation of one charge conjugated state of the kaon guarantees (tags) the existence of the other one in the opposite direction. A pure $`K_S`$ beam without using any regenerator is therefore available.
At the design luminosity, DA$`\mathrm{\Phi }`$NE will produce, in one HEP physics year ($`10^7`$ s), $`8.5\times 10^9`$ $`K_S`$$`K_L`$, $`1.2\times 10^{11}`$ $`K^+K^{}`$, $`2.5\times 10^8`$ $`\eta `$, $`2.5\times 10^6`$ $`\eta ^{}`$ and $``$ $`10^5`$ $`a_o`$ and $`f_o`$ .
A detector, KLOE, has been built to exploit the following main physics goals:
* the measurement of $`CP/CPT`$ parameters from interferometry and the double ratio R, with a sensitivity of $`10^4`$;
* the measurement of the kaon form factors, $`K_S`$ rare decays and the $`K_S`$ semileptonic asymmetry (not measured so far);
* the study of the radiative $`\varphi `$ decays, investigating the nature of the $`a_o`$, $`f_o`$ mesons and providing a precise determination of the relative branching ratio between $`\varphi \eta \gamma `$ and $`\varphi \eta ^{}\gamma `$.
With the first 100 pb<sup>-1</sup> of collected data we already expect to measure $`\mathrm{}(ฯต^{}/ฯต)`$ with an error of $``$ 0.1%, to improve the measurement of the kaon form factors and to carry out most of the program on the radiative $`\varphi `$ decays.
Since the double ratio (R) is related to $`\mathrm{}(ฯต^{}/ฯต)`$ by:
$$R=\frac{N(K_L\pi ^+\pi ^{})/N(K_S\pi ^+\pi ^{})}{N(K_L\pi ^o\pi ^o)/N(K_L\pi ^o\pi ^o)}=1+6\mathrm{}(ฯต^{}/ฯต)$$
(1)
the final goal of the experiment of measuring $`\mathrm{}(ฯต^{}/ฯต)`$ with a precision of $`10^4`$, translates into a global error on R of few $`10^4`$. While the production rate of the kaon pairs and the tagging efficiency for the $`K_S`$ and the $`K_L`$ cancel out identically in the double ratio, this is not true for the detection efficiencies, which for this reason have been kept as high as possible in the detector design stage.
In addition, KLOE is a self-calibrating experiment: the abundance of events such as $`K_L`$ $`\pi ^+\pi ^{}\pi ^0`$ ($`K_L^{\pi \pi \pi }`$) or $`K^\pm \pi ^\pm \pi ^o`$ which produce in the final state both charged and neutral pions, will allow the determination of the efficiencies from the data themselves.
Since at these energies the decay length ($`\lambda `$) of the $`K_S`$ and $`K_L`$ is 0.6 cm and 345 cm, event counting will have to be performed in two different fiducial volumes. While essentially all the $`K_S`$ decay in vacuum inside a beam pipe of radius 10 cm, which approximates an infinite decay volume, the difference between the boundaries of charged and neutral decay volumes of the $`K_L`$ has to be known with a precision of $``$ 0.5 mm. Again, the use of $`K_L^{\pi \pi \pi }`$ and $`K^\pm \pi ^\pm \pi ^o`$ events will help in surveying and correcting for this difference.
Background subtraction is finally another source of systematics. Main backgrounds for the CP violating $`K_L`$ decays are $`K_L\pi ^0\pi ^0\pi ^0`$ ($`K_L^{000}`$) and $`K_L\pi \mu \nu `$ ($`K_{\mu 3}`$) which have signal to noise ratios of 1/240 and 1/135 respectively, so that rejection factors of $`8(6)\times 10^4`$ are needed.
Statistics is the other half of the game: in order to reach the desired accuracy in $`\mathrm{}(ฯต^{}/ฯต)`$ we need to collect $`4\times 10^6K_L\pi ^o\pi ^o`$, which can be done in two HEP years ($`2\times 10^7\mathrm{s}`$) at DA$`\mathrm{\Phi }`$NE design luminosity $`L_05\times 10^{32}\mathrm{cm}^2\mathrm{s}^1`$.
## 2 DA$`\mathrm{\Phi }`$NE and KLOE: the pilot run
In order to reach this ambitious goal, DA$`\mathrm{\Phi }`$NE has been designed to operate at a conservative single bunch luminosity ($`L_0L_{VEPP2M}=4\times 10^{30}cm^2s^1`$) and with a large number of bunches (up to 120). The DA$`\mathrm{\Phi }`$NE complex consists of two coplanar rings, a Linac accelerating $`e^+`$ up to 550 MeV and $`e^{}`$ to 800 MeV and a 510 MeV accumulator ring for fast topping-up. The two beams ($`\sigma _x,\sigma _y,\sigma _z=20\mu \mathrm{m},2\mathrm{mm},3\mathrm{cm}`$) intersect with an half angle of $``$ 12 mrad to reduce parasitic interactions between ingoing and outgoing bunches. Crossing frequency is up to 368 MHz (with 120 bunches). Typical number of particles/bunch is $`10^{11}`$, corresponding to $``$ 43 mA/bunch, with a total current per beam of 5.2 A.
A coupling $`\beta _y/\beta _x`$ ($``$ 1%) with a $`\beta _y=0.045\mathrm{m}`$ , together with a rather aggressive value for the tune-shift ($`\xi =0.04`$) have been chosen. Before KLOE rolled-in, a single bunch luminosity $`L1.5\times 10^{29}cm^2s^1`$ was achieved with 20 mA/beam and soon after $`L10^{31}cm^2s^1`$, with 13 bunches and $``$ 200 mA/beam, was also reached.
KLOE insertion has introduced a large perturbation in the rings optics, due to the large Bdl of the solenoid (2.4 T/m, to be compared with a beam rigidity $`B\rho `$ 1.7 T/m). If uncompensated, the KLOE field would rotate the beam of $`40`$ degrees at the interaction point.
Compensation has to be obtained on both rings through the careful tuning of solenoidal compensators, while the quadrupoles of the low $`\beta `$ insertion need to be rotated, in order to reduce transverse coupling. Only a first order pass has been made so far, which still cannot achieve neither a satisfactory coupling. nor a good matching of the $`\beta ^{}`$ on both rings.
At the moment of this conference, single bunch luminosities $`L_02\times 10^{29}cm^2s^1`$ with 10 mA/bunch and multi-bunch $`L2\times 10^{30}\mathrm{cm}^2\mathrm{s}^1`$ with ten bunches have been reached.
The KLOE detector was fully operational since July 1998 and a lot of experience on the detector calibration, data acquisition and triggering has been gained running with cosmic rays. The detector was installed in DA$`\mathrm{\Phi }`$NE early in 1999 and was fully operational in March 1999. First collisions were observed and recorded on the 23rd April 1999.
Another few weeks were later dedicated to the study of the $`\varphi `$ line shape running in single bunch mode. Once the beam energy was determined, multi-bunch operation started, for a total of 10 days of stable operation. The integrated luminosity of this pilot run was $`200\text{ nb}\text{-1}`$, which allowed us to tune-up the calibration and reconstruction procedures.
## 3 Overview of KLOE and its performances
The KLOE detector is a general purpose $`e^+e^{}`$ detector of respectable dimension ($``$ 7 m diameter, 6 m length). Going outward (fig. 1), a 0.5 mm thick cylindrical beryllium beam pipe with a spherical shape of 10 cm radius (16 $`\lambda _s`$) surrounds the interaction point. The inner permanent quadrupole region is instrumented by two lead-scintillator tiles calorimeters of $``$ 5 cm thickness, each organized in 16 wedges. Their major function is to improve the rejection efficiency for $`K_L^{000}`$.
Altough their reduced shower containement results in a poor energy resolution, $`\gamma `$ detection efficiency, timing and position resolution ($`ฯต98\%`$, $`\sigma _x5\mathrm{cm}`$, $`\sigma _t1\mathrm{ns}`$) are enough to achieve the desired rejection power.
A large drift chamber, DCH , of 4 m diameter and 3.5 m length, follows surrounded by an electromagnetic calorimeter, EMC, which covers almost hermetically the solid angle. Everything is embedded inside the coil cryostat and a specially shaped iron yoke. The superconducting coil generates a field of 6 kG. The trigger and data acquisition, DAQ , systems allow to collect a very high event rate. At the highest luminosity we expect to have $``$ 2.5 KHz of $`\varphi `$, $``$ 50KHz Bhabha, 2.5 KHz cosmics and few KHz of machine background. DAQ is able to handle a data throughput well above 50 MB/sec, through two levels of data concentration via dedicated hardware processors and online event building in parallel farms. The two levels trigger is based on EMC energy deposit and DCH hit multiplicity. Its main task is to reduce the event rate to the one allowed by DAQ without loosing any relevant $`\varphi `$ decay, while retaining some portions of Bhabha and cosmics for calibration purposes.
The main task of the EMC is to reconstruct the $`K_L^{00}`$ decay vertex and to efficiently reject the $`K_L^{000}`$ background. The vertex reconstruction of the decay $`K_L\pi ^0\pi ^0`$ is performed by accurately measuring the arrival time to the calorimeter of all the photons . Using the flight direction of the $`K_S\pi ^+\pi ^{}`$, one can easily demonstrate that a single photon is sufficient to determine the $`K_L`$ decay vertex. These requirements correspond to the following design specifications:
(1) $`\sigma _E/E`$ $``$ $`5\%/\sqrt{E\text{ (GeV)}}`$; (2) $`\sigma _T70\text{ ps}/\sqrt{E\text{ (GeV)}}`$;
(3) full efficiency for $`\gamma `$โs in the range 20โ280 MeV; (4) hermeticity.
The KLOE EmC is a fine sampling lead-scintillating fiber calorimeter with photomultiplier (PM) read-out. The central part (barrel) approximating a cylindrical shell of 4 m inner diameter, 4.3 m active length and 23 cm thickness ($`15`$ $`X_0`$), consists of 24 modules. Two end-caps, consisting of 32 โCโ shaped modules, close hermetically the barrel. The modules are read out on the two sides with a granularity of$`4.4\times 4.4`$ cm<sup>2</sup> by fine mesh PMโs, for a total number of $`4880`$ channels. The basic calorimeter structure consists in an alternating stack of 1 mm scintillating fiber layers glued between thin grooved lead foils. The final composite has a fiber:lead:glue volume ratio of approximately 48:42:10, a density of $`5`$ g/cm<sup>3</sup> and a $`X_0`$ of $`1.6`$ cm.
The EmC modules and the front-end electronics were fully installed in KLOE at the beginning of 1998. The calorimeter has been fully operational since then. After the installation, a first calibration of the calorimeter with cosmic rays has been performed, to obtain the minimum ionizing peak for each calorimeter channels with an accuracy better than 1%. The responses of all PMs is equalized to a few percent, in order not to bias the trigger response. The times measured at the two ends of a cell allow to obtain the arrival time of the particles and the coordinate along the fiber.
Fitting the measured time pattern for the fired cells in high energy cosmic ray events, the time offset of each channel is determined with a precision of $``$ 10 ps. A check of this calibration consists in measuring the velocity of cosmic rays as a function of their momentum, as reconstructed by the drift chamber. A fit to the measured distribution of $`\beta =p/\sqrt{p^2+m^2}`$, leaving the mass $`m`$ as the only free parameter, yields a value in good agreement with the muon mass, as shown in fig. 2-left.
The Bhabha and $`e^+e^{}\gamma \gamma `$ events allow then to set the energy and time scales. The measured energy resolution at 510 MeV is $``$ 7.8 % , corresponding to $`\sigma /E5.5\%/\sqrt{E(\text{GeV})}`$. As a check of the energy scale calibration, we reconstructed the invariant mass of the photon pair from $`\varphi \pi ^+\pi ^{}\pi ^0`$ events. The peak in the invariant mass distribution is within 1 MeV from the $`\pi ^0`$ mass, as shown in fig. 2-right.
In order to evaluate the time performances of the calorimeter, the quantity $`\mathrm{\Delta }t=t_{clu1}R_1/c(t_{clu2}R_2/c)`$ is measured for the two photons of $`e^+e^{}\gamma \gamma `$ events; $`t_{clu}`$ is the arrival time of the particle as measured by the calorimeter cluster, $`R`$ is the position of the cluster centroid. We observe $`\sigma _{\mathrm{\Delta }T}150`$ps, corresponding to $`\sigma _t`$ $`75\text{ ps}/\sqrt{E\text{ (GeV)}}`$;
The DCH should provide 3-D tracking with a resolution of $``$ 200 $`\mu \mathrm{m}`$ in the bending plane and a z-resolution of $``$ 1 mm on the decay vertices over the whole sensitive volume. In order to efficiently reject the $`K_{\mu 3}`$ background, a momentum resolution of 0.5% for low momentum tracks is required; the active volume should be as transparent as possible to minimize multiple scattering. Since the $`K_L`$โs decay almost uniformly all over the whole chamber, the tracking reconstruction efficiency should have minimal dependence on the position in the DCH. The design of the tracker is also driven by the requirement of having a light and homogeneous active volume to prevent $`K_L`$ regeneration and photon conversions; for the same reason low-mass walls are necessary. The adopted solution is a cylindrical drift chamber whose supporting structure is entirely made of carbon-fiber and filled with an ultra-light gas mixture (90% He-10% iC<sub>4</sub>H<sub>10</sub>). The requirement of a 3-D reconstruction and of uniform tracking forces the choice of square cells arranged in layers with alternating stereo angles. In order to keep a constant drop, the stereo angle varies with increasing radius from $`\pm `$60 mrad to $`\pm `$ 150 mrad. The cells are organized in 12 inner layers of smaller cells ($`2\times 2\mathrm{cm}^2`$) and 46 outer layers of bigger cells ($`3\times 3\mathrm{cm}^2`$).
Good efficiency and space resolution, together with negligible ageing effects, are obtained running at a gas gain of $``$ 10<sup>5</sup>. This requirement is met with 25 $`\mu \mathrm{m}`$ W(Au) sense wires and 80 $`\mu \mathrm{m}`$ field Al(Ag) wires kept at $``$ 1.9 KV. The DCH has a total of 52000 wires with a ratio 3:1 between the number of field to sense wires.
The DCH was moved inside KLOE on April 1998 and has been kept in operation for more than 1 year with a very low number of dead/hot channels/wires (below 0.1%). The calibration of the drift chamber proceeds, after subtracting the time offsets of each single wire, via an iterative procedure which minimizes the fit residuals by redefining the $`st`$ relations. In this gas mixture we do not expect to have a saturated drift velocity. The $`st`$ relations are therefore not linear and they are parametrized with a 5th order Chebychev polynomial. The minimization procedure stops whenever the mean of the fit residuals is below 100 $`\mu m`$. The behaviour of the space resolution along the cell shows the usual dependence on primary ionization and longitudinal diffusion. The average value of the resolution is $``$ 150 $`\mu m`$. The whole chamber calibration can be performed in $``$ 4 hours, using cosmic rays.
Using Bhabha events and $`K_s\pi ^+\pi ^{}`$, we can then evaluate the momentum resolution and the momentum scale. As shown in fig. 3-left, we obtain a momentum resolution better than 0.4% for polar angles greater than 45. The invariant mass of the $`K_S`$ is also shown in fig. 3-right; fitting it with a gaussian, we obtain M$`{}_{K_s}{}^{}496`$ MeV with $`\sigma `$ of $``$ 1 MeV.
## 4 CONCLUSIONS
KLOE has started operations on the beam.
DA$`\mathrm{\Phi }`$NE is presently delivering a luminosity L of $`2\times 10^{30}\mathrm{cm}^2\mathrm{s}^1`$.
Using the 200 $`\mathrm{nb}^1`$ of delivered luminosity so far, the detector has shown to be fully operational in all its hardware and software components, reaching or being very close already to its design specifications. KLOE is ready for real data. |
warning/0002/astro-ph0002394.html | ar5iv | text | # Innermost stable circular orbits around strange stars and kHz QPOs in low-mass X-ray binaries
## 1 Introduction
Observations of quasi periodic oscillations (QPOs) in the X-ray fluxes from low-mass X-ray binaries (LMXB), which are believed to be due to the orbital motion of matter in an accretion disk, raised hopes concerning observational constraints on the equation of state (EOS) of matter at supranuclear densities (Kaaret et al. 1997, Kluลบniak 1998, Zhang et al. 1998, Miller et al. 1998, Thampan et al. 1999, Schaab & Weigel 1999). General relativity predicts the existence of the marginally stable (MS) orbit, within which no stable circular motion is possible. This implies the existence of the innermost stable circular orbit (ISCO) around neutron stars. The frequency of the ISCO is an upper bound on the frequency of stable orbital motion around neutron stars. Whether the ISCO is separated from neutron star surface by a gap, or its radius coincides with stellar equatorial radius, depends on the star mass and on the EOS of neutron star matter. On the other hand, accreting neutron stars in LMXBs are expected to be rotating, and this influences both neutron star structure and the ISCO. Therefore, in order to attempt to use observed frequencies of QPOs to constrain the EOS of dense matter, one has to calculate the ISCO as a function of stellar mass and stellar rotation frequency. Such a procedure is based on the assumption that the observed upper QPO frequency is due to orbital motion, and that the effects of magnetic field, accretion, and radiation drag on the matter flow can be neglected.
A basic assumption of the present paper is that the frequency of the upper kHz QPO is the orbital frequency of the inner edge of an accretion disk surrounding the compact object, which will be identified with the ISCO. This is the leading interpretation of the QPOs. However, alternative models of the kHz QPOs were also proposed. In a model of Alpar & Yilmaz (1997) the kHz QPOs are explained in terms of wave packets of sound waves in the inner disk. In a series of papers, Titarchuk and collaborators propose a model in which the QPOs result from radial oscillations of the plasma in the boundary layer, i.e. in the region between the ISCO and stellar surface (see Titarchuk & Osherovich 1999 and references therein). These alternative models will not be considered in our study.
In the present paper we describe results of exact calculations of the ISCO, under the assumption that the compact object is not a neutron star, but a strange star. A strange star is composed of self-bound quark matter, which at zero pressure would constitute a real ground state of matter (strange matter), with energy per unit baryon number lower than that of $`{}_{}{}^{56}\mathrm{Fe}`$ crystal (Witten 1984, Farhi & Jaffe 1984, Haensel et al. 1986, Alcock et al. 1986; for a recent review of physics and astrophysics of strange matter, see Madsen 1999). Recently, strange stars were invoked by several authors in the context of modeling of observational properties of some X-ray and gamma-ray sources (Bombaci 1997, Cheng et al. 1998, Dai & Lu 1998, Li et al. 1999). First study of possibility of existence of strange stars in LMXBs exhibiting kHz QPOs was restricted to slow-rotation approximation for the ISCO, neglected the effect of rotation on the strange star structure, and used simplified EOS of strange matter, with massless, non-interacting quarks (Bulik et al. 1999 ). Very recently, the ISCOs around bare strange stars were calculated, assuming a simplified EOS of strange matter, for the limiting case of rotation at Keplerian frequency (Stergioulas et al. 1999). In both these studies the possible presence of the solid crust on the strange star surface was not taken into account.
In principle, a strange star could be covered by a thin crust of normal matter, a possibility which is particularly natural in the case of LMXB. The problem of formation and structure of a crust on an accreting strange star was studied by Haensel & Zdunik (1991)(see also Miralda-Escudรฉ et al. 1990). Because of its low mass, typically $`10^5\mathrm{M}_{}`$, the effect of the crust on the exterior spacetime is negligible. However, it determines the location of the star surface, due to its finite thickness of $`200300`$m. The matter distribution within the strange core, relevant for the exterior metric of rotating strange star, is characterized by a very flat density profile: for a massive strange star, density at the stellar center is typically only 2-3 times larger than that at the outer edge of the strange core. This has to be contrasted with the density distribution within a massive neutron star, which decreases continuously from $`10^{15}\mathrm{g}\mathrm{cm}^3`$ at the center to a few $`\mathrm{g}\mathrm{cm}^3`$ at the surface. The differences between the density profiles of a strange star and a neutron star result from the basic difference in the EOS of their interiors. In the case of a rapidly rotating compact object (situation relevant to LMXB), the differences between matter distributions within neutron star and strange star may be expected to imply differences in the spacetime exterior to the compact object, and in particular, differences in the properties of the ISCO.
## 2 ISCOs around strange stars for standard MIT Bag Model of strange matter
Our EOS of strange matter, composed of massless u, d quarks, and massive s quarks, is based on the MIT Bag Model. It involves three basic parameters: the bag constant, $`B`$, the mass of the strange quarks, $`m_\mathrm{s}`$, and the QCD coupling constant, $`\alpha _\mathrm{c}`$ (Farhi & Jaffe 1984, Haensel et al. 1986, Alcock et al. 1986). Our basic EOS corresponds to standard values of the Bag Model parameters for strange matter: $`B=56\mathrm{MeV}/\mathrm{fm}^3`$, $`m_\mathrm{s}=200\mathrm{MeV}/\mathrm{c}^2`$, and $`\alpha _\mathrm{c}=0.2`$ (Farhi & Jaffe 1984, Haensel et al. 1986, Alcock et al. 1986). This EOS of strange quark matter will be hereafter referred to as SQM1. It yields energy per unit baryon number at zero pressure $`E_0=918.8\mathrm{MeV}<E(^{56}\mathrm{Fe})=930.4`$MeV. For the SQM1 EOS maximum allowable mass for static strange star models is $`M_{\mathrm{max}}^{\mathrm{stat}}=1.8\mathrm{M}_{}`$.
The general relativistic models of stationary rotating strange stars have been calculated by means of the multi-domain spectral method, developed recently by Bonazzola et al. (1998). Details of the calculation method, specifically adapted for rotating strange stars, may be found in Gourgoulhon et al. (1999). Having calculated a particular stationary rotating strange star model, and its exterior spacetime, we determine the frequency of a particle in stable circular orbit in the equatorial plane, $`\nu _{\mathrm{orb}}(r)`$, where $`r`$ is the radial coordinate of the orbit. By testing the stability of orbital motion, we determine the radius of the innermost, marginally stable orbit, $`R_{\mathrm{ms}}`$, and its frequency $`\nu _{\mathrm{ms}}`$ (see, e.g., Datta et al. 1998, for the equations to be solved). Our numerical code calculating the ISCO has been successfully tested by comparing our results for the polytropic $`\gamma =2`$ EOS with those obtained by Cook et al. (1994a). Let us notice that the high precision of our numerical method makes it particularly suitable for the determination of $`R_{\mathrm{ms}}`$, which requires calculation of second derivatives of metric functions: these latter are better evaluated by the spectral method we employ than by means of finite differences.
No orbital motion is possible for $`r<R_{\mathrm{ms}}`$. The values of $`R_{\mathrm{ms}}`$ and $`\nu _{\mathrm{ms}}`$ for particles corotating with strange star differ from those for counterrotating ones. In the present paper we restricted ourselves to the corotating case, relevant for the LMXB. We neglect the effect of magnetic field, accretion, and radiation drag on the location of the ISCO, which is justified for $`B10^8`$G and $`\dot{M}\dot{M}_{\mathrm{Edd}}`$.
Let us consider a strange star, rotating at a frequency $`\nu _{\mathrm{rot}}`$, with equatorial radius $`R_{\mathrm{eq}}`$. If $`R_{\mathrm{ms}}>R_{\mathrm{eq}}`$, then stable orbits exist for $`r>R_{\mathrm{ms}}`$; the ISCO has then the radius $`R_{\mathrm{ms}}`$ and the frequency $`\nu _{\mathrm{ms}}`$, and there is a gap of width $`R_{\mathrm{ms}}R_{\mathrm{eq}}`$ between the ISCO and the strange star surface. However, if $`R_{\mathrm{ms}}<R_{\mathrm{eq}}`$, then $`R_{\mathrm{ISCO}}=R_{\mathrm{eq}}`$, $`\nu _{\mathrm{ISCO}}=\nu _{\mathrm{orb}}(R_{\mathrm{eq}})`$; and the accretion disk extends then down to the strange star surface (or, more precisely, joins stellar surface via a boundary layer).
While the exterior spacetime is practically not influenced by the presence of a solid crust on the strange star surface, the value of $`R_{\mathrm{eq}}`$ is affected by it. Neutrons are absorbed by strange matter, and therefore the density at the bottom of the crust, $`\rho _{\mathrm{bott}.\mathrm{cr}.}`$, cannot be higher than $`\rho _{\mathrm{n}\mathrm{drip}}4\times 10^{11}\mathrm{g}\mathrm{cm}^3`$ (lower values of $`\rho _{\mathrm{bott}.\mathrm{cr}.}`$ were discussed by Huang & Lu 1997). The equatorial thickness of the crust, $`t_{\mathrm{eq}}`$, which we calculate, corresponds to $`\rho _{\mathrm{bott}.\mathrm{cr}.}=\rho _{\mathrm{n}\mathrm{drip}}`$, and is therefore an upper bound on $`t_{\mathrm{eq}}`$. At a fixed baryon number, rotation increases $`t_{\mathrm{eq}}`$, as compared to the value for a static strange star, $`t_0`$. Dependence of $`t_{\mathrm{eq}}`$ on $`\nu _{\mathrm{rot}}`$ is well described by a formula $`t_{\mathrm{eq}}(\nu _{\mathrm{rot}})=t_0[1+0.7(\nu _{\mathrm{rot}}/\nu _\mathrm{K})^2]`$, where $`\nu _\mathrm{K}`$ is the Keplerian (mass shedding) frequency of strange star. For rotating strange stars our formula for $`t_{\mathrm{eq}}`$ reproduces numerical results of Glendenning & Weber (1992) within better than 2% in all cases considered by these authors. It is obvious that the bare strange star rotating at Keplerian limit would be unstable if we added a crust of nonzero thickness due to the increase of the radius, leaving the mass and the angular momentum practically unaltered. Thus at a fixed baryon mass of rotating strange star, the presence of the crust implies a decrease of the Keplerian frequency. Knowing the dependence of the radius of the strange core and rotational frequency on the stellar angular momentum one can estimate the point of the Keplerian instability for a strange star with crust.
Let us consider first rotating strange stars for the SQM1 EOS, which corresponds to the โstandard setโ of the Bag Model parameters for strange matter. A sample of our results for sequences of rotating strange star models with fixed baryon number are presented in Fig. 1.
The form of Fig. 1 is analogous to that constructed by Miller et al. (1998) for neutron stars, and therefore is suitable for discussion of the differences between neutron stars and strange stars. For strange stars with static mass $`M1.4\mathrm{M}_{}`$, we have always $`R_{\mathrm{ms}}>R_{\mathrm{eq}}`$, for any $`\nu _{\mathrm{rot}}`$. So, for $`M1.4\mathrm{M}_{}`$ the gap between strange star surface (with or without solid crust) and the ISCO exists at any strange star rotation rate. Even for lower $`M`$, the gap, which disappears at moderate rotation rates, reappears at $``$ 1 kHz frequency of rotation; this is visualized by the $`M=1.2\mathrm{M}_{}`$ case in Fig. 1. At $`\nu _{\mathrm{rot}}=\nu _\mathrm{K}`$, the ISCO is always separated from the strange star surface by a gap. Clearly, these features of the ISCO around strange stars are quite different from those obtained by Miller et al. (1998) for neutron stars with the FPS EOS (see their Fig. 1). Note that the existence of a gap ($`R_{\mathrm{ms}}>R_{\mathrm{eq}}`$) is expected to lead to a qualitatively different spectrum of X-ray radiation from LMXB, compared to the no-gap case (Kluลบniak et al. 1990).
Constraints on EOS of dense matter, resulting from the kHz QPOs observations, were initially derived within the slow-rotation approximation, in which $`R_{\mathrm{ms}}R_{\mathrm{ms}}^{\mathrm{s}.\mathrm{r}.}=6GM/c^2[1(2/3)^{3/2}j]`$, $`jJc/GM^2`$, and $`J`$ is stellar angular momentum (Kluลบniak & Wagoner 1985). However as pointed out by Shibata & Sasaki (1998), the mass quadrupole moment is as important as the angular momentum in determining the ISCO. Indeed as we can see in Fig. 1, slow rotation approximation yields $`R_{\mathrm{ms}}^{\mathrm{s}.\mathrm{r}.}`$, which in the case of $`M=1.4\mathrm{M}_{}`$ diverges from exact $`R_{\mathrm{ms}}`$ for $`\nu _{\mathrm{rot}}500`$Hz. Moreover, for rotating strange stars $`R_{\mathrm{ms}}^{\mathrm{s}.\mathrm{r}.}`$ leads always to disappearance of the gap at sufficiently high $`\nu _{\mathrm{rot}}`$, in contrast to exact calculation.
A quantity of particular interest in the context of the interpretation of observed kHz QPOs in LMXB, is the maximum frequency of the stable circular orbit at a given $`\nu _{\mathrm{rot}}`$, which we identify here with that of the ISCO. In principle, both $`\nu _{\mathrm{ISCO}}`$ and $`\nu _{\mathrm{rot}}`$ are observable (measurable) quantities, which can be thus used for confronting stellar models with observations. In Fig. 2 we present curves $`\nu _{\mathrm{ISCO}}(\nu _{\mathrm{rot}})`$ for the SQM1 EOS. As in Fig. 1, baryon masses are fixed along each curve, while labels correspond to gravitational mass of non-rotating strange star.
In all cases, displayed in Fig. 2, gap between stellar surface and ISCO exists, and therefore $`\nu _{\mathrm{ISCO}}=\nu _{\mathrm{ms}}`$. The dash-dotted line was calculated for a simplified EOS, with massless, non-interacting quarks ($`m_\mathrm{s}=0`$, $`\alpha _\mathrm{c}=0`$, $`B=56\mathrm{MeV}/\mathrm{fm}^3`$), hereafter referred to as SQM0 (such a type of the EOS of strange quark matter was used in Bulik et al. 1999). For $`\nu _{\mathrm{rot}}500`$Hz, neglecting strange quark mass (and, to a smaller extent, neglecting QCD interactions) leads to a rather severe underestimate of $`\nu _{\mathrm{ISCO}}`$ for rapidly rotating strange stars (by 200 Hz at $`\nu _{\mathrm{rot}}=1`$kHz). As we stressed before, presence of the solid crust does not influence the space-time outside rotating strange strange star. However, solid crust decreases (by about 10%) the value of $`J_{\mathrm{max}}`$ of strange stars of a given baryon mass. Complete curves in Fig. 2 correspond to bare strange stars. Rotating configurations with crust terminate at filled dots, corresponding to the Keplerian limit in the presence of the crust. The effect of the presence of the crust on the value of $`\nu _{\mathrm{ISCO}}`$ at $`\nu _\mathrm{K}`$ turns out to be significant, which is due to the steepness of the $`\nu _{\mathrm{ISCO}}(\nu _{\mathrm{rot}})`$ curve for bare strange stars at $`JJ_{\mathrm{max}}`$. For $`J`$ approaching $`J_{\mathrm{max}}`$ bare strange star undergoes strong deformation with increasing $`J`$. This deformation in turn implies strong decrease of $`\nu _{\mathrm{ISCO}}`$ with increasing rotation frequency. Consequently, the values of $`\nu _{\mathrm{ISCO}}`$ for maximally rotating strange stars with crust is about hundred Hz higher than for bare strange stars. This effect increases with decreasing strange star mass. At fixed $`B`$, the effect is stronger for the EOS which produces less compact strange stars of a given mass. Therefore, it is strongest for the SQM0 EOS with massless, noninteracting quarks, where maximally rotating configurations of $`1.4\mathrm{M}_{}`$ with crust have the ISCO frequency of 1.1 kHz, to be compared with less than 1 kHz for maximally rotating bare strange stars.
The problem of an appropriate parametrization of the one-parameter family of rotating strange stars with fixed baryon mass deserves a comment. These configurations may be labeled by the value of the total angular momentum $`J`$, which changes from $`J=0`$ in the static case to $`J_{\mathrm{max}}`$ at the Keplerian limit. As one can see in Fig.2, for bare strange stars, Keplerian configuration is not that with maximum $`\nu _{\mathrm{rot}}`$. The reason is that for very rapidly rotating strange stars the increase of the total angular momentum results mainly in the oblateness of the configurations leading to the significant increase of the equatorial radius without an increase of $`\nu _{\mathrm{rot}}`$ (or even with a decrease of $`\nu _{\mathrm{rot}}`$ very close to $`J_{\mathrm{max}}`$). As a consequence at fixed baryon mass the Keplerian configurations is reached not due to the increase of $`\nu _{\mathrm{rot}}`$ but because of the increase of the equatorial radius related to the deformation of the star. It is worth noticing that the difference between $`\nu _{\mathrm{rot},\mathrm{max}}`$ and $`\nu _\mathrm{K}`$ is of the order of one percent. The existence of this difference implies that for $`JJ_{\mathrm{max}}`$ it is in principle possible to spin up the strange star by the angular momentum loss. Such a situation was previously discussed in the case of supramassive neutron stars (Cook et al. 1994b) and supramassive strange stars (Gourgoulhon et al. 1999).
## 3 Confronting the standard MIT Bag Model of strange matter with QPO observations
Let us pass now to the confrontation of our results for strange stars with observations of the QPOs. Nearly twenty LMXBs, exhibiting QPOs, have been observed (van der Klis 2000). The upper-peak frequency, $`\nu _{\mathrm{QPO}}^{\mathrm{u}.\mathrm{p}.}`$, is usually interpreted as the frequency of the orbital motion around a neutron star. The most general observational constraint on a neutron star in LMXB is thus $`\nu _{\mathrm{QPO}}^{\mathrm{u}.\mathrm{p}.}\nu _{\mathrm{ISCO}}`$. Highest observed $`\nu _{\mathrm{QPO}}^{\mathrm{u}.\mathrm{p}.}`$ is $`1329\pm 4`$ Hz in 4U 0614+09 (van Straaten et al. Straaten (2000)). Condition $`\nu _{\mathrm{ISCO}}1.33`$ kHz is satisfied by nearly all strange star models displayed in Fig. 2 (except those rotating very close to the Keplerian frequency and slowly rotating maximum mass model). In particular for the spin frequency of the star $`\nu _{\mathrm{spin}}=312`$ Hz (Ford et al. Ford97 (1997), van Straaten et al. Straaten (2000)) all stellar configurations for SQM1 model of strange matter are allowed.
For neutron stars, condition $`\nu _{\mathrm{ISCO}}1.2`$ kHz (considered by Thampan et al. 1999 as the highest $`\nu _{\mathrm{QPO}}^{\mathrm{u}.\mathrm{p}.}`$) eliminates stellar masses below some limit, ranging from $`0.6\mathrm{M}_{}`$ for stiff EOS to $`1.4\mathrm{M}_{}`$ for soft EOS (Thampan et al. 1999). In this case the innermost allowed orbit is defined by the radius of the star and corresponds to the Keplerian frequency at the surface $`\nu _\mathrm{K}`$. This conclusion would be stronger in the case of $`\nu _{\mathrm{QPO}}^{\mathrm{u}.\mathrm{p}.}=1.33`$ kHz excluding the softest EOS and shifting the above mass limits to a little higher values (see Fig 1. in the paper by Thampan et al. 1999). Such a constraint does not apply to bare strange stars, for which in the limit of $`M\mathrm{M}_{}`$ one gets $`\nu _\mathrm{K}(G\rho _{\mathrm{sm}}/3\pi )^{1/2}=0.841(\rho _{\mathrm{sm},14})^{1/2}`$ kHz, where $`\rho _{\mathrm{sm},14}`$ is the density of strange matter at zero pressure, in the units of $`10^{14}\mathrm{g}/\mathrm{cm}^3`$. For reasonably high values of $`\rho _{\mathrm{sm}}`$, in particular for those considered in the present paper, one gets $`\nu _\mathrm{K}>1.33`$ kHz for low-mass, slowly rotating bare strange stars. In the case of strange stars with crust, condition $`\nu _\mathrm{K}>1.33`$ kHz turns out to be violated for $`M0.4\mathrm{M}_{}`$.
The behavior of QPOs in 4U 1820-30 has been interpreted as evidence for $`\nu _{\mathrm{QPO}}^{\mathrm{u}.\mathrm{p}.}=\nu _{\mathrm{ms}}`$ in this LMXB (Kaaret et al. 1999, and references therein). Accepting such an interpretation of $`\nu _{\mathrm{QPO}}^{\mathrm{u}.\mathrm{p}}=1.07`$ kHz implies strong constraints on neutron star model, and therefore, on the neutron star EOS (Kluลบniak 1998, Miller et al. 1998). Only a few existing EOS of neutron star matter allow simultaneously for $`R_{\mathrm{eq}}<R_{\mathrm{ms}}`$ and $`\nu _{\mathrm{ms}}=1.07`$kHz. It is clear from Fig. 2 that SQM1 models of strange stars cannot give $`\nu _{\mathrm{ISCO}}`$ as low as 1.07 kHz at slow rotation rates. In the case of the SQM0 EOS one is able to get such low $`\nu _{\mathrm{ISCO}}`$ for bare strange stars of $`M1.4\mathrm{M}_{}`$ rotating close to Keplerian frequency; we confirm in this way result of Stergioulas et al. (1999). However, as we see in Fig. 2, passing to an EOS which at the same value of $`B`$ includes effects of strange quark mass and of lowest order QCD interaction increases the values of $`\nu _{\mathrm{ISCO}}`$ of bare strange stars at high rotation rates to such extent, that the value of $`1.07`$ kHz cannot be reproduced.The presence of solid crust on rotating strange stars described by the โstandard strange matter EOSโ SQM1 excludes $`\nu _{\mathrm{ISCO}}(\nu _\mathrm{K})`$ lower than 1.2 kHz. In the case of the simplest SQM0 EOS the value of $`\nu _{\mathrm{ISCO}}(\nu _\mathrm{K})`$ is increased by the presence of the crust a little above 1.07 kHz. Generally, the presence of the crust on rotating strange star with standard strange matter EOS, such as SQM1 (or SQM0) excludes possibility of getting $`\nu _{\mathrm{ISCO}}`$ as low as 1.07 kHz for any possible rotation rates.
## 4 MIT Bag Model of strange matter consistent with QPO observations
In order to get $`\nu _{\mathrm{ISCO}}`$ as low as 1.07 kHz at slow and moderate rotation rates, one has to consider a specific set of the MIT Bag Model parameters, characterized by significantly lower values of both $`B`$ and $`m_\mathrm{s}`$, and higher value of $`\alpha _\mathrm{c}`$, than those characteristic of the SQM1 model. In this way one is able to increase significantly the value of $`M_{\mathrm{max}}^{\mathrm{stat}}`$, and get ISCO frequencies as low as 1 kHz at slow rotation rates. For such a choice of EOS it is also relatively easy to get $`\nu _{\mathrm{ISCO}}1`$ kHz for a broad range of masses of configurations rotating close to the Keplerian limit (see below). An example of such an EOS, hereafter referred to as the SQM2, was obtained assuming $`B=40\mathrm{MeV}/\mathrm{fm}^3`$, $`m_\mathrm{s}=100\mathrm{MeV}/\mathrm{c}^2`$, and $`\alpha _\mathrm{c}=0.6`$. At zero pressure, the SQM2 model yields energy per unit baryon number $`E_0=874.2`$ MeV. Let us stress that despite the relatively low value of $`B`$, the standard condition that neutrons do not fuse (coagulate) spontaneously into strangelets (droplets of quark matter), is satisfied by this model. Maximum mass of static strange stars for the SQM2 EOS is $`2.3\mathrm{M}_{}`$.
Our SQM2 model is a rather extreme one, as far as the values of the $`B`$, $`m_\mathrm{s}`$, and $`\alpha _\mathrm{c}`$ parameters are concerned. Canonical value of $`B`$, resulting from fitting hadronic masses, is $`59\mathrm{MeV}/\mathrm{fm}^3`$ (De Grand et al. 1975), significantly higher than $`B=40\mathrm{MeV}/\mathrm{fm}^3`$ used in the SQM2 model. On the other hand, $`m_\mathrm{s}=100\mathrm{MeV}/\mathrm{c}^2`$ of the SQM2 model is on the lower side of usually considered $`m_\mathrm{s}`$ values (Farhi & Jaffe 1984, Madsen 1999). Finally, $`\alpha _\mathrm{c}=0.6`$ is on the upper side of the interval of the $`\alpha _\mathrm{c}`$ values considered in the strange matter calculations (Farhi & Jaffe 1984).
Our results for the SQM2 EOS, analogous to those displayed in Fig. 1 and Fig. 2 for the SQM1 EOS, are shown in Fig. 3 and Fig. 4. The main differences between these models can be explained by the scaling laws with the bag constant, discussed in Sect. 5. The features of $`R_{\mathrm{ms}}`$ and radius of the rotating strange star of given mass for SQM1 model corresponds to the star SQM2 with the mass larger by the factor $`(B_1/B_2)^{1/2}`$. As one can see in Fig. 4, slowly rotating strange stars can have ISCO frequencies as low as $`11.1`$ kHz, provided their mass is sufficiently high, $`M2.22.3\mathrm{M}_{}`$ just because maximum allowable mass for static strange stars is sufficiently high. Moreover, for bare strange stars, the ISCO frequency below 1 kHz can also be reached for very rapid rotation close to the Keplerian limit. Less massive is bare strange star, lower is the ISCO frequency reached at the Keplerian limit. The presence of the solid crust makes the window (subset) of rapidly rotating configurations allowing for $`\nu _{\mathrm{ISCO}}=1.07`$ kHz significantly narrower. These configurations are very close to the Keplerian ones. Notice that the SQM2 EOS is simultaneously consistent with $`\nu _{\mathrm{ISCO}}1.33`$ kHz, provided the strange star mass $`M1.8\mathrm{M}_{}`$.
## 5 Discussion and conclusion
The features of ISCOs around rapidly rotating strange stars, described in the present paper for a particular choice of strange matter EOS, are actually generic. The MIT Bag Model EOS of strange quark matter depends on $`B`$, $`m_\mathrm{s}`$, and $`\alpha _\mathrm{c}`$ in a way, which implies specific scaling properties with respect to change of $`B`$ (Haensel et al. 1986, Zdunik & Haensel 1990). As a consequence, the global parameters of rotating strange stars scale with some power of $`B`$, which allows one to determine the values of $`M`$, $`R_{\mathrm{eq}}`$, $`R_{\mathrm{ms}}`$, etc., for $`B`$, $`\alpha _\mathrm{c}`$, and $`m_\mathrm{s}`$, from those calculated for $`B_0`$, $`\alpha _\mathrm{c}`$ and strange quark mass $`m_\mathrm{s}(B_0/B)^{1/4}`$. All length-type quantities (stellar radius, thickness of the crust and radius of the ISCO) scale as $`B^{1/2}`$, and all frequencies ($`\nu _{\mathrm{ISCO}}`$, $`\nu _{\mathrm{rot}}`$) scale as $`B^{1/2}`$, e.g., $`\nu _{\mathrm{ISCO}}[B]=\nu _{\mathrm{ISCO}}[B_0](B/B_0)^{1/2}`$ and $`R_{\mathrm{ms}}[B]=R_{\mathrm{ms}}[B_0](B/B_0)^{1/2}`$. Thus, for other values of $`B`$ the patterns of lines in Figs. 1-4 do not change, provided one rescales the axes and stellar masses.
Our calculations show that the properties of the ISCOs around strange stars differ from those around neutron stars. A generic property is the existence of the gap between the ISCO and the stellar surface, for both slowly and rapidly rotating strange stars.
The highest observed QPO frequency of 1.33 kHz in 4U 0614+91 can be easily interpreted as an orbital frequency around strange star based on the standard SQM1 EOS of strange matter, with no significant constraint on strange star mass and rotation rate. In the case of the SQM2 EOS, the orbital origin of the 1.33 kHz QPO implies $`M1.8\mathrm{M}_{}`$ at rotation frequencies $`300`$ Hz, while frequencies close to the mass shedding limit are excluded.
The value of the ISCO frequency at the Keplerian limit is significantly influenced by the presence of a crust on the strange star surface, which increases this frequency by about hundred Hz compared to the value for a bare strange star of the same mass. As one expects the presence of a crust on a strange star in a LMXB, we conclude that only slowly rotating strange stars with mass above $`2.2\mathrm{M}_{}`$ seem to be consistent with $`\nu _{\mathrm{ISCO}}=1.07`$ kHz. This excludes EOS of strange matter corresponding to the standard bag model parameters, and can be satisfied only by choosing a set of parameters quite different from the standard one.
The numerical results discussed in the present paper show that consistency of the ISCOs around slowly rotating strange stars with orbital-motion interpretation of QPOs in LMXBs can be achieved only with a substantial tuning of the MIT Bag Model parameters of strange matter. Our SQM2 EOS is a result of such a tuning. For this EOS, the condition $`\nu _{\mathrm{ISCO}}1`$ kHz is satisfied not only for slowly rotating massive models with $`M2.2\mathrm{M}_{}`$, but also for a broad range of masses of configurations close to the Keplerian limit. In contrast to $`\nu _{\mathrm{ISCO}}`$ at low rotation rate, which decreases with increasing baryon mass, the ISCO frequency at the Keplerian limit decreases with decreasing baryon mass of rotating strange star.
###### Acknowledgements.
During his stay at DARC, Observatoire de Paris, P. Haensel was supported by the PAST professorship of French MENRT. This research was partially supported by the KBN grants No. 2P03D.014.13, 2P03D.021.17. The numerical calculations have been performed on computers purchased thanks to a special grant from the SPM and SDU departments of CNRS. We are very grateful to the referee, P. Kaaret, for helpful comments and suggestions, which influenced the final version of the present paper. |
warning/0002/astro-ph0002227.html | ar5iv | text | # A Single Circumbinary Disk in the HD 98800 Quadruple System
## 1 Introduction
The evolution of circumstellar dust around young stars is traced by a time-dependent signature in excess infrared emission. The evidence lies primarily in the spectral distribution of radiation as it chronicles infall from a protostellar envelope, viscous accretion in a gas-rich circumstellar disk, and dispersal in a dusty โdebris diskโ that survives as a last vestige of planet formation (Adams, Lada, & Shu 1987; Backman & Paresce 1993). Imaging has dramatically confirmed this interpretation, providing support for a standard model of circumstellar evolution and elucidating the role of circumstellar disks throughout the process (Beckwith & Sargent 1996; Holland et al. 1998; Koerner 1997; Koerner et al. 1998). The co-existence of disks with stellar companions is attested by comparison of high-resolution binary surveys (Ghez, Neugebauer, & Matthews 1993; Leinert et al. 1993) with the results of imaging and long-wavelength flux measurements (Jensen et al. 1996a,b; Mathieu et al. 2000). Disks are found to be reduced in mass for binaries with separations in the 10โ100 AU range, similar to the typical disk size. However, circumstellar disks in binaries wider than 100 AU (Beckwith et al. 1990; Osterloh & Beckwith 1995; Jensen et al. 1996a), and circumbinary disks around spectroscopic binaries (Jensen & Mathieu 1997) are not obviously different from disks around single stars with respect to either their global properties or frequency of occurrence. These results argue strongly for the possibility of an abundant and diverse population of extra-solar planets.
Among preโmain-sequence spectroscopic binaries with separations of 1 AU or less, there is growing evidence that circumbinary disks are common. Massive circumbinary disks have been found in a handful of cases, demonstrating unequivocally that the presence of a small-separation binary is not an impediment to the formation of a protoplanetary disk (Jensen & Mathieu 1997). Examples include GW Ori (Mathieu et al. 1991, 1995), UZ Tau E (Jensen et al. 1996b; Mathieu et al. 1996), and DQ Tau (Mathieu et al. 1997), all with projected orbital separations of order 1 AU or less and disk masses that are comparable to or greater than that estimated for the minimum mass solar nebula.
The degree of complexity possible for multiple star-disk systems is perhaps nowhere better illustrated than in the case of the post-T Tauri quadruple system, HD 98800. It is composed of a pair of low-mass spectroscopic binaries, each with a K-dwarf primary, that have a projected separation of 0.8<sup>โฒโฒ</sup> (37.6 AU at the 47 pc distance determined by Hipparcos) and estimated ages of $`10`$ Myr (Soderblom et al. 1998). Despite the presence of many stellar components, HD 98800 is associated with an unusually strong IRAS signature of dust emission with a temperature similar to the solar zodiacal dust bands (Walker & Wolstencroft 1988; Zuckerman & Becklin 1993; Sylvester et al. 1996) and with evidence for silicate emission from dust grains (Skinner, Barlow, & Justtanont 1992). Until recently, there were no observations that provided a hint as to how this dust was distributed among the stellar components of the system. N-band imaging that marginally resolved the binary has now shown that most of the dust is associated with the optical secondary and spectroscopic binary HD 98800B (Gehrz et al. 1999). Here we present sub-arcsecond images from 5 to 25 $`\mu `$m that fully resolve the 0.8<sup>โฒโฒ</sup> binary components of the HD 98800 system.
## 2 Observations and Results
HD 98800 was observed with JPLโs mid-infrared camera MIRLIN at the f/40 bent-Cassegrain focus of the Keck II telescope on UT 14 March 1998. MIRLIN employs a Boeing 128$`\times `$128 pixel, high-flux Si:As BIB detector with a plate scale at Keck II of 0.137<sup>โฒโฒ</sup> per pixel and 17.5<sup>โฒโฒ</sup> field of view. Background subtraction was carried out by chopping the secondary mirror at a $``$4 Hz rate with 8<sup>โฒโฒ</sup> throw in the north-south direction, and by nodding the telescope a similar distance east-west after coadding a few hundred chop pairs. Images of the source on the double-differenced frames were shifted and added to make the final 32 $`\times `$ 32 (4.4<sup>โฒโฒ</sup> $`\times `$ 4.4<sup>โฒโฒ</sup>) images. Observations were carried out at wavelengths from 4.7 to 24.5 $`\mu `$m in the spectral bands listed in Table I. Small dither steps were taken between chop-nod cycles. Infrared standards $`\beta `$ Leo (A3 V) and $`\alpha `$ Hya (K3 III) were observed in the same way at similar airmasses.
The resulting images of HD 98800 are displayed in Fig. 1. Since the half-maximum width of the point spread function (PSF) is between 0.3<sup>โฒโฒ</sup> and 0.55<sup>โฒโฒ</sup> over the full wavelength range, it is possible to identify unambiguously the relative flux densities of the mid-infrared emission for the first time. Two point sources, separated by 0.81$`{}_{}{}^{\prime \prime }\pm 0.02^{\prime \prime }`$, are detected at wavelengths up to $`\lambda `$ = 12.5 $`\mu `$m with an orientation that corresponds to the binary optical components with A to the south and B to the north (cf. Soderblom et al. 1998). Only a single point source is detected at the longest wavelengths. It is immediately apparent from the images that this emission arises predominantly from the northern source, corresponding to the optical secondary HD 98800B. In contrast, emission from the optical primary decreases steadily towards longer wavelengths. Separate-component flux densities were derived by fitting a measured PSF to each component and using the resulting flux component ratio to decompose the total flux into values for HD 98800A and B. Results are listed in Table I and plotted as a spectral energy distribution in Fig. 2 together with measurements from HST, IRAS, and the JCMT (Sylvester et al. 1996; Soderblom et al. 1998). Mid-infrared flux densities measured for the total system are in excellent agreement with earlier values published in the literature (Zuckerman & Becklin 1993; Sylvester et al. 1996).
The distribution of mid-infrared flux between the components HD 98800A and B clearly indicates that the total infrared excess of the system is dominated by the contribution from HD 98800B. Values for the flux density of HD 98800B at $`\lambda `$ = 12.5 and 24.5 $`\mu `$m agree very well with 12 and 25 $`\mu `$m IRAS fluxes measured for the whole system. In contrast, flux densities for HD 98800A decrease approximately as $`\lambda ^2`$ between $`\lambda `$ = 7.9 and 12.5 $`\mu `$m, consistent with origin in a stellar photosphere. At 12.5 $`\mu `$m, emission from HD 98800A contributes less than 4% of the total emission. At 24.5 $`\mu `$m, an upper limit to its contribution comprises only 2% and an estimated photospheric contribution only 0.2% of the total flux. It is thus a good approximation to ascribe all the emission from unresolved measurements at $`\lambda >25\mu `$m to HD 98800B and neglect any contribution from HD 98800A. This result is largely in agreement with the conclusion of Gehrz et al. (1999), who nevertheless attributed some of the mid-infrared excess emission to HD 98800A on the basis of lower resolution imaging (1<sup>โฒโฒ</sup> at $`\lambda `$ = 9.8 $`\mu `$m) which only marginally resolved the 0.8<sup>โฒโฒ</sup> separation of components A and B.
A spectral signature of silicate emission at $`\lambda `$ $``$ 10 $`\mu `$m is evident in the flux measurements of HD 98800B plotted in Fig. 2. It is displayed in more detail in Fig. 3, where the measurements at 7.9 and 12.5 $`\mu `$m have been assumed to represent featureless thermal continuum emission, and a simple linear extrapolation between the two points has been subtracted off. The spectrum was then scaled to give the 7.9 and 12.5 um points a value of one to facilitate comparison with other data from the literature (see Hanner, Lynch, & Russell 1994 for comparison of different continuum removal techniques). Silicate features from comets and interstellar dust are plotted in Fig. 3 for comparison. It is readily apparent that the circumstellar dust feature resembles that from comets more than that from the interstellar medium. The feature is broader and does not show a single narrow peak between 9 and 10 $`\mu `$m as seen for interstellar grains in the Trapezium. For comets, this broadened line-shape has been interpreted as diagnostic of a mixture of amorphous and crystalline silicates that radiate predominantly at 9.8 and 11.2 $`\mu `$m, respectively (Hanner, Lynch, & Russell 1994).
## 3 Modeling and Discussion
To better interpret the emission from HD 98800, we fit model emission from stellar photospheres to optical (WFPC2) and near infrared (NICMOS) HST imaging that resolved components A and B (Soderblom et al. 1998; Low et al. 1999). These were matched by reddened model atmospheres from Kurucz by varying only the stellar luminosity, as described by Jensen & Mathieu (1997). Discrepant NICMOS measurements at roughly the same wavelength were averaged and weighted as a single point in the fit. Stellar effective temperatures were adopted from Case C of Soderblom et al. (1998) where the single value $`T_{\mathrm{eff}}=4350`$ K was given for the spectroscopic binary HD 98800A, and $`T_{\mathrm{eff}}=4250`$ and 3700 K were reported for the two stars in the double-lined spectroscopic binary HD 98800B. Soderblom et al. (1998) reported $`A_\mathrm{V}=0.44`$ mag for HD 98800B, but gave no $`A_\mathrm{V}`$ value for HD 98800A. We assumed $`A_\mathrm{V}=0`$ for HD 98800A and used a standard interstellar extinction law with $`A_\mathrm{V}=3.1E_{(\mathrm{B}\mathrm{V})}`$ to redden the model for HD 98800B. The luminosity ratio of the two components was fixed at 2.7 based on the absolute V magnitudes given by Soderblom et al. (1998) and bolometric corrections from Kenyon & Hartmann (1995). The best-fit models gave $`L_{\mathrm{star}}`$ = 0.78 $`L_{}`$ and 0.56 $`L_{}`$ and are plotted as a dotted and dashed line in Fig. 2 for the A and B components, respectively.
An average dust temperature of 150 K was derived by fitting a Planck function to the excess continuum emission from HD 98800B, omitting points associated with the silicate feature. For grains 1-10 $`\mu `$m in size, this temperature corresponds to a 4-12 AU distance from a single star of luminosity 0.56 $`L_{}`$. Given the 1 AU orbital separation estimated for the components of HD 98800B (Soderblom et al. 1998), it implies that most of the dust is located in a circumbinary configuration around the spectroscopic binary. To estimate the radial extent of the dust, we also fit the spectral energy distribution with a model of a disk around a single star of luminosity 0.56 $`L_{}`$. The model parametrization and fitting method are described in Koerner et al. (1998). Five parameters were varied in the fit, including inner radius $`R_{in}`$, radial extent $`\mathrm{\Delta }R`$, effective particle size $`\lambda _0`$, and the radial power-law index, $`\gamma `$, of the optical depth, $`\tau (r)=\tau _0(r/r_0)^\gamma `$. The optical depth scaling, $`\tau _0`$, was derived after varying the area-integrated optical depth, $`\sigma _{total}=_{R_{in}}^{R_{in}+\mathrm{\Delta }R}\tau _0(r/r_0)^\gamma 2\pi r๐r`$, an indicator of the total cross-sectional area of the grains. Parameter ranges considered were 0โ9 AU for $`R_{in}`$, 1โ25 AU for $`\mathrm{\Delta }R`$, $`10^1`$$`10^3\mu `$m for $`\lambda _0`$, -4.0โ4.0 for $`\gamma `$ and 5โ50 AU<sup>2</sup> for $`\sigma _{total}`$. A disk model with most-likely values of these parameters is displayed in Fig. 2; these are $`R_{in}=5.0\pm 2.5`$ AU, $`\mathrm{\Delta }R=13\pm 8`$ AU, $`\lambda _0=2_{1.5}^{+4}\mu `$m, $`\gamma =0\pm 2.5`$, and $`\sigma _{total}=16\pm 3`$ AU<sup>2</sup>, where the values quoted are central within a range of probabilities enclosing the 68% confidence level. The probability distribution is fairly flat within these ranges and peak values are not always central but lie at $`R_{in}=3.0`$ AU, $`\mathrm{\Delta }R=22`$ AU, $`\lambda _0=1.8\mu `$m, $`\gamma =1`$, and $`\sigma _{total}=16`$ AU<sup>2</sup>.
Many of these values are not narrowly constrained by the flux measurements alone, largely because the temperature dependence on both particle size and radial distance from the star makes it impossible to determine them uniquely. However, taken over the whole range of parameter space, there is greater than a 90% probability that the dust is distributed in a circumbinary disk, with $`\mathrm{\Delta }R/R_{in}>1`$, rather than a narrow ring like that around HR 4796A ($`\mathrm{\Delta }R/R_{in}<0.25`$; Koerner et al. 1998; Schneider et al. 1999). We emphasize the caveat that these estimates apply only under the assumptions of this particular disk model. An estimate of the true inner radius, for example, should take into account radiation from the two stellar components, and some temperature broadening may be due to a range of emissivities inherent in an unknown particle-size distribution. However, these effects are unlikely to alter our general conclusion about the disk vs ring-like nature of the dust.
The total cross section for dust grains around HD98800B, $`\sigma _{total}=16\pm 3`$ AU<sup>2</sup>, is 2-3 orders of magnitude smaller than for several other debris disks (e.g., $`\beta `$ Pic, HR4796A, and 49 Cet). Thus, from the standpoint of circumstellar mass, the disk around HD98800B is not as remarkable as suggested by the infrared excess alone (cf. Zuckerman & Becklin 1993). The relatively high fractional luminosity is, instead, a consequence of dust location close to the star where grains intercept a greater fraction of the stellar radiation. Assuming a range of plausible grain densities, $`\rho `$ = 1.0-3.0 g cm<sup>-3</sup>, values of $`\sigma _{total}`$ and $`\lambda _0`$ (grain radius $`a`$ = $`\lambda _0`$/1.5; cf. Backman et al. 1992) imply a disk mass in the range of 0.001โ0.1 lunar masses.
Models that incorporate a circumbinary disk surrounding an optically thin region of warmer dust have served to explain the spectral energy distributions of younger T-Tauri spectroscopic binaries (Jensen & Mathieu 1997). It is likely that HD 98800B is a similar system in a later phase of evolution. Modeling of the circumbinary dust emission indicates location of the dust in a radial zone associated with planet building early in the life of our own solar system. Consequently, it may well represent the telltale signature of planet formation in a hierarchically ordered multiple star system. If so, we can expect our picture of the plenitude and diversity of extra-solar planetary systems to become increasingly rich as it is revealed by impending surveys with high-resolution techniques now under development.
We gratefully acknowledge support of the NSFโs โLife in Extreme Environmentsโ program through grant AST 9714246. Data presented herein were obtained at the W.M. Keck Observatory (WMKO), which is operated as a scientific partnership among the California Institute of Technology, the University of California and the National Aeronautics and Space Administration. The Observatory was made possible by the generous financial support of the W.M. Keck Foundation. We wish to thank an anonymous referee for useful comments. A great debt is due, also, to Robert Goodrich and the WMKO summit staff for their many hours of assistance in adapting MIRLIN to the Keck II visitor instrument port. |
warning/0002/cond-mat0002034.html | ar5iv | text | # Hole-burning experiments within solvable glassy models
\[
## Abstract
We reproduce the results of non-resonant spectral hole-burning experiments with fully-connected (equivalently infinite-dimensional) glassy models that are generalizations of the mode-coupling approach to nonequilibrium situations. We show that an ac-field modifies the integrated linear response and the correlation function in a way that depends on the amplitude and frequency of the pumping field. We study the effect of the waiting and recovery-times and the number of oscillations applied. This calculation will help descriminating which results can and which cannot be attributed to dynamic heterogeneities in real systems.
preprint: LPTENS-97/18
PACS Numbers: 64.70.Pf, 75.10Nr \]
One of the most interesting questions in glassy physics is whether localized spatial heterogenities are generated in supercooled liquids and glasses.
In most supercooled liquids, the linear response to small external perturbations is nonexponential in the time-difference $`\tau `$. Within the โheteregeneous scenarioโ, the stretching is due to the existence of dynamically distinguishable entities in the sample, each of them relaxing exponentially with its own characteristic time. A different interpretation is that the macroscopic response is intrinsically nonexponential. In the glass phase, the relaxation is nonstationary and the dependence in $`\tau `$ is also much slower than exponential.
The heterogeneous regions, if they exist, are expected to be nanoscopic. The development of experimental techniques capable of giving evidence for the existence of such distinguishable spatial regions has been a challenge for experimentalists.
With non-resonant spectral hole-burning (NSHB) techniques one expects to probe, selectively, the microscopic responses. The method is based on a wait, pump, recovery and probe scheme depicted in Fig. 1. The amplitude of the ac perturbation is sufficiently large to pump energy in the sample, modifying the response as a linear function of the absorbed energy. The step-like perturbation $`\delta `$ is very weak and serves as a probe to measure the integrated linear response of the full system. The large ac and small dc fields can be magnetic, electric, or any other perturbation relevant for the sample studied. The idea behind the method is that the comparison of the modified (perturbed by the oscillation) and unmodified (unperturbed) integrated responses yield information about the microscopic structure of the sample. On the one hand, a spatially homogeneous sample will absorb energy uniformly and its modified integrated response is expected to be a simple translation towards shorter time-differences $`\tau `$ of the unmodified one. On the other hand, in a heterogeneous sample, the degrees of freedom that respond near the pump frequency $`\mathrm{\Omega }`$ are expected to absorb an important amount of energy and a maximum difference in the relaxation (equivalently, a spectral hole) is expected to generate around $`t1/\mathrm{\Omega }`$.
The NSHB technique has been first applied to the study of supercooled liquids. The polarization response of dielectric samples, glycerol and propylene carbonate, was measured after being modified by an ac electric field. More recently, ion-conducting glasses like CKN , relaxor ferroelectrics (90PMN-10PT ceramics) and spin-glasses (5% Au:Fe) were studied with similar methods. The results have been interpreted as evidence for the existence of spatial heterogeneities. We show here that their main features can be reproduced by a system with no spatial structure. We use one model, out of a family, that captures many of the experimentally observed features of super-cooled liquids and glasses as, for instance, a two-step equilibrium relaxation close and above $`T_c`$ , aging effects below $`T_c`$ , etc. The model is the $`p`$ spherical spin-glass , that is intimately related to the $`F_{p1}`$ mode-coupling model . It can be interpreted as a system of $`N`$ fully-connected continuous spins or as a model of a particles in an infinite dimensional random environment. In both cases, no reference to a geometry in real space nor any identification of spatially distinguishable regions can be made.
In the presence of a uniform field, the model is
$$H_J[๐ฌ]=\underset{i_1\mathrm{}i_p}{}J_{i_1\mathrm{}i_p}s_{i_1}\mathrm{}s_{i_p}+h\underset{i=1}{\overset{N}{}}s_i.$$
(1)
The interactions $`J_{i_1\mathrm{}i_p}`$ are quenched independent random variables taken from a Gaussian distribution with zero mean and variance $`[J_{i_1\mathrm{}i_p}^2]_J=\stackrel{~}{J}^2p!/(2N^{p1})`$. $`p`$ is a parameter and we take $`p=3`$. Hereafter $`[]_J`$ represents an average over $`P[J]`$ and $`\stackrel{~}{J}=1`$. The continuous variables $`s_i`$ are constrained spherically $`_{i=1}^Ns_i^2=N`$. A stochastic evolution is given to $`๐`$, $`\dot{s}_i(t)=\delta _{s_i(t)}H_J[๐]+\xi _i(t)`$ with $`\xi _i`$ a white noise with $`\xi _i=0`$ and $`\xi _i(t)\xi _i(t^{})=2T\delta (tt^{})`$. When $`N\mathrm{}`$, standard techniques lead to a set of coupled integro-differential equations for the autocorrelation $`NC(t,t^{})_{i=1}^N[s_i(t)s_j(t^{})]_J`$ and the linear response $`R(t,t^{})_{i=1}^N\delta [s_i(t)]_J/\delta \delta _i(t^{})|_{\delta =0}`$, with $`\delta _i(t^{})`$ an infinitesimal perturbation modifying the energy at time $`t^{}`$ according to $`HH_i\delta _is_i`$. The dynamic equations read
$`_tC(t,t^{})=z(t)C(t,t^{})+{\displaystyle \frac{p}{2}}{\displaystyle _0^t^{}}๐t^{\prime \prime }C^{p1}(t,t^{\prime \prime })R(t^{},t^{\prime \prime })`$ (2)
$`+{\displaystyle \frac{p(p1)}{2}}{\displaystyle _0^t}๐t^{\prime \prime }C^{p2}(t,t^{\prime \prime })R(t,t^{\prime \prime })C(t^{\prime \prime },t^{})`$ (3)
$`+2TR(t^{},t)+h(t){\displaystyle _0^t^{}}๐t^{\prime \prime }h(t^{\prime \prime })R(t^{},t^{\prime \prime }),`$ (4)
$`_tR(t,t^{})=z(t)R(t,t^{})`$ (5)
$`+{\displaystyle \frac{p(p1)}{2}}{\displaystyle _t^{}^t}๐t^{\prime \prime }C^{p2}(t,t^{\prime \prime })R(t,t^{\prime \prime })R(t^{\prime \prime },t^{}),`$ (6)
The Lagrange multiplier $`z(t)`$ enforces the spherical constraint and an integral equation for it follows from Eq. (4) and the condition $`C(t,t)=1`$. In deriving these equations, a random initial condition at $`t_0=0`$ has been used. It corresponds to an infinitely fast quench from equilibrium at $`T=\mathrm{}`$ to the working temperature $`T`$. The evolution continues in isothermal conditions.
In the absence of energy pumping, these models have a dynamic phase transition at a ($`p`$-dependent) critical temperature $`T_c`$, $`T_c0.61`$ for $`p=3`$. When an external ac-field is applied, it drives the system out-of-equilibrium and stationarity and FDT do not necessarily hold at any temperature. The question as to whether the clearcut dynamic transition survives under an oscillatory field is open and we do not address it here. We simply study the dynamics close to the critical temperature in the absence of the field by constructing a numerical solution to Eqs. (4) and (6) with a constant grid algorithm of spacing $`ฯต`$. We present data for small spacings, typically $`ฯต=0.02`$, to minimize the numerical errors. Due to the fact that Eqs. (4) and (6) include integrals ranging from $`t_0=0`$ to present time $`t`$, the algorithm is limited to a maximum number of iterations of the order of $`8000`$ that imposes a lower limit $`\mathrm{\Omega }2\pi /(8000ฯต)0.1`$ to the frequencies we use.
A word of caution concerning the scheme in Fig. 1 and the times involved is in order. For the purpose of collecting the data for each reference unmodified integrated response, the sample is prepared at the working temperature $`T`$ at $`t_0=0`$ and let freely evolve during a total waiting time $`t_w+t_1+t_r`$. Depending on $`T`$, this interval may or may not be enough to equilibrate the sample. ($`t_1`$ is chosen as $`t_1=2\pi n_c/\mathrm{\Omega }`$ with $`\mathrm{\Omega }`$ the angular velocity of the field that will be used to record the modified curve.) A constant infinitesimal probe $`\delta `$ is applied after $`t_w+t_1+t_r`$ to measure
$`\mathrm{\Phi }(\tau ){\displaystyle _0^\tau }๐\tau ^{}R(t_w+t_1+t_r+\tau ,t_w+t_1+t_r+\tau ^{}).`$
As an abuse of notation we explicitate only the $`\tau `$ dependence and eliminate the possible $`t_w+t_1+t_r`$ dependence. The modified integrated response $`\mathrm{\Phi }^{}`$ is measured after waiting $`t_w`$, applying $`n_c`$ oscillations of duration $`t_1=2\pi n_c/\mathrm{\Omega }`$, further waiting $`t_r`$, and only then applying the probe $`\delta `$. The effect of the ac perturbation is then quantified by studying the difference:
$$\mathrm{\Delta }\mathrm{\Phi }\mathrm{\Phi }^{}\mathrm{\Phi }.$$
(7)
We have examined $`\mathrm{\Delta }\mathrm{\Phi }`$ at $`T=0.8>T_c`$ and $`T=0.59<T_c`$. We pump one oscillation with $`h_F=0.1`$ and later check that this field is small enough to provoke a spectral modification that is linear in the absorbed energy (see Fig. 5 below). For simplicity, we start by choosing $`t_w=t_r=0`$. In Fig. 3 we show $`\mathrm{\Delta }\mathrm{\Phi }`$ against $`\mathrm{log}\tau `$ for different $`\mathrm{\Omega }`$ at $`T=0.8`$. All the curves are bell-shaped and vanish both at short and long times. In panel a, the $`\mathrm{\Omega }`$s are larger than a threshold value $`\mathrm{\Omega }_c1`$. The height of the peak $`\mathrm{\Delta }\mathrm{\Phi }_m\mathrm{max}(\mathrm{\Delta }\mathrm{\Phi })`$ decreases with increasing frequency reaching the limit $`\mathrm{\Delta }\mathrm{\Phi }_m=0`$ for $`\mathrm{\Omega }\mathrm{}`$. In addition, the location of the peak $`t_m`$ moves towards longer times when $`\mathrm{\Omega }`$ decreases. In panel b, $`\mathrm{\Omega }<\mathrm{\Omega }_c`$ and the behaviour of the height of the peak is the opposite, it decreases when $`\mathrm{\Omega }`$ decreases and, within numerically errors, its position is either independent of $`\mathrm{\Omega }`$ or it very smoothly moves towards shorter times for increasing $`\mathrm{\Omega }`$. The nonmonotonic behaviour of $`\mathrm{\Phi }_m`$ with $`\mathrm{\Omega }`$ is a consequence of the interplay between $`t_\alpha `$, the $`\alpha `$ relaxation time, and $`2\pi /\mathrm{\Omega }`$ the period of the oscillation. The term $`_0^{\mathrm{min}(t^{},t_1)}๐t^{\prime \prime }h(t^{\prime \prime })R(t^{},t^{\prime \prime })`$ in Eq. (4) controls the effect of the field and, clearly, vanishes in the limits $`\mathrm{\Omega }\mathrm{}`$ and $`\mathrm{\Omega }0`$. The inversion then occurs at a frequency $`\mathrm{\Omega }_c`$ that is of the order of $`2\pi /t_\alpha `$. These results qualitatively coincide with the measurements of the electric relaxation in CKN at $`T<T_g`$ in Fig 1 a and b of Ref. . In Fig. 3 we show $`\mathrm{\Delta }\mathrm{\Phi }`$ against $`\mathrm{log}\tau `$ for different $`\mathrm{\Omega }`$ at $`T=0.59`$. For all $`\mathrm{\Omega }`$ we reproduced the situation of panel a in Fig. 3, as if $`\mathrm{\Omega }>\mathrm{\Omega }_c`$. We have not found a threshold $`\mathrm{\Omega }_c`$, that has gone below the minimum $`\mathrm{\Omega }`$ reachable with the algorithm.
The maximum modification of the relaxation $`\mathrm{\Delta }\mathrm{\Phi }_m`$ increases quadratically with the square of the amplitude of the pumping field $`h_F`$, and hence linearly in the absorbed energy, as long as $`h_F1`$. In Fig. 5 we display the relation $`\mathrm{\Delta }\mathrm{\Phi }_mh_F^2`$ in a log-log scale for the two temperatures explored. The amplitude $`h_F=0.1`$ used in Figs. 3 and 3 is in the linear regime.
The effect of the pump diminishes with increasing recovery time $`t_r`$. A convenient way of displaying this result is to plot the normalized maximum deviation $`\mathrm{\Delta }\mathrm{\Phi }_m(t_r)/\mathrm{\Delta }\mathrm{\Phi }_m(0)`$ vs $`\mathrm{\Omega }t_r`$. Using several frequencies and recovery times, we verified that this scaling holds for $`T=0.59`$ but does not hold for $`T=0.8`$, as shown in Fig. 5. This simple saling holds very nicely in the relaxor ferroelectric and in the spin-glass but it is very different from the $`\mathrm{\Omega }`$-independence of the propylene carbonate .
Up to now, the effect of a single cycle of different frequencies has been studied. Another procedure can be envisaged. Since $`t_1=2\pi n_c/\mathrm{\Omega }`$, we can change $`t_1`$ by applying different numbers of cycles $`n_c`$ while keeping $`\mathrm{\Omega }`$ fixed. In Fig. 7 we show the distortion due to $`n_c=10,2,1`$ cycles with $`\mathrm{\Omega }=10`$ at $`T=0.8`$. The qualitative dependence on $`n_c`$ is indeed the same as the dependence on $`1/\mathrm{\Omega }`$: the peaks are displaced towards longer times with increasing $`n_c`$ (longer $`t_1`$). This behaviour is similar to the results obtained for propylene carbonate in Fig. 11 of Ref. b. though we do not reach the expected saturation within our accessible time window.
Below $`T_c`$ the nonperturbed model never equilibrates and the relaxation depends on $`t_w`$. Indeed, $`t_\alpha `$ is an approximately linear function of $`t_w`$ and the distortion might depend on $`t_w`$. We compare $`\mathrm{\Delta }\mathrm{\Phi }`$ vs $`\mathrm{log}\tau `$ for two $`t_w`$โs in Fig. 7.
Finally, we checked that the effect of one or many pump oscillations on the difference $`\mathrm{\Delta }CC^{}(t_w+t_1+t_r+\tau ,t_w+t_1+t_r)C(t_w+t_1+t_r+\tau ,t_w+t_1+t_r)`$ is very similar to the one observed in $`\mathrm{\Delta }\mathrm{\Phi }`$. This observation is interesting since it is easier to compute numerically correlations than responses. Figure 8 shows the modification observed at $`T=0.8`$ and $`\mathrm{\Omega }>\mathrm{\Omega }_c`$ (to be compared to Fig. 3).
We conclude by stressing that we do not claim that spatial heterogeneities do not exist in real glassy systems. We just wish to stress that the ambiguities in the interpretation of experimental results have to be eliminated in order to have unequivocal evidence for them. The detailed comparison of the experimental measurements to the behaviour of glassy models with and without space will certainly help us refine the experimental techniques. Numerical simulations can play an important role in this respect.
LFC and JLI thank the Dept. of Phys. (UNMDP) and LPTHE (Jussieu) for hospitality, and ECOS-Sud, CONICET and UNMDP for financial support. We thank R. Bรถhmer, H. Cummins, G. Diezemann, M. Ediger, J. Kurchan and G. Mc Kenna for very useful discussions and T. Grigera, N. Israeloff and E. Vidal-Russel for introducing us to the hole-burning experiments. |
warning/0002/cond-mat0002420.html | ar5iv | text | # Conformal Dynamics of Fractal Growth Patterns Without Randomness
## I Introduction
In this paper we introduce a new class of fractal growth patterns in two dimensions, constructed in terms of the conformal maps from the exterior of the unit circle to the exterior of the growing cluster. Until now most of the interesting fractal growth models included randomness as an essential aspect of the growth algorithms. Foremost in such models has been the diffusion limited aggregation (DLA) model that was introduced in 1981 by T. Witten and L. Sander . This model has been shown to underlie many pattern forming processes including dielectric breakdown , two-fluid flow , and electro-chemical deposition . The algorithm begins with fixing one particle at the center of coordinates in $`d`$-dimensions, and follows the creation of a cluster by releasing random walkers from infinity, allowing them to walk around until they hit any particle belonging to the cluster. Upon hitting they are attached to the growing cluster. The growth probability for a random walker to hit the interface is known as the โharmonic measureโ, being the solution of the harmonic (Laplace) equation with the appropriate boundary conditions. The DLA model was generalized to a family of models known collectively as Dielectric Breakdown Models, in which the density of growth probability is the density of the harmonic measure raised to a power $`\eta `$ . For $`\eta =1`$ one regains the DLA model; the interval $`\eta (0,\mathrm{})`$ generates a family of growth patterns from compact to a single needle. For $`\eta =0`$ one obtains a growth probability that is uniform for all boundary points. This is known as the Eden model that was introduced originally to describe the growth of cancer cells .
The fundamental difficulty of all these models is that their mathematical description calls for solving equations with boundary conditions on a complex, evolving interface. It is therefore advantageous to swap for a simple boundary, like the unit circle, and to delegate the complexity to the dynamics of the conformal map from the exterior of the unit circle to the exterior of the growing cluster. For continuous time processes this method had been around for decades , and had been used extensively. For discrete particle growth such a language was developed recently , showing that DLA in two dimensions can be grown by iterating stochastic conformal maps. In this paper we employ this new language to define models in which the stochasticity is eliminated altogether, to create deterministic iterations of conformal maps with very interesting fractal growth properties. It is stressed below that these new models and their interesting properties are natural extensions of the discrete conformal dynamics; it may be very difficult to study such models with the traditional techniques in physical space.
A central thesis of this paper is that the growth models introduced below are simpler to understand than DLA, even though the fractal geometry exhibited does not seem simpler. Indeed, we present below some tools and concepts that allow us to explain why the growing cluster is fractal. We present a scaling theory of the growing clusters, and identify the exponent that determines the fractal dimension. In sect. 2 we review the basic ideas of conformal dynamics as a method to grow DLA and related growth patterns. In Sec.3 we make the point that within this framework randomness can be eliminated from the discussion without changing the properties of the fractal growth: one can have deterministic growth rules with clusters that are indistinguishable from DLA. In Sec.4 we introduce fractal growth patterns that are obtained from quasiperiodic itineraries of iterated conformal maps. These itineraries are characterized by a winding number $`W`$. The growing clusters have complex geometries and a difference appearance for every $`W`$. We propose nevertheless that all the quadratic irrationals belong to the same universality class, and that the dimensions of their clusters are the same. In Sec.5 we consider rational approximants $`P/Q`$ to the quadratic irrational winding numbers $`W`$โs. With rational approximants the growth patterns crossover from a fractal phase of growth to a 1-dimensional star-like growth pattern. We argue that the analysis of the crossover as a function of $`Q`$ provides us with a scaling theory, allowing the introduction of universality classes and the achievement of data collapse. In Sec.6 we elucidate the mechanism for crossover from fractal to 1-dimensional growth, and identify the exponent that determines the fractal dimension. In Sec.7 we summarize and offer final remarks regarding the availability of a renormalization group treatment and of the road ahead.
## II Discrete Conformal Dynamics for Fractal Growth Patterns
The basic idea is to follow the evolution of the conformal mapping $`\mathrm{\Phi }^{(n)}(w)`$ which maps the exterior of the unit circle $`e^{i\theta }`$ in the mathematical $`w`$โplane onto the complement of the (simply-connected) cluster of $`n`$ particles in the physical $`z`$โplane . The unit circle is mapped to the boundary of the cluster which is parametrized by the arc length $`s`$, $`z(s)=\mathrm{\Phi }^{(n)}(e^{i\theta })`$. This map $`\mathrm{\Phi }^{(n)}(w)`$ is made from compositions of elementary maps $`\varphi _{\lambda ,\theta }`$,
$$\mathrm{\Phi }^{(n)}(w)=\mathrm{\Phi }^{(n1)}(\varphi _{\lambda _n,\theta _n}(w)),$$
(1)
where the elementary map $`\varphi _{\lambda ,\theta }`$ transforms the unit circle to a circle with a โbumpโ of linear size $`\sqrt{\lambda }`$ around the point $`w=e^{i\theta }`$. Accordingly the map $`\mathrm{\Phi }^{(n)}(w)`$ adds on a new bump to the image of the unit circle under $`\mathrm{\Phi }^{(n1)}(w)`$. The bumps in the $`z`$-plane simulate the accreted particles in the physical space formulation of the growth process. The main idea in this construction is to choose the positions of the bumps $`\theta _n`$ and their sizes $`\sqrt{\lambda _n}`$ such as to achieve accretion of fixed linear size bumps on the boundary of the growing cluster according to the growth rules appropriate for the particular growth model that we discuss.
As an example consider DLA. In $`z`$-space we want to accrete particles according to the harmonic measure. This means that the probability for the $`n`$th particle to hit a boundary element $`ds`$ equals $`P(s)ds`$, where $`P(s)`$ (the density of the harmonic measure ) and $`ds`$ are:
$`P(s)`$ $`=`$ $`{\displaystyle \frac{1}{|\mathrm{\Phi }_{}^{(n1)}{}_{}{}^{^{}}(e^{i\theta })|}},`$ (2)
$`ds`$ $`=`$ $`|\mathrm{\Phi }_{}^{(n1)}{}_{}{}^{^{}}(e^{i\theta })|d\theta .`$ (3)
Here $`e^{i\theta }`$ is the preimage of $`z(s)`$. Accordingly the probability to grow on an interval $`d\theta `$ is uniform (independent of $`\theta `$). Thus to grow a DLA we have to choose random positions $`\theta _n`$, and $`\lambda _n`$ in Eq.(1) according to
$$\lambda _n=\frac{\lambda _0}{|\mathrm{\Phi }_{}^{(n1)}{}_{}{}^{}(e^{i\theta _n})|^2}.$$
(4)
This way we accrete fixed size bumps in the physical plane according to the harmonic measure. The elementary map $`\varphi _{\lambda ,\theta }`$ is chosen as
$`\varphi _{\lambda ,0}(w)=w^{1a}\{{\displaystyle \frac{(1+\lambda )}{2w}}(1+w)`$ (5)
$`\times [1+w+w(1+{\displaystyle \frac{1}{w^2}}{\displaystyle \frac{2}{w}}{\displaystyle \frac{1\lambda }{1+\lambda }})^{1/2}]1\}^a`$ (6)
$`\varphi _{\lambda ,\theta }(w)=e^{i\theta }\varphi _{\lambda ,0}(e^{i\theta }w),`$ (7)
The parameter $`a`$ is confined in the range $`0<a<1`$, determining the shape of the bump. We employ $`a=2/3`$ which is consistent with semicircular bumps. The recursive dynamics can be represented as iterations of the map $`\varphi _{\lambda _n,\theta _n}(w)`$,
$$\mathrm{\Phi }^{(n)}(w)=\varphi _{\lambda _1,\theta _1}\varphi _{\lambda _2,\theta _2}\mathrm{}\varphi _{\lambda _n,\theta _n}(\omega ).$$
(8)
The DLA cluster is fully determined by the stochastic itinerary $`\{\theta _i\}_{i=1}^n`$. In Fig. 1 we present a typical DLA cluster grown by this method to size $`n=10^5`$. The main point of this paper is that the same method can be now used to grow a large variety of interesting fractal shapes, but without any randomness in the growth algorithm.
## III DLA-like clusters without randomness
As a first example of a new model we will remove the stochasticity of DLA, leaving the growth characteristics unchanged. To this aim consider an itinerary
$$\theta _{n+1}=2\theta _n\mathrm{mod2}\pi ,$$
(9)
together with Eq.(4). Such an itinerary, although deterministic, is chaotic (in fact Bernoulli, Kolmogorov and ergodic), covering the unit circle uniformly, with $`\delta `$-function correlation between consecutive $`\theta `$ values. Accordingly, we expect the growing cluster to be indistinguishable from a DLA, as is indeed the case, cf. Fig. 2.
One advantage of the present formalism is that such a statement can be made quantitatively, not by eyeball. The function $`\mathrm{\Phi }^{(n)}(w)`$ and $`\varphi _{\lambda ,\theta }(w)`$ can be expanded in a Laurent series in which the highest power is $`w`$ :
$$\mathrm{\Phi }^{(n)}(w)=F_1^{(n)}w+F_0^{(n)}+F_1^{(n)}w^1+F_2^{(n)}w^2+\mathrm{}$$
(10)
The recursion equations for the Laurent coefficients of $`\mathrm{\Phi }^{(n)}(w)`$ can be obtained analytically, and in particular one shows that
$$F_1^{(n)}=\underset{k=1}{\overset{n}{}}[1+\lambda _k]^a.$$
(11)
The importance of this lies in the fact that $`F_1^{(n)}`$ determines that fractal dimension of the cluster. Defining $`R_n`$ as the minimal radius of all circles in $`z`$ that contain the $`n`$-cluster, one can prove that
$$R_n4F_1^{(n)}.$$
(12)
Accordingly one expects that
$$F_1^{(n)}n^{1/D}\sqrt{\lambda _0},$$
(13)
as $`\sqrt{\lambda _0}`$ is the only length scale in the problem. We can thus present, as an example, plots of $`F_1^{(n)}`$ for our deterministic model (9) together with $`F_1^{(n)}`$ in any stochastic DLA growth, see Fig. 3. Another comparison is furnished by the statistics of $`\lambda _n`$. For the DLA case it was shown in that
$$\lambda _n=\frac{1}{aDn},$$
(14)
where the average is taken over the harmonic measure. This is in agreement with the โelectrostatic relationโ derived by Halsey, . In the Bernoulli itinerary there is no randomness and no probability measure, but we may still define a โrunning averageโ by, say, the last $`M`$ iterations
$$\lambda _n_M\frac{1}{M}\underset{k=nM}{\overset{n}{}}\lambda _k.$$
(15)
In Fig. 4 we show a related quantity, $`(_{k=nM}^Mk\lambda _k)/M`$ for $`M=1000`$ and $`M=10000`$. We see that up to the expected fluctuations it settles down very quickly on the appropriate value of the DLA cluster, i.e. $`1/aD=.877..`$. Any other quantitative comparison that one can think of leads to the same conclusion, i.e the Bernoulli itinerary is a bona fide generic DLA. Of course, this is not surprising: the correlation properties of successive values of $`\theta _n`$ in (9) are indistinguishable from random numbers on the interval $`[0,2\pi ]`$. Nevertheless, our point is that the present growth algorithm gives us freedom to choose deterministic itineraries resulting in DLA or other growth patterns, and we next exploit this freedom to explore new geometries.
## IV Fractal Growth with quasi-periodic itineraries
### A Winding numbers and geometry
A new class of models is obtained by using a quasi-periodic itinerary. Consider a simple map of the circle with a winding number W:
$$\theta _{i+1}=\theta _i+2\pi W.$$
(16)
If we choose $`W`$ rational, $`W=P/Q`$, then after a cross-over time the cluster grown is locked into a 1-dimensional object made of rays. In the next subsection we present an extensive discussion of the cross over time and of the properties of the 1-dimensional phase of growth. As an example consider in Fig. 5 the cluster resulting from (16) with $`W=233/144`$. On the other hand, for an irrational winding $`W`$ the itinerary is ergodic and the cluster grown is geometrically non-trivial. As a first example we present the case $`W=\rho `$ where $`\rho `$ is the Golden Mean, $`\rho =(\sqrt{5}+1)/2`$. The fractal cluster that is associated with this rule is shown in Fig.6. The cluster has a fractal dimension $`D=1.86\pm 0.03`$, as determined from the scaling of $`F_1^{(n)}`$. This is considerably higher than DLA (for which $`D1.71`$).
The golden mean is best approximated by the continued fraction representation
$$\rho =\frac{1}{1+\frac{1}{1+\frac{1}{1+\mathrm{}}}}.$$
(17)
Such a continued fraction is denoted below as $`[0,\overline{1}]`$. It is known that the golden mean is special in presenting the slowest converging continued fraction. Other quadratic irrationals also have periodic continued fractions that converge faster. In Figs. 7-11 we show the clusters grown with $`W=\sqrt{2},\sqrt{3},(1+\sqrt{10})/3,(\sqrt{13}1)/2`$ and $`\sqrt{7}`$ respectively. The continued fraction representations of these winding numbers are $`[1,\overline{2}],[1,\overline{1,2}],[\overline{1,2,1}],[1,\overline{3}],[2,\overline{1,1,1,4}]`$ respectively. In choosing these examples we picked quadratic irrationals whose representations converge relatively slowly. This facilitates the exposition of scaling theory presented below.
We note that the clusters shown have very complicated geometry. Consider for example the cases $`W=(\sqrt{13}1)/2`$ and $`W=\sqrt{7}`$ shown in Figs. 10, 11 respectively. They exhibit thin spiral growth patterns at their root, and then become bushy and thin in an apparently oscillatory fashion. Accordingly, it becomes unclear whether the different quadratic irrational winding numbers result in the same overall fractal dimension. This question warrants some extra analysis. We will argue below that in spite of the difference appearance and the oscillations in the โbushinessโ, the clusters grown by quadratic irrational winding numbers have the same fractal dimension $`D`$.
### B Different growth rules: period doubling itinerary
Clearly, one can come up with an arbitrary number of different growth rules. In this paper we will consider only one additional itinerary, to underline the fact that quadratic irrational windings lead to a class of their own. This itinerary is constructed such that after every $`2^n`$ iterations the points $`\theta _k`$ chosen on the circle are equidistributed without repetitions. The order of visitation is determined by the following rule:
$`\theta _i`$ $`=`$ $`2\pi x_i,`$ (18)
$`x_{i+1}`$ $`=`$ $`x_i+{\displaystyle \frac{3}{2^{k_i+1}}}1`$ (19)
$`k_i`$ $`=`$ $`[\mathrm{log}_2(1x_i)],`$ (20)
where $`[\mathrm{}]`$ stands for the integer value. We refer to this itinerary below as the โperiod doublingโ algorithm. The cluster grown with this rule is shown in Fig. 12 . The dimension of this cluster is $`D=1.77\pm 0.02`$. In contrast to the quadratic irrationals in this case a comparison of $`F_1^{(n)}`$ of this cluster to $`F_1^{(n)}`$ of the golden mean itinerary shows a different scaling dependence on $`n`$ (cf. Fig. 13 ).
### C Universality classes?
In the previous section we noted that the geometry of some of the clusters with quadratic irrational winding exhibit oscillations. It is thus not clear whether they have the same fractal dimension $`D`$. In this subsection we provide numerical test of the claim that the quadratic irrationals belong to the same universality class. In the the following sections we address this question using additional tools.
To study quantitatively the oscillatory fractal geometry we consider the dependence of $`F_1^{(n)}`$ on $`n`$. In Fig. 14 panel a we present compensated plots of $`F_1^{(n)}(\sqrt{2})`$ vs.$`F_1^{(n)}(\sqrt{3})`$ and vs. $`F_1^{(n)}`$ of a DLA as a function of $`n`$.
It appears that although this ratio exhibits oscillations, these are bounded and decreasing in amplitude, at least up to $`n=10^5`$. For comparison we show in panel b of Fig. 14 a plot of $`F_1^{(n)}(\sqrt{2})/F_1^{(n)}(DLA)`$. Here we see the clear difference in dimension as seen in the ratio approaching zero as a power law in $`n`$. In Fig. 15 we show compensated plots of $`F_1^{(n)}`$ of the clusters in Figs. 7-11 versus $`F_1^{(n)}`$ of the golden mean growth. We see oscillations on the logarithmic scale, but again these are bounded, and we propose that this points towards the possibility that all quadratic irrationals winding numbers lead to the same overall dimension of the cluster. In the next section we address the issue of universality classes using additional tools.
## V Towards a scaling theory: winding with rational approximants
To gain understanding of the geometry of the clusters grown with quadratic irrational winding numbers we will make use now of the well known fact that these irrationals can be systematically approximated by rational approximants. Thus, having a cluster constructed with a golden mean itinerary, a natural question is what happens to the growth pattern when $`\rho `$ is replaced by ratios of successive Fibonacci numbers which are defined by the recursion relation $`F_{m+1}=F_m+F_{m1}`$, $`F_0=0`$, $`F_1=1`$. Using rational approximants $`\rho _m=F_{m1}/F_m`$, the itinerary becomes periodic on the unit circle with period $`F_m`$, and it is observed in simulations (see Fig. 5 ) that while for small clusters $`nn_c(F_m)`$ the cluster appears fractal, for $`nn_c(F_m)`$ the cluster consists of a set of $`F_m`$ rays, sometimes fused into a smaller set of 1-dimensional rays whose number is extremely sensitive to the initial conditions (here controlled by the value of $`\lambda _0`$).
### A The 1-dimensional phase
The properties of the 1-dimensional phase are important for developing a scaling theory. As an example of the interesting behaviour seen as a function of $`\lambda _0`$ consider Fig. 16 in which clusters with $`W=144/89`$ are grown with 4 values of $`\lambda _0`$ which are 0.11, 0.22, 0.44, and 0.88. Evidently the cross over from fractal to 1-dimensional behaviour depends on $`\lambda _0`$. We also note that the number of rays in the 1-dimensional phase has a nonmonotonic dependence on $`\lambda _0`$. This indicates high sensitivity of the number of rays to changes in the initial conditions. Obviously, the radius of the cluster in the 1-dimensional case is inversely proportional to the number of rays. On the other hand, we have found a surprising invariant: $`F_1^{(n)}`$ is asymptotically invariant to the number of rays (i.e. to initial conditions) being always equal to $`n\sqrt{\lambda _0}/Q`$, up to a constant of proportionality depending on the microscopic parameter $`a`$ only. The numerical evidence is shown in Fig. 17. Note the convergence to the golden mean in panel a, and to $`\sqrt{2}`$ in panel b (which is the value of the ratios of $`\sqrt{\lambda _0}`$). This finding puts strict bounds on the number of possible rays. The upper bound is obviously $`Q`$. The lower bound stems from the inequality $`R_n4F_1^{(n)}`$, meaning that the number of rays must be larger than $`Q/4`$. This invariance also indicates that the geometry of the rays is not arbitrary, and that the angles between them are arranged to agree with an invariant $`F_1^{(n)}`$.
### B Scaling Function
The crossover in fractal shape is a general result for any periodic itinerary with $`W=P/Q`$, and suggests the existence of a scaling for $`F_1^{(n)}`$ of the form
$$F_1^{(n)}=n^{1/D}\sqrt{\lambda _0}f(n^{1/\alpha }/Q),$$
(21)
where we have assumed that the crossover cluster size scales as
$$n_c(Q)Q^\alpha .$$
(22)
The asymptotic forms of $`F_1^{(n)}`$ obey $`F_1^{(n)}n^{1/D}\sqrt{\lambda _0}`$ for $`nn_c(Q)`$, while $`F_1^{(n)}(n/Q)\sqrt{\lambda _0}`$ for $`nn_c(Q)`$. In the first asymptote we expect $`D`$ to be the same for all values of rational approximants to $`\rho `$, including the limiting fractal cluster. The growing cluster cannot distinguish between the rational approximant and the limiting irrational as long as the fractal phase is observed. The second asymptote is demonstrated in the previous subsection. Thus we require that the asymptotic forms of the scaling function obey
$`f(u)`$ $``$ $`f(0)\text{ as }u0,`$ (23)
$`f(u)`$ $``$ $`u\text{ as }u\mathrm{}.`$ (24)
The second asymptote (24) determines the scaling relation
$$\alpha =D/(D1).$$
(25)
For the Golden Mean fractal $`D1.86`$ and consequently in this case $`\alpha 2.16`$.
In Fig. 18 $`F_1^{(n)}/(n^{1/D}\sqrt{\lambda _0})`$ is plotted against the scaling variable $`u=n^{1/\alpha }/Q`$ for six different clusters with different values of $`W`$ and $`\lambda _0`$. The best data collapse was obtained using the value $`\alpha =2.15`$ The data collapse achieved is readily apparent with the scaling function $`f(u)`$ predicted by the theory.
## VI The crossover and the estimate of the dimension
In this section we discuss the properties of the conformal map $`\varphi _{\lambda ,\theta }`$ which determine the cross over from fractal to 1-dimensional growth. In other words, we will attempt to provide an independent estimate of $`n_c`$ as a function of the winding number $`W`$. If we succeeded to estimate the exponent $`\alpha `$ in (22) independently from Eq.(25), we would have an equation for the dimension.
To understand the crossover, we note that the reason for the fractal growth phase with rational winding is that after every event of growth the interface $`z(e^{i\theta })`$ is non-locally reparametrized in addition to the local growth event. Accordingly, a periodic orbit on the unit circle is not necessarily mapped to a periodic orbit in $`z`$. The region in the unit circle which is significantly affected by growing the $`n`$th bump has a scale $`\sqrt{\lambda _n}`$ centered around $`\theta _n`$ . Accordingly we can estimate when reparametrization will cause a โmissโ in the mapped orbit: as long as
$$\sqrt{\lambda _n}\frac{2\pi }{Q},$$
(26)
the growth will remain fractal. We can therefore expect a crossover to 1-dimensional growth when this condition is violated, something that is bound to happen since typical values of $`\lambda _n`$ are expected to decreases with $`n`$, cf. Eq.(14) and the discussion below.
What remains is to estimate $`\lambda _n`$ as a function of $`n`$ in the cross over region that is defined by
$$\sqrt{\lambda _{n_c}}2\pi /Q.$$
(27)
In the fractal region $`\lambda _n`$ is a highly erratic function of $`n`$. Even though we do not have here randomness in the sense of DLA, it is natural to consider, in a fashion similar to Eq.(15), the distribution of $`\lambda _k`$ over $`Q`$ successive steps of growth. For $`Q`$ large enough such distributions have well defined moments. In particular consider the first moment
$$\lambda _n_Q\frac{1}{Q}\underset{k=nQ}{\overset{n}{}}\lambda _k$$
(28)
The power law dependence of $`F_1^{(n)}`$ and Eq.(11) imply that this moment has to be
$$\lambda _n_Q=\frac{1}{anD}.$$
(29)
If we estimate $`\lambda _{n_c}`$ in Eq. (27) by its mean (29), we would write
$$\lambda _{n_c}1/n_cn_CQ^2.$$
(30)
Thus $`D/(D1)=2`$ or $`D=2`$. Even though we get an overestimate, this is a good indication that we are on the right track. The reason for the overestimate is that we neglected the fluctuations that sometime lead to $`\lambda _n`$ much larger than the mean. We expect a cross over to occur when the largest $`\sqrt{\lambda _k}`$ are smaller than $`2\pi /Q^2`$, since it is enough to have a few large $`\lambda _k`$ to cause a reparametrization that will ruin a potential periodic orbit. We thus seek a condition
$$\lambda _n^{\mathrm{max}}\mathrm{max}\{\lambda _k\}_{k=nQ}^n\frac{4\pi ^2}{Q^2}.$$
(31)
We note that $`\lambda _k`$ is an erratic function of $`k`$, and therefore the condition (31) can be met more than once in a given series $`\lambda _k`$. In Fig. 19 we show two log-log plots of $`n_c`$ computed from the value of $`n`$ for which $`\sqrt{\lambda _n^{\mathrm{max}}}=2\pi /Q`$, plotted as a function of $`Q=F_m`$. The cross over value $`n_c`$ was computed in two different ways. In circles we exhibit the values obtained from measuring when $`\sqrt{\lambda _n^{\mathrm{max}}}=2\pi /Q`$ for the first time, and in squares we exhibit the values obtained from $`\sqrt{\lambda _n^{\mathrm{max}}}=2\pi /Q`$ for the last time. Computing the slopes by linear regression and averaging between them we find the scaling law
$$n_cQ^{2.17\pm 0.03}.$$
(32)
Comparing with Eqs.(22), (25) we get an estimate for $`D=1.86\pm 0.03`$, in excellent agreement with the determination of the dimension by $`F_1^{(n)}`$.
We note in passing that $`\lambda _n^{\mathrm{max}}`$ can be assigned a generalized dimension $`D_{\mathrm{}}`$ in the language of Hentschel and Procaccia . Define
$$\lambda _n^q_Q\frac{1}{Q}\underset{k=nQ}{\overset{n}{}}\lambda _k^q.$$
(33)
From the precise scaling law is
$$\lambda _n^{\mathrm{max}}=\underset{q\mathrm{}}{lim}\lambda _n^q_Q^{1/q}n^{2D_{\mathrm{}}/D}.$$
(34)
Comparing with (25) we conclude that in this case there exists a scaling relation
$$D_{\mathrm{}}=D1.$$
(35)
Such a scaling relation was conjectured by Turkevich and Scher for DLA (of course with a different $`D`$ and $`D_{\mathrm{}}`$). While there are severe doubts about the correctness of this conjecture for DLA, we point out that in our case it follows directly from elementary considerations.
### A The period doubling itinerary
Even though the period doubling itinerary leads to a cluster whose fractal dimension differs from the quadratic irrational windings, we show here that the ideas presented above pertain equally to this growth pattern. Instead of rational approximants we use here, naturally, $`2^n`$-periodic orbits which are obtained by cutting the itinerary (20) after $`2^n`$ iterations and repeating it periodically. The crossover from fractal to 1-dimensional growth is seen also in this case, and we can use it in a very similar way to identify the crucial exponent that determines the dimension of the asymptotic cluster. Indeed, the whole set of ideas developed above repeats verbatim by changing $`Q`$ with $`2^n`$. What remains is to find $`\lambda _n^{\mathrm{max}}`$ as a function of $`n`$. In Fig .20 we show the data collapse obtained as in Fig. 18 for the quasiperiodic analog. We show nine different data sets with periodic itineraries of periods 32, 64 and 128 and $`\sqrt{\lambda _0}`$ values of 0.22, 0.44 and 0.88. The scaling function for these data sets is plotted as a function of $`u=n^{.44}/Q`$, where the exponent is computed from $`D=1.78`$. It is noteworthy that the scaling function obtained appears identical to the scaling function $`f(u)`$ for the quasiperiodic family. For comparison we added in Fig .20 also one curve from the quasiperiodic class, and it appears indistinguishable from the rest.
The conclusion from this data collapse is that the mechanism governing the crossover from fractal to 1-dimensional growth phases here is the same as the one discussed above for the quasiperiodic itineraries. The difference between the dimensions of the period doubling cluster and the quasiperiodic cluster must lie in the different numerical value of the exponent characterizing $`\lambda _n^{\mathrm{max}}`$ as a function of $`n`$. In this case the natural averaging cycles are of length $`Q=2^n`$. Fig. 21 is the analog of Fig. 19 for the period doubling itinerary, where the critical value $`n_c`$ was estimated from the first time that $`\sqrt{\lambda _n^{\mathrm{max}}}`$ became smaller than $`2\pi /2^n`$.
The linear regression provides us with the the scaling law
$$n_cQ^{2.33\pm 0.1}.$$
(36)
Computing $`D`$ we find $`D=1.75\pm 0.05`$ in good agreement with the numerical estimate from $`F_1^{(n)}`$.
## VII Summary and the road ahead
The main points of this paper are as follows:
* The iterated conformal maps algorithm for fractal growth patterns offers a convenient way to introduce a large number of deterministic growth models with highly non-trivial fractal geometry.
* Itineraries with irrational winding numbers generate fractal growth patterns. We proposed that all the quadratic irrationals produce clusters of the same fractal dimension, in spite of different appearance.
* By considering a series of rational approximants we could produce a scaling theory of the growing clusters, achieving data collapse for all values of $`n,\lambda _0`$ and $`P/Q`$.
* Identifying the mechanism for the cross over from fractal to 1-dimensional growth phases we could pinpoint the exponent that determines the fractal dimension $`D`$. This exponent characterizes the $`n`$ dependence of the extremal values of $`\lambda _n`$.
* The mechanism appears general; itineraries leading to different cluster dimensions, like the period doubling itinerary (20) and its truncated versions, can be understood in the same way. The scaling function (21) and the scaling relation (25) are general, but the exponent $`\alpha `$ changes. Its determination by the scaling of $`\lambda _n^{\mathrm{max}}`$ Eq.(31) is however general.
We note that all the numerical tests point out in favour of this scenario, and in our opinion rule out a value $`D=2`$ for the clusters discussed above. The only way to get 2-dimensional growth, as shown above, is if the distribution of $`\lambda _n`$ does not multiscale, i.e. all $`D_q`$ are the same, and the scaling of $`\lambda _n^{\mathrm{max}}`$ identifies with the scaling of the average of $`\lambda _n`$.
Nevertheless, we point out that the crucial step in our scenario, the determination of the exponent $`\alpha `$ in Eq.(22), was achieved numerically. The scaling theory presented above has a strong flavour of a renormalization group approach. It appears that such an underlying theory may have a low codimension, maybe with 1 important exponent, the one characterizing the rate of crossover of the rational approximants to the irrational limit. The search of such a theory appears to be an important task for the near future.
###### Acknowledgements.
We benefitted from discussions with T.C. Halsey, C. Tresser and L. Peliti. This work has been supported in part by the European Commission under the TMR program and the Naftali and Anna Backenroth-Bronicki Fund for Research in Chaos and Complexity. |
warning/0002/math-ph0002003.html | ar5iv | text | # On the complete ionization of a periodically perturbed quantum system
## 1. Introduction and results
The ionization of atoms<sup>1</sup><sup>1</sup>footnotetext: Also Department of Physics.
To appear in Proceedings of the CRM meeting โNonlinear Analysis and Renormalization Groupโ subjected to external time dependent perturbations is an issue of central importance in quantum mechanics which has attracted substantial theoretical and experimental interest , . There exists by now a variety of theoretical methods, and a vast amount of literature, devoted to the subject. Beyond the celebrated Fermiโs golden rule, approaches include higher order perturbation theory, semi-classical phase-space analysis, Floquet theory, complex dilation, some exact results for small fields and bounds for large fields and numerical integration of the time dependent Schrรถdinger equation -. Nevertheless there is apparently no complete analysis of the ionization of any periodically perturbed model with no restrictions on the amplitudes and frequencies of the perturbing field. This is not so surprising considering the very complex behavior we find in even the most elementary of such systems.
In the present paper we show rigorously the full ionization, in all ranges of amplitudes and frequencies, of one of the simplest models with spatial structure which, with a different perturbing potential, is however frequently used as a model system , , . The unperturbed Hamiltonian we consider is
(1.1)
$$_0=\frac{\mathrm{}^2}{2m}\frac{\mathrm{d}^2}{\mathrm{d}x^2}g\delta (x),g>0,\mathrm{}<x<\mathrm{}.$$
$`_0`$ has a single bound state $`u_b(x)=\sqrt{p_0}e^{p_0|x|},p_0=\frac{m}{\mathrm{}^2}g`$ with energy $`\mathrm{}\omega _0=\mathrm{}^2p_0^2/2m`$ and a continuous uniform spectrum on the positive real line, with generalized eigenfunctions
$$u(k,x)=\frac{1}{\sqrt{2\pi }}\left(e^{ikx}\frac{p_0}{p_0+i|k|}e^{i|kx|}\right),\mathrm{}<k<\mathrm{}$$
and energies $`\mathrm{}^2k^2/2m`$.
Beginning at $`t=0`$, we apply a parametric perturbing potential, i.e. for $`t0`$ we have
(1.2)
$$(t)=_0g\eta (t)\delta (x)$$
and solve the time dependent Schrรถdinger equation for $`\psi (x,t)`$,
$`\psi (x,t)=\theta (t)u_b(x)e^{i\omega _0t}`$
(1.3) $`+{\displaystyle _{\mathrm{}}^{\mathrm{}}}\mathrm{\Theta }(k,t)u(k,x)e^{i\frac{\mathrm{}k^2}{2m}t}๐k(t0)`$
with initial values $`\theta (0)=1,\mathrm{\Theta }(k,0)=0`$. This gives the survival probability $`|\theta (t)|^2`$, as well as the fraction of ejected electrons $`|\mathrm{\Theta }(k,t)|^2dk`$ with (quasi-) momentum in the interval $`dk`$.
This problem can be reduced to the solution of an integral equation . Setting
(1.4) $`\theta (t)=1+2i{\displaystyle _0^t}Y(s)๐s`$
(1.5) $`\mathrm{\Theta }(k,t)=2|k|/\left[\sqrt{2\pi }(1i|k|)\right]{\displaystyle _0^t}Y(s)e^{i(1+k^2)s}๐s`$
$`Y(t)`$ satisfies the integral equation
(1.6)
$$Y(t)=\eta (t)\left\{1+_0^t[2i+M(tt^{})]Y(t^{})๐t^{}\right\}=\eta (t)\left(1+(2i+M)Y\right)$$
where $`\mathrm{},2m`$ and $`\frac{g}{2}`$ have been set equal to $`1`$ (implying $`p_0=1`$, $`\omega _0=1`$),
$$M(s)=\frac{2i}{\pi }_0^{\mathrm{}}\frac{u^2e^{is(1+u^2)}}{1+u^2}๐u=\frac{1+i}{2\sqrt{2}\pi }_s^{\mathrm{}}\frac{e^{iu}}{u^{3/2}}๐u$$
and
$$fg=_0^tf(s)g(ts)๐s$$
###### Theorem 1.
When $`\eta (t)=r\mathrm{sin}\omega t`$ the survival probability $`|\theta (t)|^2`$ tends to zero as $`t\mathrm{}`$, for all $`\omega >0`$ and $`r0`$.
Note: For definiteness we assume in the following that $`r>0`$.
The method of proof relies on the properties of the Laplace transform of $`Y`$, $`y(p)=Y(p)=_0^{\mathrm{}}e^{pt^{}}Y(t^{})๐t^{}`$ (note that $`y(p)=\frac{i}{2}(1p\theta )`$). In particular we need to show that $`y(p)`$ is bounded in the closed right half of the complex $`p`$ plane. Before the proof we describe briefly some additional results on this model system, cf. .
### 1.1. Further results not proven in the present paper
(1) Theorem 1 generalizes to the case when $`\eta (t)`$ is a trigonometric polynomial:
(1.7)
$$\eta (t)=\underset{j=1}{\overset{K}{}}[A_j\mathrm{sin}(j\omega t)+B_j\mathrm{cos}(j\omega t)],$$
where we assume $`|A_K|+|B_K|0`$.
(2) The detailed behavior of the system as a function of $`t`$, $`\omega `$, and $`r`$ is obtained from the singularities of $`y(p)`$ in the complex $`p`$-plane. We summarize them for small $`r`$; below $`\frac{1}{2}<\delta <1`$.
At $`p=\{in\omega i:n\}`$, $`y`$ has square root branch points and $`y`$ is analytic in the right half plane and also in an open neighborhood $`๐ฉ`$ of the imaginary axis with cuts through the branch points. As $`|\mathrm{}(p)|\mathrm{}`$ in $`๐ฉ`$ we have $`|y(p)|=O(r\omega |p|^2)`$. If $`|\omega \frac{1}{n}|>\mathrm{const}_nO(r^{2\delta }),n^+`$, then for small $`r`$ the function $`y`$ has a unique pole $`p_m`$ in each of the strips $`m\omega >\mathrm{}(p)+1\pm O(r^{2\delta })>m\omega \omega ,m`$. $`\mathrm{}(p_m)`$ is strictly independent of $`m`$ and gives the exponential decay of $`\theta `$. After suitable contour deformation of the inverse Laplace transform, $`\theta `$ can be (uniquely) written in the form
(1.8) $`\theta (t)=e^{\gamma (r;\omega )t}F_\omega (t)+{\displaystyle \underset{m=\mathrm{}}{\overset{\mathrm{}}{}}}e^{(mi\omega i)t}h_m(t)`$
where $`F_\omega `$ is periodic of period $`2\pi \omega ^1`$ and
$$h_m(t)\underset{j=0}{\overset{\mathrm{}}{}}c_{m,j}t^{3/2j}\text{as }t\mathrm{},\mathrm{arg}(t)(\frac{\pi }{2}ฯต,\frac{\pi }{2}+ฯต)$$
Not too close to resonances, i.e. when $`|\omega n^1|>O(r^{2\delta })`$, for all integer $`n`$, $`|F_\omega (t)|=1\pm O(r^2)`$ and its Fourier coefficients decay faster than $`r^{|2m|}|m|^{|m|/2}`$. Also, the sum in (1.8) does not exceed $`O(r^2t^{3/2})`$ for large $`t`$, and the $`h_m`$ decrease with $`m`$ faster than $`r^{|m|}`$.
(3) By (1.8), for times of order $`1/\mathrm{\Gamma }`$ where $`\mathrm{\Gamma }=2\mathrm{}(\gamma )`$, the survival probability for $`\omega `$ not close to a resonance decays as $`\mathrm{exp}(\mathrm{\Gamma }t)`$. This is illustrated in Figure 1 where it is seen that for $`r1/2`$ the exponential decay holds up to times at which the survival probability is extremely small, after which $`|\theta (t)|^2=O(t^3)`$ with many oscillations as described by (1.8). Note the slow decay for $`\omega =.8`$, when ionization requires the absorption of two photons.
(4) When $`r`$ is larger the polynomial-oscillatory behavior starts sooner. Since the amplitude of the late asymptotic terms is $`O(r^2)`$ for small $`r`$, increased $`r`$ yields higher late time survival probability. This phenomenon, sometimes referred to as atomic stabilization , , can be associated with the perturbation-induced probability of back-transitions to the well.
(5) Using the continued fraction representation (2.12) $`\mathrm{\Gamma }`$ can be calculated convergently for any $`\omega `$ and $`r`$.
The limiting behavior for small $`r`$ of the exponent $`\mathrm{\Gamma }`$ is described as follows. Let $`n`$ be the integer part of $`\omega ^1+1`$ and assume $`\omega ^1`$. Then we have, for $`T>0`$ ($`t=r^{2n}T`$),
(1.9)
$$\widehat{\mathrm{\Gamma }}=T^1\underset{r0}{lim}\mathrm{ln}\left|\theta (r^{2n}T)\right|^2=\frac{2^{2n+2}\sqrt{n\omega 1}}{n\omega {\displaystyle \underset{m<n}{}}(1\sqrt{1m\omega })^2}$$
(6) The behavior of $`\mathrm{\Gamma }`$ is different at the resonances $`\omega ^1`$. For instance, whereas if $`\omega `$ is not close to $`1`$, the scaling of $`\mathrm{\Gamma }`$ implied by (1.9) is $`r^2`$ when $`\omega >1`$ and $`r^4`$ when $`\frac{1}{2}<\omega <1`$, by taking $`\omega 1=r^2/\sqrt{2}`$ we find
$$T^1\underset{\begin{array}{c}r0\\ \omega =1+r^2/\sqrt{2}\end{array}}{lim}\mathrm{ln}\left|\theta (r^3T)\right|^2=\frac{2^{1/4}}{8}\frac{2^{3/4}}{16}$$
## 2. Proofs of Theorem 1
###### Lemma 2.
(i) $`Y`$ exists and is analytic in the right half plane $`=\{p:\mathrm{}(p)>0\}`$. Furthermore, $`y(p)0`$ as $`\mathrm{}(p)\pm \mathrm{}`$ in $``$.
(ii) The function $`y(p)`$ satisfies (and is determined by) the functional equation
(2.1) $`y=r\left(T^{}T^+\right)\left(h_1+h_2y\right)`$
with
$$\left(T^\pm f\right)(p)=f(p\pm i\omega ),h_1(p)=\frac{i}{2p}\text{ and }h_2(p)=\frac{1}{2p}\left(1+\sqrt{1ip}\right)$$
and by the boundary condition $`y(p)0`$ as $`\mathrm{}(p)\pm \mathrm{}`$ in $``$.
The branch of the square root is such that for $`p`$, the real part of $`\sqrt{1ip}`$ is nonnegative and the imaginary part nonpositive.
###### Proof.
(i) The time evolution of $`\psi `$ is unitary and thus $`|\psi |u_b|=|\theta (t)|1`$. The stated analyticity is an immediate consequence of the elementary properties of the Laplace transform <sup>2</sup><sup>2</sup>2See also the appendix for a proof of analyticity for $`\mathrm{}(p)>p_0`$ (all that is required in the subsequent analysis), relying only on the properties of the convolution equation.. The asymptotic behavior follows then from the Riemann-Lebesgue lemma.
(ii) We have in $``$,
(2.2) $`M=\underset{a0}{lim}{\displaystyle \frac{2i}{\pi }}{\displaystyle _0^{\mathrm{}}}dxe^{px}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{u^2e^{i(xia)(1+u^2)}}{1+u^2}}du`$
(2.3) $`={\displaystyle \frac{i}{\pi }}{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \frac{u^2}{(1+u^2)(p+i(1+u^2))}}du`$
For $`\mathrm{}(p)>0`$ we push the integration contour through the upper half plane. At the two poles in the upper half plane $`u^2+1`$ equals $`0`$ and $`ip`$ respectively, so that
(2.4)
$$\begin{array}{c}\frac{i}{\pi }_{\mathrm{}}^{\mathrm{}}\frac{u^2}{(1+u^2)(p+i(1+u^2))}du\hfill \\ \hfill =\frac{i}{\pi }\left(\frac{(1)}{(2i)(p)}\frac{ds}{s}+\frac{u_0^2}{(ip)(2iu_0)}\frac{ds}{s}\right)=\frac{i}{p}+\frac{u_0}{p}\end{array}$$
where $`u_0`$ is the root of $`p+i(1+u^2)=0`$ in the upper half plane. Thus
(2.5)
$$M=\frac{i}{p}+\frac{i\sqrt{1ip}}{p}$$
with the branch satisfying $`\sqrt{1ip}1`$ as $`p0`$ in $``$. As $`p`$ varies in $``$, $`1ip`$ belongs to the lower half plane $`i`$ and then $`\sqrt{1ip}`$ varies in the fourth quadrant.
For $`\mathrm{}(p)>0,\omega >0`$ we have
(2.6) $`\left(e^{\pm i\omega }M\right)={\displaystyle \frac{i}{pi\omega }}+{\displaystyle \frac{i\sqrt{1ip\omega }}{pi\omega }}`$
$`(\text{with }\sqrt{1ip\omega }=i\sqrt{\omega 1+ip}\text{ if }\omega >1)`$
and relation (2.1) follows. โ
After the substitution $`y(p)=2(\sqrt{1ip}1)e^{\frac{\pi p}{2\omega }}v(p)`$ we get
(2.7)
$$v(pi\omega )+v(p+i\omega )=\frac{2}{r}(\sqrt{1ip}1)v(p)+\frac{i\omega }{\omega ^2+p^2}$$
###### Remark 3.
It is clear that the functional equation (2.7) only links the points on one dimensional lattice $`\{p+i\omega \}`$. It is convenient to take $`p_0`$ such that $`p=p_0+in\omega `$ with $`\mathrm{}(p_0)=\mathrm{}(p)`$ and
(2.8)
$$\mathrm{}(p_0)[0,\omega )$$
and write $`v(p)=v(p_0+in\omega )=v_n`$ which transforms (2.7) to a recurrence relation:
(2.9)
$$v_{n+1}+v_{n1}=\frac{2}{r}(\sqrt{1ip_0+n\omega }1)v_n+\frac{i\omega }{\omega ^2+(p_0+in\omega )^2}$$
where $`v_n`$ depends parametrically on $`p_0`$. It will be seen that the asymptotic conditions as well as analyticity in $`p_0`$ determine the solution of (2.9) uniquely.
###### Remark 4.
The approach is based on a discrete analog of the Wronskian technique. The regularity of the bounded solution of (2.9) will be a consequence of the absence of a bounded solution of the homogeneous equation
(2.10)
$$v_{n+1}+v_{n1}=\frac{2}{r}(\sqrt{1ip_0+n\omega }1)v_n=D_nv_n$$
a problem which we analyze first.
###### Proposition 5.
For $`p_0`$ satisfying (2.8) and $`\mathrm{}(p_0)0`$ (actually for any $`p_0\overline{}=i`$) there is no nonzero solution of (2.10) such that $`vl_2()`$.
###### Proof.
To get a contradiction, assume $`v0`$ is an $`l_2()`$ solution of (2.10). Multiplying (2.10) by $`\overline{v_n}`$, and summing with respect to $`n`$ from $`\mathrm{}`$ to $`+\mathrm{}`$ we get
(2.11)
$$\begin{array}{c}\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}v_{n+1}\overline{v}_n+\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}v_{n1}\overline{v}_n\hfill \\ \hfill =2\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}\mathrm{}(v_n\overline{v}_{n+1})=\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}\frac{2}{r}\left(\sqrt{1ip_0+n\omega }1\right)|v|_n^2\end{array}$$
For $`p_0\overline{}`$ the imaginary part of $`\sqrt{1ip_0+n\omega }`$ is nonpositive, by Lemma 2, and is strictly negative for $`n<0`$ large enough. Thus if for some such $`n`$, $`v_n`$ is nonzero then the last sum in (2.11) has a strictly negative imaginary part, which is impossible since the left side is real. If on the other hand $`v_n`$ is zero when $`n`$ is large negative, then solving (2.10) for $`v_{n+1}`$ in terms of the $`v_n,v_{n1}`$ it would follow inductively that $`v0`$, contradicting the assumption.
###### Lemma 6.
(i) There is, up to multiplicative constants, a unique pair of solutions $`v^+`$ and $`v^{}`$ of (2.10) such that $`v_n^\pm 0`$ as $`n\pm \mathrm{}`$ ( respectively). These solutions are related to convergent continued fractions representations:
(2.12)
$$v_{n1}^\pm /v_n^\pm =:\frac{1}{\rho _n^\pm }=D_n\frac{1}{D_{n\pm 1}{\displaystyle \frac{1}{D_{n\pm 2}}}\mathrm{}}$$
(ii) We have the following estimates
(2.13)
$$\frac{1}{\rho _n^\pm }=\frac{1}{\stackrel{~}{\rho }_n^\pm }+O(n^{3/2})(n\pm \mathrm{})$$
where
(2.14) $`{\displaystyle \frac{1}{\stackrel{~}{\rho }_n^+}}={\displaystyle \frac{2}{r}}\sqrt{n\omega }{\displaystyle \frac{2}{r}}{\displaystyle \frac{r^22+2ip_0}{2r\sqrt{n\omega }}}{\displaystyle \frac{r}{2\omega n}}(n>0)`$
$`{\displaystyle \frac{1}{\stackrel{~}{\rho }_n^{}}}={\displaystyle \frac{2i}{r}}\sqrt{|n|\omega }{\displaystyle \frac{2}{r}}+{\displaystyle \frac{(2r^2)i+2p_0}{2r\sqrt{|n|\omega }}}+{\displaystyle \frac{r}{2\omega |n|}}(n<0)`$
Let $`\stackrel{~}{v}_n^\pm `$ be solutions of the one step recurrences $`\stackrel{~}{v}_n^\pm =\stackrel{~}{v}_{n1}^\pm \stackrel{~}{\rho }_n^\pm `$. Then
(2.15)
$$\begin{array}{c}\mathrm{ln}\stackrel{~}{v}_n^+=\frac{1}{2}n\mathrm{ln}n+n\mathrm{ln}\left(\frac{r}{2}\sqrt{\frac{e}{\omega }}\right)\hfill \\ \hfill +2\sqrt{\frac{n}{\omega }}+\left(\frac{2ip_0+r^2+\omega }{4\omega }\right)\mathrm{ln}n+o(1)(n\mathrm{})\end{array}$$
and
(2.16)
$$\begin{array}{c}\mathrm{ln}(\stackrel{~}{v}_n^{})=\frac{1}{2}|n|\mathrm{ln}|n|+|n|\mathrm{ln}\left(\frac{r}{2}\sqrt{\frac{e}{\omega }}\right)+i\pi |n|2i\sqrt{|n|/\omega }\hfill \\ \hfill +\left(\frac{2ip_0+r^2+\omega }{4\omega }\right)\mathrm{ln}|n|+o(1)(n\mathrm{})\end{array}$$
and, for some constants $`K^\pm `$,
(2.17)
$$\mathrm{ln}(v_n^\pm )=\mathrm{ln}(\stackrel{~}{v}_n^\pm )+K^\pm +o(1)$$
($`v_n^\pm `$ decay roughly as $`1/\sqrt{|n|!}`$ for $`n\pm \mathrm{}`$, respectively).
(iii) Two special solutions of (2.10), $`v^+`$ and $`v^{}`$, are well defined by:
(2.18)
$$v_n^+=\stackrel{~}{v}_n^+\underset{jn+1}{}\frac{\stackrel{~}{\rho }_j^+}{\rho _j^+}\text{ for }n>N,\text{ and }v_n^{}=\stackrel{~}{v}_n^{}\underset{jn1}{}\frac{\stackrel{~}{\rho }_j^{}}{\rho _j^{}}\text{ for }n<N$$
if $`N`$ is sufficiently large (this amounts to making a convenient choice of the free multiplicative constant in (i)). These functions do not depend on $`N`$. $`v^+`$ and $`v^{}`$ are linearly independent for $`p_0\overline{}`$: their discrete Wronskian, defined by $`W(v^+,v^{})_n=v_n^+v_{n+1}^{}v_n^{}v_{n+1}^+`$, satisfies
(2.19)
$$W(v^+,v^{})=const0$$
As functions of parameters, $`v^\pm `$ and $`W(v^+,v^{})`$ are analytic in $`p_0`$. If $`\omega \{0,n^1:n\}`$ then $`v^\pm `$ and $`W(v^+,v^{})`$ are analytic in some neighborhood of $`p_0=0`$ as well. For any $`\omega >0`$, $`v_n^\pm `$ are Lipschitz continuous of exponent at least $`1/2`$ in $`p_0`$, for $`p_0`$.
###### Proof.
(i) We look at $`v^+`$, the case of $`v^{}`$ being similar. Dropping the <sup>+</sup> superscript we have from (2.10)
(2.20)
$$\rho _n=\frac{1}{D_n\rho _{n+1}}$$
To find the analytic properties of the solution $`\rho _n`$ it is convenient to regard (2.20) as a contractive equation in the space $`\mathrm{}^{\mathrm{}}(S_N)`$ of sequences $`\{\rho _j\}_{j>N}`$ in the norm $`\rho _{\mathrm{}}=sup_{j>N}|\rho _j|`$. Let $`N`$ be large. The map $`J:S_NS_N`$ defined by
(2.21)
$$J(\rho )_n=\frac{1}{D_n\rho _{n+1}}$$
depends analytically on $`p_0`$ and is Lipschitz continuous of exponent at least $`1/2`$ if $`\mathrm{}(p_0)0`$. In addition, if $`\rho _j_{\mathrm{}}1`$ we have for sufficiently large $`N=N(p_0,\omega ,r)`$
(2.22)
$$J(\rho )_{\mathrm{}}\frac{1}{\frac{2}{r}\left(\sqrt{|N\omega |1|p_0|}1\right)1}<\frac{|r|}{|\omega |^{1/2}}\frac{1}{\sqrt{N}}$$
Similarly,
(2.23)
$$J(\rho )J(\rho ^{})_{\mathrm{}}\frac{\rho \rho ^{}_{\mathrm{}}}{\left[\frac{2}{r}\left(\sqrt{|N\omega |1|p_0|}1\right)1\right]^2}<\frac{|r|^2}{N|\omega |}\rho \rho ^{}_{\mathrm{}}$$
for sufficiently large $`N`$ which shows that $`J`$ is contractive in the unit ball in $`\mathrm{}^{\mathrm{}}(S_N)`$. Thus, equation (2.21) has a unique solution in $`S_N`$, which depends analytically on $`p_0`$ and is Lipschitz continuous of exponent at least $`1/2`$ if $`\mathrm{}(p_0)0`$. This also implies the convergence of (2.12).
Note that given $`๐ฆ_1\overline{}`$ and $`๐ฆ_2^+`$ both compact, $`N`$ can be chosen the same for all $`p_0๐ฆ_1`$ and $`r๐ฆ_2`$.
(ii) From (2.22) it is seen that $`|\rho _j|=O(j^{1/2})`$ for large $`j`$. Thus, we may write, for large $`j`$,
(2.24)
$$\frac{1}{\rho _j^+}=D_j\frac{1}{D_{j+1}{\displaystyle \frac{1}{D_{j+2}+O(j^{1/2})}}}$$
The estimates (2.13) now follow by a straightforward calculation. Since $`\mathrm{ln}v_n^+=\mathrm{ln}v_N^++_{j=N+1}^n\mathrm{ln}\rho _j^+`$, the estimates follow from (2.24) and the Euler-Maclaurin summation formula.
(iii) As before, we only need to look at $`v^+`$. We take two compact sets $`๐ฆ_1`$ and $`๐ฆ_2`$, and choose $`N`$ as in the note at the end of the proof of (i). Taking the log in the definition (2.18), the infinite sums are absolutely convergent. By standard measure theory, $`v_n^+`$ has the same analyticity properties in the interior of $`๐ฆ_1\times ๐ฆ_2`$ and Lipschitz continuity in $`๐ฆ_1\times ๐ฆ_2`$ as those of $`\rho ^+`$, when $`n>N`$. Now, (2.10) easily implies that the same is true for $`nN`$ as well.
If $`f_n`$ and $`g_n`$ are solutions of (2.10) then $`(g_{n+1}+g_{n1})f_n(f_{n+1}+f_{n1})g_n=0`$ and thus $`W_n(f,g)=f_ng_{n+1}g_nf_{n+1}=const.`$ Thus, if $`W_n(f,g)=0`$ for some $`n`$ then $`W_n0`$ and $`fconstg`$. The smoothness properties follow from the proof of (iii).
###### Proposition 7.
There exists a unique solution of (2.7) which is bounded as $`\mathrm{}(p)\pm \mathrm{}`$ in $``$. This solution is analytic in $`p`$, and $`v(p)=O(p^2)`$ as $`\mathrm{}(p)\pm \mathrm{},p`$.
###### Proof.
By analyticity and continuity $`W(v^+,v^{})`$ does not vanish for any $`p`$ and $`r>0,\omega >0`$. By Lemma 11 the function $`v`$ defined through $`v(p_0+in\omega )=f_n`$, where
(2.25)
$$f_n:=W(v^+,v^{})^1\left(v_n^+\underset{l=\mathrm{}}{\overset{n1}{}}v_l^{}H_l+v_n^{}\underset{l=n}{\overset{\mathrm{}}{}}v_l^+H_l\right)$$
and
(2.26)
$$H_n=\frac{i\omega \mathrm{exp}\left(\frac{\pi p_n}{2\omega }\right)}{p_n^2+\omega ^2}$$
has the required properties. Since no solution of the homogeneous equation is bounded on $``$, $`v`$ is the unique solution with the desired properties. โ
###### Note 8.
The link between $`y`$ and $`f_n`$ is
(2.27)
$$y(p)=2(\sqrt{1ip}1)e^{\frac{\pi p}{2\omega }}f_n;\text{ }\text{for }p=p_0+in\omega $$
###### Proposition 9.
The function $`y(p)`$ is analytic in the right half plane, Lipschitz continuous of exponent at least $`1/2`$ on the imaginary axis and $`lim_{p0}y(p)=i/2`$.
###### Proof.
Since $`W`$ is analytic in $``$, continuous and nonzero in $`\overline{}`$, $`W`$ is bounded below in compact sets in $`\overline{}`$. Then, the smoothness properties of $`y`$ derive easily from those of $`q_n:=W(v^+,v^{})f_n`$ on which we concentrate now.
(a) For $`n2`$ we write, using (2.26),
(2.28)
$$\begin{array}{c}q_n=v_n^+\underset{\begin{array}{c}l=\mathrm{}\\ l\pm 1\end{array}}{\overset{n1}{}}v_l^{}H_l+v_n^{}\underset{l=n}{\overset{\mathrm{}}{}}v_l^+H_l\hfill \\ \hfill +v_n^+i\omega e^{\frac{\pi p_0}{2\omega }}\left(\frac{iv_1^{}}{p_0(p_02i\omega )}+\frac{iv_1^{}}{p_0(p_0+2i\omega )}\right)\end{array}$$
The last term in parenthesis can be rewritten, using also (2.10), as
(2.29)
$$\begin{array}{c}\frac{i(v_1^{}v_1^{})}{p_0^2+4\omega ^2}+\frac{2\omega }{p_0^2+4\omega ^2}\left(\frac{v_1^{}+v_1^{}}{p_0}\right)\hfill \\ \hfill =\frac{i(v_1^{}v_1^{})}{p_0^2+4\omega ^2}+\frac{4\omega }{r(p_0^2+4\omega ^2)}\frac{\sqrt{1ip_0}1}{p_0}v_0^{}\end{array}$$
Thus we see that $`q_n`$ is continuous as $`\mathrm{}(p_0)0`$ and $`\mathrm{}(p_0)[0,\omega )`$ \[cf. (2.8)\], if $`n2`$. A very similar calculation shows the continuity of $`q_n`$ if $`n1`$.
(b) By part (a), $`y(p)`$ is continuous as $`\mathrm{}(p)0`$ with $`\mathrm{}(p)2`$ or $`\mathrm{}(p)<0`$. Now, (2.1) written in the form
(2.30)
$$\begin{array}{c}rph_2(p)y(p)=rp\left(h_1(p+2i\omega )h_1(p)\right)\hfill \\ \hfill +rp\left(y(p+i\omega )+h_2(p+2i\omega )y(p+2i\omega )\right)\end{array}$$
shows that $`y(p)`$ is Lipschitz continuous as $`\mathrm{}(p)0`$ if $`\mathrm{}(p)>2`$ thus for all $`\mathrm{}(p)`$. The value of $`y(0)`$ is easily calculated using (2.30). โ
###### Proposition 10.
$`1+2ilim_x\mathrm{}_0^xY(s)๐s=0`$.
###### Proof.
Indeed,
(2.31)
$$\begin{array}{c}2\pi i_0^{\mathrm{}}Y(s)๐s\hfill \\ \hfill =\underset{x\mathrm{}}{lim}\underset{\delta 0^+}{lim}\left(_i\mathrm{}^{i\delta }+_{i\delta }^i\mathrm{}\right)\frac{e^{xp}}{p}(i/2+(y(p)i/2))dp=\pi \end{array}$$
## 3. Appendix
###### Lemma 11.
Equation (1.6) has a unique solution $`YL_{loc}^1(^+)`$ and $`|Y(x)|<Ke^{Cx}`$ for some $`K^+`$ and $`C`$.
###### Proof.
Consider $`L_{loc}^1[0,A]`$ endowed with the norm $`F_\nu :=_0^A|F(s)|e^{\nu s}๐s`$, where $`\nu >0`$. If $`f`$ is continuous and $`F,GL_{loc}^1[0,A]`$, a straightforward calculation shows that
(3.1) $`fF_\nu <F_\nu \underset{[0,A]}{sup}|f|`$
(3.2) $`FG_\nu <F_\nu G_\nu `$
(3.3) $`F_\nu 0\text{ as }\nu \mathrm{}`$
where the last relation follows from the Riemann-Lebesgue lemma.
The integral equation (1.6) can be written as
(3.4)
$$Y=r\eta +๐ฅY\text{ where }๐ฅF:=r\eta (2i+M)F$$
Since $`M`$ is locally in $`L^1`$ and bounded for large $`x`$ it is clear that for large enough $`C_2`$, and for any $`A`$, (1.6) is contractive if $`\nu >C_2`$. โ
Acknowledgments. The authors would like to thank A. Soffer and M. Weinstein for interesting discussions and suggestions. Work of O. C. was supported by NSF Grant 9704968, that of J. L. L. and A. R. by AFOSR Grant F49620-98-1-0207 and NSF Grant DMR-9813268. |
warning/0002/hep-ph0002289.html | ar5iv | text | # RIGHT-HANDED VECTOR ๐ AND AXIAL ๐ด COUPLINGS IN WEAK INTERACTIONS
## 1 Introduction
The present theory of weak interactions (the Standard Model of electroweak interactions ) describes only what has been measured so far. These are, most of all, the measurements of nuclear observables and of observables for massive leptons. It means the measurement of the electron helicity , the indirect measurement of the neutrino helicity , the asymmetry in the distribution of the electrons from $`\beta `$-decay , the experiment with muon decay confirming parity violation . Basing their inferences on these results, among others, Feynman, Gell-Mann and Sudarshan, Marshak established that only left-handed vector $`V`$, axial $`A`$ couplings are involved in weak interactions because this yields the maximum symmetry breaking under space inversion, under charge conjugation; the two-component neutrino theory of negative helicity; the conservation of the combined symmetry $`CP`$ and of the lepton number. The Fermi hamiltonian, being low-energy approximation of the Salam-Weinberg model, has thus a vector-axial $`(VA)`$ structure, and the three remaining scalar $`S`$, tensor $`T`$, pseudoscalar $`P`$ couplings are eliminated by the assumption that only left-handed states can take part. The $`VA`$ theory has a chiral symmetry.
The investigation of the completeness of the Lorentz structure and of the handedness structure of weak interactions at low energies can be reduced to two main scenarios. The first conception assumes the participation of the scalar $`S`$, tensor $`T`$ and pseudoscalar $`P`$ couplings in addition to the standard vector $`V`$ and axial $`A`$ couplings. S. C. Wu indicated explicitly that possibility. According to her, both left-handed $`(V,A)_L`$ couplings and exotic right-handed $`(S,T,P)_R`$ couplings may be responsible for the negative electron helicity observed in $`\beta `$-decay. K. Mursula et al. analyzed all the available data on the charged leptonic weak interactions while testing different models which admit the participation of additional $`(S,T,P)`$ couplings beside the standard $`(V,A)`$ couplings.
The other conception assumes that the right-handed vector $`V_R`$ and axial $`A_R`$ couplings participate in weak interactions beside left-handed standard $`(V,A)_L`$ couplings. This scenario is studied and analyzed in this work. There are many theoretical and experimental papers devoted to this problem. The models with $`SU(2)_L\times SU(2)_R\times U(1)`$ as the gauge group emerged first in the framework of a class of grand unified theories (GUT) . The manifest left-right symmetry model predicts the existence of the additional heavy vector bosons of the masses much larger than the masses of the bosons of the Standard Model. Bรฉg considered the bounds on the admixtures of the right-handed currents obtained from the measurements of the lepton polarization in semileptonic decays and by the determination of the parameters characterizing the spectrum in muon decay. Herczeg model is the generalization of the manifest left-right symmetric model, in which the fermions couple to distinct charged gauge-boson fields $`W_L`$ and $`W_R`$ with the different coupling constants $`g_L`$ and $`g_R`$, respectively . The effective hamiltonian has the structure of the four-fermion point interaction, for both muon decay and semileptonic processes. There are many experimental constraints on the possible mass of the right-handed vector bosons $`W_R`$ obtained from weak interaction processes at low energy and from high energy collider experiments. The lower mass limits for the $`W_R`$ received at the Tevatron collider are $`M_R652GeV(95\%CL`$, CDF-collaboration ) and $`M_R720GeV(95\%CL`$, D0-collaboration ), respectively. The lower bound obtained from $`K_0\overline{K_0}`$ mass difference is $`M_R1.6TeV`$ . A. Jodidio et al. measured the positron spectrum from muon decay, which allowed them to give the lower bound of $`432GeV(90\%CL)`$ on the possible mass of a new vector boson. J. Maalampi et al. explored the structure of the charged leptonic weak currents in the framework of the $`SU(2)_L\times SU(2)_R\times U(1)`$ models. They fitted the parameters of this model to experimental results obtained from pseudoscalar mezon decay, muon decay, nuclear $`\beta `$-decay and inverse muon decay. It allowed them to determine the values of the mass ratio of the charged gauge bosons and the mixing angle. N. B. Shulโgina introduced the admixtures of the right-handed $`(V,A)_R`$ currents into the interaction lagrangian, which helped to explain, e. g., the neutron paradox. Recent measurements of the longitudinal polarization of the positrons emitted by the polarized $`{}_{}{}^{107}In`$ and $`{}_{}{}^{12}N`$ nuclei gave the lower limit of $`306GeV(90\%CL)`$ on the mass of the right-handed gauge boson . The recommended lower constraint on the mass of the additional gauge boson is $`M_R>549GeV`$ . M. Zraลek et al. considered the possibility of the existence of neutrino magnetic moments in the framework of the left-right symmetry models. That could be especially interesting in the context of the solar neutrino deficit. However, all these limits are model-dependent and they can be considerably weakened. The stringent bound $`M_R1.6TeV`$ can be relaxed to the $`300GeV`$ range if one assumes that the Cabbibo-Kobayashi-Maskawa matrix elements for the right-handed quarks and for the left-handed quarks are not identical, and also that the $`SU(2)_{L,R}`$ gauge coupling constants are distinct, $`g_Lg_R`$. Therefore one should give the bounds for the nonstandard $`(V,A)_R`$ couplings without model assumptions . There are the present limits on all possible coupling constants obtained from normal muon decay and inverse muon decay .
However, to verify uniquely a scenario admitting the possible participation of the right-handed $`(V,A)_R`$ couplings beside the standard $`(V,A)_L`$ couplings in weak interactions at low energies, one proposes neutrino observables in the process of $`\mu `$-capture by proton. Both $`T`$-odd and $`T`$-even transverse components of the neutrino polarization are taken into account. Only in the quantities of this type the interference terms between the standard $`(V,A)_L`$ and nonstandard $`(V,A)_R`$ couplings appear and they do not depend explicite on the neutrino mass. It is analogical to the situation analysed by T.D. Lee and C.N. Yang in $`\beta `$-decay . They proposed an observable which was pseudoscalar under space inversion to determine uniquely if parity is violated. Only in this quantity interferences between couplings for the parity-conserving interactions $`C`$ and the parity-nonconserving interactions $`C^{}`$ can appear. Using the current data from $`\mu `$-decay and inverse $`\mu `$-decay, the magnitude of effects coming from the transverse components of the neutrino polarization can be determined. At the end, we can derive the lower limit on the mass of the right-handed gauge bosons (assuming, for example, manifest left-right symmetry).
Recently J. Sromicki measured the $`CP`$-odd transverse electron polarization in $`{}_{}{}^{8}Li`$ $`\beta `$-decay. The final results indicated the compatibility with the Standard Model prediction and $`CP`$-conservation in $`\beta `$-decay. Armbruster et al. measured the energy spectrum of electron neutrinos $`\nu _e`$ from $`\mu `$-decay at rest in the KARMEN experiment using the reaction $`{}_{}{}^{12}C(\nu _e,e^{})^{12}N_{g.s.}`$. They determined the upper limit of $`|g_{RL}^S+2g_{RL}^T|0.78(90\%CL)`$ on the possible interference term between scalar $`S`$ and tensor $`T`$ couplings. M. Abe et al. searched for T-odd transverse components of the muon polarization in $`K^+\pi ^0+\mu ^++\nu _\mu `$ decay at rest. They pointed out that the contribution to this observable from the Standard Model is of the order of $`10^7`$, so nonzero values of this quantity would indicate the beginning of new physics beyond the Standard Model. In our case, there is no contribution to the transverse neutrino polarization from the Standard Model (massless Dirac neutrinos), so nonzero values of such observable would be the unique proof of the participation of the $`(V,A)_R`$ couplings and of the production of the right-handed neutrinos. These last experiments made at high precision show that the problem of the completeness of the Lorentz structure and of the handedness structure of weak interactions is still explored.
The purpose of this paper is motivated by the desire to test how right-handed vector $`V_R`$ and axial $`A_R`$ couplings with the participation of left-handed standard $`(V,A)_L`$ couplings enter different observables such as: longitudinal and transverse neutron polarization, longitudinal neutrino polarization and, most of all, transverse neutrino polarization.
The structure of the work is as follows: Sect. 2 concentrates on the qualitative description of muon capture and on the assumptions concerning the calculations. In Sect. 3 the results obtained for the transverse, longitudinal neutron polarization and longitudinal neutrino polarization, among others, are presented. In Sect. 4 the results for the transverse neutrino polarization are dealt with. Sect. 5 gives the conclusions. In these considerations the system of natural units with $`\mathrm{}=c=1`$, Dirac hermitian matrices $`\gamma _\lambda `$ and the four-plus metric are used .
## 2 Muon capture by proton
The research is based on the reaction of the muon capture by proton $`\mu ^{}+pn+\nu _\mu `$. In the Standard Model it is a coherent low-energy process at the lepton-quark level. The typical energy transfer is of the order of $`1MeV`$ and therefore the space-time area within the interactions coincides with the size of the muonic atom, because of that both hadrons and leptons participating in this process are point objects. In the light of the above muon capture is considered at the level of the Fermi theory, whose hamiltonian describes the local, derivative-free, lepton-number-conserving, four-fermion point (contact) interaction. Right-handed $`(V,A)_R`$ couplings are assumed to take part in muon capture in addition to left-handed standard $`(V,A)_L`$ couplings. The coupling constants are denoted as $`C_V^L`$, $`C_A^L`$ and $`C_V^R`$, $`C_A^R`$ respectively to the neutrino handedness.
$`H_\mu ^{}`$ $`=`$ $`C_V^L(\overline{\mathrm{\Psi }}_\nu \gamma _\lambda (1+\gamma _5)\mathrm{\Psi }_\mu )(\overline{\mathrm{\Psi }}_n\gamma _\lambda \mathrm{\Psi }_p)`$
$`+C_A^L(\overline{\mathrm{\Psi }}_\nu i\gamma _5\gamma _\lambda (1+\gamma _5)\mathrm{\Psi }_\mu )(\overline{\mathrm{\Psi }}_ni\gamma _5\gamma _\lambda \mathrm{\Psi }_p)`$
$`+C_V^R(\overline{\mathrm{\Psi }}_\nu \gamma _\lambda (1\gamma _5)\mathrm{\Psi }_\mu )(\overline{\mathrm{\Psi }}_n\gamma _\lambda \mathrm{\Psi }_p)`$
$`+C_A^R(\overline{\mathrm{\Psi }}_\nu i\gamma _5\gamma _\lambda (1\gamma _5)\mathrm{\Psi }_\mu )(\overline{\mathrm{\Psi }}_ni\gamma _5\gamma _\lambda \mathrm{\Psi }_p)`$
where $`\mathrm{\Psi }_\mu ,\mathrm{\Psi }_\nu ,\mathrm{\Psi }_p,\mathrm{\Psi }_n`$ \- Dirac bispinors for the muon, muonic neutrino, proton and neutron. The above hamiltonian can be derived from the one by T. D. Lee and C. N. Yang , when the following expression is put $`C_V+C_V^{}\gamma _5=C_V^L(1+\gamma _5)+C_V^R(1\gamma _5),C_A+C_A^{}\gamma _5=C_A^L(1+\gamma _5)+C_A^R(1\gamma _5)`$, where $`C_V^L=(C_V+C_V^{})/2,C_V^R=(C_VC_V^{})/2,C_A^L=(C_A+C_A^{})/2,C_A^R=(C_AC_A^{})/2,C_V,C_A,C_V^{},C_A^{}`$ \- the vector and axial couplings for the parity-conserving interactions and the parity-nonconserving interactions. The remaining couplings are omitted $`C_S=C_S^{}=C_P=C_P^{}=C_T=C_T^{}=0`$. The Fermi hamiltonian, Eq. (2), can be modified if one puts $`C_V^LC_A^L\gamma _5=(C_V^LC_A^L)(1+\gamma _5)/2+(C_V^L+C_A^L)(1\gamma _5)/2,C_V^R+C_A^R\gamma _5=(C_V^R+C_A^R)(1+\gamma _5)/2+(C_V^RC_A^R)(1\gamma _5)/2`$, and then, one obtains the effective hamiltonian of the form:
$`H_\mu ^{}`$ $`=`$ $`{\displaystyle \frac{C_V^LC_A^L}{2}}(\overline{\mathrm{\Psi }}_\nu \gamma _\lambda (1+\gamma _5)\mathrm{\Psi }_\mu )(\overline{\mathrm{\Psi }}_n\gamma _\lambda (1+\gamma _5)\mathrm{\Psi }_p)`$
$`+{\displaystyle \frac{C_V^L+C_A^L}{2}}(\overline{\mathrm{\Psi }}_\nu \gamma _\lambda (1+\gamma _5)\mathrm{\Psi }_\mu )(\overline{\mathrm{\Psi }}_n\gamma _\lambda (1\gamma _5)\mathrm{\Psi }_p)`$
$`+{\displaystyle \frac{C_V^R+C_A^R}{2}}(\overline{\mathrm{\Psi }}_\nu \gamma _\lambda (1\gamma _5)\mathrm{\Psi }_\mu )(\overline{\mathrm{\Psi }}_n\gamma _\lambda (1+\gamma _5)\mathrm{\Psi }_p)`$
$`+{\displaystyle \frac{C_V^RC_A^R}{2}}(\overline{\mathrm{\Psi }}_\nu \gamma _\lambda (1\gamma _5)\mathrm{\Psi }_\mu )(\overline{\mathrm{\Psi }}_n\gamma _\lambda (1\gamma _5)\mathrm{\Psi }_p)`$
Muonic neutrinos are assumed to be massive and of Dirac nature. The nonrelativistic approximation is applied both for the nucleons and the muon in the $`1s`$ state. To describe muon capture the following observables are used: $`\stackrel{}{P}_\mu `$ \- the initial muon polarization in the $`1s`$ state, $`\stackrel{}{S}_\nu `$ \- the operator of the neutrino spin, $`\stackrel{}{q}`$ \- the momentum of the outgoing neutrino, $`\stackrel{}{J}_n`$ \- the operator of the neutron spin. $`\stackrel{}{P}_\mu `$ and $`\stackrel{}{q}`$ are assumed to be perpendicular to each other. To illustrate the calculation method the definition of the $`T`$-even components of the transverse neutrino polarization is given $`<\stackrel{}{S}_\nu \widehat{\stackrel{}{P}}_\mu >_fTr[\stackrel{}{S}_\nu \widehat{\stackrel{}{P}}_\mu \rho _f]`$, where $`\rho _f`$ \- the density operator of the final state (neutrino-neutron), $`\widehat{\stackrel{}{P}}_\mu `$ \- the direction of the muon polarization in the $`1s`$ state. The calculations are made with the Reduce computer program.
## 3 Longitudinal and transverse neutron non-polarization and longitudinal neutrino non-polarization
In this section the results for the longitudinal and transverse neutron polarization and for the longitudinal neutrino polarization are presented. The final results are as follows:
$`<\widehat{\stackrel{}{q}}\stackrel{}{J}_n>_f`$ $`=`$ $`{\displaystyle \frac{\varphi _\mu (0)^2}{4\pi }}\{{\displaystyle \frac{q}{E}}(C_A^L^2C_A^R^2)`$
$`+Re[(2{\displaystyle \frac{q}{M}}+{\displaystyle \frac{q}{E}})(C_V^RC_A^RC_V^LC_A^L)`$
$`+{\displaystyle \frac{m_\nu }{E}}{\displaystyle \frac{q}{M}}(C_V^RC_A^LC_V^LC_A^R)]\}`$
$`<\stackrel{}{J}_n(\widehat{\stackrel{}{P}_\mu }\times \widehat{\stackrel{}{q}})>_f`$ $`=`$ $`{\displaystyle \frac{\varphi _\mu (0)^2}{4\pi }}|\stackrel{}{P_\mu }|Im\{({\displaystyle \frac{q}{E}}+{\displaystyle \frac{q}{M}})(C_V^RC_A^RC_V^LC_A^L)`$
$`+{\displaystyle \frac{m_\nu }{E}}{\displaystyle \frac{q}{M}}(C_V^RC_A^LC_V^LC_A^RC_A^LC_A^RC_V^LC_V^R)\}`$
$`<\stackrel{}{S}_\nu \widehat{\stackrel{}{q}}>_f`$ $`=`$ $`{\displaystyle \frac{\varphi _\mu (0)^2}{4\pi }}\{Re[{\displaystyle \frac{q}{M}}(C_V^LC_A^LC_V^RC_A^R)`$
$`+{\displaystyle \frac{m_\nu }{E}}{\displaystyle \frac{q}{M}}(C_V^LC_A^RC_V^RC_A^L)]`$
$`+({\displaystyle \frac{1}{2}}{\displaystyle \frac{q}{E}}+{\displaystyle \frac{q}{2M}})(C_V^R^2C_V^L^2)`$
$`+({\displaystyle \frac{3}{2}}{\displaystyle \frac{q}{E}}+{\displaystyle \frac{q}{2M}})(C_A^R^2C_A^L^2)\}`$
where, $`\varphi _\mu (0)`$ \- the value of the large radial component of the muon Dirac bispinor for $`r=0`$, $`\widehat{\stackrel{}{q}}`$ \- the direction of the neutrino momentum, $`|\stackrel{}{P_\mu }|`$ \- the value of the muon polarization in the $`1s`$ state, $`q/2M`$ \- the momentum corrections, $`q,E,m_\nu ,M`$ \- the value of the neutrino momentum, its energy, its mass and the nucleon mass, respectively.
It can be noticed that in these observables the occurrence of the interferences between the right-handed $`(V,A)_R`$ and left-handed $`(V,A)_L`$ couplings depends explicite on the muonic neutrino mass. Thus, the so-called โconspiracyโ of interference terms appears here. This โconspiracyโ makes the measurement of the relative phase between these two coupling types impossible because a very small mass of the neutrino $`(m_\nu <0.17MeVCL=90\%`$ ) at its high energy $`(E_\nu 100MeV)`$ suppresses such an interference in practice. The additional interference attenuation is further caused by the momentum corrections. The factor $`(m_\nu /E)(q/M)1710^5`$ is very small. We can see that new interference contributions can not be detected at the present level of experimental precision. Longitudinal neutrino polarization behaves as a typical nuclear quantity. The โconspiracyโ of interference terms caused by the neutrino mass occurs here. The measurement of this observable would not allow the unique determination of the possible participation of the $`(V,A)_R`$ couplings. Therefore, the observables in which such difficulties do not appear are proposed.
## 4 Why transverse components of neutrino polarization?
In this section the results for two $`T`$-odd neutrino observables and for one $`T`$-even neutrino quantity are given. All these three cases concern the transverse neutrino polarization. In practice, that would mean the measurements of the components of the neutrino polarization perpendicular to its direction of momentum. The final results are as follows:
$`<\stackrel{}{S}_\nu (\widehat{\stackrel{}{P}}_\mu \times \widehat{\stackrel{}{q}})>_f`$
$`=`$ $`{\displaystyle \frac{\varphi _\mu (0)^2}{4\pi }}|\stackrel{}{P_\mu }|Im\{({\displaystyle \frac{q}{E}}+{\displaystyle \frac{q}{M}})(C_A^LC_A^RC_V^LC_V^R)\}`$
$`<\stackrel{}{S}_\nu (\stackrel{}{J}_n\times \widehat{\stackrel{}{q}})>_f`$ $`=`$ $`{\displaystyle \frac{\varphi _\mu (0)^2}{4\pi }}Im\{({\displaystyle \frac{q}{E}}+{\displaystyle \frac{q}{M}})(C_V^RC_A^LC_V^LC_A^R)`$
$`+{\displaystyle \frac{q}{M}}C_V^LC_V^R+(2{\displaystyle \frac{q}{E}}+{\displaystyle \frac{q}{M}})C_A^LC_A^R`$
$`+{\displaystyle \frac{m_\nu }{E}}{\displaystyle \frac{q}{M}}(C_V^RC_A^RC_V^LC_A^L)\}`$
$`<\stackrel{}{S}_\nu \widehat{\stackrel{}{P}}_\mu >_f`$ $`=`$ $`{\displaystyle \frac{\varphi _\mu (0)^2}{4\pi }}|\stackrel{}{P_\mu }|Re\{(1+{\displaystyle \frac{q}{E}}{\displaystyle \frac{q}{M}})(C_V^LC_V^RC_A^LC_A^R)`$
$`+{\displaystyle \frac{1}{2}}{\displaystyle \frac{m_\nu }{E}}(C_V^R^2C_A^R^2+C_V^L^2C_A^L^2)\}`$
The obtained result, Eq. (4), consists exclusively of the interference terms between the left-handed standard $`(V,A)_L`$ couplings and the right-handed nonstandard $`(V,A)_R`$ couplings, whose occurrence does not depend explicite on the muonic neutrino mass. It can be understood as the interference between the neutrino waves of negative and positive chirality. It creates the possibility of measuring the relative phase between the two coupling types. In the next observables, Eq. (4) and (4), we have the additional dependence on the neutrino mass which occurs only at the terms of the type: $`C_V^R^2,C_A^R^2,C_V^L^2,C_A^L^2,C_V^RC_A^R,C_V^LC_A^L`$. It gives a very small contribution in the relation to the main one coming from the interferences between the $`(V,A)_L`$ and $`(V,A)_R`$ couplings. There is no contribution to these observables from the Standard Model, in which neutrinos are massless.
Now, we will express our coupling constants $`C_{V,A}^{L,R}`$ by Fetscherโs couplings $`g_{ฯต\mu }^\gamma `$ assuming the universality of weak interactions. The induced couplings generated by the dressing of hadrons are neglected as their presence does not change qualitatively the conclusions about transverse neutrino polarization. Here, $`\gamma =S,V,T`$ indicates a scalar, vector, tensor interaction; $`ฯต,\mu =L,R`$ indicate the chirality of the electron or muon and the neutrino chiralities are uniquely determined for given $`\gamma ,ฯต,\mu `$. We get the following relations: $`C_V^L=A(g_{LL}^V+g_{RL}^V),C_A^L=A(g_{LL}^Vg_{RL}^V),C_V^R=A(g_{LR}^V+g_{RR}^V),C_A^R=A(g_{LR}^Vg_{RR}^V)`$, where $`A(4G_F/\sqrt{2})cos\theta _c,G_F=1.16639(1)\times 10^5GeV^2`$ is the Fermi coupling constant, $`\theta _c`$ is the Cabbibo angle ($`cos\theta _c=0.9740\pm 0.0010`$ ). We can derive the contributions coming from the $`C_{V,A}^{L,R}`$ coupling constants in $`\mu `$-capture, using the current data : $`|C_V^L|>0.850A,|C_A^L|>1.070A,|C_V^R|<0.093A,|C_A^R|<0.027A`$. From the above, we can see that the effects of the transverse neutrino polarization connected with the right-handed $`(V,A)_R`$ couplings may be of the order of $`7\%`$. This value is model-independent. In this way, we can give the lower limit of $`M_R305GeV`$ on the mass of the right-handed vector boson $`W_R`$ (for manifest left-right symmetry, $`m_\nu =0`$, without $`W_LW_R`$ mixing). It is compatible with the current bounds on the mass of the $`W_R`$ received from the weak interaction processes at low energy .
When the neutrinos are massive and only standard $`(V,A)_L`$ couplings participate in muon capture, the transverse components of the neutrino polarization could be observed in both cases:
$$<\stackrel{}{S}_\nu (\stackrel{}{J}_n\times \widehat{\stackrel{}{q}})>_f=\frac{\varphi _\mu (0)^2}{4\pi }\frac{m_\nu }{E}\frac{q}{M}Im(C_V^LC_A^L)$$
(9)
$`<\stackrel{}{S}_\nu \widehat{\stackrel{}{P}}_\mu >_f`$ $`=`$ $`{\displaystyle \frac{\varphi _\mu (0)^2}{8\pi }}|\stackrel{}{P_\mu }|{\displaystyle \frac{m_\nu }{E}}(C_V^L^2C_A^L^2)`$ (10)
However, we can see that the eventual effect of the nonzero transverse components of the neutrino polarization connected with the neutrino mass would be much weaker than the one coming from the nonstandard couplings. In such a scenario, the observable determined by the Eq. (4) always equals zero, $`<\stackrel{}{S}_\nu (\widehat{\stackrel{}{P}}_\mu \times \widehat{\stackrel{}{q}})>_f=0`$. Because the direct measurement of the transverse neutrino polarization in $`\mu ^{}`$-capture by proton is very difficult now, one proposes to use the muonic neutrino-electron elastic scattering $`\nu _\mu +e^{}e^{}+\nu _\mu `$ to detect the effects of the transverse neutrino polarization connected with the nonstandard couplings. From the differential cross section for this process the possible neutrino-electron correlations in terms of the longitudinal and transverse neutrino polarization will be seen (in preparation). Observing the change in the distribution of the electrons in relation to the distribution with standard neutrino, one would obtain a clear signal of the nonzero values of the transverse neutrino polarization. Maalampi et al. considered the presence of muonic neutrinos of a given initial polarization in the inverse muon decay. In the future the experimental verification of the hypothesis concerning the transverse components of the neutrino polarization could be carried out by the Fermi laboratory. Currently at Fermilab, the BooNE experiment (The Booster Neutrino Experiment) with the intense neutrino source is designed to search for the muonic neutrino oscillations, the mass difference, the mixing angle, the CP violation in the lepton sector, the muon-neutrino disappearance signal, the neutrino magnetic moment, and the helicity structure of the weak neutral current. This experiment will also look for the non-oscillation neutrino physics using, among others, the neutrino-electron elastic scattering. At Fermilab, essentially all of muon-neutrinos come from $`\pi ^+`$-decay. However, on the quark level, in $`\mu ^{}`$-capture and $`\pi ^+`$-decay, there is the same semileptonic interaction. Therefore the conclusions regarding transverse neutrino polarization in $`\mu ^{}`$-capture are also correct for the muon-neutrinos coming from $`\pi ^+`$-decay.
## 5 Conclusions
The measurements of the $`T`$-odd, Eq. (4), transverse components of the neutrino polarization could verify the possibility of the right-handed $`(V,A)_R`$ couplings participation in weak interactions. They would also be the proof of the production of the right-handed neutrinos. As far as the $`T`$-odd components of the transverse neutrino polarization are concerned, one would also obtain a proof of the symmetry breaking under time inversion $`T(CP)`$ in a semileptonic process. The measurement of longitudinal neutrino polarization does not offer such possibilities because of the suppressing of interferences between the $`(V,A)_L`$ and $`(V,A)_R`$ couplings caused by the neutrino mass. In this way, that will always lead to the compatibility with the Standard Model. The similar regularity can be observed in nuclear observables: longitudinal, transverse neutron polarization and also probability of muon capture, and the quantities of only this type are measured today.
The BooNE experiment, which is now being constructed, will be able to measure the nonstandard neutrino-electron correlations using the $`\nu _\mu +e^{}e^{}+\nu _\mu `$ process. This experiment will be started in the year 2001.
I am greatly indebted to Prof. S. Ciechanowicz for many useful and helpful discussions and his interest in my research. I owe much to Prof. M. Zraลek for interesting critical remarks. |
warning/0002/cond-mat0002150.html | ar5iv | text | # Semiflexible polymer on an anisotropic Bethe lattice
## I INTRODUCTION
Chain polymers are often approximated as self- and mutually avoiding walks (SAWโs) on a lattice, and much information about the behavior of polymers both in a melt or in solution has been understood theoretically through this model . One of the characterizations of the conformations of a walk is through its mean-square end-to-end distance $`R^2`$, where the mean is taken over all configurations of the $`N`$-step walk on the lattice. In the limit $`N\mathrm{}`$ a scaling behavior $`R^2N^{2\nu }`$ is observed, where the exponent $`\nu `$ exhibits universal behavior, with a mean-field value $`\nu =1/2`$ for random walks, which correspond to ideal chains, and $`\nu =3/4`$ for SAWโs on two-dimensional lattices , for example.
An interesting question arises if the chains are not considered to be totally flexible, an energy being associated to bends of the chain. This is often observed for real polymers. Let us, for simplicity, restrict ourselves to SAWโS on hypercubic lattices. In this case, consecutive steps of the walk are either in the same direction or perpendicular. So, a Boltzmann factor $`z`$ may be associated to each pair of perpendicular consecutive steps of walk. This problem of semiflexible polymers (also called persistent or biased walks) has been studied for some time, and there occurs a crossover in the behavior of the walk between a rodlike behavior $`\nu _r=1`$ for $`z=0`$, where the polymer is totally stiff, and the usual behavior with a different exponent $`\nu `$ for nonzero values of $`z`$. Stating this point more precisely, the mean square end-to-end distance displays scaling behavior in the limit $`N\mathrm{};z0;N^\psi z=constant`$, which is given by
$$R^2N^{2\nu _r}F(zN^\psi ),$$
(1)
the observed values being $`\nu _r=1`$ and $`\psi =1`$.
The scaling function has the behavior $`F(x)x^{(2\nu 2)/\psi }`$ in the limit $`x\mathrm{}`$. This scaling form has been verified through several techniques, although in three dimensions a mean field exponent $`\nu =1/2`$ was found for intermediate values of the number of steps $`N`$, the crossover to the three-dimensional value occurring at rather high values of $`N`$.
In this paper, we consider the problem of a semiflexible polymer on a Bethe lattice, calculating exactly the mean square end-to-end distance of walks on the Cayley tree which start at the central site and have $`N`$ steps, supposing that the walks will never reach the surface of the Cayley tree, thus remaining in its core. We calculate also the mean square end-to-end distance in the case when the lattice is considered anisotropic, that is, when the edges of the lattice are not equivalent with respect to their occupation by a polymer bond. The definition of the distance between two sites of the Cayley tree is not obvious, and some possibilities exploring the fact that the tree may be embedded in a hypersurface of a non-Euclidean space have been given. In this paper, however, we used a simpler definition, considering the Cayley tree in the thermodynamic limit to be embedded in an infinite-dimensional Euclidean space. The result for $`R^2(N,z)`$, for the isotropic case has the scaling form of Eq. 1. Not surprisingly the scaling function $`F(x)`$ is equal to the one obtained for random walks with no immediate return on an hypercubic lattice with the same coordination number of the Bethe lattice considered. This might be expected since Bethe lattice calculations lead to mean-field critical exponents. Also, in the limit $`N\mathrm{}`$ for nonzero values of $`z`$ the scaling behavior $`R^2N^{2\nu }`$ with the classical value $`\nu =1/2`$ is verified in the expression for $`R^2(N,z)`$. It should be mentioned that our proposal of defining the Euclidean distance between two points of the Cayley tree is similar to earlier results in the literature relating this distance to the chemical distance, measured along the chain . However, the distinction between the chemical and the Euclidean distances is not always properly considered in the literature, and this may lead to contradicting results , as we will discuss in more detail in the conclusion.
In section II we define the model and calculate the mean square end-to-end distance recursively on the anisotropic Bethe lattice. The problem is then reduced to finding the general term of a linear mapping in six dimensions. In the particular case of an isotropic lattice, we find a closed expression for $`R^2`$. In section III the asymptotic behavior is studied for the general case, based on the mapping. In section IV final comments and discussions may be found. Finally, we present in the appendices a combinatorial calculation for $`R^2`$ in the isotropic case.
## II DEFINITION OF THE MODEL AND SOLUTION FOR THE ISOTROPIC LATTICE
We consider a Cayley tree of coordination number $`q`$ and place a chain on the tree starting at the central site. Each bond of the tree is supposed to be of unit length. Figure 1 shows a tree with $`q=4`$ and a polymer with $`N=2`$ steps placed on it. Since we want the Cayley tree to be an approximation of a hypercubic lattice in $`d`$ dimensions, we will restrict ourselves to even coordination numbers $`q=2d`$. As in the hypercubic lattice, the bonds incident on any site of the tree are in $`d`$ directions, orthogonal to each other. As may be seen in Figure 1, the central site of the tree is connected to $`q`$ other sites, which belong to the first generation of sites. Each of the sites of the first generation is connected to $`(q1)`$ sites of the second generation, and this process continues until the surface of the tree is reached, after a number of steps equal to the number of generations in the tree. Upon reaching a site of the $`i`$โth generation coming from a site belonging to generation $`(i1)`$, there are $`(q1)`$ possibilities for the next step of the walk towards a site of generation $`(i+1)`$. One of them will be in the same direction as the previous step, while the remaining $`(q2)`$ will be in directions orthogonal to all previous steps. In the second case, a statistical weight $`z`$ is associated to the elementary bend in the walk. Therefore, we admit that the $`(q2)`$ bonds which are orthogonal to the last step are also orthogonal to all bonds of the lattice in earlier generations. Let us stress two consequences of this supposition: (i) A tree of coordination number $`q`$ with $`N_g`$ generations will be embedded in a space of dimension
$$D=q/2+(N_g1)(q/21).$$
(2)
The sites of the Cayley tree will all be sites of a hypercubical lattice in $`D`$ dimensions. This may be seen in Figure 1 where the sites of a tree with $`q=4`$ and $`N_g=2`$ are sites of a cubic $`(D=3)`$ lattice. (ii) By construction, there will never be loops in the tree, a property which is true for any Cayley tree. It is well known that it may be shown by other means that the Cayley tree is a infinite-dimensional lattice in the thermodynamic limit $`N_g\mathrm{}`$. Finally, the anisotropy is introduced into the model considering that bonds of the chain in $`s`$ of the $`q/2`$ directions (we will call them special) at each lattice site contribute with a factor $`y`$ to the partition function, while no additional contribution comes from bonds in the remaining $`t=q/2s`$ directions at each lattice site.
Usually , the calculation of thermodynamic properties of models defined on the Bethe lattice is done in a recursive manner, so we will follow a similar procedure in the calculation of the mean square end-to-end distance. We define a generalized partition (or generating) function for $`N`$-step chains
$$g_N=z^my^{N_e}p^{R^2},$$
(3)
where the sum is over all configurations of the chain, $`z`$ is the statistical weight of an elementary bend in the chain, $`y`$ is the statistical weight of bonds in special directions and $`p`$ is a parameter associated to the square of the end-to-end distance of the chain. At the end of the calculation, we will take $`p=1`$. The numbers of elementary bends, bonds in special directions, and the square end-to-end distance of each chain are $`m`$, $`N_e`$, and $`R^2`$, respectively. The mean square end-to-end distance may then be calculated through
$$R^2_N=\frac{1}{g_N}\left(p\frac{g_N}{p}\right)|_{p=1}.$$
(4)
The partition function may then be calculated in a recursive way if we define partial partition functions $`a_N^l`$ and $`b_N^l`$ such that the first ones include all $`N`$-bond chains whose last $`l`$ bonds are collinear and in one of the special directions (there is necessarily a bend before the $`l`$ bonds, if $`l<N`$), while the last $`l`$ bonds of the chains contributing to $`b_N^l`$ are collinear and in one of the non-special directions. The partition function may then be written as
$$g_N=\underset{l=1}{\overset{N}{}}(a_N^l+b_N^l).$$
(5)
Due to the fact that there are no closed loops on the Cayley tree, it is quite easy to write down recursion relations for the partial partition functions
$`a_{N+1}^1`$ $`=`$ $`2syzp{\displaystyle \underset{l=1}{\overset{N}{}}}b_N^l+2(s1)yzp{\displaystyle \underset{l=1}{\overset{N}{}}}a_N^l,`$ (7)
$`a_{N+1}^{l+1}`$ $`=`$ $`yp^{2l+1}a_N^l,`$ (8)
$`b_{N+1}^1`$ $`=`$ $`2tzp{\displaystyle \underset{l=1}{\overset{N}{}}}a_N^l+2(t1)zp{\displaystyle \underset{l=1}{\overset{N}{}}}b_N^l,`$ (9)
$`b_{N+1}^l`$ $`=`$ $`p^{2l+1}b_N^l,`$ (10)
with the initial conditions
$`a_1^1`$ $`=`$ $`2syp,`$ (12)
$`b_1^1`$ $`=`$ $`2tp.`$ (13)
For example, in the first expression above, the new bond may be preceded by a bond in a special direction, with $`2s`$ possibilities, or by a bond in a non-special direction, with $`2(s1)`$ possibilities. In both cases, a factor $`p`$ is present since the bond added is in a direction perpendicular to all previous ones, and thus $`R^2`$ is increased by one unit. Finally, the inclusion of the new bond introduces one bend in the chain, thus explaining the factor $`z`$, and since the bond is in a special direction the factor $`y`$ is justified. In the second expression, it should be mentioned that $`R^2`$ is increased by $`(l+1)^2l^2`$, thus explaining the exponent of $`p`$. If we now define
$`a_N={\displaystyle \underset{l=1}{\overset{N}{}}}a_N^l,`$ (15)
$`b_N={\displaystyle \underset{l=1}{\overset{N}{}}}b_N^l,`$ (16)
the mean square end-to-end distance will be
$$R^2_N=\frac{1}{a_N+b_N}\left[p\frac{}{p}(a_N+b_N)\right]|_{p=1}=\frac{c_N+d_N}{a_N+b_N}.$$
(17)
The recursion relations for $`a_N`$ and $`b_N`$, as well as the ones for the new variables $`c_N`$ and $`d_N`$ may be written, for $`p=1`$, as
$`a_{N+1}`$ $`=`$ $`2z{\displaystyle \underset{l=0}{\overset{N}{}}}y^{l+1}\left[(s1)a_{Nl}+sb_{Nl}\right],`$ (19)
$`b_{N+1}`$ $`=`$ $`2z{\displaystyle \underset{l=0}{\overset{N}{}}}\left[ta_{Nl}+sb_{Nl}\right],`$ (20)
$`c_{N+1}`$ $`=`$ $`2z{\displaystyle \underset{l=0}{\overset{N}{}}}y^{l+1}\left[(s1)(l+1)^2a_{Nl}+s(l+1)^2b_{Nl}+(s1)c_{Nl}+sd_{Nl}\right],`$ (21)
$`d_{N+1}`$ $`=`$ $`2z{\displaystyle \underset{l=0}{\overset{N}{}}}\left[t(l+1)^2a_{Nl}+(t1)(l+1)^2b_{Nl}+tc_{Nl}+(t1)d_{Nl}\right],`$ (22)
with the initial conditions
$`a_0`$ $`=`$ $`{\displaystyle \frac{2s}{z(q2)}},`$ (24)
$`b_0`$ $`=`$ $`{\displaystyle \frac{2t}{z(q2)}},`$ (25)
$`c_0`$ $`=`$ $`d_0=0`$ (26)
An undesirable feature of the recursion relations Eqs. 17 is that the new values of the iterating variables depend on all previous values. This dependence, however, is rather simple, and it is possible, introducing two more variables $`e_N`$ and $`f_N`$, to rewrite the recursion relations as a mapping involving only one previous value of each variable, valid for $`N1`$
$`a_{N+1}`$ $`=`$ $`ya_N+2zy\left[(s1)a_N+sb_N\right],`$ (28)
$`b_{N+1}`$ $`=`$ $`b_N+2z\left[ta_N+(t1)b_N\right],`$ (29)
$`c_{N+1}`$ $`=`$ $`yc_N+2zy\left[(s1)(a_N+c_N)+s(b_N+d_N)\right]+(2N+1)ya_N2ye_N,`$ (30)
$`d_{N+1}`$ $`=`$ $`d_N+2z\left[t(a_N+c_N)+(t1)(b_N+d_N)\right]+(2N+1)b_N2f_N,`$ (31)
$`e_{N+1}`$ $`=`$ $`ye_N+2Nzy\left[(s1)a_N+sb_N\right],`$ (32)
$`f_{N+1}`$ $`=`$ $`f_N+2Nz\left[ta_N+(t1)b_N\right],`$ (33)
with the initial conditions
$`a_1`$ $`=`$ $`c_1=2sy,`$ (35)
$`b_1`$ $`=`$ $`d_1=2t,`$ (36)
$`e_1`$ $`=`$ $`f_1=0.`$ (37)
The value for $`R^2`$ may be found iterating the mapping above through the Eq. 17. In principle, since the mapping is linear, it is solvable. One starts finding the general term of the first two equations, then solving the last two and finally solving the two remaining relations. A software for algebraic computing is helpful, but we realized that the general answer will be too large to be handled, and also the computer time and memory required are beyond the resources we have available. We therefore restrict ourselves to a complete solution of the isotropic case $`y=1`$ and to an exact study of the asymptotic properties of the solution for the general case. It is worthwhile to observe in the mapping Eqs. II that under a transformation
$`s^{}=t,`$ (39)
$`t^{}=s,`$ (40)
$`y^{}=1/y,`$ (41)
$`R^2`$ will be invariant, as expected.
For the isotropic case ($`y=1`$) the mapping Eqs. II is reduced to three variables
$`\alpha _N`$ $`=`$ $`a_N+b_N,`$ (43)
$`\beta _N`$ $`=`$ $`c_N+d_N,`$ (44)
$`\gamma _N`$ $`=`$ $`e_N+f_N,`$ (45)
and it may be written as
$`\alpha _{N+1}`$ $`=`$ $`\left[1+z(q2)\right]\alpha _N,`$ (47)
$`\beta _{N+1}`$ $`=`$ $`\left[1+z(q2)\right]\beta _N+\left[2N+1+z(q2)\right]\alpha _N2\gamma _N,`$ (48)
$`\gamma _{N+1}`$ $`=`$ $`\gamma _N+Nz(q2)\alpha _N.`$ (49)
The initial conditions are
$`\alpha _1`$ $`=`$ $`\beta _1=q,`$ (50)
$`\gamma _1`$ $`=`$ $`0.`$ (51)
It is easy to find the general solution for this mapping
$`\alpha _N`$ $`=`$ $`qk^{N1},`$ (53)
$`\beta _N`$ $`=`$ $`{\displaystyle \frac{q}{(k1)^2}}\left[N(k^21)k^{N1}+22k^N\right],`$ (54)
$`\gamma _N`$ $`=`$ $`{\displaystyle \frac{q}{k1}}\left[N(k^21)k^{N1}+1k^N\right],`$ (55)
where $`k=1+(q2)z`$. The substitution of this solution in Eq. 17 results in
$$R^2=\frac{2[1+a]}{a^2}\left[Na1+\frac{1}{[1+a]^N}\right]N,$$
(56)
where $`a=k1=(q2)z`$.
The properties of the mean square end-to-end distance Eq. 56 in some limiting cases show that our result has the expected behavior. First, we notice that when the bend statistical weight $`z`$ vanishes we have
$$lim_{z0}R^2=N^2,$$
(57)
for any number of steps $`N`$. This rodlike behavior is expected, since no bend will be present in the walk. In the opposite limit of infinite bending statistical weight the result is
$$lim_z\mathrm{}R^2=N,$$
(58)
which is also an expected result, since in this limit there is a bend at every internal site of the chain, so that, according to the definition of the end-to-end distance we are using, the vector $`\stackrel{}{R}`$ in this situation will have $`N`$ components, all of them being equal to 1.
In the limit of an infinite chain $`N\mathrm{}`$ we get, for nonzero $`z`$,
$$lim_N\mathrm{}R^2=\frac{(2+a)N}{a},$$
(59)
and we notice that the expected scaling behavior $`R^2N^{2\nu }`$ is obtained with the mean field exponent $`\nu =1/2`$. The asymptotic behavior of $`R^2`$ is different for zero and nonzero $`a`$, as may be appreciated comparing Eqs. 57 and 59 respectively. So we may look for the crossover between both behaviors in the limit of Eq. 1, getting the result
$$lim_{N\mathrm{};a0;aN=x}R^2=N^2F(x),$$
(60)
with a scaling function
$$F(x)=\frac{2(x1+\mathrm{exp}(x))}{x^2}.$$
(61)
It should be stressed that the square end-to-end distance given in Eq. 56 is the same obtained by adapting the general result of Flory for random walks without immediate return to hypercubic lattices. In general, it may be shown that the exact solution of statistical models with first neighbor interactions on the Bethe lattice is equivalent to the Bethe approximation on the Bravais lattice with the same coordination number . The random walk without immediate returns corresponds to the Bethe approximation of the $`n0`$ model associated to the self-avoiding walk problem , and here we show that the analogy may be extended to the mean square end-to-end distance if we define distances on the Bethe lattice as was done above. Although the results on the Bethe lattice as calculated here and the ones for ideal chains without immediate return on a hypercubic lattice with the same coordination number should have the same asymptotic behaviors, it is at first surprising that they are actually identical. However, it turns out that the mean value of the angle between successive bonds, as calculated by Flory in his original work , is actually exact for chains on the Bethe lattice as we considered.
## III ASYMPTOTIC BEHAVIOR IN THE GENERAL CASE
In this section we develop a study of the asymptotic solution of the mapping Eqs. II for $`N1`$. Let us reduce the dimension of the mapping by one defining new iteration variables
$`B_N={\displaystyle \frac{b_N}{a_N}},`$ (63)
$`C_N={\displaystyle \frac{c_N}{a_N}},`$ (64)
$`D_N={\displaystyle \frac{d_N}{a_N}},`$ (65)
$`E_N={\displaystyle \frac{e_N}{a_N}},`$ (66)
$`F_N={\displaystyle \frac{f_N}{a_N}}.`$ (67)
From Eqs. II and the initial conditions Eqs. II it is easy to write the recursion relations and initial conditions for the new iterative variables in the mapping above. In the limit of large values of $`N`$, for fixed $`z`$ and $`y`$, the following asymptotic behavior is observed
$`B_N`$ $``$ $`B^0,`$ (69)
$`C_N`$ $``$ $`C^0+C^1N,`$ (70)
$`D_N`$ $``$ $`D^0+D^1N,`$ (71)
$`E_N`$ $``$ $`E^0+E^1N,`$ (72)
$`F_N`$ $``$ $`F^0+F^1N.`$ (73)
The substitution of the these expressions into the recursion relations for the variables defined in Eqs. III, obtained from the general mapping Eqs. II, leads to the determination of the constants in the asymptotic behavior and thus we obtain
$$R^2=\frac{C_N+D_N}{1+B_N}\frac{C^1+D^1}{1+B^0}N=CN,$$
(74)
where the amplitude $`C=C^1`$ is given by
$$C=\frac{sy(B^0)^2\left[\frac{y(1+ฯต)+1}{y(1+ฯต)1}\right]+t\left[\frac{2+ฯต}{ฯต}\right]}{sy(B^0)^2+t},$$
(75)
where
$$ฯต=2z(s1+sB^0),$$
(76)
and $`B^0`$ is the positive root of
$$2zsy(B^0)^2+\left[y1+2z(syty+1)\right]B^02zt=0.$$
(77)
The amplitude of the asymptotic behavior of $`R^2`$ thus may be obtained exactly in the general case and, as may be seen in Figure 2, diverges as $`z0`$, as expected. Also, in the limit $`y\mathrm{}`$ the problem reduces to a walk on an isotropic lattice with coordination number equal to $`2s`$, and we get
$$C=\frac{2z(s1)+2}{2z(s1)},$$
(78)
which agrees with Eq. II for the isotropic case.
Now we will study the asymptotic behavior in the quasi-rigid limit $`N\mathrm{}`$, $`z0`$, and $`N(q2)z=x`$. We thus expand $`R^2`$ for small values of $`z`$
$$R^2(z,y,N)R^2(0,y,N)+\frac{R^2}{z}|_{z=0}z.$$
(79)
For $`z=0`$ the solution of the mapping Eqs. II is
$`a_N`$ $`=`$ $`2sy^N,`$ (81)
$`b_N`$ $`=`$ $`2t,`$ (82)
$`c_N`$ $`=`$ $`2sy^NN^2,`$ (83)
$`d_N`$ $`=`$ $`2tN^2,`$ (84)
$`e_N`$ $`=`$ $`f_N=0;`$ (85)
and we have $`R^2=N^2`$, as expected. From the mapping Eqs. II the recursion relations for the derivatives of the variables with respect to $`z`$ (at $`z=0`$) may be seen to be
$`a_{N+1}^{}`$ $`=`$ $`ya_N^{}+4ys\left[(s1)y^N+t\right],`$ (87)
$`b_{N+1}^{}`$ $`=`$ $`b_N^{}+4t(sy^N+t1),`$ (88)
$`c_{N+1}^{}`$ $`=`$ $`yc_N^{}+4ys(1+N^2)\left[(s1)y^N+t\right]+y(2N+1)a_N^{}2ye_N^{},`$ (89)
$`d_{N+1}^{}`$ $`=`$ $`d_N^{}+4t(1+N^2)(sy^N+t1)+(2N+1)b_N^{}2f_N^{},`$ (90)
$`e_{N+1}^{}`$ $`=`$ $`ye_N^{}+4Nys\left[(s1)y^N+t\right],`$ (91)
$`f_{N+1}^{}`$ $`=`$ $`f_N^{}+4Nt(sy^N+t1),`$ (92)
where the values for the variables (Eqs. 79) have already been substituted and the initial conditions are $`a_1^{}=b_1^{}=\mathrm{}=f_1^{}=0`$. The general solution of the recursion relations Eqs. III is not difficult to obtain with the aid of an algebra software. Considering the invariance described in Eqs. II we will restrict our discussion to the case $`y>1`$, without loss of generality. For large values of $`N`$, the dominant terms of the solution of the mapping are
$`a_N+b_N`$ $``$ $`2sy^N,`$ (94)
$`c_N+d_N`$ $``$ $`2sy^NN^2,`$ (95)
$`a_N^{}+b_N^{}`$ $``$ $`\{\begin{array}{cc}4s(s1)y^NN\hfill & \text{if }s>1\hfill \\ \frac{8ty^N}{y1}\hfill & \text{if }s=1\hfill \end{array},`$ (98)
$`c_N^{}+d_N^{}`$ $``$ $`\{\begin{array}{cc}\frac{8}{3}s(s1)y^NN^3\hfill & \text{if }s>1\hfill \\ \frac{8ty^NN^2}{y1}\hfill & \text{if }s=1\hfill \end{array}.`$ (101)
The leading term in the derivative of the mean-square end-to-end distance will be
$$\frac{R^2}{z}|_{z=0}\frac{2}{3}(s1)N^3.$$
(102)
Therefore, up to first order in $`x`$, the scaling function $`F(x)`$ in the quasi-rigid limit is found to be $`F(x)1F_1(s,t)x`$. Considering the symmetry Eq. II and the solution for the isotropic case Eq. 56, we have
$$F_1(s,t)=\{\begin{array}{ccc}\frac{2(t1)}{3(q2)}\hfill & \text{if }y<1,\hfill & \\ \frac{1}{3}\hfill & \text{if }y=1,\hfill & \\ \frac{2(s1)}{3(q2)}\hfill & \text{if }y>1.\hfill & \end{array}$$
(103)
We thus conclude that the scaling function in general displays a discontinuous derivative at $`y=1`$.
## IV CONCLUSION
We formulated the problem of the calculation of the mean square end-to-end distance of semiflexible polymers placed on a $`q`$-coordinated anisotropic Bethe lattice as a linear mapping, whose general term may in principle be obtained. In the isotropic case, the mapping may easily be solved and leads to an expression for $`R^2`$ which is identical to the one obtained for random walks without immediate return on a hypercubic lattice with the same coordination number . The identity between the two problems regarding thermodynamic properties derived from the free energy is well known , and is here extended for a thermodynamic average of a geometric property. One point which should be stressed is that the definition of the Euclidean distance between two points on the Bethe lattice is rather arbitrary. Here we defined the distance by embedding the Cayley tree in a hypercubic lattice of sufficiently high dimensionality. In the thermodynamic limit the dimensionality of this lattice diverges, as expected . Other definitions of distance may be proposed . The simple one we adopted here leads to meaningful conclusions. Since calculations on the Bethe lattice are usually done recursively, and one step in the recursion relations corresponds to adding another generation to the tree, it is tempting to define the distance between two sites on the tree as the difference between the numbers of the generations they belong to. This definition, although simple and operational, has serious drawbacks, however. This is quite clear for the particular problem we looked at here, since it implies that $`R^2`$ for any $`N`$-step chain is equal to $`N^2`$. We would thus have $`\nu =1`$, the one-dimensional value, and the identity between the results for the Bethe lattice and for walks without immediate return on hypercubic lattices would break down. This definition of distance was used recently in the exact calculation of correlation functions for a general spin-$`S`$ magnetic model , leading to $`\nu =1`$, in opposition to the generally accepted mean-field value $`\nu =1/2`$ .
The fact that all walks we considered here have their initial site located at the central site of the Cayley tree is of course convenient for the calculations and may be seen as a particular case. A closer consideration of this point, however, leads to the conclusion that our results are exact for any chains such that the assertion that at any bend the new direction is perpendicular to all previous directions of bonds holds. Thus, it is clear that if the whole chain is contained in one of the $`q`$ rooted sub-trees attached to the central site the results are still the same. If portions of the chain are located on two of these sub-trees the calculation becomes more complicated since, as may be seen in figure 1, there are bonds in the same direction in different sub-trees. However, this problem may be easily avoided by enlarging the dimension of the euclidean space in which the tree is embedded, thus assuring that any two bonds in the same direction are necessarily connected by a walk without any bend. For such a tree, our results hold for any chain, regardless of the location of its endpoints.
In the general anisotropic case, we restricted ourselves to the discussion of the asymptotic behavior of $`R^2`$, which was studied in the semiflexible case and also in the quasi-rigid limit. The expected scaling behavior was obtained in both cases, and a interesting discontinuity in the quasi-rigid limit amplitude is observed as the isotropic value $`y=1`$ is crossed.
###### Acknowledgements.
We acknowledge partial financial support from the Brazilian agencies CNPq and FINEP.
## A COMBINATORIAL SOLUTION IN THE ISOTROPIC CASE
Any $`N`$-step walk on the Cayley tree will visit a subset of sites of the D-dimensional hypercubic lattice defining a subspace whose dimensionality is between 1 and $`N`$. The limiting cases are the ones of a polymer without any bend (rod), which is one-dimensional, and a polymer where we have a bend at every internal site, and since at each bend the new bond is in a direction orthogonal to all precedent bonds of the polymer, the polymer is embedded in a $`N`$-dimensional subspace. Since the initial site of the chain is supposed to be at the central site of the tree, the end-to-end distance will be given by the modulus of the position vector of the final site, denoted by $`\stackrel{}{R}`$. For a polymer with $`m`$ bends, the number of components of this vector will be equal to $`m+1`$. For simplicity, we will admit that each bond is of unit length, so that the components of $`\stackrel{}{R}`$ will be integers. We want to compute the mean value of $`\stackrel{}{R}`$ over all polymers with $`N`$ steps
$$R^2=\frac{_{\stackrel{}{R_m^N}}z^mR^2}{_{\stackrel{}{R_m^N}}z^m},$$
(A1)
where $`m`$ is the number of bends in the walk and the sum is over all configurations $`\stackrel{}{R_m^N}`$ of polymers with $`N`$ steps. Besides the first and last components the values of the other $`m1`$ components of $`R`$ are the numbers of steps between successive bends in the walk. We should remember that there are $`q2`$ possibilities for each bend. So we may rewrite Eq. A1
$$R^2=\frac{_{m=0}^{N1}a^mB_{N,m}}{_{m=0}^{N1}a^mA_{N,m}},$$
(A2)
where $`a=(q2)z`$ embodies all dependence on coordination number and statistical weight as long as $`q4`$,
$$A_{N,m}=\underset{\stackrel{}{R}_m^N}{}1,$$
(A3)
and
$$B_{N,m}=\underset{\stackrel{}{R}_m^N}{}\underset{i=0}{\overset{m+1}{}}R_i^2.$$
(A4)
Note that the effect of the bending energy can be described by introducing an effective coordination number $`q^{}=a+2`$ for an associated totally flexible polymer. The sums in $`A_{N,m}`$ and $`B_{N,m}`$ are over all possible values for $`\stackrel{}{R}`$ with $`m+1`$ components and subjected to the constraint of the total number of steps being equal to $`N`$, that is
$$\underset{i=1}{\overset{m+1}{}}R_i=N.$$
(A5)
The sum in Eq. A3 is just the number of vectors $`\stackrel{}{R}`$ with $`m+1`$ components which obey the constraint Eq. A5. Since the minimum value of each component of $`\stackrel{}{R}`$ is equal to 1, it is convenient to define $`r_i=R_i1`$ and therefore $`A_{N,m}`$ is the number of ways to put the $`Nm1`$ remaining steps into the $`m+1`$ components of $`\stackrel{}{R}`$
$$A_{N,m}=\frac{(N1)!}{m!(Nm1)!}.$$
(A6)
The sum $`B_{N,m}`$ may then be rewritten as
$$B_{N,m}=\underset{\stackrel{}{r}_m^N}{}\underset{i=1}{\overset{m+1}{}}(1+r_i)^2,$$
(A7)
where each components $`r_i`$ assumes values between 0 and $`Nm1`$ subject to the constraint of Eq. A5
$$\underset{i=1}{\overset{m+1}{}}r_i=Nm1.$$
(A8)
The calculation of $`B_{N,m}`$ is given in Appendix B, and the result is
$$B_{N,m}=\frac{(m+1)(2Nm)N!}{(m+2)!(Nm1)!}.$$
(A9)
Performing the sum in the denominator of Eq. A2 taking Eq. A6 into account, we have
$$R^2=\frac{N}{[1+a]^{N1}}\left[2(N+1)\underset{m=0}{\overset{N1}{}}\left(\begin{array}{c}N1\\ m\end{array}\right)\frac{a^m}{m+2}\underset{m=0}{\overset{N1}{}}\left(\begin{array}{c}N1\\ m\end{array}\right)a^m\right].$$
(A10)
The first sum may be calculated by noting that
$$_0^Ax(1+x)^{N1}๐x=A^2\underset{m=0}{\overset{N1}{}}\left(\begin{array}{c}N1\\ m\end{array}\right)\frac{A^m}{m+2},$$
(A11)
and therefore
$$\underset{m=0}{\overset{N1}{}}\left(\begin{array}{c}N1\\ m\end{array}\right)\frac{a^m}{m+2}=\frac{[1+a]^N[aN1]+1}{N(N+1)a^2}.$$
(A12)
Substituting this result in Eq. A10 and performing the second sum we finally get the expression
$$R^2=\frac{2[1+a]}{a^2}\left[Na1+\frac{1}{[1+a]^N}\right]N.$$
(A13)
## B DERIVATION OF $`B_{N,m}`$
In this appendix we want to derive Eq. A9 for $`B_{N,m}`$. Using Eq. A8 and defining for convenience $`๐ฉ=Nm1`$ the equation Eq. A7 is rewritten as
$$B_{N,m}=(m+1)\underset{j=0}{\overset{๐ฉ+1}{}}\frac{(๐ฉ+mj)!j^2}{(๐ฉ+1j)!(m1)!},$$
(B1)
Redefining the summation variable with $`i=๐ฉ+1j`$ this equation turns out to be
$$B_{N,m}=(m+1)\underset{i=0}{\overset{๐ฉ}{}}(๐ฉ+1i)^2\frac{(i+m1)!}{i!(m1)!}.$$
(B2)
Using the equality
$$\underset{i=0}{\overset{N}{}}\frac{(m+i)!}{m!i!}=\frac{(m+N+1)!}{(m+1)!N!},$$
(B3)
we get after some manipulations
$`B_{N,m}`$ $`=`$ $`(m+1)\{(๐ฉ+1){\displaystyle \frac{(๐ฉ+1+m)!}{๐ฉ!m!}}2(๐ฉ+1){\displaystyle \frac{m(๐ฉ+1+m)!}{๐ฉ!(m+1)!}}`$ (B5)
$`+{\displaystyle \underset{i=1}{\overset{๐ฉ+1}{}}}{\displaystyle \frac{(i+m1)!i}{(i1)!(m1)!}}\}.`$
The last summation to be dealt with is just
$$\underset{i=1}{\overset{๐ฉ+1}{}}\frac{(i+m1)!i}{(i1)!(m1)!}.$$
(B6)
Defining $`j=i1`$ it follows that
$$\underset{j=0}{\overset{๐ฉ}{}}\frac{(j+m)!(j+1)}{j!(m1)!}=\frac{m(๐ฉ+m+1)!}{๐ฉ!(m+1)!}+m(m+1)\frac{(๐ฉ+m+1)!}{(๐ฉ1)!(m+2)!}.$$
(B7)
Substitution in Eq. B5 leads to
$$B_{N,m}=\frac{(m+1)(๐ฉ+m+1)![2๐ฉ+m+2]}{(m+2)!๐ฉ!}.$$
(B8)
Substituting $`๐ฉ=Nm1`$, we get
$$B_{N,m}=\frac{(m+1)(2Nm)N!}{(m+2)!(Nm1)!}.$$
(B9) |
warning/0002/hep-ph0002225.html | ar5iv | text | # Exact and Broken Symmetries in Particle Physics
## 1 Introduction
All experimental evidence points to the strong, weak and electromagnetic interactions of hadrons (strongly interacting particles) and of leptons as being described by a gauge theory, based on the group
$$G_{\mathrm{SM}}=SU(3)\times SU(2)\times U(1).$$
(1)
The strong interaction theoryโQCDโhas as fundamental fermionic entities a triplet of quarks, which feel the $`SU(3)`$ gauge interactions. Both the quarks and the leptons appear in nature in a repetitive fashion, in three distinct families of doublets under the $`SU(2)\times U(1)`$ electroweak group. Although $`G_{\mathrm{SM}}`$ correctly describes the symmetry of the fundamental interactions among quarks and leptons, only $`SU(3)`$ is an exact symmetry of the theory. The electroweak group, in fact, suffers a spontaneous breakdown to $`U(1)_{\mathrm{em}}`$:
$$SU(2)\times U(1)U(1)_{\mathrm{em}}.$$
(2)
In these lectures we will describe the fundamental concepts upon which the theory for these interactions is built upon. These are related to the way in which symmetries are realized in nature and to the role of gauge fields in rendering theories invariant under local transformations. A crucial notion is that of a spontaneously broken symmetry and the effect that this spontaneous breakdown has for the spectrum of excitations in the theory.
## 2 Global Symmetries in Field Theory
The natural language for elementary particle physics is that of a quantum field theory, where to each fundamental excitation one assigns a corresponding quantum field. $`^\mathrm{?}`$ Symmetries of nature are incorporated by constructing Lagrangian densities, made up of these quantum fields, which have an action
$$W=d^4x$$
(3)
explicitly invariant under the symmetry in question:
$$WW^{}=W.$$
(4)
In what follows, I will consider only continuous symmetry transformations based on some Lie group $`G`$$`^\mathrm{?}`$ Let me denote a generic quantum field by $`\chi _\alpha (x)`$, which $`x`$ being the space-time location of the quantum field and $`\alpha `$ being an (internal) index which runs over the possible components of $`\chi `$. \[For instance, for a quark field which is a triplet of $`SU(3)`$ one would have $`q_\alpha (x)`$, with $`\alpha =1,2,3`$.\] If $`a`$ is one of the operations of the symmetry group $`G`$ of transformations, and if the quantum fields $`\chi _\alpha `$ are members of an (irreducible) multiplet, then under this operation one has
$$\chi _\alpha (x)\stackrel{a}{}\chi _\alpha ^{}(x)=_{\alpha \beta }(a)\chi _\beta (x).$$
(5)
That is, under the transformation the field $`\chi `$ goes into a new field $`\chi ^{}`$ whose components are linear combinations of the old components.
Because, by assumption, the quantum fields $`\chi _\alpha `$ are members of a multiplet under $`G`$, the matrices $`(a)`$ constitute a representation matrix for the group $`G`$ and obey a characteristic composition property. This follows from comparing the sequence of transformations
$$\chi _\alpha (x)\stackrel{a}{}\chi _\alpha ^{}(x)\stackrel{a^{}}{}\chi _\alpha ^{\prime \prime }(x)$$
(6)
to the direct transformation
$$\chi _\alpha (x)\stackrel{a^{\prime \prime }}{}\chi _\alpha ^{\prime \prime }(x).$$
(7)
Hence, one finds
$$_{\alpha \beta }(a^{})_{\beta \gamma }(a)=_{\alpha \gamma }(a^{\prime \prime })$$
(8)
In the Hilbert space of the quantum field $`\chi _\alpha (x)`$ the transformation (5) is induced by a unitary operator $`U(a)`$, so that
$$U^1(a)\chi _\alpha (x)U(a)=\chi _\alpha ^{}(x)=_{\alpha \beta }(a)\chi _\beta (x).$$
(9)
It is easy to see that the composition property (8) has its counterpart in terms of the unitary operators $`U`$:
$$U(a)U(a^{})=U(a^{\prime \prime }).$$
(10)
Since we are considering continuous symmetry transformations, it suffices to focus only on infinitesimal transformations $`\delta a`$, since finite transformations can always be built up via (10) by (infinite) compounding. A given Lie group is characterized by the number of parameters associated with these infinitesimal transformations and, more specifically, by the algebra obeyed by the operators connected to the distinct infinitesimal parameters.
Let us write for an infinitesimal transformation
$$U(\delta a)=1+i\delta a_iG_i$$
(11)
where the index $`i`$ runs over all the independent infinitesimal parameters of the Lie group G \[e.g. for the rotation group in 3 dimensions $`O(3)`$, $`\delta a_i`$ would describe the three independent rotations about the $`x,y`$ and $`z`$ axis\]. The operators $`G_i`$ are called the group generators and the composition property (10) implies a group algebra for the generators. Without loss of generality the parameters $`\delta a_i`$ can be taken as real, so that the $`G_i`$ are Hermitian. They obey the Lie algebra:
$$[G_i,G_j]=ic_{ijk}G_k.$$
(12)
The structure constants $`c_{ijk}`$ characterize the group $`G`$ and can be chosen so as to be totally antisymmetric in $`i,j`$ and $`k`$.
Just as $`U(\delta a)`$ can be expanded in terms of the generators $`G_i`$, so can the representation matrices $`_{\alpha \beta }(\delta a)`$. One has, for an infinitesimal transformation
$$_{\alpha \beta }(\delta a)=\delta _{\alpha \beta }+i\delta a_i(g_i)_{\alpha \beta }.$$
(13)
It is easy to show that the matrices $`g_i`$ furnish a representation for the generators $`G_i`$ and so obey themselves Eq. (12). To see this let us use (13) and (11) in the defining equation (9). One has
$$(1i\delta a_iG_i)\chi _\alpha (x)(1+i\delta a_iG_i)=\chi _\alpha (x)+i\delta a_i(g_i)_{\alpha \beta }\chi _\beta $$
(14)
which implies
$$[G_i,\chi _\alpha (x)]=(g_i)_{\alpha \beta }\chi _\beta (x).$$
(15)
This equation embodies succinctly how the quantum fields $`\chi _\alpha `$ transform under the group $`G`$, and will be repeatedly used in what follows. By using (15) in the Jacobi identity
$$[G_i,[G_j,\chi _\alpha ]]+[\chi _\alpha ,[G_i,G_j]]+[G_j,[\chi _\alpha ,G_i]]=0$$
(16)
one readily sees that the matrices $`g_i`$ obey Eq. (12).
Let us explore the consequences of having a theory built out of the quantum fields $`\chi _\alpha `$ which is invariant under the transformations of the group $`G`$. As we shall see, the invariance of the action under $`G`$ implies the existence of conserved currents and a set of constants of the motion, which are nothing else but the generators $`G_i`$ of the group! Since the Lagrangian density $``$ depends in general on $`\chi _\alpha `$ and its space-time derivatives $`_\mu \chi _\alpha `$, the invariance statement (4) implies
$$d^4x(\chi _\alpha ,_\mu \chi _\alpha )=d^4x(\chi _\alpha ^{},_\mu \chi _\alpha ^{}).$$
(17)
For $`\chi _\alpha ^{}`$ infinitesimally different from $`\chi _\alpha `$, the stationarity of the action implies
$`0=\delta W`$ $`=`$ $`{\displaystyle d^4x\left[\frac{}{\chi _\alpha }\delta \chi _\alpha +\frac{}{_\mu \chi _\alpha }\delta _\mu \chi _\alpha \right]}`$
$`=`$ $`{\displaystyle }d^4x[\{{\displaystyle \frac{}{\chi _\alpha }}_\mu \left({\displaystyle \frac{}{_\mu \chi _\alpha }}\right)\}\delta \chi _\alpha `$
$`+_\mu \left[{\displaystyle \frac{}{_\mu \chi _\alpha }}\delta \chi _\alpha \right]].`$
The first term above in the curly brackets vanishes because of the Euler-Lagrange equations of motion. The second can be rewritten in terms of the generator matrices $`g_i`$, since
$$\delta \chi _\alpha =\chi _\alpha ^{}\chi _\alpha =i\delta a_i(g_i)_{\alpha \beta }\chi _\beta .$$
(19)
Hence
$$0=\delta W=d^4x\delta a_i_\mu \left[\frac{}{_\mu \chi _\alpha }\frac{1}{i}(g_i)_{\alpha \beta }\chi _\beta \right].$$
(20)
Since the parameters $`\delta a_i`$ are independent, it follows that the currents
$$J_i^\mu (x)=\frac{}{_\mu \chi _\alpha (x)}\frac{1}{i}(g_i)_{\alpha \beta }\chi _\beta (x),$$
(21)
as a result of the symmetry, are conserved
$$_\mu J_i^\mu (x)=0.$$
(22)
Because of (22)โ if one assumes that the fields $`\chi _\alpha `$ drop off sufficiently fast at spatial infinityโ there exists a set of constants of the motion, given by the space integral of the $`J_i^o`$. One has
$$Q_i=d^3xJ_i^o(x)$$
(23)
with
$$\frac{d}{dt}Q_i=0.$$
(24)
It is easy to checkโ and we shall do so belowโ that the operators $`Q_i`$ are precisely the generators $`G_i`$. That is, they obey both Eqs. (12) and (15). If $`H`$ is the Hamiltonian of the theory, then Heisenbergโs equation of motion imply
$$[H,G_i]=0$$
(25)
which may be a more familiar way to express the invariance of the theory under the transformations of the group $`G`$ (e.g. rotational invariance is expressed via the vanishing of the commutator $`[H,L_i]=0`$).
Let us verify that indeed
$$G_iQ_i=d^3xJ_i^o=d^3x\left[\frac{}{_o\chi _\alpha }\frac{1}{i}(g_i)_{\alpha \beta }\chi _\beta \right]$$
(26)
acts as a generator is supposed to do. For that, remark that the canonical momentum conjugate to $`\chi _\alpha `$ is precisely $`^\mathrm{?}`$
$$\pi _\alpha (x)=\frac{}{_o\chi _\alpha (x)}$$
(27)
and that (for bosonic fields) one has the equal time commutation relations
$`[\pi _\alpha (x),\chi _\beta (y)]|_{x^o=y^o}`$ $`=`$ $`{\displaystyle \frac{1}{i}}\delta ^3(\stackrel{}{x}\stackrel{}{y})\delta _{\alpha \beta }`$ (28)
$`[\pi _\alpha (x),\pi _\beta (y)]|_{x^o=y^o}`$ $`=`$ $`[\chi _\alpha (x),\chi _\beta (y)]|_{x^o=y^o}=0.`$
Then
$$G_i=d^3x\pi _\alpha (x)\frac{1}{i}(g_i)_{\alpha \beta }\chi _\beta (x).$$
(29)
Since $`G_i`$ is time-independent, in computing the commutator of $`G_i`$ with $`\chi _\gamma (y)`$ one can set the time $`x^o`$ in (29) equal to $`y^o`$. Using (28) it is then trivial to check that
$`[G_i,\chi _\gamma (y)]`$ $`=`$ $`{\displaystyle d^3x[\pi _\alpha (x)\frac{1}{i}(g_i)_{\alpha \beta }\chi _\beta (x),\chi _\gamma (y)]_{x^o=y^o}}`$ (30)
$`=`$ $`(g_i)_{\gamma \beta }\chi _\beta (y)`$
and
$`[G_i,G_j]`$ $`=`$ $`{\displaystyle d^3xd^3y[\pi _\alpha (x)\frac{1}{i}(g_i)_{\alpha \beta }\chi _\beta (x),\pi _\gamma (y)\frac{1}{i}(g_j)_{\gamma \delta }\chi _\delta (y)]_{x^o=y^o}}`$ (31)
$`=`$ $`{\displaystyle d^3x\pi _\alpha (x)\left(\frac{1}{i}[g_i,g_j]\right)_{\alpha \beta }\chi _\beta (x)}`$
$`=`$ $`ic_{ijk}{\displaystyle d^3x\pi _\alpha (x)\frac{1}{i}(g_k)_{\alpha \beta }\chi _\beta (x)}=ic_{ijk}G_k.`$
Up to now in the discussion of symmetries I focussed on the transformation properties of the quantum fields $`\chi _\alpha (x)`$. What equation (9) says is that under a group transformation the component fields $`\chi _\alpha `$ transform in a well-defined way. The correspondence between quantum fields and particles makes it natural to suppose that the quantum states associated with the fields $`\chi _\alpha (x)`$ will transform in an analogous way. Let me denote the one-particle state associated with the field $`\chi _\alpha `$ by $`|p;\alpha `$, where $`p^\mu `$ is the 4-momentum of the state and, since these states are supposed to describe particles of a given mass, $`p^2=m_\alpha ^2`$. Then, corresponding to Eq. (9), one has
$$U^1(a)|p;\alpha =_{\alpha \beta }(a)|p;\beta .$$
(32)
This equation can be used to deduce that all states of the multiplet $`|p;\alpha `$ have the same mass.
Let $`|p;\alpha _{\mathrm{rest}}`$ denote the state corresponding to 4-momentum $`p^\mu =(\stackrel{}{0},m_\alpha )`$. Then, by definition, the action of the Hamiltonian on this state is just
$$H|p;\alpha _{\mathrm{rest}}=m_\alpha |p;\alpha _{\mathrm{rest}}.$$
(33)
However, if the theory is invariant under the group $`G`$, so that $`H`$ commutes with all the generators ( c. f. Eq. (25)) it follows also that
$$[H,U^1(a)]=0.$$
(34)
Applying this equation on the rest state proves our contention, since
$`0`$ $`=`$ $`[H,U^1(a)]|p;\alpha _{\mathrm{rest}}=(HU^1(a)U^1(a)H)|p;\alpha _{\mathrm{rest}}`$ (35)
$`=`$ $`_{\alpha \beta }(a)(m_\beta m_\alpha )|p;\beta _{\mathrm{rest}}.`$
Because $`_{\alpha \beta }(a)`$ is arbitrary, it follows that $`m_\alpha =m_\beta `$.
One says that a symmetry is realized in a Wigner-Weyl way if the invariance of the action under $`G`$ leads to the appearance in nature of particle multiplets with the same mass. $`^\mathrm{?}`$ A well known example of an (approximate) Wigner-Weyl symmetry is strong isospin. This approximate global $`SU(2)`$ symmetry of the strong interaction leads to a nearly degenerate nucleon doublet $`(m_pm_n)`$ and a pion triplet $`(m_{\pi ^+}=m_\pi ^{}m_{\pi ^o})`$. Remarkably, however, the Wigner-Weyl way is not the only way in which a symmetry can be realized in nature!
## The Nambu-Goldstone Realization
It is possible that the action is invariant under a symmetry group $`G`$ but that the physical states of the theory show no trace of this symmetry. This happens in the case in which, although
$$[H,U^1(a)]=0,$$
(36)
the vacuum state is not invariant under $`G`$. Such symmetries are called spontaneously broken, or realized in a Nambu-Goldstone way. $`^\mathrm{?}`$
Eq. (32), which lead to the deduction that all states in a multiplet $`|p,\alpha `$ have the same mass, can be derived from the transformation properties of the quantum fields $`\chi _\alpha `$, provided one assumes that the vacuum state is $`G`$ invariant:
$$U(a)|0=|0.$$
(37)
The one particle states $`|p,\alpha `$ are constructed by the action of the (asymptotic) creation operators for the field $`\chi _\alpha `$$`^\mathrm{?}`$ For a scalar field $`\chi _\alpha (x)`$ one writes in the usual way
$$\chi _\alpha (x)=\frac{d^3p}{(2\pi )^32p^o}[e^{ipx}a_\alpha (p,t)+e^{ipx}a_\alpha ^{}(p,t)].$$
(38)
Then, one has
$`|p;\alpha `$ $`=`$ $`\underset{t\pm \mathrm{}}{lim}a_\alpha ^{}(p,t)|0`$ (39)
$`=`$ $`\underset{x^o\pm \mathrm{}}{lim}{\displaystyle d^3xe^{ipx}\frac{1}{i}}\stackrel{}{}_o\chi _\alpha (x)|0,`$
where
$$A\stackrel{}{}_oB=A_oB(_oA)B.$$
(40)
Consider then, as in Eq. (32), the action of $`U^1(a)`$ on the state $`|p;\alpha `$
$$U^1(a)|p,\alpha =\underset{x^o\pm \mathrm{}}{lim}d^3xe^{ipx}\frac{1}{i}\stackrel{}{}_oU^1(a)\chi _\alpha (x)|0.$$
(41)
If (37) holds, one can write
$`U^1(a)\chi _\alpha (x)|0`$ $`=`$ $`U^1(a)\chi _\alpha (x)U(a)|0`$ (42)
$`=`$ $`_{\alpha \beta }(a)\chi _\beta (x)|0`$
which immediately establishes (32). However, if the vacuum is not left invariant by a $`G`$-transformationโ i.e. if the vacuum state is degenerate or not uniqueโ then even though the fields $`\chi _\alpha `$ transform according to some irreducible representation, there are no longer degenerate multiplets in the spectrum.
When a symmetry is realized in a Nambu-Goldstone way, instead of having multiplets of particles with the same mass, there appear in the theory massless excitationsโ the so-called Goldstone bosons. To see how these ensue consider again the fields $`\chi _\alpha `$ and take the vacuum expectation value of Eq. (15)
$$0|[G_i,\chi _\alpha (x)]|0=(g_i)_{\alpha \beta }0|\chi _\beta (x)|0.$$
(43)
If the vacuum is invariant under $`G`$ transformations it follows from Eq. (37) that
$$G_i|0=0.$$
(44)
It is immediate from (43) then that the vacuum expectation values of the fields $`\chi _\alpha `$ must vanish. However, if (44) does not hold, and $`\chi _\alpha `$ are scalar fields, there is no argument why one cannot have
$$0|\chi _\alpha (x)|00.$$
(45)
\[If $`\chi _\alpha `$ correspond to fields with spin then the equivalent of Eq. (43) for Lorentz transformations, along with the invariance of the vacuum under these transformations, informs one that the vacuum expectation value of these fields must vanish.\]
A symmetry is realized in a Nambu-Goldstone way if there exist some scalar field (which may not necessarily be elementary) with non-zero vacuum expectation value. Imagine that this is so in Eq. (43). Then using the definition of the generators $`G_i`$ (Eq. (26)) one has
$$0(g_i)_{\alpha \beta }0|\chi _\beta (x)|0=d^3y0|J_i^o(y)\chi _\alpha (x)\chi _\alpha (x)J_i^o(y)|0.$$
(46)
This equation can be written in a more interesting way by inserting a complete set of states $`|n`$ and making use of translational invariance on the currents $`J_i^o(y)`$
$$J_i^o(y)=e^{iPy}J_i^o(0)e^{iPy}.$$
(47)
Then the RHS of Eq. (46) reads
$`\mathrm{RHS}`$ $`=`$ $`{\displaystyle \underset{n}{}}{\displaystyle }d^3y\{0|e^{iPy}J_i^o(0)e^{iPy}|nn|\chi _\alpha (x)|0`$ (48)
$`0|\chi _\alpha (x)|nn|e^{iPy}J_i^o(0)e^{iPy}|0\}`$
$`=`$ $`{\displaystyle \underset{n}{}}{\displaystyle d^3ye^{iP_ny}0|J_i^o(0)|nn|\chi _\alpha (x)|0}`$
$`{\displaystyle \underset{n}{}}{\displaystyle d^3ye^{iP_ny}0|\chi _\alpha (x)|nn|J_i^o(0)|0}`$
$`=`$ $`{\displaystyle \underset{n}{}}(2\pi )^3\delta ^3(\stackrel{}{p}_n)\{e^{iP_n^oy^o}0|J_i^o(0)|nn|\chi _\alpha (x)|0`$
$`e^{+iP_n^oy^o}0|\chi _\alpha (x)|nn|J_i^o(0)|0\}.`$
By assumption this expression does not vanish and, furthermore, since the LHS is independent of $`y^o`$ it must also be independent of $`y^o`$. Clearly this can only happen if in the theory there exist some massless one-particle states $`|n`$ and only these states contribute to the sum in (48). These zero mass states are the Goldstone bosons.
It is not difficult to convince oneself that for each generator $`G_i`$ that does not annihilate the vacuum there is a corresponding Goldstone boson (after all the action of $`G_i`$ on the vacuum must give some stateโand these states are associated with the Goldstone bosons!). Let us write the Goldstone boson states as $`|p;j`$, where $`p^2=0`$. Then it follows that the matrix element of the currents associated with the broken generators between the vacuum and these states are non-vanishing:
$$0|J_i^\mu (0)|p;j=if_j\delta _{ij}p^\mu $$
(49)
where $`f_j`$ are some non-vanishing constants, which are related to the vacuum expectation values of the fields $`\chi _\alpha `$. Indeed, remembering that for a one-particle state
$$\underset{n}{}\frac{d^3p_n}{(2\pi )^32p_n^o}$$
(50)
it follows from Eqs. (46) and (48) that
$$i(g_i)_{\alpha \beta }0|\chi _\beta (0)|0=\underset{p^\mu 0}{lim}\frac{1}{2}\left[f_ip;i|\chi _\alpha (0)|0+f_i^{}0|\chi _\alpha (0)|p;i\right].$$
(51)
Because the Nambu-Goldstone realization of a symmetry is so much less familiar, it is instructive to illustrate it with a very simple example. For these purposes consider the following Lagrangian density describing the interaction of a complex scalar field $`\varphi `$ with itself
$$=_\mu \varphi ^{}^\mu \varphi \lambda \left(\varphi ^{}\varphi \frac{1}{2}f\right)^2.$$
(52)
Obviously this theory is invariant under a $`U(1)`$ transformation (phase transformation)
$`\varphi (x)`$ $``$ $`\varphi ^{}(x)=e^{i\alpha }\varphi (x)`$
$`\varphi ^{}(x)`$ $``$ $`\varphi ^{}(x)=e^{i\alpha }\varphi ^{}(x).`$ (53)
The conserved current associated with this symmetry is easily constructed from our general formula (21)
$$J^\mu =\frac{}{_\mu \varphi }\frac{1}{i}(1)\varphi +\frac{}{_\mu \varphi ^{}}\frac{1}{i}(1)\varphi ^{}=i\left[(^\mu \varphi ^{})\varphi (^\mu \varphi )\varphi ^{}\right].$$
(54)
The corresponding generator
$$G=d^3xJ^o=id^3x\left[(^o\varphi ^{})\varphi (^o\varphi )\varphi ^{}\right]$$
(55)
obeys the commutation relations (15)
$`[G,\varphi (x)]`$ $`=`$ $`\varphi (x)`$ (56)
$`[G,\varphi ^{}(x)]`$ $`=`$ $`+\varphi ^{}(x).`$
In a classical sense, the second term in the Lagrangian correspond to a potential for the fields $`\varphi ,\varphi ^{}`$:
$$V(\varphi ,\varphi ^{})=\lambda \left(\varphi ^{}\varphi \frac{1}{2}f\right)^2.$$
(57)
Obviously, to guarantee the positivity of the theory, one needs that $`\lambda >0`$. However, the physics is very different depending on the sign of $`f`$. If $`f<0`$ the potential has a unique minimum at $`\varphi =\varphi ^{}=0`$ and the theory is realized in a Wigner-Weyl way, leading to a degenerate multiplet of massive states. If $`f>0`$, on the other hand, the potential has an infinity of minima characterized by the condition $`\varphi ^{}\varphi =\frac{1}{2}f`$. The theory is realized in a Nambu-Goldstone way and there is both a massless and a massive state in the theory.
Quantum mechanically, if $`f<0`$, it is sensible to expand the potential about $`\varphi =0`$, since this is the minimum of the potential. One has
$$V=\lambda \left(\varphi ^{}\varphi \frac{1}{2}f\right)^2=\frac{1}{4}\lambda f^2\lambda f\varphi ^{}\varphi +\lambda (\varphi ^{}\varphi )^2.$$
(58)
The quadratic term $`\lambda f\varphi ^{}\varphi `$, since $`f<0`$, is a perfectly good mass term for the fields $`\varphi `$ and $`\varphi ^{}`$ and one identifies
$$m_\varphi ^2=m_\varphi ^{}^2=\lambda f>0.$$
(59)
In this case, one has a degenerate multiplet of two charge-conjugate particles interacting via the $`\lambda (\varphi ^{}\varphi )^2`$ term.
If $`f>0`$, on the other hand, an expansion about $`\varphi =0`$ makes no sense as the potential has a local maximum. The only sensible point to expand the potential is about its minimum value which occurs at $`\varphi _{\mathrm{min}}=\sqrt{\frac{f}{2}}e^{i\theta }`$. In fact since $`f>0`$ there is no way that the quadratic term in $`\varphi ^{}\varphi `$ can represent a mass term.
Quantum mechanically the non-zero value of $`\varphi _{\mathrm{min}}`$ implies that $`\varphi `$ has a non-vanishing vacuum expectation value
$$\theta |\varphi (x)|\theta =\sqrt{\frac{f}{2}}e^{i\theta }.$$
(60)
The phase $`\theta `$, characterizing the vacuum state $`|\theta `$, is in fact irrelevant and can be rotated away. It is a reflection of the non-uniqueness of the vacuum state of the theory. Since under a $`U(1)`$ transformation
$$U^1(\alpha )\varphi )(x)U(\alpha )=e^{i\alpha }\varphi (x)$$
(61)
it is clear that the expectation of $`\varphi (x)`$ between the states $`U(\theta )|\theta `$ is purely real
$$\theta |U^1(\theta )\varphi (x)U(\theta )|\theta =e^{i\theta }e^{i\theta }\sqrt{\frac{f}{2}}=\sqrt{\frac{f}{2}}.$$
(62)
Obviously $`U(\theta )|\theta |0`$ is just as good a vacuum as $`|\theta `$.
Without loss of generality we can set $`\theta =0`$ and expand $`\varphi `$ as
$$\varphi =\sqrt{\frac{f}{2}}+\chi $$
(63)
where the quantum field $`\chi `$, by assumption, has a vanishing vacuum expectation value. The potential in terms of $`\chi `$ reads
$`V`$ $`=`$ $`\lambda \left(\varphi ^{}\varphi {\displaystyle \frac{f}{2}}\right)^2=\lambda \left(\chi ^{}\chi +\sqrt{{\displaystyle \frac{f}{2}}}(\chi +\chi ^{})\right)^2`$
$`=`$ $`{\displaystyle \frac{\lambda f}{2}}(\chi +\chi ^{})^2+\sqrt{2f}\lambda (\chi +\chi ^{})\chi ^{}\chi +\lambda ^2(\chi ^{}\chi )^2.`$
Obviously, it appears that a linear combination of $`\chi `$ and $`\chi ^{}`$ has a mass, while its orthogonal combination is massless. Let us write
$$\chi _+=\frac{1}{\sqrt{2}}(\chi +\chi ^{});\chi _{}=\frac{i}{\sqrt{2}}(\chi ^{}\chi ).$$
(65)
Then
$$m_+^2=2\lambda f>0;m_{}^2=0.$$
(66)
Even though the Langragian (52) is $`U(1)`$ symmetric, this symmetry is not reflected in the spectrum, when the theory is realized in the Nambu-Goldstone manner!
The above identification of $`\chi `$ as the Goldstone boson field also follows directly from the commutators (56). Since $`f`$ is real by assumption, one has
$$\chi _{}=\frac{i}{\sqrt{2}}(\chi ^{}\chi )=\frac{i}{\sqrt{2}}(\varphi ^{}\varphi )$$
(67)
and hence
$$[G,\chi _{}]=\frac{i}{\sqrt{2}}(\varphi ^{}+\varphi )=i\left[\sqrt{f}+\chi _+\right].$$
(68)
Whence, taking expectation values, one obtains
$$0|[G,\chi _{}]|0=i\sqrt{f}.$$
(69)
This equation clearly singles out $`\chi _{}`$ as the Goldstone boson field.
If $`|p`$ is the state corresponding to this Goldstone boson then, neglecting non-linearities, one expects
$$0|\chi _{}(0)|p=1$$
(70)
Eq. (69) then gives, in the same approximation,
$$0|J^\mu (0)|p=i\sqrt{f}p^\mu .$$
(71)
The decay constant $`f_i`$ of Eq. (49) here is just $`\sqrt{f}`$ and is related to the vacuum expectation value of $`\varphi `$, as expected from Eq. (51). There is an alternative way to accomplish this identification by using directly the current $`J^\mu `$ and rewriting it in terms of the fields $`\chi _+`$ and $`\chi _{}`$. One has
$`J^\mu `$ $`=`$ $`i[(^\mu \varphi ^{})\varphi (^\mu \varphi )\varphi ^{}]`$
$`=`$ $`i\left[(^\mu \chi ^{})\left(\sqrt{{\displaystyle \frac{f}{2}}}+\chi \right)(^\mu \chi )\left(\sqrt{{\displaystyle \frac{f}{2}}}+\chi ^{}\right)\right]`$
$`=`$ $`i\sqrt{f}{\displaystyle \frac{1}{\sqrt{2}}}^\mu (\chi ^{}\chi )+i[(^\mu \chi ^{})\chi (^\mu \chi )\chi ^{}]`$
$`=`$ $`\sqrt{f}^\mu \chi _{}+\text{non-linear terms}`$
which directly implies (71).
To summarize, there are two ways in which symmetries $`([H,U]=0)`$ can be realized in nature. If the vacuum state is unique $`(U|0=|0)`$, then we have a Wigner-Weyl realization with degenerate particle multiplets. If, on the other hand, the vacuum state is not unique $`(U|0|0)`$, then we have a Nambu-Goldstone realization with a number of massless excitations, one for each of the generators of the group which does not annihilate the vacuum. In this latter case one often refers to the phenomena as spontaneous symmetry breaking because, although the symmetry exists, it is not reflected in the spectrum of the states of the theory.
## 3 Local Symmetries in Field Theory
In all the preceding discussion I have talked implicitly only about global symmetry transformations. That is the parameters $`\delta a_i`$ were assumed to be independent of space-time. Clearly in this case fields at different space-time points are transformed all in the same way. One may well ask what happens if the group parameters are space-time dependent. In this case the fields $`\chi _\alpha (x)`$ and $`\chi _\alpha (x^{})`$ would be rotated in a different way by the group transformation. Transformations where this happens are called local symmetries, to distinguish them from the global case when $`\delta a_i`$ is $`x`$-independent.
Under a local transformation one has
$$\chi _\alpha (x)\stackrel{a(x)}{}\chi _\alpha ^{}(x)=_{\alpha \beta }(a(x))\chi _\beta (x).$$
(73)
Because $``$ is now space-time dependent, even though the action
$$W=d^4x(_\mu \chi _\alpha ,\chi _\alpha )$$
(74)
was invariant under global $`G`$ transformations, this action will fail to be invariant under local $`G`$ transformations. Because of the kinetic energy terms, which depends on $`_\mu \chi _\alpha `$, there will be pieces in $`W`$ which are no longer invariant. Indeed, it is easy to identify what destroys the possibility of local invariance of the action. Consider the transformation of the derivative term $`_\mu \chi _\alpha `$ under local transformations. One has
$`_\mu \chi _\alpha (x)\stackrel{a(x)}{}_\mu \chi _\alpha ^{}(x)`$ $`=`$ $`_\mu [_{\alpha \beta }(a(x))\chi _\beta (x)]`$ (75)
$`=`$ $`_{\alpha \beta }(a(x))_\mu \chi _\beta (x)+_\mu _{\alpha \beta }(a(x))\chi _\beta (x).`$
The presence of the second term above destroys the local invariance of the action. However, one can compensate for the appearance of this term by adding to the, globally invariant, Lagrangian additional fields (gauge fields) which cancel this contribution. It is clear that to make a Lagrangian locally invariant necessarily involves the introduction of more degrees of freedom in the theory.
Before giving a general prescription of how to make a globally invariant Lagrangian locally invariant, it is useful to illustrate this procedure with a simple example. Consider a free Dirac field with Lagrangian density
$$=\overline{\psi }(x)\left(\gamma ^\mu \frac{1}{i}_\mu +m\right)\psi (x).$$
(76)
Clearly $``$ is invariant under the $`U(1)`$ transformation
$`\psi (x)`$ $``$ $`\psi ^{}(x)=e^{i\alpha }\psi (x)`$
$`\overline{\psi }(x)`$ $``$ $`\overline{\psi }^{}(x)=e^{i\alpha }\overline{\psi }(x),`$ (77)
which leads to the associated current:
$$J^\mu (x)=\frac{}{_\mu \psi (x)}\frac{1}{i}(1)\psi (x)=\overline{\psi }(x)\gamma ^\mu \psi (x).$$
(78)
It is clear, however, that if $`\alpha =\alpha (x)`$ the Lagrangian (76) ceases to be invariant, since
$$_\mu \psi (x)_\mu \psi ^{}(x)=e^{i\alpha (x)}_\mu \psi (x)+i(_\mu \alpha (x))\psi (x)e^{i\alpha (x)}.$$
(79)
Thus
$`(x)\stackrel{a(x)}{}^{}(x)`$ $`=`$ $`(x)(_\mu \alpha (x))\overline{\psi }(x)\gamma ^\mu \psi (x)`$ (80)
$`=`$ $`(x)J^\mu (x)_\mu \alpha (x).`$
One may get rid of the additional contribution in (80) by augmenting the Lagrangian (76) by an additional term
$$_{\mathrm{extra}}=eA^\mu (x)J_\mu (x)$$
(81)
involving a vector field $`A^\mu (x)`$, which under a local $`U(1)`$ transformation translates by an amount $`_\mu \alpha (x)`$:
$$A^\mu (x)\stackrel{a(x)}{}A^\mu (x)=A^\mu (x)+\frac{1}{e}(^\mu \alpha (x)).$$
(82)
Of course, if this field $`A^\mu (x)`$ is to have a dynamical role, and one wants to preserve the local invariance, the kinetic energy term for $`A^\mu (x)`$ should also be invariant under (82). This is easily accomplished by introducing the field strengths:
$$F^{\mu \nu }(x)=^\mu A^\nu (x)^\nu A^\mu (x)$$
(83)
which are clearly invariant under (82). Hence, the total Lagrangian
$``$ $`=`$ $`\overline{\psi }(x)\left(\gamma ^\mu {\displaystyle \frac{1}{i}}_\mu +m\right)\psi (x)+eA^\mu (x)\overline{\psi }(x)\gamma _\mu \psi (x)`$ (84)
$`{\displaystyle \frac{1}{4}}F^{\mu \nu }(x)F_{\mu \nu }(x)`$
involving the additional gauge field $`A^\mu `$ is locally $`U(1)`$ invariant:
$$(x)\stackrel{a(x)}{}^{}(x)=(x)$$
(85)
when
$`\psi (x)`$ $`\stackrel{a(x)}{}`$ $`\psi ^{}(x)=e^{i\alpha (x)}\psi (x)`$
$`A^\mu (x)`$ $`\stackrel{a(x)}{}`$ $`A^\mu (x)=A^\mu (x)+{\displaystyle \frac{1}{e}}^\mu \alpha (x).`$ (86)
Note that to make the Lagrangian (76) locally $`U(1)`$ invariant it was necessary to introduce an interaction term between the gauge fields $`A^\mu `$ and the globally conserved $`U(1)`$ current $`J_\mu `$. There is a more geometrical way to see how the interaction (81) is necessary to guarantee local invariance. As (79) demonstrates, the reason that the original Lagrangian (76) is not locally invariant is because the derivative of the $`\psi `$ field transforms inhomogeneously under a local $`U(1)`$ rotation. If one could construct a modified derivative, $`D_\mu \psi `$, which under local transformations transformed in the same way that $`_\mu \psi `$ transforms under global transformations, then the original Lagrangian could be trivially made locally invariant by the replacement
$$(_\mu \psi ,\psi )(D_\mu \psi ,\psi ).$$
(87)
Using Eq. (82), it is clear that for the case in question this modified derivativeโ a, so called, covariant derivativeโ is
$$D_\mu \psi =_\mu \psi ieA_\mu \psi ,$$
(88)
since
$`D_\mu \psi \stackrel{a(x)}{}D_\mu ^{}\psi ^{}`$ $`=`$ $`e^{i\alpha }_\mu \psi +i(_\mu \alpha )\psi e^{i\alpha }ieA_\mu \psi e^{i\alpha }`$ (89)
$`i(_\mu \alpha )\psi e^{i\alpha }`$
$`=`$ $`e^{i\alpha (x)}[_\mu \psi ieA_\mu \psi ]=e^{i\alpha (x)}D_\mu \psi .`$
Obviously
$$=\overline{\psi }(x)\left(\gamma ^\mu \frac{1}{i}D_\mu +m\right)\psi (x)\frac{1}{4}F^{\mu \nu }F_{\mu \nu }$$
(90)
is locally $`U(1)`$ invariant and coincides with the expression (84).
Viewed from this perspective, the demand of local invariance of a Lagrangian is a marvelous prescription to fix the interactions of the globally invariant fields with the gauge fields. Furthermore, the gauge transformation (82) does not allow the introduction of a mass term for the $`A^\mu `$ field, since
$$_{\mathrm{mass}}=\frac{1}{2}m_A^2A^\mu (x)A_\mu (x)$$
(91)
breaks the local $`U(1)`$ transformation. So local invariance of a theory severly restricts the dynamics. In the example in question, it will be recognized that the demand that a Dirac field be described by a Lagrangian that is locally $`U(1)`$ invariant has produced the QED Lagrangian! To guarantee local $`U(1)`$ transformations it is necessary to introduce a massless gauge field $`A^\mu (x)`$โ the photon fieldโ interacting with strength eโ the electric chargeโ with the conserved current $`J^\mu `$.
The above simple example can be generalized to theories where the global symmetry group is bigger than the $`U(1)`$ phase symmetry, $`^\mathrm{?}`$ where the structure constants vanish (Abelian group). For these purposes, consider again a Lagrangian density $`(_\mu \chi _\alpha ,\chi _\alpha )`$ composed of fields which transform irreducibly under a non-Abelian group $`G`$ (a group where the structure constants $`c_{ijk}0)`$. Under global $`G`$ transformations, one has
$`\chi _\alpha (x)`$ $`\stackrel{a}{}`$ $`\chi _\alpha ^{}(x)=_{\alpha \beta }(a)\chi _\beta (x)`$ (92)
$`_\mu \chi _\alpha (x)`$ $`\stackrel{a}{}`$ $`_\mu \chi _\alpha ^{}(x)=_{\alpha \beta }(a)_\mu \chi _\beta (x),`$
If this Lagrangian density is invariant under these transformations then
$$(_\mu \chi _\alpha ,\chi _\alpha )\stackrel{a}{}^{}(_\mu \chi _\alpha ^{},\chi _\alpha ^{})=(_\mu \chi _\alpha ,\chi _\alpha ).$$
(93)
Suppose one were able to introduce appropriate gauges fields to construct a covariant derivative, $`D_\mu \chi _\alpha (x)`$, which under local $`G`$ transformations transformed as $`_\mu \chi _\alpha (x)`$ does under global transformations. That is,
$$D_\mu \chi _\alpha (x)\stackrel{a(x)}{}D_\mu ^{}\chi _\alpha ^{}(x)=_{\alpha \beta }(a(x))D_\mu \chi _\beta (x).$$
(94)
Then, clearly, the Lagrangian $`(D_\mu \chi _\alpha ,\chi _\alpha )`$ would be locally $`G`$ invariant
$$(D_\mu \chi _\alpha ,\chi _\alpha )\stackrel{a(x)}{}^{}(D_\mu ^{}\chi _\alpha ^{},\chi _\alpha ^{})=(D_\mu \chi _\alpha ,\chi _\alpha ).$$
(95)
For the theory to be physical, in addition, of course, one must also provide appropriate locally invariant field strengths for the gauge fields entering in the covariant derivatives $`D_\mu \chi _\alpha `$.
By assumption, the covariant derivatives required must transform under local transformations as $`_\mu \chi _\alpha `$ does in Eq. (92). In analogy to what was done for the simple $`U(1)`$ example, it is suggestive to introduce one gauge field $`A_i^\mu `$ for each of the parameters $`\delta a_i`$ of the group $`G`$. After all, the gauge fields are supposed to compensate for the local variations of the fields $`\chi _\alpha `$, and so there should be a gauge field for each of the parameters $`\delta a_i(x)`$ of the Lie group G. Taking the field $`\chi _\alpha (x)`$ transformations under $`G`$ to be those of Eq. (15)
$$[G_i,\chi _\alpha (x)]=(g_i)_{\alpha \beta }\chi _\beta $$
(96)
the $`U(1)`$ example suggest writing for the covariant derivative $`D_\mu \chi _\alpha `$ the expression
$$D_\mu \chi _\alpha (x)=\left[\delta _{\alpha \beta }_\mu ig(g_i)_{\alpha \beta }A_{\mu i}(x)\right]\chi _\beta (x),$$
(97)
where $`g`$ is some coupling constant.
For Eq. (92) to be satisfied for $`D_\mu \chi _\alpha `$, the gauge fields must respond appropriately under local transformations. To determine what this behavior should be, let us compute $`D_\mu ^{}\chi _\alpha ^{}`$ and compare it to what we expect from (92). One has
$`D_\mu ^{}\chi _\alpha ^{}(x)`$ $`=`$ $`_\mu \chi _\alpha ^{}(x)ig(g_i)_{\alpha \beta }A_{\mu i}^{}(x)\chi _\beta ^{}(x)`$ (98)
$`=`$ $`_\mu [_{\alpha \beta }(a(x))\chi _\beta (x)]ig(g_i)_{\alpha \beta }A_{\mu i}^{}(x)_{\beta \gamma }(a(x))\chi _\gamma (x)`$
$`=`$ $`_{\alpha \beta }(a(x))_\mu \chi _\beta (x)+(_\mu _{\alpha \gamma }(a(x))\chi _\gamma (x)`$
$`ig(g_i)_{\alpha \beta }A_{\mu i}^{}(x)_{\beta \gamma }(a(x))\chi _\gamma (x).`$
By definition we want
$`D_\mu ^{}\chi _\alpha ^{}(x)`$ $`=`$ $`_{\alpha \beta }(a(x))D_\mu \chi _\beta (x)`$ (99)
$`=`$ $`_{\alpha \beta }(a(x))_\mu \chi _\beta ig_{\alpha \beta }(a(x))(g_i)_{\beta \gamma }A_{\mu i}(x)\chi _\gamma (x).`$
It follows, therefore, that one must require that
$$ig(g_i)_{\alpha \beta }A_{\mu i}^{}(x)_{\beta \gamma }(a(x))+_\mu _{\alpha \gamma }(a(x))=ig_{\alpha \beta }(a(x))(g_i)_{\beta \gamma }A_{\mu i}(x).$$
(100)
Multiplying the above by $`^1`$ finally gives the transformation required for the gauge field:
$`(g_i)_{\alpha \beta }A_{\mu i}^{}(x)`$ $`=`$ $`{\displaystyle \frac{1}{ig}}[_\mu _{\alpha \gamma }(a(x))][^1(a(x))]_{\gamma \beta }`$ (101)
$`+_{\alpha \gamma }(a(x))(g_i)_{\gamma \delta }[^1(a(x))]_{\delta \beta }A_{\mu i}(x).`$
It is easy to check that this formula agrees with Eq. (82) in the Abelian $`U(1)`$ case when $`=e^{i\alpha },g_i=1`$ and $`g=e`$. In principle, however, Eq. (101) has a very troublesome aspect, since it appears that the transformation properties of the gauge fields $`A_{\mu i}^{}`$ depend on how the field $`\chi _\alpha `$ transforms under $`G`$. If this were to be really the case it would be disastrous, because to obtain a locally invariant theory one would need to introduce a separate compensating gauge field for each matter field in the theory. Fortunately, although (101) as written appears to depend on $``$ explicitly, this dependence is in fact illusory. The transformation properties of gauge fields depend only on the group $`G`$ and not on how the matter fields transform.
To prove this very important point, it is useful to consider Eq. (101) for infinitesimal transformations, where
$$_{\alpha \beta }(\delta a(x))=\delta _{\alpha \beta }+i\delta a_i(g_i)_{\alpha \beta }.$$
(102)
Using the above in (101), and employing an obvious matrix notation, one has
$`g_kA_{\mu k}^{}(x)`$ $`=`$ $`{\displaystyle \frac{1}{ig}}[_\mu (1+i\delta a_k(x)g_k)][1i\delta a_i(x)g_i]`$ (103)
$`+[1+i\delta a_j(x)g_j]g_i[1i\delta a_k(x)g_k]A_{\mu i}(x)`$
$``$ $`g_kA_{\mu k}(x)+i\delta a_j(x)[g_j,g_i]A_{\mu i}(x)+{\displaystyle \frac{1}{g}}[_\mu \delta a_k(x)]g_k.`$
Using the commutation relations for the matrices $`g_i`$
$$[g_j,g_i]=ic_{jik}g_k=ic_{ijk}g_k$$
(104)
it is easy to see that the RHS of (103) is simply proportional to $`g_k`$
$$\mathrm{RHS}=g_k\left[A_{\mu k}(x)+c_{ijk}\delta a_j(x)A_{\mu i}(x)+\frac{1}{g}[_\mu \delta a_k(x)]\right].$$
(105)
Thus, as anticipated, the transformation properties of the gauge fields are independent of the representation matrices $`g_k`$ associated with the fields $`\chi _\alpha (x)`$ and depend only on the structure constants of the group $`c_{ijk}`$:
$$A_{\mu k}^{}(x)=A_{\mu k}(x)+\delta a_j(x)c_{ijk}A_{\mu i}(x)+\frac{1}{g}_\mu (\delta a_k(x)),$$
(106)
For global transformations, where the parameters $`\delta a_k`$ are $`x`$-independent, the last term in (106) does not contribute and the transformatiion of the gauge fields can be written in the standard form one expects for a quantum field:
$$A_{\mu k}^{}(x)=A_{\mu k}(x)+i\delta a_j(\stackrel{~}{g}_j)_{ki}A_{\mu i}(x).$$
(107)
Here the โgeneratorโ matrices appropriate for the gauge fields, $`\stackrel{~}{g}`$, are expressible in terms of the structure constants of the group
$$(\stackrel{~}{g}_j)_{ki}=ic_{ijk}=ic_{jki}.$$
(108)
It is not hard to show (by using the Jacobi identity for $`\stackrel{~}{g}_i,\stackrel{~}{g}_j,`$ and $`\stackrel{~}{g}_k`$) that the matrices $`\stackrel{~}{g}`$ in Eq. (108) indeed obey the group algebra of $`G`$
$$[\stackrel{~}{g}_i,\stackrel{~}{g}_j]=ic_{ijk}\stackrel{~}{g}_k.$$
(109)
The above discussion makes it clear that the gauge fields $`A_i^\mu `$ introduced in the covariant derivative (97) transform according to a special representation of the group $`G`$, the adjoint representation. If $`G`$ has $`n`$ parameters, then the matrices $`\stackrel{~}{g}_i`$ are $`n\times n`$ matrices, whose elements are related to the structure constants $`c_{ijk}`$. Their transformation has no connection with how the matter fields $`\chi _\alpha `$ transform, but is intimately connected with the key parameters of $`G`$, its structure constants.
Having made $``$ locally invariant through the replacement of derivatives by covariant derivatives, it remains to construct the field strengths for the fields $`A_i^\mu `$, so as to be able to incorporate into the theory the kinetic energy terms for the gauge fields. It is easy to check that the naive generalization of the Abelian example
$$\stackrel{~}{F}_k^{\mu \nu }=^\mu A_k^\nu ^\nu A_k^\mu $$
(110)
will not work, since its transformation will still contain derivatives of the parameters $`\delta a_i`$. Indeed, using Eq. (106) one sees that
$`\stackrel{~}{F}_k^{\mu \nu }`$ $`=`$ $`^\mu A_k^\nu ^\nu A_k^\mu =\stackrel{~}{F}_k^{\mu \nu }+\delta a_{\mathrm{}}c_{i\mathrm{}k}\stackrel{~}{F}_i^{\mu \nu }`$ (111)
$`+c_{ijk}\left[(^\mu \delta a_j)A_i^\nu (^\nu \delta a_j)A_i^\mu \right].`$
What one wants to do to obtain the correct field strengths is to augment (110) so as to eliminate altogether the last term in (111). Since this term contains both $`^\mu \delta a_j`$ and $`A_i^\nu `$ in an antisymmetric fashion, one is led, after a bit of reflection, to try the following ansatz for the non-Abelian field strengths:
$$F_k^{\mu \nu }(x)=^\mu A_k^\nu (x)^\nu A_k^\mu (x)+gc_{kij}A_i^\mu (x)A_j^\nu (x).$$
(112)
Let us check that Eq. (112) has the right properties. Using (106), the third term in (112) transforms as
$`gc_{kij}A_i^\mu (x)A_j^\nu (x)`$ $``$ $`gc_{kij}A_i^\mu (x)A_j^\nu (x)`$ (113)
$`=`$ $`gc_{kij}[A_i^\mu (x)+\delta a_{\mathrm{}}(x)c_{m\mathrm{}i}A_m^\mu (x)+{\displaystyle \frac{1}{g}}^\mu \delta a_i(x)]`$
$`\left[A_j^\nu +\delta a_{\mathrm{}}(x)c_{m\mathrm{}j}A_m^\nu (x)+{\displaystyle \frac{1}{g}}^\nu \delta a_j(x)\right]`$
$``$ $`gc_{kij}A_i^\mu (x)A_j^\nu (x)+c_{kij}\left[(^\mu \delta a_i)A_j^\nu +(^\nu \delta a_j)A_i^\mu \right]`$
$`+\delta a_{\mathrm{}}(x)\left[gc_{kij}c_{m\mathrm{}i}A_m^\mu A_j^\nu +gc_{kij}c_{m\mathrm{}j}A_i^\mu A_m^\nu \right].`$
However, making use of the antisymmetry of the structure constants, one has:
$`c_{kij}(^\mu \delta a_i)A_j^\nu `$ $`=`$ $`c_{kji}(^\mu \delta a_j)A_i^\nu =c_{ijk}(^\mu \delta a_j)A_i^\nu `$
$`c_{kij}(^\nu \delta a_j)A_i^\mu `$ $`=`$ $`c_{ijk}(^\nu \delta a_j)A_i^\mu ,`$ (114)
and one sees that the last term in (111) precisely cancels the second term in (113).
It is also not hard to check that the last term in (113) can be written in a much more interesting form by making use of (109). Relabeling dummy indices and using (109) one obtains
$`3\mathrm{r}\mathrm{d}\mathrm{t}\mathrm{e}\mathrm{r}\mathrm{m}`$ $`=`$ $`g\delta a_{\mathrm{}}\left[c_{kij}c_{m\mathrm{}i}A_m^\mu A_j^\nu +c_{kmi}c_{j\mathrm{}i}A_m^\mu A_j^\nu \right]`$ (115)
$`=`$ $`g\delta a_{\mathrm{}}\left[c_{jki}c_{mi\mathrm{}}+c_{mki}c_{ji\mathrm{}}\right]A_m^\mu A_j^\nu `$
$`=`$ $`g\delta a_{\mathrm{}}[\stackrel{~}{g}_j,\stackrel{~}{g}_m]_k\mathrm{}A_m^\mu A_j^\nu `$
$`=`$ $`ig\delta a_{\mathrm{}}c_{jmp}[\stackrel{~}{g}_p]_k\mathrm{}A_m^\mu A_j^\nu `$
$`=`$ $`g\delta a_{\mathrm{}}c_{jmp}c_{pk\mathrm{}}A_m^\mu A_j^\nu `$
$`=`$ $`\delta a_{\mathrm{}}c_{i\mathrm{}k}\left[gc_{imj}A_m^\mu A_j^\nu \right].`$
Using the above one sees that what remains of (113) transforms in precisely the same way as the second term of $`\stackrel{~}{F}_k^{\mu \nu }`$ \[cf. Eq. (111)\].
Putting everything together, one sees that under a local transformation the field strength $`F_k^{\mu \nu }`$ transforms as
$$F_k^{\mu \nu }(x)\stackrel{\delta a(x)}{}F_k^{\mu \nu }(x)=F_k^{\mu \nu }(x)+\delta a_j(x)c_{ijk}F_i^{\mu \nu }(x).$$
(116)
The above is the desired result. Namely, that under local transformations the field strengths should transform as a quantum field which belongs to the adjoint representation of the group. In view of (116), it is easy to show that $`F_k^{\mu \nu }F_{k\mu \nu }`$ is $`G`$-invariant. One has
$`F_k^{\mu \nu }F_{k\mu \nu }F_k^{\mu \nu }F_{k\mu \nu }^{}`$ $`=`$ $`\left(F_k^{\mu \nu }+\delta a_jc_{ijk}F_i^{\mu \nu }\right)\left(F_{k\mu \nu }+\delta a_jc_{ijk}F_{i\mu \nu }\right)`$ (117)
$`=`$ $`F_k^{\mu \nu }F_{k\mu \nu }+\delta a_j\left(c_{ijk}F_i^{\mu \nu }F_{k\mu \nu }+c_{ijk}F_k^{\mu \nu }F_{i\mu \nu }\right)`$
$`=`$ $`F_k^{\mu \nu }F_{k\mu \nu },`$
since the 2nd term vanishes because of the antisymmetry of $`c_{ijk}`$: $`c_{ijk}=c_{jik}`$.
Let us recapitulate our results. The Lagrangian density $`(_\mu \chi _\alpha ,\chi _\alpha )`$โ assumed to be invariant under global $`G`$ transformationsโ can be made locally invariant by introducing gauge fields $`A_i^\mu `$, which enter in the covariant derivatives $`D_\mu \chi _\alpha `$ and the field strengths $`F_i^{\mu \nu }`$. The locally invariant Lagrangian density is simply:
$$_{\mathrm{local}}=(D_\mu \chi _\alpha ,\chi _\alpha )\frac{1}{4}F_i^{\mu \nu }F_{\mu \nu i}$$
(118)
and is completely determined from a knowledge of the global invariant Lagrangian $``$
Three remarks are in order:
i) Again, as in the Abelian case, no mass term for the gauge fields $`A_i^\mu `$ are allowed if one wants to preserve the local invariance (106).
ii) The pure gauge Lagrangian
$$=\frac{1}{4}F_i^{\mu \nu }F_{\mu \nu i}$$
(119)
which contains the kinetic energy terms for the gauge fields $`A_i^\mu `$ is already a nonlinear field theory, since $`F_i^{\mu \nu }`$ contains terms quadratic in the gauge fields $`A_i^\mu `$. For the Abelian case, where the structure constants vanish, these nonlinear terms are absent.
iii) Because the gauge fields transform nontrivially under the group $`G`$, as far as global transformations go, the symmetry currents of the full theory given by Eq. (118) now also get a contribution from the gauge fields. That is, one has
$$J_i^\mu =\frac{}{_\mu \chi _\alpha }\frac{1}{i}(g_i)_{\alpha \beta }\chi _\beta +\frac{}{_\mu A_j^\nu }\frac{1}{i}(\stackrel{~}{g}_i)_{jk}A_{\nu k}.$$
(120)
## 4 The Higgs Mechanism
We saw earlier that in the case of global symmetries, these symmetries could be realized either in a Wigner-Weyl or Nambu-Goldstone way, depending on whether the vacuum state was left, or not left, invariant by the group transformations. It is clearly of interest to know what happens in each of these cases when the global symmetry is made local, via the introduction of gauge fields. For the Wigner-Weyl case, nothing very much happens. Besides the various degenerate multiplets of particles of the global symmetry there is now also a degenerate zero mass multiplet of gauge field excitations. In the Nambu-Goldstone case, however, some remarkable things happen. When the global symmetry is gauged, the Goldstone bosons associated with the broken generators disappear and the corresponding gauge fields acquire a mass! This is the celebrated Higgs mechanism. $`^\mathrm{?}`$
To understand this phenomena, it is useful to return to the simple $`U(1)`$ model discussed earlier and see what obtains when one tries to make the $`U(1)`$ global symmetry also a local symmetry of the Lagrangian. Recall that the Lagrangian density of the model was
$$=_\mu \varphi ^{}^\mu \varphi \lambda \left(\varphi ^{}\varphi \frac{1}{2}f\right)^2,$$
(121)
and that the sign of $`f`$ determined whether one had a Wigner-Weyl realization $`(f<0)`$ or a Nambu-Goldstone realization $`(f>0)`$. To make the above Lagrangian locally $`U(1)`$ invariant it suffices to replace $`_\mu \varphi `$ by a covariant derivative $`D_\mu \varphi `$ involving a gauge field $`A_\mu `$, and include in the theory a kinetic energy term for this gauge field.
If under local $`U(1)`$ transformations one assumes that
$`\varphi (x)`$ $``$ $`\varphi ^{}(x)=e^{i\alpha (x)}\varphi (x)`$ (122)
$`A_\mu (x)`$ $``$ $`A_\mu ^{}(x)=A_\mu (x)+{\displaystyle \frac{1}{g}}_\mu \alpha (x),`$
then the covariant derivative
$$D_\mu \varphi (x)=(_\mu igA_\mu )\varphi $$
(123)
clearly transforms just like $`\varphi `$ does
$$D_\mu \varphi (x)D_\mu ^{}\varphi ^{}(x)=e^{i\alpha (x)}(D_\mu \varphi (x)).$$
(124)
Whence the augmented Lagrangian
$$=(D_\mu \varphi )^{}(D^\mu \varphi )\lambda \left(\varphi ^{}\varphi \frac{1}{2}f\right)^2\frac{1}{4}F^{\mu \nu }F_{\mu \nu }$$
(125)
with
$$F_{\mu \nu }=_\mu A_\nu _\nu A_\mu $$
(126)
is clearly locally $`U(1)`$ invariant.
If $`f<0`$, so that the global symmetry is Wigner-Weyl realized, the above Lagrangian is suitable for computation as is. It describes the interaction of a degenerate multiplet of scalar fields ($`\varphi `$ and $`\varphi ^{}`$) both with themselves and with a massless gauge field $`A_\mu `$. These latter interactionsโsince the $`\varphi `$โs are scalar fields and hence have quadratic kinetic energiesโcontain both a linear term in the gauge fields:
$$_{\mathrm{int}}^{(1)}=gA_\mu \left[i(^\mu \varphi ^{})\varphi i\varphi ^{}^\mu \varphi \right]=gA_\mu J^\mu $$
(127)
as well as a quadraticโso called โsea-gullโ termโcontribution:
$$_{\mathrm{sea}\mathrm{gull}}^{(1)}=g^2A^\mu A_\mu \varphi ^{}\varphi .$$
(128)
These interactions follow directly from the gauge invariant replacement $`_\mu \varphi D_\mu \varphi =(_\mu igA_\mu )\varphi `$.
If $`f>0`$, on the other hand, so that the global $`U(1)`$ symmetry is realized in a Nambu-Goldstone way, one must reparametrize the theory in terms of fields with vanishing expectation value (c.f. Eq. (63)). This reparametrization is such that one is computing oscillations around the minimum of the potential $`V(\varphi )`$. That is, one replaces
$$\varphi ^{}\varphi =\frac{f}{2}+\text{quantum fields}.$$
(129)
This necessary shift implies that the seagull term of Eq. (128) gives rise to a mass term for the $`A_\mu `$ field!
$$_{\mathrm{mass}}=g^2\frac{f}{2}A^\mu A_\mu \frac{1}{2}m_A^2A^\mu A_\mu .$$
(130)
If the gauge field acquires mass, it follows that it cannot be purely transverse (like the photon) but must also have a longitudinal polarization component. This extra degree of freedom must come from somewhere. It is not difficult to show that it arises from the dissapearance of the Nambu-Goldstone excitation, which would ordinarily arise from the spontaneous breakdown of the global $`U(1)`$ symmetry.
To check this assertion, it is convenient to reparametrize the field $`\varphi `$, in the case $`f>0`$, in a somewhat different way than that chosen before. \[The physics of the theory is, in fact, independent of the parametrization one chooses, but certain parametrizations are more directly physical. Different choices for $`\varphi `$ are akin to choosing different gauges for $`A_\mu `$.\] Let us write $`\varphi `$ in the following exponential parametrization:
$$\varphi (x)=\frac{1}{\sqrt{2}}[\sqrt{f}+\rho (x)]\mathrm{exp}\left[i\frac{\xi (x)}{\sqrt{f}}\right].$$
(131)
Here $`\rho (x)`$ and $`\xi (x)`$ are real fields, with $`\xi (x)`$โthe phase fieldโbeing connected to the Goldstone boson. This last assertion is easy to understand since $`\xi (x)`$ vanishes altogether from the potential $`V`$, and so obviously cannot have any mass term. One has simply
$$V=\lambda \left(\varphi ^{}\varphi \frac{f}{2}\right)^2=\lambda \left(\frac{\rho ^2}{2}+\sqrt{f}\rho \right)^2,$$
(132)
so that the $`\rho `$ field has a mass
$$m_\rho ^2=2\lambda f$$
(133)
in agreement with the value obtained earlier(cf. (66)).
It is easy to check that the phase field $`\xi `$ enters in the covariant derivative in a trivial way, so that it can also be eliminated from the kinetic energy term by an appropriate gauge choice. Thus, as advertised, the Nambu-Goldstone boson plays no role in the local theory. It is โeatenโ to give mass to the gauge fields. To prove this assertion, let us consider $`D_\mu \varphi `$ when $`\varphi `$ is parametrized as in Eq. (131):
$`D_\mu \varphi `$ $`=`$ $`(_\mu igA_\mu )\varphi =(_\mu igA_\mu ){\displaystyle \frac{1}{\sqrt{2}}}(\sqrt{f}+\rho )\mathrm{exp}\left[i{\displaystyle \frac{\xi }{\sqrt{f}}}\right]`$
$`=`$ $`{\displaystyle \frac{\mathrm{exp}\left[i\frac{\xi }{\sqrt{f}}\right]}{\sqrt{2}}}\left[_\mu \rho ig(\sqrt{f}+\rho )\left\{A_\mu {\displaystyle \frac{1}{g\sqrt{f}}}_\mu \xi \right\}\right].`$
Obviously the factor in front of the $`[]`$ bracket in (134) involving $`\mathrm{exp}\left[i\xi /\sqrt{f}\right]`$ will not appear in the Lagrangian (125), since the Lagrangian involves $`(D_\mu \varphi )^{}(D^\mu \varphi )`$. Furthermore the $`\xi `$ dependence in the curly bracket is also spurious, since it can be eliminated via a gauge transformation of the field $`A^\mu `$
$$A^\mu B^\mu =A^\mu \frac{1}{g}^\mu \frac{\xi }{\sqrt{f}}.$$
(135)
If the $`U(1)`$ global symmetry is spontaneously broken $`(f>0)`$ the Lagrangian (125) can be rewritten entirely in terms of a massive vector field $`B^\mu `$ and a massive real scalar field $`\rho `$. The resulting Lagrangian
$``$ $`=`$ $`{\displaystyle \frac{1}{2}}_\mu \rho ^\mu \rho {\displaystyle \frac{1}{2}}m_\rho ^2\rho ^2{\displaystyle \frac{1}{4}}F^{\mu \nu }F_{\mu \nu }{\displaystyle \frac{1}{2}}m_A^2B^\mu B_\mu `$
$`g^2\left(\sqrt{f}\rho +{\displaystyle \frac{1}{2}}\rho ^2\right)B^\mu B_\mu \lambda \left(\sqrt{f}\rho ^3+{\displaystyle \frac{1}{4}}\rho ^4\right)`$
where
$$m_\rho ^2=2\lambda f;m_A^2=g^2f$$
(137)
shows no explicit traces of the original $`U(1)`$ symmetry, except that certain parameters in the interactions have particular interrelations. I remark that, although we demonstrated the absorption of the Goldstone boson to produce a massive gauge field only in the Abelian case, this same phenomenon also occurs in the non-Abelian case.
Let me close this section by discussing the two versions of the model \[Wigner-Weyl $`f<0`$; Nambu-Goldstone $`f>0`$\] in terms of the degrees of freedom present in the theory. In the Wigner-Weyl case the theory has a complex scalar field $`\varphi `$ (2 degrees of freedom) plus a massless gauge field $`A^\mu `$ (2 degrees of freedom, corresponding to the two transverse polarizations). In the Nambu-Goldstone case in the theory there is a real scalar field $`\rho `$ (1 degree of freedom) plus a massive spin 1 field $`B_\mu `$ (3 degrees of freedom). Clearly both versions of the theory have the same number of degrees of freedom. However the spectrum of the excitations is completely different!
## 5 The Structure of Quantum Chromodynamics
As a first illustration, I want to describe very briefly the structure of Quantum Chromodynamics (QCD), the theory that describes the strong interactions. $`^\mathrm{?}`$ As we shall see, although QCD is a local gauge theory realized in a Wigner-Weyl way, it possesses also a set of approximate global symmetries. It turns out that some of these global symmetries are realized in a Wigner-Weyl way, while some others are realized in a Nambu-Goldstone manner. Hence, QCD provides a nice practical example of the more formal considerations we have discussed up to now.
We know in nature of the existence of six different typesโflavorsโof quarks: $`u,d,s,c,b`$, and $`t`$. Each flavor of quark is actually a triplet of fields, since the quarks transform irreducibly under the $`SU(3)`$ symmetry group that characterizes QCD. This $`SU(3)`$ symmetry is a local symmetry, so besides quarks in QCD one must introduce the $`3^21=8`$ gauge fields which are associated with the local $`SU(3)`$ symmetry. These 8 gauge fields are known as gluons, since they help bind quarks into hadronsโ like protons and $`\pi `$-mesons.
Let $`q_\alpha ^f(x)`$ stand for a quark field, with the index $`f`$ denoting the various flavors $`f=\{u,d,s,c,b,t\}`$ and $`\alpha =\{1,2,3\}`$ being an $`SU(3)`$ index. Under local infinitesimal $`SU(3)`$ transformation then one has:
$$q_\alpha ^f(x)q_\alpha ^f(x)=\left[\delta _{\alpha \beta }+i\delta a_i(x)\left(\frac{\lambda _i}{2}\right)_{\alpha \beta }\right]q_\beta ^f(x).$$
(138)
In the above, the $`\lambda _i`$ matrices $`i=1,\mathrm{},8`$ are the $`3\times 3`$ Gell-Mann matrices $`^\mathrm{?}`$ transforming as the 3 representation of $`SU(3)`$. The $`SU(3)`$ structure constantsโdenoted here by $`f_{ijk}`$โare easily found by using the explicit form of the $`\lambda `$-matrices given below:
$$\begin{array}{ccc}\lambda _1=\left[\begin{array}{ccc}0& 1& 0\\ 1& 0& 0\\ 0& 0& 0\end{array}\right];& \lambda _2=\left[\begin{array}{ccc}0& i& 0\\ i& 0& 0\\ 0& 0& 0\end{array}\right];& \lambda _3=\left[\begin{array}{ccc}1& 0& 0\\ 0& 1& 0\\ 0& 0& 0\end{array}\right];\\ & & \\ \lambda _4=\left[\begin{array}{ccc}0& 0& 1\\ 0& 0& 0\\ 1& 0& 0\end{array}\right];& \lambda _5=\left[\begin{array}{ccc}0& 0& i\\ 0& 0& 0\\ i& 0& 0\end{array}\right];& \lambda _6=\left[\begin{array}{ccc}0& 0& 0\\ 0& 0& 1\\ 0& 1& 0\end{array}\right];\\ & & \\ \lambda _7=\left[\begin{array}{ccc}0& 0& 0\\ 0& 0& i\\ 0& i& 0\end{array}\right];& \lambda _8=\frac{1}{\sqrt{3}}\left[\begin{array}{ccc}1& 0& 0\\ 0& 1& 0\\ 0& 0& 2\end{array}\right].& \end{array}$$
(139)
One has
$$[\frac{\lambda _i}{2},\frac{\lambda _j}{2}]=if_{ijk}\frac{\lambda _k}{2}.$$
(140)
The gauge fields $`A_k^\mu (x)`$ under a local infinitesimal $`SU(3)`$ transformation rotate into each other with coefficients proportional to the structure functions $`f_{ijk}`$ and shift by the gradient of the $`SU(3)`$ parameters $`\delta a_k(x)`$:
$$A_k^\mu (x)A_k^\mu (x)=A_k^\mu (x)+\delta a_jf_{ijk}A_i^\mu (x)+\frac{1}{g}^\mu \delta a_k(x).$$
(141)
The field strengths
$$F_i^{\mu \nu }=^\mu A_i^\nu ^\nu A_i^\mu +gf_{ijk}A_j^\mu A_k^\nu $$
(142)
transform in the same way as the $`A_k^\mu `$ fields do but have no inhomogeneous contribution proportional to derivatives of $`\delta a_k(x)`$. Finally, the covariant derivatives of the quark fields
$$D_{\alpha \beta }^\mu q_\beta ^f=\left[^\mu \delta _{\alpha \beta }ig\left(\frac{\lambda _i}{2}\right)_{\alpha \beta }A_i^\mu \right]q_\beta ^f$$
(143)
transform under local $`SU(3)`$ transformations precisely as the quark fields themselves do.
Using the above equations, it is easy to see that the QCD Lagrangian
$$_{\mathrm{QCD}}=\underset{f}{}\overline{q}_\alpha ^f\left[\gamma ^\mu \frac{1}{i}(D_\mu )_{\alpha \beta }+m_f\delta _{\alpha \beta }\right]q_\beta ^f\frac{1}{4}F_i^{\mu \nu }F_{i\mu \nu }$$
(144)
is locally $`SU(3)`$ invariant. In the above, the parameters $`m_f`$ are mass terms for each flavor $`f`$ of quarks. If these terms were absent, that is if one could set $`m_f0`$, it is clear that the QCD Lagrangian has a large global symmetry in which quarks of one flavor are changed into quarks of another flavor. For six flavors of quarks, it is not difficult to show that, in the limit $`m_f0`$, the QCD Lagrangian is invariant under a $`U(6)\times U(6)`$ group of global transformations.
Physically, it turns out that whether one can, or one cannot, approximately neglect the quark mass terms $`m_f`$ depends on whether the mass $`m_f`$ is much smaller, or much greater, than the dynamical scale, $`\mathrm{\Lambda }_{\mathrm{QCD}}`$, associated with QCD. This latter scale is of order 300 MeV which is, in fact, much greater than both the $`u`$\- and $`d`$-quark masses. Although these, so called light quarks have masses much smaller than $`\mathrm{\Lambda }_{\mathrm{QCD}}`$:
$$m_{u,d}\mathrm{\Lambda }_{\mathrm{QCD}},$$
(145)
it turns out that $`m_s\mathrm{\Lambda }_{\mathrm{QCD}}`$, while $`\mathrm{\Lambda }_{\mathrm{QCD}}`$ is much smaller than the masses of the $`c`$-, $`b`$\- and $`t`$-quarks. For this reason, in what follows, I will consider only the QCD piece of the Lagrangian involving the $`u`$\- and $`d`$-light quarks. This, of course, is particularly interesting since these quarks are the ones that make up ordinary hadrons, like the proton, neutron and the pions.
The QCD Lagrangian for this 2-flavor case, if we neglect for the moment altogether $`m_u`$ and $`m_d`$, reads
$$_{\mathrm{QCD}}^{2\mathrm{flavor}}|_{m_u=m_d=0}=\overline{u}_\alpha \gamma ^\mu \frac{1}{i}[D_\mu ]_{\alpha \beta }u_\beta \overline{d}_\alpha \gamma ^\mu \frac{1}{i}[D_\mu ]_{\alpha \beta }d_\beta \frac{1}{4}F_i^{\mu \nu }F_{i\mu \nu }.$$
(146)
Let us organize the $`u`$\- and $`d`$-quarks into a doublet
$$Q_\alpha =\left(\begin{array}{c}u\\ d\end{array}\right)_\alpha ,$$
(147)
then the above Lagrangian can be written simply as
$$_{\mathrm{QCD}}^{2\mathrm{flavor}}|_{m_u=m_d=0}=\overline{Q}_\alpha \gamma ^\mu \frac{1}{i}[D_\mu ]_{\alpha \beta }Q_\beta \frac{1}{4}F_i^{\mu \nu }F_{i\mu \nu }.$$
(148)
It is easy to check that this Lagrangian is invariant under a global $`U(2)_V\times U(2)_A`$ symmetry in which
$`Q\stackrel{V}{}Q^{}=\mathrm{exp}\left[i\alpha _i^VT_i\right]Q`$ ; $`\overline{Q}\stackrel{V}{}\overline{Q}^{}=\overline{Q}\mathrm{exp}\left[i\alpha _i^VT_i\right]`$
$`Q\stackrel{A}{}Q^{}=\mathrm{exp}\left[i\alpha _i^AT_i\gamma _5\right]Q`$ ; $`\overline{Q}^{}\stackrel{A}{}\overline{Q}^{}=\overline{Q}\mathrm{exp}\left[i\alpha _i^AT_i\gamma _5\right],`$ (149)
where the four $`2\times 2`$ matrices $`T_i`$ are just
$$T_i=\{\tau _i,1\},$$
(150)
with $`\tau _i`$ the Pauli matrices. The invariance under $`U(2)_V`$ is trivial to see. That under $`U(2)_A`$ follows once one realizes that the Dirac matrix $`\gamma _5`$ anticommutes with all the $`\gamma `$-matrices: $`\{\gamma _5,\gamma ^\mu \}=0`$.
I should remark that even in the case when one restores the light quark masses in the Lagrangian, $`m_u0,m_d0`$, the QCD Lagrangian is invariant under a common phase transformation of the $`u`$\- and $`d`$-quark fields. This invariance just corresponds to the phase invariance associated with the overall quark number, or baryon number, with $`U(1)_VU(2)_V`$. If $`m_u=m_d`$, then it is easy to show that also the remaining $`SU(2)`$ subgroup in $`U(2)_V`$, $`SU(2)_VU(2)_V`$, is conserved. This subgroup is just the usual isospin, well known from nuclear physics. $`^\mathrm{?}`$ Notice, however, that isospin is an approximate symmetry of QCD even if $`m_um_d`$, provided that the absolute value of these masses is much less than $`\mathrm{\Lambda }_{\mathrm{QCD}}`$. In this case, to a good approximation one can neglect both $`m_u`$ and $`m_d`$ (even if they are not equal!) and the strong interactions are then invariant under the $`SU(2)`$ isospin group.
The approximate $`U(2)_V\times U(2)_A`$ invariance of the strong interactions was discovered in the 1960โs even before QCD was put forth as the theory of the strong interactions. $`^\mathrm{?}`$ It was realized then, however, that while the $`U(2)_V`$ global symmetry appeared to be realized in nature as a Wigner-Weyl symmetry, the $`U(2)_A`$ symmetry was realized in a Nambu-Goldstone way. Indeed, if $`U(2)_A`$ were an approximate Wigner-Weyl symmetry, not only would one expect a degenerate neutron-proton doublet but also one should have another doublet of states, of opposite parity, approximately degenerate with the neutron-proton doublet. Because this additional degenerate doublet was not seen in the spectrum of baryons, this approximate $`U(2)_A`$ symmetry must be spontaneously broken. In this case, one would expect some (nearly) massless Nambu-Goldstone states to appear in the theory. The triplet of pions $`(\pi ^+,\pi ^{},\pi ^o)`$, which are much lighter than any other meson states, were suggested as the likely candidate for these approximately Nambu-Goldstone states. Indeed, one can show that, dynamically, these states really behave as approximate Nambu-Goldstone states should. For instance, at low energy their couplings vanish linearly with energy.
Matters were clarified further with the advent of QCD, since one was able to understand better both the origin of the approximate symmetry and the mechanism which causes the breakdown of $`U(2)_V\times U(2)_A`$. Let me briefly comment on this last point. In QCD, because of the same strong forces that confine quarks into hadrons, condensates of $`u`$\- and $`d`$-quarks can form. These condensates are nothing but non-zero expectation values of quark bilinears in the QCD vacuum. Clearly if
$$0|\overline{u}(0)u(0)|0=0|\overline{d}(0)d(0)|00,$$
(151)
although $`U(2)_V\times U(2)_A`$ is an (approximate) symmetry of the QCD Lagrangian, only $`U(2)_V`$ remains as a true symmetry of the spectrum. That is, the above condensates breaks
$$U(2)_V\times U(2)_AU(2)_V$$
(152)
Naively one would expect as a result of the above spontaneous breakdown that four Nambu-Goldstone bosons should appear in the theory. In fact, the $`U(1)_A`$ subgroup of the $`U(2)_A`$ group, although it is a symmetry at the Lagrangian level, can be shown not to be a real quantum symmetry of QCD. $`^\mathrm{?}`$ Radiative effects cause the divergence of the $`U(1)_A`$ current not to vanish. Unfortunately, the argument why the $`U(1)_A`$ symmetry acquires an anomalous divergenceโ a, so called, chiral anomaly $`^\mathrm{?}`$โ is too complex to enter upon here. Nevertheless, taking this result at face value, one expects that the formation of the condensates above should produce 3 Nambu-Goldstone bosons, associated with the breakdown of the $`SU(2)_A`$ symmetry. These states are the pions. Indeed, one can show that the pion mass attains a finite value once one turns on the $`u`$\- and $`d`$-quark masses, but vanishes in the limit as $`m_u`$, $`m_d0`$$`^\mathrm{?}`$ I will not pursue this point further here, but note only how simply one can understand the approximate symmetry properties of the strong interactions, deduced in the 1960s after much hard work, $`^\mathrm{?}`$ directly from the QCD Lagrangian and a few dynamical assumptions, Eqs. (145) and (151).
## 6 The Structure of the $`SU(2)\times U(1)`$ Theory
The ideas we have just discussed of a spontaneously broken gauge theory have found a spectacular application in the $`SU(2)\times U(1)`$ model of the electroweak interactions of Glashow, Salam and Weinberg. $`^\mathrm{?}`$ At first sight, it appears that weak and electromagnetic interactions have little in common, so that a combined gauge model of these forces does not appear very natural. However, there were at least two phenomenological similarities which hinted at a common link, and which helped in the formulation of the $`SU(2)\times U(1)`$ model.
The first of these similarities is that in both weak and electromagnetic interactions currents are involved. In the electromagnetic case the interaction Lagrangian
$$_{\mathrm{em}}=eA^\mu J_\mu ^{\mathrm{em}}$$
(153)
gives rise to long-range forces between charged particles due to the exchange of a massless photon field. The $`1/r`$ potential between charged particles follows from the $`1/q^2`$ propagator for the photon field. The effective action among charged particles due to (153) is simply
$`W_{\mathrm{em}}^{\mathrm{eff}}`$ $`=`$ $`{\displaystyle \frac{i}{2}}{\displaystyle d^4xeJ_{\mathrm{em}}^\mu (x)T(A_\mu (x)A_\nu (y))d^4yeJ_{\mathrm{em}}^\nu (y)}`$
$`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle d^4xeJ_{\mathrm{em}}^\mu (x)D_{\mu \nu }(xy)d^4yeJ_{\mathrm{em}}^\nu },`$
where $`D_{\mu \nu }`$ is the photon propagator. Since the currents $`J_\mu ^{\mathrm{em}}`$ are conserved, one can take effectively
$$D_{\mu \nu }(xy)=\eta _{\mu \nu }\frac{d^4q}{(2\pi )^4}e^{iq(xy)}\frac{1}{q^2iฯต},$$
(155)
where
$$\eta _{\mu \nu }=\left(\begin{array}{cccc}1& & & \\ & 1& & \\ & & 1& \\ & & & 1\end{array}\right)$$
is the metric tensor. Hence
$`W_{\mathrm{eff}}^{\mathrm{em}}`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle d^4xd^4y\frac{d^4q}{(2\pi )^4}eJ_{\mathrm{em}}^\mu (x)e^{iq(xy)}\frac{1}{q^2iฯต}eJ_\mu ^{\mathrm{em}}(y)}`$
$`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \frac{d^4q}{(2\pi )^4}\left[eJ_{\mathrm{em}}^\mu (q)\frac{1}{q^2iฯต}eJ_\mu ^{\mathrm{em}}(q)\right]}.`$
Thus, in momentum space, one has simply
$$_{\mathrm{eff}}^{\mathrm{em}}(q)=\frac{1}{2}\left[eJ_{\mathrm{em}}^\mu (q)\frac{1}{q^2}eJ_\mu ^{\mathrm{em}}(q)\right].$$
(157)
For the charged current weak interactions, which are responsible for the rather long lived nuclear disintegrations, like neutron $`\beta `$ decay, one has known for a long time that they could be described by an effective current-current theory, the Fermi theory: $`^\mathrm{?}`$
$$_{\mathrm{Fermi}}=\frac{G_F}{\sqrt{2}}J_+^\mu (x)J_\mu (x).$$
(158)
Here $`G_F`$โthe Fermi constantโhas dimensions of $`(\mathrm{mass})^2`$ and $`G_F10^5`$ (GeV)<sup>-2</sup>. In momentum space (159) looks like the e.m. case, except that the photon propagator $`1/q^2`$ is replaced by the constant $`G_F/\sqrt{2}`$. In momentum space, one has
$$_{\mathrm{eff}}^{\mathrm{cc}}(q)=\frac{G_F}{\sqrt{2}}\left[J_+^\mu (q)J_\mu (q)\right].$$
(159)
This phenomenological resemblance can be sharpened by imagining that the contact nature of the charged current weak interactions is due to the exchange of a very heavy โweak bosonโ. For low momentum transfer processes, the propagator of the weak boson would be effectively constant
$$\frac{1}{q^2+M_W^2}\stackrel{q^2M_W^2}{}\frac{1}{M_W^2}.$$
(160)
So Eq. (159) could arise from an interaction Lagrangian very similar to that of electromagnetism: $`^\mathrm{?}`$
$$_{\mathrm{weak}}=\stackrel{~}{g}[J_+^\mu (x)W_\mu (x)+J_{}^\mu (x)W_{+\mu }(x)]$$
(161)
involving some spin one bosons $`W_\pm ^\mu `$. Then one could obtain, for $`q^2M_W^2`$, $`_{\mathrm{eff}}^{\mathrm{cc}}`$ from the exchange of these massive fields.
$$_{\mathrm{eff}}^{\mathrm{cc}}(q)\stackrel{q^2M_W^2}{}\frac{\stackrel{~}{g}^2}{M_W^2}\left[J_+^\mu (q)J_\mu (q)\right],$$
(162)
which identifies the Fermi constant as
$$\frac{G_F}{\sqrt{2}}=\frac{\stackrel{~}{g}^2}{M_W^2}.$$
(163)
Note that if $`\stackrel{~}{g}^2e^2`$ then from the value of $`G_F`$ one infers that the masses of the weak bosons are really heavy: $`M_W100`$ GeV!
The second similarity between weak and electromagnetic processes is that the charged currents that enter in weak decays appear to be related to the electromagnetic currentโat least as far as the strongly interacting particles go. This interrelation was discussed long ago by Feynman and Gell-Mann, and by Marshak and Sudarshan. $`^\mathrm{?}`$ The vector piece of the $`J_\pm ^\mu `$ currents are identical to the $`1i2`$ components of the strong isospin current. In turn the isovector piece of the electromagnetic current is the 3rd component of this same strong isospin current.
Although the above two points hint at a possible common origin of weak and electromagnetic interactions, they are not per se compelling. The dominant reason for attempting to treat both interactions on the same footing is theoretical. The Fermi theory (159) is actually a very sick theory as it stands, since in higher order in perturbation theory one encounters divergences which one cannot eliminate from the theory. These divergences occur because of the very singular nature of the contact interaction (159) which, in contrast to what happens in QED, is not being damped at all for large $`q^2`$.
It turns out that matters are not ameliorated even if the Fermi theory is replaced by an interaction like (162), involving mediating heavy vector bosons $`W_\pm ^\mu `$. This is because the propagator for such a massive boson contains in the numerator a propagator factor, characteristic of a spin one object, which is badly behaved at large $`q^2`$:
$`\mathrm{\Delta }_{\mu \nu }(q)`$ $`=`$ $`{\displaystyle \frac{1}{q^2+M_W^2}}\left(\eta _{\mu \nu }+{\displaystyle \frac{q_\mu q_\nu }{M_W^2}}\right)`$
$`\stackrel{q^\mu \mathrm{large}}{}`$ $`O(1).`$
Thus it is not possible to add โby handโ an interaction like (162) and hope to obtain a sensible weak interaction theory. If, however, the interaction (162) resulted from making a global symmetry localโ so that the $`W_\pm ^\mu `$ are gauge fields which are massive because of the Higgs mechanism โ then the situation is vastly improved. It turns out that the gauge invariance of the theory allows one to calculate higher order corrections with propagators for the $`W`$-fields which have only the $`\eta _{\mu \nu }`$ term. These theories, as first shown by โt Hooft, $`^\mathrm{?}`$ have the same good asymptotic behavior as QED. They are renormalizable.
The above argues for a theory of the weak interactions based on some symmetry group $`G`$ which spontaneously breaks down. Two of the currents associated with $`G`$ must include $`J_+^\mu `$ and $`J_{}^\mu `$. However, the generator algebra must close and so one expects naturally also some neutral current. This current, in general, will be related to the electromagnetic current. Thus we see that renormalizability has lead us directly to contemplate models in which at the Lagrangian level, weak and electromagnetic currents enter on the same footing!
The simplest unified model of the electroweak interactions, which contains $`J_+^\mu ,J_{}^\mu `$ and $`J_{\mathrm{em}}^\mu `$ is based on the group $`O(3)`$$`^\mathrm{?}`$ However, the discovery of weak neutral current processes experimentally argued for at least a 4-parameter group. The suggestion of Glashow, Salam and Weinberg, made well before the discovery of these neutral currents processes, was that the electroweak interactions are based on an $`SU(2)\times U(1)`$ gauge theory, which suffers spontaneous breakdown to $`U(1)_{\mathrm{em}}`$. This theory has three massive gauge bosons, associated with the broken generators, and a massless gauge field, associated with the photon. The model was built to reproduce the known structure of the charged current weak interactions. It then predicted particular neutral current interactions, whose experimental verification provided a direct test of the model. Furthermore, the model also predicts the masses of the gauge fields associated with the spontaneous breakdown. The observation at CERN of the $`W^\pm `$ and $`Z^o`$ bosons, $`^\mathrm{?}`$ with the masses predicted by the model, provided the final experimental confirmation of the validity of the $`SU(2)\times U(1)`$ theory.
To detail the structure of the GSW model, one has to specify how the matter degrees of freedom transform under the $`SU(2)\times U(1)`$ group. This could be deduced from the form of the charged currents $`J_\pm ^\mu `$, which a long series of experiments in the 1950โs and 1960โs showed to have a (V-A) form. $`^\mathrm{?}`$ That is, only the left-handed projection of the fermionic fields appear to participate in these interactions. For instance, from a study of $`\beta `$-decay for the muon one established that the current $`J_+^\mu `$ had both $`\mu \nu _\mu `$ and $`e\nu _e`$ terms, in which only the left-handed neutrino fields entered:
$$J_+^\mu =\overline{e}\gamma ^\mu (1\gamma _5)\nu _e+\overline{\mu }\gamma ^\mu (1\gamma _5)\nu _\mu +\mathrm{}.$$
(165)
Writing the projections
$$\psi =\frac{1}{2}(1\gamma _5)\psi +\frac{1}{2}(1+\gamma _5)\psi =\psi _\mathrm{L}+\psi _\mathrm{R}$$
(166)
and using the properties
$$\{\gamma _5,\gamma ^\mu \}=0;\gamma _5^2=1;\overline{\psi }=\psi ^{}\gamma ^o;\gamma _5^{}=\gamma _5,$$
(167)
one sees that
$$J_+^\mu =2\overline{e}_\mathrm{L}\gamma ^\mu \nu _{e\mathrm{L}}+2\overline{\mu }_\mathrm{L}\gamma ^\mu \nu _{\mu \mathrm{L}}+\mathrm{}.$$
(168)
That is, the charged currents only contain left-handed fields.
The structure of $`J_+^\mu `$, and its complex conjugate $`J_{}^\mu `$, suggests that under $`SU(2)`$ the $`\nu _{e\mathrm{L}}`$ and the $`e_\mathrm{L}`$ fields (and the $`\nu _{\mu \mathrm{L}}`$ and $`\mu _\mathrm{L}`$ fields) transform as a doublet. The appropriate generator matrix for an $`SU(2)`$ doublet is $`\frac{\tau _i}{2}`$, where $`\tau _i`$ are the Pauli matrices. Indeed these matrices obey the $`SU(2)`$ Lie algebra
$$[\frac{\tau _i}{2},\frac{\tau _j}{2}]=iฯต_{ijk}\frac{\tau _k}{2}.$$
(169)
Hence, if $`\left(\begin{array}{c}\nu _e\\ e\end{array}\right)_\mathrm{L}`$ transforms as a doublet, the relevant piece of the $`SU(2)`$ current involving these fields is
$$J_i^\mu =(\overline{\nu }_e\overline{e})_\mathrm{L}\gamma ^\mu \frac{\tau _i}{2}\left(\begin{array}{c}\nu _e\\ e\end{array}\right)_\mathrm{L}$$
(170)
and one sees that indeed
$`2(J_1^\mu iJ_2^\mu )`$ $`=`$ $`2(\overline{\nu }_e\overline{e})_\mathrm{L}\gamma ^\mu \left(\begin{array}{cc}0& 0\\ 1& 0\end{array}\right)\left(\begin{array}{c}\nu _e\\ e\end{array}\right)_\mathrm{L}`$
$`=`$ $`2\overline{e}_\mathrm{L}\gamma ^\mu \nu _{e\mathrm{L}}=(J_+^\mu )_{\nu _ee}`$
The fundamental matter entities presently known are quarks and leptons, which appear in a repetitive pattern as far as the $`SU(2)\times U(1)`$ interactions are concerned \[cf. the $`\nu _ee`$ and $`\nu _\mu \mu `$ terms in $`J_+^\mu `$ of Eq. (166)\]. To date we know of the existence of three generations of quarks and leptons: the electron family: $`(\nu _e,e;u,d)`$; the muon family $`(\nu _\mu ,\mu ;c,s)`$ and the $`\tau `$-lepton family $`(\nu _\tau ,\tau ;t,b)`$, where to each lepton doublet there are associated a pair of quarks. The quarks in the pair actually are comprised each of three states, since each quark carries a color index $`\alpha =1,2,3`$. As we just discussed, these color degrees of freedom are associated with the strong interactions of quarks, which are based on an $`SU(3)`$ gauge theory realized in a Wigner-Weyl wayโ QCD.
Because all the three families transform in the same way under $`SU(2)\times U(1)`$, I will only describe the $`SU(2)\times U(1)`$ properties of the electron family. In view of the preceding discussion, it is clear that $`\left(\begin{array}{c}\nu _e\\ e\end{array}\right)_\mathrm{L}`$ transforms as an $`SU(2)`$ doublet. So does the quark pair $`\left(\begin{array}{c}u\\ d\end{array}\right)_\mathrm{L}`$, as an analysis of beta decay of nuclei indicates. Furthermore, since only left-handed fields enter in the weak charged currents, it must be that the right-handed components of the electron family are $`SU(2)`$ singlets. Since the $`SU(2)\times U(1)`$ group must eventually break down to $`U(1)_{\mathrm{em}}`$, it follows that the electromagnetic charge must be a linear combination of the $`U(1)`$ generator and of the neutral $`T_3`$ generator of $`SU(2)`$, which is diagonal. Thus one can write
$$Q=T_3+Y,$$
(176)
with $`Y`$ being the $`U(1)`$ generator. Hence the $`U(1)`$ quantum numbers of the fields in the electron family follow from their known charges. These considerations allow us to build the following table for the transformation properties of $`\nu _e,e,u`$ and $`d`$ under $`SU(2)\times U(1)`$. The right-handed neutrino field $`\nu _\mathrm{R}`$ in Table 1 is usually not included as a real excitation, since it is a total $`SU(2)\times U(1)`$ singlet and so does not participate in these interactions.
Given the transformation properties of the quarks and leptons under $`SU(2)\times U(1)`$, we may now immediately write down the locally $`SU(2)\times U(1)`$ invariant Lagrangian which describes their interactions. For that purpose we need only to replace in the free Dirac Lagrangian for the fermion fields the ordinary derivatives $`_\mu \psi `$ by the appropriate $`SU(2)\times U(1)`$ covariant derivatives $`D_\mu \psi `$ and add the gauge field interactions. Using Table 1, it is trivial to write down these covariant derivatives. One has
$`D_\mu \left(\begin{array}{c}\nu _e\\ e\end{array}\right)_\mathrm{L}`$ $`=`$ $`\left(_\mu ig{\displaystyle \frac{\tau _i}{2}}W_{\mu i}+ig^{}{\displaystyle \frac{1}{2}}Y_\mu \right)\left(\begin{array}{c}\nu _e\\ e\end{array}\right)_\mathrm{L}`$ (181)
$`D_\mu \left(\begin{array}{c}u\\ d\end{array}\right)_\mathrm{L}`$ $`=`$ $`\left(_\mu ig{\displaystyle \frac{\tau _i}{2}}W_{\mu i}ig^{}{\displaystyle \frac{1}{6}}Y_\mu \right)\left(\begin{array}{c}u\\ d\end{array}\right)_\mathrm{L}`$ (186)
$`D_\mu \nu _\mathrm{R}`$ $`=`$ $`(_\mu )\nu _\mathrm{R}`$ (187)
$`D_\mu e_\mathrm{R}`$ $`=`$ $`(_\mu +ig^{}Y_\mu )e_\mathrm{R}`$ (188)
$`D_\mu u_\mathrm{R}`$ $`=`$ $`(_\mu ig^{}{\displaystyle \frac{2}{3}}Y_\mu )u_\mathrm{R}`$ (189)
$`D_\mu d_\mathrm{R}`$ $`=`$ $`(_\mu +ig^{}{\displaystyle \frac{1}{3}}Y_\mu )d_\mathrm{R}.`$ (190)
Here $`g,g^{}`$ are the $`SU(2)`$ and $`U(1)`$ coupling constants, respectively, while $`W_{\mu i}`$ and $`Y_\mu `$ are the $`SU(2)`$ and $`U(1)`$ gauge fields, respectively.
The Lagrangian for the $`SU(2)\times U(1)`$ model of Glashow, Salam and Weinbergโ as far as the interactions among the fermions of the electron family and the gauge fields goโ is then simply
$`_{\mathrm{FG}}=`$ $``$ $`(\overline{\nu }_e\overline{e})_\mathrm{L}\gamma ^\mu {\displaystyle \frac{1}{i}}D_\mu \left(\begin{array}{c}\nu _e\\ e\end{array}\right)_\mathrm{L}(\overline{u}\overline{d})_\mathrm{L}\gamma ^\mu {\displaystyle \frac{1}{i}}D_\mu \left(\begin{array}{c}u\\ d\end{array}\right)_\mathrm{L}`$
$``$ $`\overline{e}_\mathrm{R}\gamma ^\mu {\displaystyle \frac{1}{i}}D_\mu e_\mathrm{R}\overline{u}_\mathrm{R}\gamma ^\mu {\displaystyle \frac{1}{i}}D_\mu e_\mathrm{R}\overline{d}_\mathrm{R}\gamma ^\mu {\displaystyle \frac{1}{i}}D_\mu d_\mathrm{R}`$
$``$ $`{\displaystyle \frac{1}{4}}W_i^{\mu \nu }W_{\mu \nu i}{\displaystyle \frac{1}{4}}Y^{\mu \nu }Y_{\mu \nu },`$
where the field strengths $`W_i^{\mu \nu }`$ and $`Y^{\mu \nu }`$ are given by
$`W_i^{\mu \nu }`$ $`=`$ $`^\mu W_i^\nu ^\nu W_i^\mu +gฯต_{ijk}W_j^\mu W_k^\nu `$ (196)
$`Y^{\mu \nu }`$ $`=`$ $`^\mu Y^\nu ^\nu Y^\mu .`$ (197)
Note that the Lagrangians (180) contains no mass terms for the fermion fields. Mass terms involve a left-right transition
$$_{\mathrm{mass}}=m\overline{\psi }\psi =m(\overline{\psi }_\mathrm{L}\psi _\mathrm{R}+\overline{\psi }_\mathrm{R}\psi _\mathrm{L}).$$
(198)
Since under $`SU(2)`$ $`\psi _\mathrm{L}2`$ and $`\psi _\mathrm{R}1`$, clearly the $`SU(2)\times U(1)`$ symmetry permits no fermion mass terms. As I will show later, however, masses can be generated when $`SU(2)\times U(1)`$ is spontaneously broken down.
Before we discuss the breakdown of $`SU(2)\times U(1)`$ it is useful to organize a bit the interaction terms which emerge from the Lagrangian (180). These take the simple form
$$_{\mathrm{int}}=gW_i^\mu J_{\mu i}+g^{}Y^\mu J_{\mu Y}$$
(199)
where the $`SU(2)`$ and $`U(1)`$ currents, $`J_i^\mu `$ and $`J_Y^\mu `$ are readily seen to be
$`J_i^\mu `$ $`=`$ $`(\overline{\nu }_e\overline{e})_\mathrm{L}\gamma ^\mu {\displaystyle \frac{\tau _i}{2}}\left(\begin{array}{c}\nu _e\\ e\end{array}\right)_\mathrm{L}+(\overline{u}\overline{d})_\mathrm{L}\gamma ^\mu {\displaystyle \frac{\tau _i}{2}}\left(\begin{array}{c}u\\ d\end{array}\right)_\mathrm{L}`$ (204)
$`J_Y^\mu `$ $`=`$ $`{\displaystyle \frac{1}{2}}(\overline{\nu }_e\overline{e})_\mathrm{L}\gamma ^\mu \left(\begin{array}{c}\nu _e\\ e\end{array}\right)_\mathrm{L}+{\displaystyle \frac{1}{6}}(\overline{u}\overline{d})_\mathrm{L}\gamma ^\mu \left(\begin{array}{c}u\\ d\end{array}\right)_\mathrm{L}`$
$`\overline{e}_\mathrm{R}\gamma ^\mu e_\mathrm{R}+{\displaystyle \frac{2}{3}}\overline{u}_\mathrm{R}\gamma ^\mu u_\mathrm{R}{\displaystyle \frac{1}{3}}\overline{d}_\mathrm{R}\gamma ^\mu d_\mathrm{R}.`$
I note that since in the model the electromagnetic current is given by \[cf. Eq. (173)\]
$$J_{\mathrm{em}}^\mu =J_3^\mu +J_Y^\mu $$
(210)
the phenomenological observation mentioned earlier, that the vector piece of the weak charged currents and the isovector piece of $`J_{\mathrm{em}}^\mu `$ are related, is built in already in (187).
It is convenient to rewrite (184) in terms of physical fields. If the model is to reproduce the weak interactions, the $`SU(2)\times U(1)`$ symmetry must suffer a spontaneous breakdown to $`U(1)_{\mathrm{em}}`$. This means that of the four gauge fields $`W_i^\mu ,Y^\mu `$, three must acquire a mass and one will remain massless. Now, in general, $`U(1)_{\mathrm{em}}`$ is a linear combination of an $`U(1)SU(2)`$ and $`U(1)_Y`$, so that one expects the photon fields to be a linear combination of $`W_3^\mu `$ and $`Y^\mu `$. The orthogonal combination then corresponds to a massive neutral fieldโ the $`Z^o`$ boson. It has become conventional to parametrize these linear combinations in terms of an angle $`\theta _W`$โ the Weinberg angle.
$`W_3^\mu `$ $`=`$ $`\mathrm{cos}\theta _WZ^\mu +\mathrm{sin}\theta _WA^\mu `$
$`Y^\mu `$ $`=`$ $`\mathrm{sin}\theta _WZ^\mu +\mathrm{cos}\theta _WA^\mu `$ (211)
It proves useful also to rewrite $`W_1^\mu `$ and $`W_2^\mu `$ in terms of fields of definite charge
$$W_\pm ^\mu =\frac{1}{\sqrt{2}}(W_1^\mu iW_2^\mu )$$
(212)
and use the charged currents $`J_\pm ^\mu `$, which enter in the Fermi theory \[cf. Eq. (159)\]
$$J_\pm ^\mu =2(J_1^\mu iJ_2^\mu ).$$
(213)
With all these definitions the interaction Lagrangian of Eq. (184) becomes
$`_{\mathrm{int}}`$ $`=`$ $`{\displaystyle \frac{g}{2\sqrt{2}}}[W_+^\mu J_\mu +W_{}^\mu J_{+\mu }]`$ (214)
$`\left\{(g\mathrm{cos}\theta _W+g^{}\mathrm{sin}\theta _W)J_3^\mu g^{}\mathrm{sin}\theta _WJ_{\mathrm{em}}^\mu \right\}Z_\mu `$
$`+\left\{g^{}\mathrm{cos}\theta _WJ_{\mathrm{em}}^\mu +(g^{}\mathrm{cos}\theta _Wg\mathrm{sin}\theta _W)J_3^\mu \right\}A_\mu .`$
In the above, I have made use of (187) to eliminate altogether $`J_Y^\mu `$ in favor of $`J_{\mathrm{em}}^\mu `$.
The above interaction is supposed to reproduce both the electromagnetic interaction (153) and the charged current weak interaction (162). It predicts as well a new neutral current weak interaction involving the $`Z^o`$ boson. Since the photon field is supposed to only interact with $`J_{\mathrm{em}}^\mu `$ with strength $`e`$, one sees that one must require the Weinberg angle to obey the unification condition
$$g^{}\mathrm{cos}\theta _W=g\mathrm{sin}\theta _W=e.$$
(215)
Using this information to eliminate $`g`$ and $`g^{}`$ in terms of $`\theta _W`$ and $`e`$ allows one to write for the interaction Lagrangian the expression:
$`_{\mathrm{int}}`$ $`=`$ $`{\displaystyle \frac{e}{2\sqrt{2}\mathrm{sin}\theta _W}}(W_+^\mu J_\mu +W_{}^\mu J_{\mu +})+eJ_{\mathrm{em}}^\mu A_\mu `$ (216)
$`+{\displaystyle \frac{e}{2\mathrm{cos}\theta _W\mathrm{sin}\theta _W}}J_{\mathrm{NC}}^\mu Z_\mu .`$
Here the neutral current $`J_{\mathrm{NC}}^\mu `$ which interacts with the $`Z_\mu `$ field is
$$J_{\mathrm{NC}}^\mu =2[J_3^\mu \mathrm{sin}^2\theta _WJ_{\mathrm{em}}^\mu ].$$
(217)
Comparing this result with our earlier discussion, the coupling $`\stackrel{~}{g}`$ of Eq. (162) is seen to be
$$\stackrel{~}{g}=\frac{e}{2\sqrt{2}\mathrm{sin}\theta _W}.$$
(218)
Hence, the comparison with the Fermi theory \[cf. Eq. (164)\] gives for the Fermi constant the expression
$$\frac{G_F}{\sqrt{2}}=\frac{\stackrel{~}{g}^2}{M_W^2}=\frac{e^2}{8\mathrm{sin}^2\theta _WM_W^2}.$$
(219)
One sees that a knowledge of the Weinberg angleโwhich enters in the neutral currentโgives direct information on the mass of the heavy weak boson which mediates the charged current weak interactions. One finds experimentally that $`\mathrm{sin}^2\theta _W1/4`$$`^\mathrm{?}`$ which predicts for $`M_W`$ a value of around 80 GeV. This prediction has been spectacularly confirmed by the discovery at the CERN Collider of a particle of this mass with all the characteristic of the $`W`$ boson.
Just as charged current interactions, for processes where the momentum transfer $`q^2M_W^2`$, can be described by the Fermi theory, one can arrive at a similar structure for neutral current interactions. In the same approximation, $`q^2M_Z^2`$, one has
$$_{\mathrm{Fermi}}^{\mathrm{NC}}\frac{1}{2}\left[\frac{e}{2\mathrm{sin}\theta _W\mathrm{cos}\theta _W}\right]^2\frac{1}{M_Z^2}J_{\mathrm{NC}}^\mu J_{\mu \mathrm{NC}}.$$
(220)
Using the identification (195) of the Fermi constant, one has
$`_{\mathrm{Fermi}}^{\mathrm{NC}}`$ $`=`$ $`{\displaystyle \frac{G_F}{\sqrt{2}}}\left[{\displaystyle \frac{M_W^2}{M_Z^2\mathrm{cos}^2\theta _W}}\right]J_{\mathrm{NC}}^\mu J_{\mu \mathrm{NC}}`$ (221)
$`=`$ $`{\displaystyle \frac{G_F}{\sqrt{2}}}\rho J_{\mathrm{NC}}^\mu J_{\mu \mathrm{NC}},`$
where the ratio
$$\rho =\frac{M_W^2}{M_Z^2\mathrm{cos}^2\theta _W}$$
(222)
gives the relative strength of neutral to charged current weak processes.
To summarize, the weak interactions in the Glashow Salam Weinberg model, in the limit in which $`q^2M_W^2,M_Z^2`$ can be written in a current-current form
$$_{\mathrm{Weak}}^{\mathrm{eff}}=\frac{G_F}{\sqrt{2}}[J_+^\mu J_\mu +\rho J_{\mathrm{NC}}^\mu J_{\mu \mathrm{NC}}].$$
(223)
The charged current weak interactions by construction agree with experiment. Neutral current weak interactions test the model, since all experiments must be describable by the only two free parameters $`\rho `$ and $`\mathrm{sin}^2\theta _W`$, which enters in the definition of $`J_{\mathrm{NC}}^\mu `$, present in (194). All neutral current experiments, indeed, can be fitted with a common value of $`\mathrm{sin}^2\theta _W1/4`$ and of $`\rho 1`$$`^\mathrm{?}`$ thereby providing strong confirmation of the validity of the GSW model. Furthermore, given $`\rho `$ and $`\mathrm{sin}^2\theta _W`$, one can determine the mass of the $`Z^o`$ and $`W^\pm `$ bosons from Eqs. (196) and (199). The discovery at the CERN collider of the $`W^\pm `$ bosons and, soon thereafter, of a neutral heavy particle of mass around 90 GeV, in agreement with the value predicted by the GSW model, provided a splendid confirmation of the model. $`^\mathrm{?}`$
To complete the GSW model, it is necessary to describe briefly the mechanism by which the $`W^\pm `$ and $`Z^o`$ bosons get mass. The idea here is very much like that described in the last section, when I discussed the Higgs mechanism for the Abelian $`U(1)`$ model. Namely, one introduces some scalar field whose self interactions cause the $`SU(2)\times U(1)`$ symmetry to break down. Since we want $`SU(2)`$ to break down, the scalar field introduced into the theory must carry $`SU(2)`$ quantum numbers. The simplest possibility is afforded by an $`SU(2)`$ doublet. Furthermore, since we want also to break the $`U(1)`$ symmetry, this doublet must be complex. Thus the simplest agent to carry through the desired breakdown is the complex doublet.
$$\mathrm{\Phi }=\left(\begin{array}{c}\varphi ^o\\ \varphi ^{}\end{array}\right),$$
(224)
where $`\varphi ^o`$ and $`\varphi ^{}`$ are complex fields, and the charge assignments identify $`Y_\varphi =1/2`$.
To accomplish the breakdown we consider a potential analogous to that in (58). In addition we must introduce an appropriately $`SU(2)\times U(1)`$ covariant kinetic energy term for the field $`\mathrm{\Phi }`$, using the covariant derivative
$$D_\mu \mathrm{\Phi }=\left(_\mu ig\frac{\tau _i}{2}W_{\mu i}+i\frac{g^{}}{2}Y_\mu \right)\mathrm{\Phi }.$$
(225)
The interaction Lagrangian involving the scalar field $`\mathrm{\Phi }`$โthe Higgs fieldโis just then:
$$_{\mathrm{HG}}=(D_\mu \mathrm{\Phi })^{}(D^\mu \mathrm{\Phi })\lambda \left(\mathrm{\Phi }^{}\mathrm{\Phi }\frac{v^2}{2}\right)^2.$$
(226)
It is clear that the potential term in (203) will cause $`SU(2)\times U(1)`$ to break down. The choice of vacuum expectation value
$$\mathrm{\Phi }=\frac{v}{\sqrt{2}}\left(\begin{array}{c}1\\ 0\end{array}\right)$$
(227)
guarantees that $`SU(2)\times U(1)U(1)_{\mathrm{em}}`$. \[Actually, with just one doublet $`\mathrm{\Phi }`$ one can always define $`U(1)_{\mathrm{em}}`$ as the $`U(1)`$ left unbroken in $`V`$. The choice (204) is dictated by our definition of charge. Any other choice would do, but it would change what we called $`Q`$.\]
Given the vacuum expectation value (204), the mass terms for the gauge fields are read off immediately from the seagull terms:
$$_{\mathrm{seagull}}=\left[\left(g\frac{\tau _i}{2}W_{\mu i}g^{}\frac{1}{2}Y_\mu \right)\mathrm{\Phi }\right]^{}\left[\left(g\frac{\tau _i}{2}W_i^\mu g^{}\frac{1}{2}Y^\mu \right)\mathrm{\Phi }\right].$$
(228)
Using Eq. (188), the gauge field matrix in (205) is easily seen to be
$`g{\displaystyle \frac{\tau _i}{2}}W_i^\mu g^{}{\displaystyle \frac{1}{2}}Y^\mu `$ $`=`$ $`\left[\begin{array}{cc}\frac{g}{2}W_3^\mu \frac{g^{}}{2}Y^\mu & \frac{g}{\sqrt{2}}W_+^\mu \\ \frac{g}{\sqrt{2}}W_{}^\mu & \frac{g}{2}W_3^\mu \frac{g^{}}{2}Y^\mu \end{array}\right]`$ (231)
$`=`$ $`\left[\begin{array}{cc}\frac{g}{2\mathrm{cos}\theta _W}Z^\mu & \frac{g}{\sqrt{2}}W_+^\mu \\ \frac{g}{\sqrt{2}}W_{}^\mu & \frac{g}{2\mathrm{cos}\theta _W}[\mathrm{sin}^2\theta _W\mathrm{cos}^2\theta _W]Z^\mu A^\mu \end{array}\right].`$ (234)
Since the vacuum expectation value (204) only has an upper component, one sees that only $`Z^\mu `$ and not $`A^\mu `$ acquires a mass, confirming our previous identification of this latter field as the photon field. Replacing in (205) $`\mathrm{\Phi }\mathrm{\Phi }`$ gives the following mass terms for the gauge fields:
$$_{\mathrm{mass}}=\left(\frac{gv}{2}\right)^2W_+^\mu W_\mu \frac{1}{2}\left(\frac{gv}{2\mathrm{cos}\theta _W}\right)^2Z^\mu Z_\mu .$$
(235)
Hence
$$M_W^2=\frac{1}{4}(gv)^2;M_Z^2=\frac{1}{4\mathrm{cos}^2\theta _W}(gv)^2.$$
(236)
We see that the simplest choice of Higgs field to give the $`SU(2)\times U(1)U(1)_{\mathrm{em}}`$ breaking predicts that the parameter $`\rho `$ in the neutral current interactions is unity!
$$\rho =\frac{M_W^2}{M_Z^2\mathrm{cos}^2\theta _W}=1.$$
(237)
The experimental indications that $`\rho 1`$ suggest therefore that nature has chosen (again) the simplest course. Using Eq. (208) and the relation of $`M_W`$ to the Fermi constant identifies the scale parameter $`v`$ in the Higgs potential as
$$v=(\sqrt{2}\mathrm{G}_\mathrm{F})^{1/2}250\mathrm{GeV}.$$
(238)
The introduction of a doublet Higgs field $`\mathrm{\Phi }`$ into the theory has another salutary effectโit allows for the possibility of generating masses for the quarks and leptons! I will illustrate the idea with the up quark. Since $`\mathrm{\Phi }`$ carries hypercharge $`Y_\mathrm{\Phi }=1/2`$ and is an $`SU(2)`$ doublet, the interaction (Yukawa interaction) of $`\mathrm{\Phi }`$ with $`u_\mathrm{R}`$ and the $`(\overline{u}\overline{d})_\mathrm{L}`$ doublet is allowed by the $`SU(2)\times U(1)`$ symmetry:
$$_{\mathrm{Yukawa}}=h(\overline{u}\overline{d})_\mathrm{L}\mathrm{\Phi }u_\mathrm{R}h^{}\overline{u}_\mathrm{R}\mathrm{\Phi }^{}\left(\begin{array}{c}u\\ d\end{array}\right)_\mathrm{L}.$$
(239)
Obviously, when $`\mathrm{\Phi }`$ has a vacuum expectation value this interaction will generate a mass term for the $`u`$ quark. Taking $`h`$ real, one has
$$_{\mathrm{mass}}=\frac{hv}{\sqrt{2}}\overline{u}u=m_u\overline{u}u,$$
(240)
so that $`m_u`$ is also related to the breakdown parameter $`v`$. Unfortunately since $`h`$ is not known, no predictions follow. This same mass generation procedure holds for all quarks and leptons.
## Acknowledgements
I am extremely grateful to J. Tran Than Van for having invited me to lecture in the fifth Vietnam School of Physics. The extremely friendly atmosphere of the school and of all the participants made my stay in Hanoi a real pleasure. This work was supported in part by the department of energy under contract No. DE-FG03-91ER40662, Task C.
## References |
warning/0002/hep-ph0002009.html | ar5iv | text | # Stabilization of internal spaces in multidimensional cosmology
## 1 Introduction
Stabilization of additional dimensions near their present day values (dilaton/geometrical moduli stabilization) is one of the main problems for any multidimensional theory because a dynamical behavior of the internal spaces results in a variation of the fundamental physical constants. Observations show that internal spaces should be static or nearly static at least from the time of recombination (in some papers arguments are given in favor of the assumption that variations of the fundamental constants are absent from the time of primordial nucleosynthesis ).
Observations further indicated that Standard Model (SM) matter cannot propagate a large distance in extra dimensions. This allowed for two classes of model building implications:
The first class consists of models with extra spacetime dimensions compactified at scales less the Fermi length $`L_F10^{17}\text{cm}`$ as characteristic scale of the experimentally tested electroweak interaction $`L_FM_{EW}^11\text{TeV}^1`$. Up to the early 1990s this was a standard assumption in string phenomenology with string scale slightly below the 4-dimensional Planck scale $`M_s10^{16}`$ GeV . The question about a concrete mechanism for the stabilization of the compactification scales (moduli stabilization) remained open in this discussion .
The second class of models starts from the assumptions that observable SM matter is confined to a 3-brane located in a higher dimensional bulk spacetime and that gravitational interactions can propagate in the whole bulk spacetime provided that a mechanism exists which ensures usual Newtonโs $`r^2`$ law at distances $`\begin{array}{c}>\hfill \\ \hfill \end{array}1`$ cm accessable to present gravitational tests. The thickness of the 3-brane in this case should be of order of the Fermi length $`L_F`$. The additional bulk dimensions can be compactified or non-compact.
Historically, the first proposal for an interpretation of our appearantly 4-dimensional Universe as a submanifold embedded into a non-compact higher dimensional bulk space dates back to the 1983 work of Rubakov and Shaposhnikov and Akama (still without accounting for gravitational interactions) and Visserโs consideration from 1985 (studying the localization/trapping of particles via gravity to a 4-dimensional submanifold of a 5-dimensional โrealโ world).
Within the framework of superstring theory/M-theory new arguments have been given for a selfcontent embedding of the 4-dimensional $`SU(3)\times SU(2)\times U(1)`$ Standard Model of strong and electroweak interactions into a fundamentally higher dimensional spacetime manifold. For example, in Hoลava-Witten theory one starts from the strongly coupled regime of $`E_8\times E_8`$ heterotic string theory and interprets it as M-theory on an orbifold $`\text{}^{10}\times S^1/\text{}_2`$. After compactification on a Calabi-Yau three-fold one arrives at solutions which may be considered as a pair of parallel 3-branes with opposite tension, and location at the orbifold planes.
After 1995 it became clear from investigations in Type I string theory that due to compactified higher dimensions the string scale $`M_s`$ can be much smaller than the 4-dimensional Planck scale $`M_{Pl(4)}=1.22\times 10^{19}`$ GeV and that it is bounded from below only experimentally by the scale of electroweak interaction $`M_{EW}\begin{array}{c}<\hfill \\ \hfill \end{array}M_s\begin{array}{c}<\hfill \\ \hfill \end{array}M_{Pl(4)}`$ . As suggested by Arkani-Hamed et al it is even possible to lower the fundamental Planck scale $`M_{Pl(4+D^{})}`$ of the $`(4+D^{})`$dimensional theory down to the SM electroweak scale $`M_{Pl(4+D^{})}M_{EW}1`$ TeV. This allows for a solution of the hierarchy problem not relying on supersymmetry or technicolor. In this approach gravity can propagate in all multidimensional bulk space whereas ordinary SM fields are localized on a 3-brane with thickness in the extra dimensions of order the Fermi length $`L_F`$. As a result, the 4-dimensional Planck scale of the external space is connected with the electroweak scale by the relation
$$M_{Pl(4)}^2V_D^{}M_{EW}^{(2+D^{})},$$
(1.1)
where $`V_D^{}`$ is the volume of the internal spaces. Thus, the scale of the internal space compactification is of order
$$aV_D^{}^{1/D^{}}10^{\frac{32}{D^{}}17}\text{cm}.$$
(1.2)
In this model physically acceptable values correspond to $`D^{}2`$, e.g. for $`D^{}=2`$ the internal space scale of compactification is $`a10^1\text{cm}`$.
The stabilization of extra dimensions (geometrical moduli stabilization) in models with sub-millimetre internal spaces was considered in Refs. where the dynamics of the conformal excitations of the internal spaces near minima of an effective potential have been investigated. Due to the product topology of the $`(4+D^{})`$dimensional bulk spacetime constructed from Einstein spaces with scale (warp) factors depending only on the coordinates of the external 4-dimensional component, the conformal excitations have the form of massive scalar fields living in the external space. Within the framework of multidimensional cosmological models (MCMs) we investigated such excitations in and called them gravitational excitons. Later, since the submillimetre weak-scale compactification approach these geometrical moduli excitations are known as radions .
Recently Randall and Sundrum proposed an interesting construction for the solution of the hierarchy problem localizing low energy SM matter as well as low energy gravity on a 3-brane in a slice of anti-de Sitter space $`AdS_5`$. Subsequently, it has been shown that such a localization of low energy physics can be also achieved at the intersection of a system of $`(n+2)`$branes in $`AdS_{4+n}`$ allowing for an interpretation of our observable universe e.g. as a defect in a higher dimensional brane crystal . But the RS proposal and its generalizations are not considered in the present paper.
The main goal of our present comments consists in a clarification of conditions which ensure the stabilization of the internal spaces in multidimensional models with a minimal coupled scalar field as a matter source (section 2). A general method for the solution of such problems in models with an arbitrary number of internal spaces was proposed in Ref. . There it was shown that the problem of the internal space stabilization can be solved most easily in the Einstein frame (although it is clear that if stabilization takes place in the Einstein frame it will also take place in the Brans-Dicke frame and vice versa). Our investigations (see also ) show that inflation of the external space which was maintained for some models in earlier Refs. (see e.g. ) is destroyed by a required stabilization of the internal spaces.
On the other hand there are also papers devoted to inflation where stabilization of the internal spaces was supposed a priori (see e.g. ). We would like to stress here that it is necessary to be rather careful in this case because stabilization can destroy inflation. For example, the appearance of a negative effective cosmological constant, which in some models is a necessary condition for the internal spaces stabilization, can either destroy inflation at all or make problematic its succesfull completion. This situation occures e.g. in the simple toy model which we consider in section 3 of the present paper. We use this model in order to show exactly under which conditions stabilization takes place in multidimensional cosmological models with a minimal coupled scalar field and to discuss briefly a possibility for inflation in these models.
In the present paper most of the calculations are performed in the electroweak fundamental scale approach. In the conclusion section 4 we compare the corresponding results with those for the Planck fundamental scale approach and show that the transition from one approach to the other results in a rescaling of the effective cosmological constant $`\mathrm{\Lambda }_{eff}`$ as well as of gravitational exciton masses $`m_i`$. The corresponding rescaling prefactors which appear due to the transition (see eq. (4.3)) lead to a different functional dependence of $`\mathrm{\Lambda }_{eff}`$ and $`m_i`$ on the compactification sizes of the internal spaces in the two approaches. As result, in the Planck fundamental scale approach the values of $`\mathrm{\Lambda }_{eff}`$ and $`m_i`$ can be much smaller than in the electroweak approach. Finally, we discuss some bounds on the parameters of the model which follow from observable cosmological data. These bounds strongly depend on the details of the behavior of the inflaton and gravitational exciton fields after inflation, e.g. on the times of their reheating and decay.
## 2 Stabilization of the internal spaces
We consider a cosmological model with metric
$$g=g^{(0)}+\underset{i=1}{\overset{n}{}}e^{2\beta ^i(x)}g^{(i)},$$
(2.1)
which is defined on a manifold with product topology
$$M=M_0\times M_1\times \mathrm{}\times M_n,$$
(2.2)
where $`x`$ are some coordinates of the $`D_0=(d_0+1)`$ \- dimensional manifold $`M_0`$ and
$$g^{(0)}=g_{\mu \nu }^{(0)}(x)dx^\mu dx^\nu .$$
(2.3)
Let manifolds $`M_i`$ be $`d_i`$ \- dimensional Einstein spaces with metric $`g^{(i)},`$ i.e.
$$R_{mn}\left[g^{(i)}\right]=\lambda ^ig_{mn}^{(i)},m,n=1,\mathrm{},d_i$$
(2.4)
and
$$R\left[g^{(i)}\right]=\lambda ^id_iR_i.$$
(2.5)
In the case of constant curvature spaces parameters $`\lambda ^i`$ are normalized as $`\lambda ^i=k_i(d_i1)`$ with $`k_i=\pm 1,0`$. Later on we shall not specify the structure of the spaces $`M_i`$. We require only $`M_i`$ to be compact spaces with arbitrary sign of curvature.
With total dimension $`D=D_0+_{i=1}^nd_i`$, $`\kappa _D^2`$ a $`D`$dimensional gravitational constant, $`\mathrm{\Lambda }`$ \- a $`D`$dimensional cosmological constant and $`S_{YGH}`$ the standard York - Gibbons - Hawking boundary term , we consider an action of the form
$$S=\frac{1}{2\kappa _D^2}\underset{M}{}d^Dx\sqrt{|g|}\left\{R[g]2\mathrm{\Lambda }\right\}\frac{1}{2}\underset{M}{}d^Dx\sqrt{|g|}\left(g^{MN}_M\mathrm{\Phi }_N\mathrm{\Phi }+2U(\mathrm{\Phi })\right)+S_{YGH},$$
(2.6)
where the minimal coupled scalar field $`\mathrm{\Phi }`$ with an arbitrary potential $`U(\mathrm{\Phi })`$ depends on the external coordinates $`x`$ only. This field can be understood as a zero mode of a bulk field. Such a scalar field can naturally originate also in non-linear $`D`$dimensional theories where metric ansatz (2.1) ensures its dependence on $`x`$ only.
Let $`\beta _0^i`$ be the scale of compactification of the internal spaces at the present time and
$$V_D^{}V_I\times v_0\underset{i=1}{\overset{n}{}}\underset{M_i}{}d^{d_i}y\sqrt{|g^{(i)}|}\times \underset{i=1}{\overset{n}{}}e^{d_i\beta _0^i}$$
(2.7)
the corresponding total volume of the internal spaces ($`[V_D^{}]=\text{c}m^D^{}`$, $`[V_I]=1`$, where $`D^{}=DD_0`$ is the number of extra dimensions). Instead of $`\beta ^i`$ it is convenient to introduce a shifted quantity:
$$\stackrel{~}{\beta }^i=\beta ^i\beta _0^i.$$
(2.8)
Then, after dimensional reduction action (2.6) reads
$`S`$ $`=`$ $`{\displaystyle \frac{1}{2\kappa _0^2}}{\displaystyle \underset{M_0}{}}d^{D_0}x\sqrt{|g^{(0)}|}{\displaystyle \underset{i=1}{\overset{n}{}}}e^{d_i\stackrel{~}{\beta }^i}\{R\left[g^{(0)}\right]G_{ij}g^{(0)\mu \nu }_\mu \stackrel{~}{\beta }^i_\nu \stackrel{~}{\beta }^j+`$ (2.9)
$`+`$ $`{\displaystyle \underset{i=1}{\overset{n}{}}}\stackrel{~}{R}_ie^{2\stackrel{~}{\beta }^i}2\mathrm{\Lambda }g^{(0)\mu \nu }\kappa _D^2_\mu \mathrm{\Phi }_\nu \mathrm{\Phi }2\kappa _D^2U(\mathrm{\Phi })\},`$
where $`\stackrel{~}{R}_i:=R_ie^{2\beta _0^i}`$, $`G_{ij}=d_i\delta _{ij}d_id_j(i,j=1,\mathrm{},n)`$ is the midisuperspace metric and
$$\kappa _0^2:=\frac{\kappa _D^2}{V_D^{}}$$
(2.10)
is the $`D_0`$dimensional (4-dimensional) gravitational constant. If we take the electroweek scale $`M_{EW}`$ and the Planck scale $`M_{Pl}`$ as fundamental ones for $`D`$dimensional and 4-dimensional space-times respectively:
$`\kappa _D^2`$ $`=`$ $`{\displaystyle \frac{8\pi }{M_{EW}^{2+D^{}}}},`$ (2.11)
$`\kappa _0^2`$ $`=`$ $`{\displaystyle \frac{8\pi }{M_{Pl}^2}},`$
then we reproduce eqs. (1.1) and (1.2).
Action (2.9) describes a generalized $`\sigma `$model with target space metric $`G_{ij}`$ where the scale factors $`\beta ^i`$ play the role of scalar fields. The problem of the internal space stabilization is reduced now to the investigation of the dynamics of these fields. Most easily this can be done in the Einstein frame. For this purpose we perform a conformal transformation
$$g_{\mu \nu }^{(0)}=\mathrm{\Omega }^2\stackrel{~}{g}_{\mu \nu }^{(0)}:=\left(\underset{i=1}{\overset{n}{}}e^{d_i\stackrel{~}{\beta }^i}\right)^{\frac{2}{D_02}}\stackrel{~}{g}_{\mu \nu }^{(0)}$$
(2.12)
which yields
$$S=\frac{1}{2\kappa _0^2}\underset{M_0}{}d^{D_0}x\sqrt{|\stackrel{~}{g}^{(0)}|}\left\{\stackrel{~}{R}\left[\stackrel{~}{g}^{(0)}\right]\overline{G}_{ij}\stackrel{~}{g}^{(0)\mu \nu }_\mu \stackrel{~}{\beta }^i_\nu \stackrel{~}{\beta }^j\stackrel{~}{g}^{(0)\mu \nu }\kappa _D^2_\mu \mathrm{\Phi }_\nu \mathrm{\Phi }2U_{eff}\right\},$$
(2.13)
where $`\overline{G}_{ij}=d_i\delta _{ij}+\frac{1}{D_02}d_id_j`$ and
$$U_{eff}[\stackrel{~}{\beta },\mathrm{\Phi }]=\left(\underset{i=1}{\overset{n}{}}e^{d_i\stackrel{~}{\beta }^i}\right)^{\frac{2}{D_02}}\left[\frac{1}{2}\underset{i=1}{\overset{n}{}}\stackrel{~}{R}_ie^{2\stackrel{~}{\beta }^i}+\mathrm{\Lambda }+\kappa _D^2U(\mathrm{\Phi })\right]$$
(2.14)
is the effective potential.
With the help of a regular coordinate transformation $`\phi =Q\beta ,\beta =Q^1\phi `$ midisuperspace metric (target space metric) $`\overline{G}`$ can be transformed to a pure Euclidean form: $`\overline{G}_{ij}d\beta ^id\beta ^j=\sigma _{ij}d\phi ^id\phi ^j=_{i=1}^nd\phi ^id\phi ^i,\sigma =\mathrm{diag}(+1,+1,\mathrm{},+1)`$. An appropriate transformation $`Q:\beta ^i\phi ^j=Q_i^j\beta ^i`$ can be found e.g. in . We note that in the case of one internal space $`(n=1)`$ this transformation is reduced to a simple redefinition
$$\phi \phi ^1:=\pm \sqrt{\frac{d_1(D2)}{D_02}}\stackrel{~}{\beta }^1$$
(2.15)
which yields
$$U_{eff}[\phi ,\mathrm{\Phi }]=e^{2\phi \sqrt{\frac{d_1}{(D2)(D_02)}}}\left[\frac{1}{2}\stackrel{~}{R}_1e^{2\phi \sqrt{\frac{D_02}{d_1(D2)}}}+\mathrm{\Lambda }+\kappa _D^2U(\mathrm{\Phi })\right].$$
(2.16)
(For definiteness we use the minus sign in eq. (2.15).)
It is clear now that stabilization of the internal spaces can be achieved iff the effective potential $`U_{eff}`$ has a minimum with respect to fields $`\stackrel{~}{\beta }^i`$ (or fields $`\phi ^i`$). In general it is possible for potential $`U_{eff}`$ to have more than one extremum. But it can be easily seen that for the model under consideration we can get one extremum only. Let us find conditions which ensure a minimum at $`\stackrel{~}{\beta }=0`$.
The extremum condition yields:
$$\frac{U_{eff}}{\stackrel{~}{\beta }^k}|_{\stackrel{~}{\beta }=0}=0\stackrel{~}{R}_k=\frac{d_k}{D_02}\left(\underset{i=1}{\overset{n}{}}\stackrel{~}{R}_i2(\mathrm{\Lambda }+\kappa _D^2U(\mathrm{\Phi }))\right).$$
(2.17)
The left-hand side of this equation is a constant but the right-hand side is a dynamical function. Thus, stabilization of the internal spaces in such type of models is possible only when the effective potential has also a minimum with respect to the scalar field $`\mathrm{\Phi }`$ (in Ref. it was proved that for this model the only possible solutions with static internal spaces correspond to the case when the minimal coupled scalar field is in its extremum position too). Let $`\mathrm{\Phi }_0`$ be the minimum position for field $`\mathrm{\Phi }`$. From the structure of the effective potential (2.14) it is clear that minimum positions of the potentials $`U_{eff}[\stackrel{~}{\beta },\mathrm{\Phi }]`$ and $`U(\mathrm{\Phi })`$ with respect to field $`\mathrm{\Phi }`$ coincide with each other:
$$\frac{U_{eff}}{\mathrm{\Phi }}|_{\mathrm{\Phi }_0}=0\frac{U(\mathrm{\Phi })}{\mathrm{\Phi }}|_{\mathrm{\Phi }_0}=0.$$
(2.18)
Hence, we should look for parameters which ensure a minimum of $`U_{eff}`$ at the point $`\stackrel{~}{\beta }^i=0,\mathrm{\Phi }=\mathrm{\Phi }_0`$. Eqs. (2.17) show that there exists a fine tuning condition for the scalar curvatures of the internal spaces:
$$\frac{\stackrel{~}{R}_k}{d_k}=\frac{\stackrel{~}{R}_i}{d_i},(i,k=1,\mathrm{},n).$$
(2.19)
Introducing the auxiliary quantity
$$\stackrel{~}{\mathrm{\Lambda }}\mathrm{\Lambda }+\kappa _D^2U(\mathrm{\Phi })|_{\mathrm{\Phi }_0},$$
(2.20)
we get the useful relations
$$\mathrm{\Lambda }_{eff}:=U_{eff}|_{\genfrac{}{}{0pt}{}{\stackrel{~}{\beta }^i=0,}{\mathrm{\Phi }=\mathrm{\Phi }_0}}=\frac{D_02}{D2}\stackrel{~}{\mathrm{\Lambda }}=\frac{D_02}{2}\frac{\stackrel{~}{R}_k}{d_k},$$
(2.21)
which show that $`\text{sign}\mathrm{\Lambda }_{eff}=\text{sign}\stackrel{~}{\mathrm{\Lambda }}=\text{sign}R_k`$. It is clear that $`\mathrm{\Lambda }_{eff}`$ plays the role of an effective cosmological constant in the external space-time. For the masses of the normal mode excitations of the internal spaces (gravitational excitons) and of the scalar field near the extremum position we obtain respectively :
$`m_1^2`$ $`=`$ $`\mathrm{}=m_n^2={\displaystyle \frac{4\mathrm{\Lambda }_{eff}}{D_02}}=2{\displaystyle \frac{\stackrel{~}{R}_k}{d_k}}>0,`$ (2.22)
$`m_\mathrm{\Phi }^2`$ $`:=`$ $`{\displaystyle \frac{^2U(\mathrm{\Phi })}{\mathrm{\Phi }^2}}|_{\mathrm{\Phi }_0}.`$
These equations show that for our specific model a global minimum can only exist in the case of compact internal spaces with negative curvature $`R_k<0(k=1,\mathrm{},n)`$. The effective cosmological constant is negative also: $`\mathrm{\Lambda }_{eff}<0`$. Obviously, in this model it is impossible to trap the internal spaces at a minimum of $`U_{eff}`$ if they are tori ($`\stackrel{~}{R}_i=0`$) because for Ricci-flat internal spaces the effective potential has no minimum at all. Eqs. (2.21) and (2.22) show also that a stabilization by trapping takes place only for $`\stackrel{~}{\mathrm{\Lambda }}<0`$<sup>1</sup><sup>1</sup>1An interesting scenario for a dynamical stabilization of the internal space was proposed in Ref. for a model with $`\stackrel{~}{\mathrm{\Lambda }}=\stackrel{~}{R}_i=0`$. If, at some stage of the Universe evolution, the inflaton field $`\mathrm{\Phi }`$ reached its zero minimum and was frozen out, then there exists a solution $`\stackrel{~}{\beta }0`$ for times $`t\mathrm{}`$ which corresponds to a dynamical stabilization of the internal space. However, the inflaton field is never frozen out completely and its dynamics can destabilize the internal space. An investigation of this problem in collaboration with Anupam Mazumdar will be presented soon in a common paper.. This means that the minimum of the scalar field potential should be negative $`U(\mathrm{\Phi }_0)<0`$ for non-negative bare cosmological constant $`\mathrm{\Lambda }0`$ or it should satisfy inequality $`\kappa _D^2U(\mathrm{\Phi }_0)<|\mathrm{\Lambda }|`$ for $`\mathrm{\Lambda }<0`$.
For small fluctuations of the normal modes in the vicinity of the minima of the effective potential action (2.13) reads
$`S`$ $`=`$ $`{\displaystyle \frac{1}{2\kappa _0^2}}{\displaystyle \underset{M_0}{}}d^{D_0}x\sqrt{|\stackrel{~}{g}^{(0)}|}\left\{\stackrel{~}{R}\left[\stackrel{~}{g}^{(0)}\right]2\mathrm{\Lambda }_{eff}\right\}`$ (2.23)
$``$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{M_0}{}}d^{D_0}x\sqrt{|\stackrel{~}{g}^{(0)}|}\left\{{\displaystyle \underset{i=1}{\overset{n}{}}}\left(\stackrel{~}{g}^{(0)\mu \nu }\psi _{,\mu }^i\psi _{,\nu }^i+m_i^2\psi ^i\psi ^i\right)+\stackrel{~}{g}^{(0)\mu \nu }\varphi _{,\mu }\varphi _{,\nu }+m_\varphi ^2\varphi \varphi \right\}.`$
(For convenience we use here the normalizations: $`\kappa _0^1\stackrel{~}{\beta }\stackrel{~}{\beta }`$ and $`\sqrt{V_D^{}}(\mathrm{\Phi }\mathrm{\Phi }_0)\varphi `$.) Thus, conformal excitations of the metric of the internal spaces behave as massive scalar fields developing on the background of the external spacetime. In analogy with excitons in solid state physics where they are excitations of the electronic subsystem of a crystal, we called the excitations of the subsystem of internal spaces gravitational excitons . Later, since these particles are also known as radions.
## 3 Inflation of the external space
In this section we discuss briefly the possibility for inflation in the external space of our model. We perform the analysis in the Einstein frame where the effective theory is described by action (2.13) and inflation depends on the form of potential (2.14).
For simplicity we consider a model with only one internal space and an effective potential given by equation (2.16). All our conclusions can be easily generalized to a model with $`n`$ internal spaces.
First, we consider region
$$e^{2\phi \sqrt{\frac{D_02}{d_1(D2)}}}|\mathrm{\Lambda }+\kappa _D^2U(\mathrm{\Phi })|,$$
(3.1)
where the effective potential reads
$$U_{eff}\frac{1}{2}|\stackrel{~}{R}_1|e^{2\phi \sqrt{\frac{D2}{d_1(D_02)}}}.$$
(3.2)
It is well known that for models with potential $`U(\phi )Ae^{\lambda \phi }`$ the scale factor behaves as $`\stackrel{~}{a}\stackrel{~}{t}^{2/\lambda ^2}`$ and power law inflation takes place if $`\lambda ^2<2`$. In our case we have
$$\lambda ^2=\frac{4(D2)}{d_1(D_02)}|_{D_0=4}=2\left(1+\frac{2}{d_1}\right)>2$$
(3.3)
and power law inflation is impossible in this region of the model. For the model with $`n`$ internal spaces the assisted inflation proposed in Ref. is impossible also in this region because of the form of the effective potential (it is impossible to split the effective potential into a sum of $`n`$ terms where each of them depends on one scalar field only).
Second, we consider the region near the minimum of the effective potential. In the scenario of assisted chaotic inflation with a sufficiently large number of scalar fields $`\psi ^i`$ inflation occurs at scales much less than Planck scale: $`|\psi ^i|1`$. In our model the effective action for these fields is given by eq. (2.23) and it has the typical form of an action allowing for this type of inflation. Therefore, it is of interest to investigate the possibility for assisted chaotic inflation here. Unfortunately, for our particular model the internal space stabilization takes place only for negative effective cosmological constant. This destroys inflation because, as it follows from eq. (2.22) $`m_i^2|\mathrm{\Lambda }_{eff}|`$, the energy density of the potential $`U_{eff}`$ is not sufficient for inflation. There is also another drawback of theories with negative cosmological constant. Even if they have a period of inflation there is a problem of succesful completion of it. We shall return to this problem in the next section.
Third, we consider the region
$$\kappa _D^2U(\mathrm{\Phi })\mathrm{\Lambda }+\frac{1}{2}|\stackrel{~}{R}_1|e^{2\phi \sqrt{\frac{D_02}{d_1(D2)}}},$$
(3.4)
where the effective potential reads
$$U_{eff}e^{2p\phi }\kappa _D^2U(\mathrm{\Phi }),p:=\sqrt{\frac{d_1}{(D2)(D_02)}}.$$
(3.5)
For models with $`n+1`$ scalar fields the slow roll conditions are :
$$ฯต\frac{1}{2U_{eff}^2}\underset{i=1}{\overset{n+1}{}}\left(\frac{U_{eff}}{\phi ^i}\right)^2$$
(3.6)
and
$$\eta _iฯต+\frac{1}{U_{eff}}\underset{j=1}{\overset{n+1}{}}\frac{^2U_{eff}}{\phi ^i\phi ^j}\left(\frac{U_{eff}}{\phi ^j}/\frac{U_{eff}}{\phi ^i}\right),i=1,\mathrm{},n+1.$$
(3.7)
Inflation is possible if these parameters are small: $`ฯต,|\eta _i|<1`$. For potential (3.5) we get:
$`ฯต`$ $``$ $`\eta _12p^2+ฯต_\mathrm{\Phi },`$
$`\eta _2`$ $``$ $`2p^2+\eta _\mathrm{\Phi },`$ (3.8)
where $`ฯต_\mathrm{\Phi }:=\frac{1}{2}\left(\frac{U^{}(\mathrm{\Phi })}{U(\mathrm{\Phi })}\right)^2`$ and $`\eta _\mathrm{\Phi }:=ฯต_\mathrm{\Phi }+\frac{U^{\prime \prime }(\mathrm{\Phi })}{U(\mathrm{\Phi })}`$. Because of
$$2p^2|_{D_0=4}=1\frac{2}{d_1+2}<1,$$
(3.9)
inflation is possible in this region if
$$ฯต_\mathrm{\Phi },\eta _\mathrm{\Phi }1.$$
(3.10)
Thus, the scalar field $`\mathrm{\Phi }`$ can act as inflaton and drive the inflation of the external space if its potential in region (3.4) satisfies conditions (3.10). It is clear that estimates (3.9) and (3.10) are rather crude and they show only the principal possibility for inflation to occur. For each particular form of $`U(\mathrm{\Phi })`$ a detailed analysis of the dynamical behavior of the fields in this region should be performed to confirm inflation. Obviously, if the inflation in our model is realized it takes place before the stabilization of the internal spaces. In the case of constant scalar field $`\mathrm{\Phi }=\text{const}`$ or its absence inflation of the external space in our model is impossible at all.
## 4 Discussion and conclusions
In the present paper we considered the possibility for stabilization of the internal space and inflation in the external space using as example a multidimensional cosmological toy model with minimal coupled scalar field as matter source. The calculations above were performed in a model with the electroweek scale $`M_{EW}`$ as fundamental scale of the $`D`$dimensional theory (see eq. (2.11)). Clearly, it is also possible to choose the Planck scale as the fundamental scale.
For this purpose we will not fix the compactification scale of the internal spaces at the present time. We consider them as free parameters of the model and demand only that $`L_{Pl}<a_{(0)i}=e^{\beta _0^i}<L_F`$. So, we shall not transform $`\beta ^i`$ to $`\stackrel{~}{\beta }^i`$. In this case, after dimensional reduction of action (2.6) the effective $`D_0`$dimensional gravitational constant $`\kappa _0^2`$ is defined as
$$\frac{1}{\kappa _0^2}=V_I\frac{\left(L_{Pl}\right)^D^{}}{\kappa _D^2}.$$
(4.1)
At the other hand there holds $`\kappa _0^2=8\pi /M_{Pl}^2`$ (for $`D_0=4`$). Thus, $`\kappa _D^2=8\pi V_I/M_{Pl}^{(2+D^{})}`$, so that the Planck scale becomes the fundamental scale of $`D`$dimensional theory. In this approach eqs. (2.9), (2.12) - (2.16) preserve their form with only substitutions $`\stackrel{~}{\beta }\beta `$ and $`\stackrel{~}{R}_iR_i`$.
The analysis of the internal space stabilization shows that the fine tuning condition (2.19) is not changed:
$$\frac{R_k}{d_k}e^{2\beta _0^k}=\frac{R_i}{d_i}e^{2\beta _0^i},i,k=1,\mathrm{},n$$
(4.2)
and the masses squared of the gravitational excitons and the effective cosmological constant are shifted by the same prefactor:
$`m_i^2`$ $``$ $`\left({\displaystyle \underset{i=1}{\overset{n}{}}}e^{d_i\beta _0^i}\right)^{\frac{2}{D_02}}m_i^2=2\left({\displaystyle \underset{i=1}{\overset{n}{}}}e^{d_i\beta _0^i}\right)^{\frac{2}{D_02}}{\displaystyle \frac{R_i}{d_i}},`$ (4.3)
$`\mathrm{\Lambda }_{eff}`$ $``$ $`\left({\displaystyle \underset{i=1}{\overset{n}{}}}e^{d_i\beta _0^i}\right)^{\frac{2}{D_02}}\mathrm{\Lambda }_{eff}={\displaystyle \frac{D_02}{2}}\left({\displaystyle \underset{i=1}{\overset{n}{}}}e^{d_i\beta _0^i}\right)^{\frac{2}{D_02}}{\displaystyle \frac{R_i}{d_i}}.`$
For example, in the one-internal-space case we get the estimate<sup>2</sup><sup>2</sup>2We use standard Planck length unit conventions with $`[m]=\text{cm}^1`$ and the corresponding shorthand,
e.g. $`m_1^2(a_{(0)1})^{(D2)}\left(\frac{a_{(0)1}}{L_{Pl}}\right)^{(D2)}L_{Pl}^2.`$ :
$$|\mathrm{\Lambda }_{eff}|m_1^2\left(e^{\beta _0^1}\right)^{2\frac{D2}{D_02}}|_{D_0=4}=\left(a_{(0)1}\right)^{(D2)}.$$
(4.4)
This expression shows that due to the power $`(2D)`$ the effective cosmological constant and the masses of the gravitational excitons can be very far from the planckian values even for scales of compactification of the internal spaces close to the Planck length.
Another important note consists in the observation that the Einstein frame metrics of the external spacetime in both approaches are equivalent to each other up to a numerical prefactor:
$$\stackrel{~}{g}_{\mu \nu }^{(0)}|_{EW}=v_0^{2/(D_02)}\stackrel{~}{g}_{\mu \nu }^{(0)}|_{Pl}.$$
(4.5)
Equation (2.12) shows that in the electroweak approach the Brans-Dicke and Einstein scales coincide with each other at the point of stabilization: $`\stackrel{~}{\beta }^i=0\mathrm{\Omega }=1`$. In the Planck fundamental scale approach this has place when the internal scale factors are equal to the Planck length: $`\beta ^i=0\mathrm{\Omega }=1`$. This does not mean that in the latter approach the stabilization of the internal space takes place at the Planck length. Depending on the concrete form of the effective potential $`U_{eff}`$ its minimum position/stabilization point $`\beta _0^i`$ can be located at much larger scales $`L_{Pl}a_{(0)i}=e^{\beta _0^i}L_{Pl}`$. Generally speaking, we should not exclude also a possibility for internal spaces to change very slowly with time. In this case $`\beta _0^i`$ is not so strictly defined as for models with the internal space stabilization in minima of the effective potential.
Let us return to the comparision of the electroweak and the Planck scale approaches. From eqs. (4.3) it is clear that the reason for the rescaling/lightening of the effective cosmological constant as well as of the gravitational exciton masses in the Planck scale approach consists in the prefactor $`\left(_{i=1}^ne^{d_i\beta _0^i}\right)^{\frac{2}{D_02}}`$. In spite of the smallness of the internal space sizes in the Planck fundamental scale approach ($`L_{Pl}<a_{(0)1}<L_F`$) in comparison with the sizes in the electroweak fundamental scale approach ($`a_{(0)1}10^1`$cm for $`D^{}=2`$ and $`a_{(0)1}10^{17}`$cm for $`D^{}\mathrm{}`$), the prefactors in eqs. (4.3) can considerably reduce the values of $`\mathrm{\Lambda }_{eff}`$ and $`m_i`$ making them much smaller then in the electroweak approach.
Let us compare now some estimates following from the electroweak as well as from the Planck fundamental scale approaches. (We use the obvious subscripts $`EW`$ and $`Pl`$ respectively.)
In the first case, the scale of the internal space compactification is given by formula (1.2). We take for definiteness the total number of dimensions $`D=6`$ and $`D=10`$ and obtain respectively the following scales of compactification: $`a_{(0)1}10^1\text{c}m`$ for $`D=6`$ and $`a_{(0)1}10^9\text{c}m`$ for $`D=10`$. Then, from eqs. (2.21) and (2.22) we get:
$$|\mathrm{\Lambda }_{eff}||_{EW}\frac{1}{a_{(0)1}^2}\{\begin{array}{ccc}\hfill 10^2\text{c}m^2& & 10^{64}\mathrm{\Lambda }_{Pl},D=6\hfill \\ \hfill 10^{18}\text{c}m^2& & 10^{48}\mathrm{\Lambda }_{Pl},D=10\hfill \end{array}$$
(4.6)
and
$$m_1|_{EW}\frac{1}{a_{(0)1}}\{\begin{array}{ccc}\hfill 10^{32}M_{Pl}& & 10^4eV,D=6\hfill \\ \hfill 10^{24}M_{Pl}& & 10^4eV,D=10\hfill \end{array}.$$
(4.7)
In the second case, the scale of compactification is not fixed, but a free parameter. We demand only that it should be smaller then the Fermi length. For definiteness let us use $`a_{(0)1}10^{18}\text{c}m`$. Then, from eq. (4.4) we obtain:
$$|\mathrm{\Lambda }_{eff}||_{Pl}a_{(0)1}^{(D2)}\{\begin{array}{ccc}\hfill 10^6\text{c}m^2& & 10^{60}\mathrm{\Lambda }_{Pl},D=6\hfill \\ \hfill 10^{54}\text{c}m^2& & 10^{120}\mathrm{\Lambda }_{Pl},D=10\hfill \end{array}$$
(4.8)
and
$$m_1|_{Pl}a_{(0)1}^{(D2)/2}\{\begin{array}{ccc}\hfill 10^{30}M_{Pl}& & 10^2eV,D=6\hfill \\ \hfill 10^{60}M_{Pl}& & 10^{32}eV,D=10\hfill \end{array}$$
(4.9)
Estimates (4.6) and (4.8) show that for the electroweak scale the effective cosmological constant is much greater than the present day observable limit $`\mathrm{\Lambda }10^{122}\mathrm{\Lambda }_{Pl}10^{57}\text{c}m^2`$ (for our model $`|\mathrm{\Lambda }_{eff}||_{EW}10^2\text{cm}^2`$), whereas in the Planck scale approach we can satisfy this limit even for very small compactification scales. For example, if we demand in accordance with observations $`|\mathrm{\Lambda }_{eff}|10^{122}\mathrm{\Lambda }_{Pl}`$ then eq. (4.4) gives a compactification scale $`a_{(0)1}10^{122/(D2)}L_{PL}`$. Thus, $`a_{(0)1}10^{15}L_{Pl}10^{18}\text{c}m`$ for $`D=10`$ and $`a_{(0)1}10^5L_{Pl}10^{28}\text{c}m`$ for $`D=26`$, which does not contradict to observations because for this approach the scales of compactification should be $`a_{(0)1}10^{17}\text{c}m`$. Assuming an estimate $`\mathrm{\Lambda }_{eff}10^{122}L_{Pl}`$, we automatically get from eq. (4.4) the value of the gravitational exciton mass: $`m_110^{61}M_{Pl}10^{33}eV10^{66}\text{g}`$ which is extremely light. Nevertheless such light particles are not in contradiction with the observable Universe, as we shall show below.
Similar to the Polonyi fields in spontaneously broken supergravity or moduli fields in the hidden sector of SUSY the gravitational excitons are WIMPs (Weakly-Interacting Massive Particles ) because their coupling to the observable matter is suppresed by powers of the Planck scale. In Ref. we show that the decay rate of the gravitational excitons with mass $`m_\phi `$ is $`\mathrm{\Gamma }m_\phi ^3/M_{Pl}^2`$ as for Polonyi and moduli fields.
Let us assume for a moment that after inflation the inflaton field $`\varphi `$ has already decayed and produced the main reheating of the Universe. For our model it may happen if $`m_\varphi m_\phi `$ and the inflaton field starts to oscillate and decay much earlier than the $`\phi `$field (coherent oscillations of field $`\varphi `$ with mass $`m_\varphi `$ usually start when the Hubble constant $`Hm_\varphi `$). The Universe is radiation dominated in this period and the Hubble constant is defined by $`HT^2/M_{Pl}`$. After the temperature is fallen to the value $`T_{in}\sqrt{m_\phi M_{Pl}}`$ the scalar field<sup>3</sup><sup>3</sup>3Here, $`\phi =\pm \frac{M_{Pl}}{\sqrt{8\pi }}\sqrt{\frac{d_1(D2)}{D_02}}\beta ^1`$ where $`\beta ^1`$ is the logarithm of the internal space scale factor: $`a_1=e^{\beta ^1}L_{Pl}`$. If stabilization occurs at $`a_{(0)1}10^nL_{Pl},(0<n<18)`$, then it corresponds to the minimum position $`\phi _0=\frac{n\mathrm{ln}10}{\sqrt{8\pi }}\sqrt{\frac{d_1(D2)}{D_02}}M_{Pl}`$. $`\phi `$ begins to oscillate coherently around the minumim and its density evolves as $`T^3`$ :
$$\rho _\phi (T)=\rho _\phi (T_{in})\left(T/T_{in}\right)^3=m_\phi ^2\phi _{in}^2(T_{in})\left(T/T_{in}\right)^3,$$
(4.10)
where $`\phi _{in}:=(\phi \phi _0)_{in}`$ is the amplitude of initial oscillations of the field $`\phi `$ near the minimum position. It is clear that for the extremely light particles we can neglect their decay ($`\mathrm{\Gamma }_\phi 0`$). Then, because the ratio $`\rho _\phi /\rho _{rad}`$ increases as $`1/T`$, at some themperature the Universe will be dominated (up to present time) by the energy density of the coherent oscillations. We can easily estimate the mass of the gravitational excitons which overclose the Universe. Assuming that at present time $`\rho _\phi \begin{array}{c}<\hfill \\ \hfill \end{array}\rho _c`$, where $`\rho _c`$ is the critical density of the present day Universe, we obtain<sup>4</sup><sup>4</sup>4See also Note added.
$$m_\phi \begin{array}{c}<\hfill \\ \hfill \end{array}10^{56}M_{Pl}\left(\frac{M_{Pl}}{\phi _{in}}\right)^4.$$
(4.11)
Usually, it is assumed that $`\phi _{in}O(M_{Pl})`$ although it depends on the form of $`U_{eff}`$ and can be considerably less than $`M_{Pl}`$. If we put $`\phi _{in}O(M_{Pl})`$ then excitons with masses $`m_\phi \begin{array}{c}<\hfill \\ \hfill \end{array}10^{28}\text{e}V`$ will not overclose the Universe . If $`\phi _{in}M_{Pl}`$ this estimate will be not so severe. We see that our mass $`m_\phi 10^{33}\text{e}V`$ satisfies the most severe estimate. It can be considered as hot dark matter which negligibly contributes to the total amount of dark matter and does not contradict to the model of cold dark matter.
Of course, as it follows from eqs. (4.4) and (4.9) the mass $`m_\phi `$ could be considerably heavier than $`10^{33}\text{e}V`$ but as result we would arrive at an effective cosmological constant greater than the observable one (see eqs. (4.8) and (4.9) for $`D=6`$ and $`a_{(0)1}10^{18}\text{c}m`$) and we would need a mechanism for its reduction to the observable value. An example for such a reduction of the cosmological constant was proposed in for SUSY breaking models with moduli masses $`m10^210^3\text{e}V`$. Such masses we get also in our model if we take for the Planck scale approach $`D=6`$ and $`a_{(0)1}10^{18}\text{c}m`$ (see (4.9)). For these particles we cannot neglect the decay rate $`\mathrm{\Gamma }_\phi `$ which results in converting of the coherent oscillations into radiation. In this case the Universe has a further reheating to the themperature
$$T_{RH}\sqrt{\frac{m_\phi ^3}{M_{Pl}}}.$$
(4.12)
For $`m_\phi 10^2`$ eV the reheating temperature $`T_{RH}10^{23}\text{ MeV}1\text{ MeV}`$ is much less than the temperature $`T1\text{M}eV`$ at which the nucleosynthesis begins. Thus, either decaying particles should have masses $`m_\phi >10^4\text{G}eV`$ to get $`T_{RH}>1\text{M}eV`$ or we should get rid off such particles before nucleosynthesis. The latter can be achieved if the decay rate becomes larger. In it was proposed that at a very early stage of the Universe evolution (after inflation) WIMPs collapse into stars (e.g. modular stars) where their field strength could be very large and leads to a substantial enhancement of the decay into ordinary particles. For example, in Ref. it is shown that gravexcitons have a coupling to photons of the form $`\frac{\phi }{M_{Pl}}F_{\mu \nu }^2`$. In the core of such stars the gravexciton amplitude $`\phi `$ might be much larger than $`M_{Pl}`$, enhancing the coupling of this field to photons and leading to explosions of these stars into bursts of photons.
As it follows from eqs. (2.22) and (4.7), in the electroweak approach gravexciton masses should satisfy the inequality $`m_\phi \begin{array}{c}>\hfill \\ \hfill \end{array}10^4\text{e}V`$. If the above mentioned mechanism of the gravexciton energy dilution due to modular star explosions or due to some other reasons<sup>5</sup><sup>5</sup>5For example, in for this purpose a short period of late inflation was proposed which should be followed by a reheating. However, it is necessary to be rather careful to avoid the generation of quantum fluctuations of gravexcitons during inflation again . does not work, the bound $`m_\phi \begin{array}{c}>\hfill \\ \hfill \end{array}10^4\text{G}eV`$ is valid and leads to the large $`D^{}`$ limit ($`D^{}30`$) with a scale of compactification $`a_{(0)1}\sqrt{D^{}}m_\phi ^1`$ (see (2.22)). Thus for $`D^{}100`$ and $`m_\phi 10^4\text{G}eV`$ we get $`a_{(0)1}10^{17}`$ cm which is not in strong contradiction with the value $`a_{(0)1}10^{16,7}\text{c}m`$ which follows from eq. (1.2). It is clear that in this approach an increasing of the mass by one order requires an increasing of the number of internal dimensions by two orders.
Above, we considered the case $`m_\varphi m_\phi `$ when the inflaton field starts to oscillate coherently much earlier than the scale factors of the internal spaces. Let us suppose now that $`m_\varphi m_\phi m`$. Thus, the inflaton $`\varphi `$ and gravexciton $`\phi `$ fields start to oscillate coherently at the same time $`t_{in}`$ with approximately the same initial amplitude $`\varphi _{in}\phi _{in}`$. by At this time the Universe becomes matter dominated with $`\rho _\phi \rho _\varphi 1/\stackrel{~}{a}^3`$ where $`\stackrel{~}{a}`$ is the scale factor of the external space-time. We assume also that the inflaton $`\varphi `$ is not a WIMP and its decay rate $`\mathrm{\Gamma }_\varphi \alpha _\varphi ^2m\mathrm{\Gamma }_\phi m^3/M_{Pl}^2`$. Thus the effective coupling $`\alpha _\varphi `$ of the inflaton field $`\varphi `$ satisfies: $`\alpha _\varphi m/M_{Pl}`$. Because $`mM_{Pl}`$ the effective coupling $`\alpha _\varphi `$ still may be much less than 1.
First, we consider the case when the gravexciton decay rate is negligibly small: $`\mathrm{\Gamma }_\phi 0`$. Let $`t_{RH}`$ be the time of reheating due to inflaton decay and let us suppose that all the inflaton energy is converted into radiation ($`\rho _\varphi (t_{RH})\rho _{rad}|_{RH}T_{RH}^4`$). It can be easily seen that for $`t>t_{RH}`$ the relative contribution of $`\phi `$ to the energy density starts to increase as
$$\frac{\rho _\phi (T)}{\rho _{rad}(T)}=\frac{T_{RH}}{T}.$$
(4.13)
Here, in the sudden decay approximation the reheating temperature, $`T_{RH}`$, is defined by equating the Hubble constant with the rate of decay: $`H(t_D)\mathrm{\Gamma }_\varphi \alpha _\varphi ^2m`$, where $`t_Dt_{RH}`$ is the decay time. Because $`H^2(t_D)M_{Pl}^2\rho _\varphi |_{t_D}T_{RH}^4/M_{Pl}^2`$ we get
$$m\frac{1}{\alpha _\varphi ^2}\frac{T_{RH}^2}{M_{Pl}}.$$
(4.14)
This formula shows that to get the temperature $`T_{RH}>1\text{M}eV`$, which is necessary for the nucleosynthesis, the mass should satisfy the inequality
$$m\begin{array}{c}>\hfill \\ \hfill \end{array}\frac{1}{\alpha _\varphi ^2}10^{16}\text{e}V.$$
(4.15)
At the other hand, at present time (which we denote by a subscript 0) the condition that gravexcitons do not overclose the Universe reads: $`\rho _\varphi |_0=\left(T_{RH}/T_0\right)\rho _{rad}|_0\begin{array}{c}<\hfill \\ \hfill \end{array}\rho _c`$ and gives a second limit for the mass:
$$m\begin{array}{c}<\hfill \\ \hfill \end{array}\frac{1}{\alpha _\varphi ^2}\left(\frac{\rho _c}{\rho _{rad}|_0}\right)^2\frac{T_0^2}{M_{Pl}}.$$
(4.16)
Inserting into this formula the present day values for the temperature $`T_0`$ and the critical energy density $`\rho _c`$ we obtain
$$m\begin{array}{c}<\hfill \\ \hfill \end{array}\frac{1}{\alpha _\varphi ^2}10^{26}\text{e}V,$$
(4.17)
which obviously is in contradiction to the previous estimate (4.15).
Second, to solve this problem we consider the possibility of a further reheating due to gravexciton decay: $`\mathrm{\Gamma }_\phi 0`$. In order to estimate the temperature at which this decay occurs we should take into account that after the first reheating (with the temperature defined by (4.14)) the Universe is matter dominated because $`\rho _\phi /\rho _{rad}=T_{RH}/T>1`$ for $`T<T_{RH}`$ and for the Hubble constant holds $`H^2\rho _\phi /M_{Pl}^2`$. Thus, equating the Hubble constant with the decay rate: $`H(T_D^{})\mathrm{\Gamma }_\phi m^3/M_{Pl}^2`$ we obtain the temperature of the gravexciton decay:
$$\left(T_D^{}\right)^3\frac{1}{\alpha _\varphi }\frac{m^{11/2}}{M_{Pl}^{5/2}}.$$
(4.18)
In this scenario the temperatures of the gravexciton decay and the reheating are denoted by a prime to distinguish them from the corresponding temperatures of inflaton decay and reheating. In the sudden decay approximation the temperature of the second reheating is obtained by equating the squared decay rate and the radiation energy density just after reheating (because $`H^2(T_D^{})\mathrm{\Gamma }_\phi ^2\rho _{rad}/M_{Pl}^2`$) which obviously leads again to eq. (4.12) (where $`T_{RH}`$ should be replaced by $`T_{RH}^{}`$). Again, for a successful nucleosynthesis with $`T_{RH}^{}>1\text{M}eV`$ the mass should be $`m\begin{array}{c}>\hfill \\ \hfill \end{array}10^4\text{G}eV`$. The reheating from $`T_D^{}`$ to $`T_{RH}^{}`$ produces an entropy increase given by
$$\mathrm{\Delta }=\left(\frac{T_{RH}^{}}{T_D^{}}\right)^3\alpha _\varphi \frac{M_{Pl}}{m}1,$$
(4.19)
which is much greater than 1 because $`\alpha _\varphi m/M_{Pl}`$. However, it may be much less (not so severe) than the usual estimate : $`\mathrm{\Delta }M_{Pl}/m`$ because $`\alpha _\varphi `$ may be much less than 1. If we require $`\mathrm{\Delta }\begin{array}{c}<\hfill \\ \hfill \end{array}10^5`$, as a maximal permissible factor for the dilution of the high-temperature baryogenesis, we obtain the bound $`m\begin{array}{c}>\hfill \\ \hfill \end{array}\alpha _\varphi 10^5M_{Pl}`$ and for $`\alpha _\varphi 1`$ this bound is not so strong as the usual one: $`m\begin{array}{c}>\hfill \\ \hfill \end{array}10^{14}`$GeV.
Summarizing the discussion we see that in models, where the coherent oscillation of gravitational excitons starts in the radiation dominated era, the gravexcitons should be either extremely light (see eq. (4.11)) or very heavy particles ($`m_\phi >10^4`$GeV for a successful nucleosynthesis; in case that the hot baryogenesis is taken into account: $`m_\phi >10^{14}`$GeV). In models, where inflaton and gravitational exciton start their coherent oscillation at the same time, extremely light excitons are forbidden. Heavier excitons with masses $`m_\phi \begin{array}{c}>\hfill \\ \hfill \end{array}\alpha _\phi 10^5M_{Pl}`$ are allowed (for a successful nucleosynthesis and high-temperature baryogenesis).
As conclusion we would like to note that in our toy model the stabilization of the internal spaces is realized only when the effective cosmological constant is negative (for both fundamental scale approaches). It is well known that for such models inflation is never succeffully completed , because in this case our (external) space has a turning point at its maximal scale factor where it stops to expand and begins to contract. If the spatial curvature of our Universe is non-negative (according to the latest observational data it is zero), then the internal scale factors cannot be freezed because solutions $`\beta ^i=\beta _0^i=\text{const}`$ correspond to a negative squared Hubble constant: $`H^2=\frac{1}{3}\left[\frac{1}{2}\overline{G}_{ij}\dot{\beta }^i\dot{\beta }^j+U_{eff}\right]\frac{1}{3}\mathrm{\Lambda }_{eff}<0`$. To describe the post-inflationary stage for such models we should extend our consideration including e.g. additional perfect fluid terms into the action functional which correspond to usual matter in the universe (see for the details of this method). Another possible generalization consists in an inclusion of additional terms which result in a positive effective cosmological constant in accordance with recent observational data . This can be achieved e.g. with the help of antisymmetric form-fields . For these models the gravexciton masses and the (positive) effective cosmological constant are defined by equations similar to (4.4). Such models can solve the following three important problems simultaneously: they yield stabilization of the internal spaces, allow for inflation of the external space, and lead to a positive observable effective cosmological constant. In these models the mechanism of lightening of the effective cosmological constant as well as the gravitational exciton masses will work also in the Planck fundamental scale approach because eqs. (4.3) are general for this type of models.
Note added
Alexander Sakharov informed us about another upper bound on $`m_\phi `$ following from isocurvature gravexciton fluctuations if $`m_\varphi m_\phi `$ because in this case gravexcitons on the stage of inflation can be considered as massless particles. These isocurvature fluctuations result in a CMBR anisotropy $`\delta T/T`$. The amplitude of these fluctuations can be estimated as $`\delta \phi H_{inf}/2\pi `$ and is connected with $`\delta T/T`$ as follows: $`\delta T/T(\rho _\phi /\rho _c)(\delta \phi /\phi _{in})(\rho _\phi /\rho _c)(H_{inf}/2\pi \phi _{in})`$, where $`H_{inf}`$ is the Hubble constant at the inflation stage. According to COBE data, $`\delta T/T\begin{array}{c}<\hfill \\ \hfill \end{array}10^5`$ and $`H_{inf}10^5M_{Pl}`$. Thus, we get following limitation on the gravexciton energy density at present time: $`\rho _\phi \begin{array}{c}<\hfill \\ \hfill \end{array}2\pi \rho _c\phi _{in}/M_{Pl}`$. Substitution of eq. (4.10) into this limitation gives
$$m_\phi \begin{array}{c}<\hfill \\ \hfill \end{array}10^{55}M_{Pl}\left(\frac{M_{Pl}}{\phi _{in}}\right)^2.$$
(4.20)
So, if $`\phi _{in}O(M_{Pl})`$ then both eqs. (4.11) and (4.20) give close limitations on $`m_\phi `$. However, for $`\phi _{in}M_{Pl}`$ we should use eq. (4.11), (4.20) correspondingly.
In the case of decaying gravexcitons the CMBR anisotropy due to gravexciton isocurvature fluctuations is washed out.
Acknowledgments
We would like to thank Valery Rubakov and Alexander Sakharov for valuable correspondence, Nemanja Kaloper for useful comments concerning his recent work and Martin Rainer for interesting discussions. A.Z. thanks H. Nicolai and the Albert Einstein Institute for kind hospitality. U.G. acknowledges financial support from DFG grant KON 1575/1999/GU 522/1. |
warning/0002/math-ph0002033.html | ar5iv | text | # On bifurcations from normal solutions for superconducting states
## 1 Introduction
### 1.1 Our model.
Following the paper by Berger-Rubinstein \[BeRu\], we would like to understand the minima (or more generally the extrema) of the following Ginzburg-Landau functional. In a bounded, connected, regular<sup>1</sup><sup>1</sup>1with $`C^{\mathrm{}}`$ boundary, open set $`\mathrm{\Omega }IR^2`$ and, for any $`\lambda >0`$ and $`\kappa >0`$, this functional $`G_{\lambda ,\kappa }`$ is defined, for $`uH^1(\mathrm{\Omega };C\text{ }\text{ })`$ and $`AH_{loc}^1(IR^2;IR^2)`$ such that $`\mathrm{rot}AL^2`$, by
$$\begin{array}{cc}G_{\lambda ,\kappa }(u,A)\hfill & =_\mathrm{\Omega }(\lambda (|u|^2+\frac{1}{2}|u|^4)+|(iA)u|^2))dx_1dx_2\hfill \\ & +\kappa ^2\lambda ^1_{IR^2}|\mathrm{rot}AH_e|^2๐x_1๐x_2.\hfill \end{array}$$
(1.1)
Here, for $`A=(A_1,A_2)`$, $`\mathrm{rot}A=_{x_1}A_2_{x_2}A_1`$, $`\mathrm{div}A=_{x_1}A_1+_{x_2}A_2`$ and $`H_e`$ is a $`C_0^{\mathrm{}}`$ function on $`IR^2`$ (or more generally some function in $`L^2(IR^2)`$). Physically $`H_e`$ represents the exterior magnetic field.
Let $`A_e`$ be a solution of
$$\begin{array}{c}\mathrm{rot}A_e=H_e\hfill \\ \mathrm{div}A_e=0.\hfill \end{array}$$
(1.2)
It is easy to verify that such a solution exists by looking for $`A_e`$ in the form $`A_e=(_{x_2}\psi _e,_{x_1}\psi _e)`$. We have then to solve $`\mathrm{\Delta }\psi _e=H_e`$ and it is known to be solvable in $`๐ฎ^{}(IR^2)C^{\mathrm{}}(IR^2)`$ (or in $`๐ฎ^{}(IR^2)H_{loc}^2(IR^2)`$ if $`H_eL^2(IR^2)`$). Of course $`A_e`$ is not unique but we shall discuss about uniqueness modulo gauge transform later and at the end this is mainly the restriction of $`A_e`$ to $`\mathrm{\Omega }`$ which will be considered.
We shall sometimes use the identification between vector fields $`A`$ and $`1`$-forms $`\omega _A`$.
When analyzing the extrema of the GL-functional, it is natural to first analyze the corresponding Euler-Lagrange equations. This is a system of two equations (with a boundary equation) :
$$\begin{array}{cc}(GL)_1\hfill & (iA)^2u+\lambda u(|u|^21)=0,\text{ in }\mathrm{\Omega },\hfill \\ (GL)_2\hfill & \mathrm{rot}^{}(\mathrm{rot}AH_e)=\lambda \kappa ^2\mathrm{Im}\left[\overline{u}(iA)u\right]1_\mathrm{\Omega },\hfill \\ (GL)_3\hfill & (iA)u\nu =0,\text{ in }\mathrm{\Omega }.\hfill \end{array}$$
(1.3)
Here $`\nu `$ is a unit exterior normal to $`\mathrm{\Omega }`$. The operator $`\mathrm{rot}^{}`$ is defined by $`\mathrm{rot}^{}f:=(_{x_2}f,_{x_1}f)`$.
Moreover, without loss of generality in our problem, we shall add the condition
$$(GL)_4\mathrm{div}A=0\text{ in }\mathrm{\Omega }.$$
(1.4)
One can also assume if necessary that the vector potential satisfies
$$A\stackrel{}{\nu }=0,$$
(1.5)
on the boundary of $`\mathrm{\Omega }`$, where $`\nu `$ is a normal unit vector to $`\mathrm{\Omega }`$.
Let us briefly recall the argument. One would like to find $`\theta `$ in $`C^{\mathrm{}}(\overline{\mathrm{\Omega }})`$ such that $`\stackrel{~}{A}=A+d\theta `$ satisfies (1.4) and (1.5). One can proceed in two steps. The first step is to find a gauge transformation such that (1.5) is satisfied. This is immediate if the boundary is regular.
We now assume this condition.
The second step consists in solving
$$\begin{array}{c}\mathrm{\Delta }\theta =\mathrm{div}A\text{ in }\mathrm{\Omega },\hfill \\ \frac{\theta }{\nu }=0,\text{ on }\mathrm{\Omega }.\hfill \end{array}$$
This is a Neumann problem, which is solvable iff the right hand-side is orthogonal to the first eigenfunction of the Neumann realization of the Laplacian, that is the constant function $`x1`$. We have only to observe that $`_\mathrm{\Omega }\mathrm{div}A๐x=0`$ if (1.5) is satisfied.
An important remark is that the pair $`(0,A_e)`$ is a solution of the system. This solution is called the normal solution. Of course, any solution of the form $`(0,A_e+\varphi )`$ with $`\varphi `$ harmonic is also a solution.
###### Remark 1.1
.
Note also that the normalization of the functional leads to the property that
$$G_{\lambda ,\kappa }(0,A_e)=0.$$
(1.6)
The first proposition is standard.
###### Proposition 1.2
.
If $`\mathrm{\Omega }`$ is bounded, the functional $`G_{\lambda ,\kappa }`$ admits a global minimizer which is a solution of the equation.
We refer to \[DGP\], for a proof together with the discussion of the next subsection.
### 1.2 Comparison with other models
Let us observe that there is another natural problem which may be considered. This is the problem of minimizing, for $`(u,A)H^1(\mathrm{\Omega },C\text{ }\text{ })\times H^1(\mathrm{\Omega },IR^2)`$, the functional $`G_\lambda ^\mathrm{\Omega }`$ defined by
$$\begin{array}{cc}G_{\lambda ,\kappa }^\mathrm{\Omega }(u,A)\hfill & =_\mathrm{\Omega }\left(\lambda (|u|^2+\frac{1}{2}|u|^4)+|(iA)u|^2\right)๐x_1๐x_2\hfill \\ & +\kappa ^2\lambda ^1_\mathrm{\Omega }|\mathrm{rot}AH_e|^2๐x_1๐x_2.\hfill \end{array}$$
(1.7)
This may lead to a different result in the case when $`\mathrm{\Omega }`$ is not simply connected. According to discussion with Akkermans, this is the first problem which is the most physical (See also the discussion in the appendix).
A comparison between $`G_{\lambda ,\kappa }`$ and $`G_{\lambda ,\kappa }^{\mathrm{\Omega },D}`$ where $`D`$ is a ball containing $`\mathrm{\Omega }`$ and $`G_{\lambda ,\kappa }^{\mathrm{\Omega },D}`$ is defined by
$$\begin{array}{cc}G_{\lambda ,\kappa }^{\mathrm{\Omega },D}(u,A)\hfill & =_\mathrm{\Omega }\left(\lambda (|u|^2+\frac{1}{2}|u|^4)+|(iA)u|^2\right)๐x_1๐x_2\hfill \\ & +\kappa ^2\lambda ^1_D|\mathrm{rot}AH_e|^2๐x_1๐x_2.\hfill \end{array}$$
(1.8)
is useful. If $`b`$ is given with support outside of the ball $`D`$, it is easy to see (assuming that $`b`$ is regular) that there exists $`a`$ with support outside $`D`$ such that $`\mathrm{rot}a=b`$. It is indeed sufficient to take the usual transversal gauge
$$a_1=x_2_0^1sb(sx)๐s,a_2=x_1_0^1sb(sx)๐s.$$
(1.9)
This shows that, for any $`D`$ containing $`\mathrm{\Omega }`$, we have
$$infG_{\lambda ,\kappa }(u,A)=infG_{\lambda ,\kappa }^{\mathrm{\Omega },D}(u,A).$$
(1.10)
In particular it is enough to consider minimizing sequences $`(u_n,A_e+a_n)`$ where $`suppa_nD`$ and $`D`$ is a ball containing $`\mathrm{\Omega }`$. The proof of the existence of minimizers is then greatly simplified.
Finally, it is natural <sup>2</sup><sup>2</sup>2This is at least clear when $`\stackrel{~}{\mathrm{\Omega }}`$ is a star-shaped domain by the previous proof. See Section 3, in the proof of Proposition 3.2 for a complementary argument. to think that one can replace $`D`$ by
$$\stackrel{~}{\mathrm{\Omega }}:=\mathrm{\Omega }_i๐ช_i,$$
(1.11)
where the $`๐ช_i`$ are the holes, that are the bounded connected components of $`IR^2\mathrm{\Omega }`$. A proof can be obtained by analyzing the Ginzburg-Landau equations satisfied by a minimizer of $`G_{\lambda ,\kappa }^{\mathrm{\Omega },D}`$. We finally get :
$$infG_{\lambda ,\kappa }(u,A)=infG_{\lambda ,\kappa }^{\mathrm{\Omega },\stackrel{~}{\mathrm{\Omega }}}(u,A).$$
(1.12)
###### Remark 1.3
.
If $`(u,A_e+a)`$ is a solution of the GL-equation then $`\mathrm{rot}a=0`$ in the unbounded component of $`IR^2\mathrm{\Omega }`$ and $`\mathrm{rot}a=\mathrm{const}.`$ in each hole (See Lemma 2.1 in \[GiPh\]). It would be interesting to discuss what the possible values of these constants are : are they equal to $`0`$ ?
### 1.3 Standard results
The second proposition which is also quite standard (See for example \[DGP\]) is
###### Proposition 1.4
.
If $`u`$ is a solution of the first GL-equation with the Neumann boundary condition then
$$|u(x)|1,x\mathrm{\Omega }.$$
(1.13)
We note for further use that the solutions of the G-L system are in $`C^{\mathrm{}}(\overline{\mathrm{\Omega }})`$ under the assumption that $`\mathrm{\Omega }`$ is regular.
## 2 Is the normal state a minimizer ?
The aim of this section is to give a proof of a result suggested in \[BeRu\] who said โ We expect the normal state to be a stable solution for small $`\lambda `$โฆโ.
Although, this result is probably known as folk theorem, we think it is useful to give a proof (following considerations by M. Dutour in a near context \[Du\]) of this property.
Note that connected results are obtained in \[GiPh\] and more recently in \[LuPa1\], \[LuPa2\].
Before stating the theorem, let us recall that we have called normal state a pair $`(u,A)`$ of the form :
$$(u,A)=(0,A_e),$$
(2.1)
where $`A_e`$ is any solution of (1.1).
We note that this is well defined up to gauge transformation. Moreover, we have :
###### Lemma 2.1
.
$`(0,A_e)`$ is a solution of the GL-system.
So it is effectively natural to ask if $`(0,A_e)`$ is a global minimum. The first result in this direction is the following easy proposition about the normal state. But let us first introduce :
###### Definition 2.2
.
We denote by $`\lambda ^{(1)}`$ the lowest eigenvalue of the Neumann realization in $`\mathrm{\Omega }`$ of
$$\mathrm{\Delta }_{A_e}:=(iA_e)^2.$$
We shall frequently use the assumption
$$\lambda ^{(1)}>0.$$
(2.2)
Note the following necessary and sufficient condition for this property (cf \[He\]).
###### Proposition 2.3
.
The condition (2.2) is satisfied if and only if one of two following conditions is satisfied :
1. $`H_e`$ is not identically zero in $`\mathrm{\Omega }`$;
2. $`H_e`$ is identically zero in $`\mathrm{\Omega }`$ but there exists a closed path $`\gamma `$ in $`\mathrm{\Omega }`$ such that $`\frac{1}{2\pi }_\gamma \omega _{A_e}ZZ`$.
Let us observe that the second case can only occur when $`\mathrm{\Omega }`$ is non simply connected.
###### Proposition 2.4
.
Under condition (2.2) and if $`\lambda ]0,\lambda ^{(1)}[`$, the pair $`(0,A_e)`$ is a non-degenerate (up to gauge transforms) local minimum of $`G_{\lambda ,\kappa }`$.
The Hessian at $`(0,A_e)`$ of the GL-functional is indeed the map
$$(\delta u,\delta a)((\mathrm{\Delta }_{A_e}\lambda )\delta u,\mathrm{rot}^{}\mathrm{rot}\delta a),$$
where we assume that $`\mathrm{div}\delta a=0`$ and $`\delta a\nu =0`$ at the boundary of $`\stackrel{~}{\mathrm{\Omega }}`$.
Note that this proof gives also :
###### Proposition 2.5
.
If $`\lambda >\lambda ^{(1)}`$, the pair $`(0,A_e)`$ is not a local minimum of $`G_{\lambda ,\kappa }`$.
We refer to \[LuPa1\] for a connected result. Proposition 2.5 does not answer completely to the question about global minimizers. The next theorem gives a complementary information.
###### Theorem 2.6
.
Under assumption (2.2), then, for any $`\kappa >0`$, there exists $`\lambda _0(\kappa )>0`$ such that, for $`\lambda ]0,\lambda _0(\kappa )]`$, $`G_{\lambda ,\kappa }`$ has only normal solutions as global minimizers.
###### Remark 2.7
.
Using a variant of the techniques used in \[Du\] in a similar context, one can actually show that, for any $`\kappa >0`$, there exists $`\lambda _1(\kappa )>0`$ such that, for $`\lambda ]0,\lambda _1(\kappa )]`$, $`G_{\lambda ,\kappa }`$ has only the normal solutions as solutions of the Ginzburg-Landau equations. This will be analyzed in Section 5.
Proof of Theorem 2.6:
Let $`(u,A):=(u,A_e+a)`$ be a minimizer of the $`(GL)`$ functional. So it is a solution<sup>3</sup><sup>3</sup>3 We actually do not use this property in the proof. of (GL) and moreover we have, using (1.6), the following property :
$$G_{\lambda ,\kappa }(u,A)0.$$
(2.3)
Using the inequality $`|u|^2\frac{1}{2}|u|^4\frac{1}{2}`$ and (2.3), we first get, with $`b=\mathrm{rot}a`$ :
$$\frac{\kappa ^2}{\lambda }_{IR^2}b^2๐x\frac{\lambda }{2}|\mathrm{\Omega }|,$$
(2.4)
where $`|\mathrm{\Omega }|`$ is the area of $`\mathrm{\Omega }`$.
We now discuss the link between $`b`$ and $`a`$ in $`\stackrel{~}{\mathrm{\Omega }}`$. So we shall only use from (2.4) :
$$\frac{\kappa ^2}{\lambda }_{\stackrel{~}{\mathrm{\Omega }}}b^2๐x\frac{\lambda }{2}|\mathrm{\Omega }|,$$
(2.5)
Let us now consider in $`\stackrel{~}{\mathrm{\Omega }}`$, $`\stackrel{~}{a}`$ the problem of finding a solution of
$$\begin{array}{cc}\mathrm{rot}\stackrel{~}{a}=b,\hfill & \mathrm{div}\stackrel{~}{a}=0,\hfill \\ \stackrel{~}{a}\nu =0,\text{ on }\stackrel{~}{\mathrm{\Omega }}.\hfill & \end{array}$$
(2.6)
We have the following standard proposition (see Lemma 2.3 in \[GiPh\]).
###### Proposition 2.8
.
The problem (2.6) admits, for any $`bL^2(\stackrel{~}{\mathrm{\Omega }})`$, a unique solution $`\stackrel{~}{a}`$ in $`H^1(\stackrel{~}{\mathrm{\Omega }})`$. Moreover, there exists a constant $`C`$ such that
$$\stackrel{~}{a}_{H^1(\stackrel{~}{\mathrm{\Omega }})}Cb_{L^2(\stackrel{~}{\mathrm{\Omega }})},bL^2.$$
(2.7)
Proof of Proposition 2.8.
Following a suggestion of F. Bethuel, we look for a solution in the form : $`\stackrel{~}{a}=\mathrm{rot}^{}\psi `$. We then solve the Dirichlet problem $`\mathrm{\Delta }\psi =b`$ in $`\stackrel{~}{\mathrm{\Omega }}`$. This gives a solution with the right regularity. For the uniqueness, we observe that $`\stackrel{~}{\mathrm{\Omega }}`$ being connected and simply connected a solution of $`\mathrm{rot}\widehat{a}=0`$ is of the form $`\widehat{a}=d\theta `$ (with $`\theta H^2(\stackrel{~}{\mathrm{\Omega }})`$, and if $`\mathrm{div}\widehat{a}=0`$ and $`\widehat{a}\nu `$ on $`\stackrel{~}{\mathrm{\Omega }}`$, we get the equations $`\mathrm{\Delta }\theta =0`$ and $`\theta \nu =0`$ on $`\stackrel{~}{\mathrm{\Omega }}`$, which implies $`\theta =\mathrm{const}.`$ and consequently $`\widehat{a}=0`$.
We can now use the Sobolev estimates in order to get
$$a_{L^4(\stackrel{~}{\mathrm{\Omega }})}C_1a_{H^1(\stackrel{~}{\mathrm{\Omega }})}.$$
(2.8)
ยฟFrom (2.5), (2.7) and (2.8), we get the existence of a constant $`C_2`$ such that
$$a_{L^4(\stackrel{~}{\mathrm{\Omega }})}C_2\frac{\lambda }{\kappa }.$$
(2.9)
The second point is to observe, that, for any $`ฯต]0,1[`$, we have the inequality
$$_\mathrm{\Omega }|(iA)u|^2๐x(1ฯต)(iA_e)u_{L^2(\mathrm{\Omega })}^2\frac{(1ฯต)}{ฯต}au_{L^2(\mathrm{\Omega })}^2.$$
(2.10)
Taking $`ฯต=\frac{1}{4}`$ and using Hรถlderโs inequality, we get
$$_\mathrm{\Omega }|(iA)u|^2๐x\frac{3}{4}(iA_e)u_{L^2(\mathrm{\Omega })}^23a_{L^4(\mathrm{\Omega })}^2u_{L^4(\mathrm{\Omega })}^2.$$
(2.11)
Using now the ellipticity of $`\mathrm{\Delta }_{A_e}`$ in the form of the existence of a constant $`C_1`$
$$u_{H^1(\mathrm{\Omega })}^2C_1\left((iA_e)u_{L^2(\mathrm{\Omega })}^2+u_{L^2(\mathrm{\Omega })}^2\right),$$
(2.12)
and again the Sobolev inequality, we then obtain the existence of a constant $`C_2`$ such that
$$_\mathrm{\Omega }|(iA)u|^2๐x\left(\frac{3}{4}C_2a_{L^4(\mathrm{\Omega })}^2\right)(iA_e)u_{L^2(\mathrm{\Omega })}^2C_2a_{L^4(\mathrm{\Omega })}^2u_{L^2(\mathrm{\Omega })}^2.$$
(2.13)
We get then from (2.3) and (2.9), and for a suitable new constant $`C`$ (depending only on $`\mathrm{\Omega }`$ and $`H_e`$),
$$\left[\frac{3}{4}\lambda ^{(1)}C\frac{\lambda ^2}{\kappa ^2}\lambda \right]u_{L^2(\mathrm{\Omega })}^20.$$
(2.14)
Using the assumption (2.2), this gives $`u=0`$ for $`\lambda `$ small enough and the proof of Theorem 2.6.
###### Remark 2.9
.
Note that with a small improvement of the method, it is possible (taking $`ฯต=\frac{1}{\kappa }`$ in (2.10) ) to show that one can choose, in the limit $`\kappa +\mathrm{}`$, $`\lambda _0(\kappa )`$ satisfying :
$$\lambda _0(\kappa )\lambda ^{(1)}๐ช(\frac{1}{\kappa }).$$
(2.15)
This will be developped in Section 4.
###### Remark 2.10
.
Observing that $`\lambda \frac{1}{\lambda }G_{\lambda ,\kappa }(u,A)`$ is monotonically decreasing, one easily obtains, that the set of $`\lambda `$โs such that $`(0,A_e)`$ is a global minimum is an interval of the form $`]0,\lambda _0^{opt}(\kappa )]`$. Inequality (2.15) implies :
$$\lambda _0^{opt}(\kappa )\lambda ^{(1)}๐ช(\frac{1}{\kappa }).$$
(2.16)
Similar arguments are used in \[Du\] for the Abrikosovโs case. We recall that, in this case, the domain $`\mathrm{\Omega }`$, is replaced by a torus $`IR^2/`$ where $``$ is the lattice generated over $`ZZ^2`$ by two independent vectors of $`IR^2`$.
Observing now that $`\kappa G_{\lambda ,\kappa }(u,A)`$ is monotonically increasing, one easily obtains that $`\kappa \lambda _0^{opt}(\kappa )`$ is increasing. Using (2.16) and Proposition 3.1, one gets that $`\kappa \lambda _0^{opt}(\kappa )`$ is increasing from $`0`$ to $`\lambda ^{(1)}`$ for $`\kappa ]0,+\mathrm{}[`$.
## 3 Estimates in the case $`\kappa `$ small.
We have already shown in Proposition 2.5 that, if $`\lambda >\lambda ^{(1)}(\kappa )`$, then the normal state is not a minimizer. In other words (see Remark 2.10), under condition (2.2), we have :
$$0<\lambda _0^{opt}(\kappa )\lambda ^{(1)}.$$
(3.1)
If we come back to the formula (2.14), one immediately obtains the following first result :
###### Proposition 3.1
.
There exist constants $`\mu _0]0,\lambda ^{(1)}]`$ and $`\alpha _0>0`$ such that, for $`\lambda ]0,\mu _0]`$ satisfying
$$\lambda \alpha _0\kappa ,$$
(3.2)
the minimizer is necessarily the normal solution.
In order to get complementary results, it is also interesting to compute the energy of the pair $`(u,A)=(1,0)`$. This will give, in some asymptotic regime, some information about the possibility for the normal solution (or later for a bifurcating solution) to correspond to a global minimum of the functional. An immediate computation gives :
$$G_{\lambda ,\kappa }(1,0)=\frac{\lambda }{2}|\mathrm{\Omega }|+\frac{\kappa ^2}{\lambda }_{IR^2}H_e^2๐x.$$
(3.3)
We see in particular that when $`\frac{\kappa }{\lambda }`$ is small, the normal solution cannot be a global minimizer of $`G_{\lambda ,\kappa }`$.
As already observed in Subsection 1.2, what is more relevant is probably the integral $`_{\stackrel{~}{\mathrm{\Omega }}}H_e^2๐x`$ instead of $`_{IR^2}H_e^2๐x`$ in (3.3). Note also that it would be quite interesting to determine the minimizers in the limit $`\kappa 0`$. We note indeed that $`(1,0)`$ is not a solution of the GL-system, unless $`H_e`$ is identically zero in $`\mathrm{\Omega }`$. Let us show the following proposition.
###### Proposition 3.2
.
If
$$\kappa <\lambda \left(\frac{|\mathrm{\Omega }|}{2_\mathrm{\Omega }H_e^2๐x}\right)^{\frac{1}{2}},$$
(3.4)
and if $`\mathrm{\Omega }`$ is simply connected then the normal solution is not a global minimum.
Proof.
Let $`\psi _n`$ a sequence of $`C^{\mathrm{}}`$ functions such that
* $`0\psi _n1`$;
* $`\psi _n=0`$ in a neighborhood of $`\overline{\mathrm{\Omega }}`$;
* $`\psi _n(x)1,x\overline{\mathrm{\Omega }}`$ ;
We observe that
$$_{IR^2}((1\psi _n)H_e)^2๐x_\mathrm{\Omega }H_e^2๐x.$$
(3.5)
We can consequently choose $`n`$ such that :
$$\kappa <\lambda \left(\frac{|\mathrm{\Omega }|}{2_{IR^2}\left((1\psi _n)H_e\right)^2๐x}\right)^{\frac{1}{2}},$$
(3.6)
We now try to find $`A_n`$ such that
* $`\mathrm{rot}A_n=\psi _nH_e`$;
* $`suppA_n\mathrm{\Omega }=\mathrm{}`$.
We have already shown how to proceed when $`\mathrm{\Omega }`$ is starshaped. In the general case, we first choose $`\stackrel{~}{A}_n`$ such that : $`\mathrm{rot}\stackrel{~}{A}_n=\psi _nH_e`$, without the condition of support (see (1.2) for the argument).
We now observe that $`\mathrm{rot}\stackrel{~}{A}_n=0`$ in $`\mathrm{\Omega }`$. Using the simple connexity, we can find $`\varphi _n`$ in $`C^{\mathrm{}}(\overline{\mathrm{\Omega }})`$ such that $`\stackrel{~}{A}_n=\varphi _n`$. We can now extend $`\varphi _n`$ outside $`\mathrm{\Omega }`$ as a compactly supported $`C^{\mathrm{}}`$ function in $`IR^2`$ $`\stackrel{~}{\varphi }_n`$. We then take $`A_n=\stackrel{~}{A}_n\stackrel{~}{\varphi }_n`$.
It remains to compute the energy of the pair $`(1,A_n)`$ (which is strictly negative) in order to achieve the proof of the proposition.
###### Remark 3.3
.
In the case when $`\mathrm{\Omega }`$ is not simply connected. Proposition 3.2 remains true, if we replace $`\mathrm{\Omega }`$ by $`\stackrel{~}{\mathrm{\Omega }}`$, where $`\stackrel{~}{\mathrm{\Omega }}`$ is the smallest simply connected open set containing $`\mathrm{\Omega }`$.
###### Remark 3.4
.
It would be interesting to see how one can use the techniques of \[AfDa\] for analyzing the properties of the zeros of the minimizers, when they are not normal solutions. The link between the two papers is given by the relation $`\lambda =(\kappa d)^2`$.
In conclusion, we have obtained, the following theorem :
###### Theorem 3.5
.
Under condition (2.2), there exists $`\alpha _0>0`$, such that :
$$\left(\frac{|\mathrm{\Omega }|}{2_{\stackrel{~}{\mathrm{\Omega }}}H_e^2๐x}\right)^{\frac{1}{2}}\frac{\lambda _0^{opt}(\kappa )}{\kappa }inf(\alpha _0,\frac{\lambda ^{(1)}}{\kappa }).$$
(3.7)
## 4 Localization of pairs with small energy, in the case $`\kappa `$ large.
When $`\kappa `$ is large and $`\lambda \lambda ^{(1)}`$ is small enough, we will show as in \[Du\] that all the solutions of non positive energy of the GL-systems are in a suitable neighborhood of $`(0,A_e)`$ independent of $`\kappa \kappa _0>0`$. This suggests that in this limiting regime these solutions of the GL-equations (if there exist and if they appear as local minima) will furnish global minimizers. Let us show this localization statement. The proof is quite similar to the proof of Theorem 2.6. We recall that we have (2.4)-(2.10). Now we add the condition that, for some $`\eta >0`$,
$$\lambda \lambda ^{(1)}+\eta .$$
(4.1)
Note that we have already solved the problem when $`\lambda \lambda ^{(1)}\frac{C}{\kappa }`$, so we are mainly interested in the $`\lambda `$โs in an interval of the form $`[\lambda ^{(1)}\frac{C}{\kappa },\lambda ^{(1)}+\eta ]`$.
The second assumption is that we consider only pairs $`(u,A)H^1(\mathrm{\Omega })\times H_{loc}^1(IR^2)`$ such that
$$G_\lambda (u,A)0.$$
(4.2)
We improve (2.10) into
$$||(iA)u||_{L^2(\mathrm{\Omega })}^2((1ฯต\frac{C}{ฯต}||a||_{L^4(\mathrm{\Omega })}^2)__+\lambda ^{(1)}\frac{C}{ฯต}||a||_{L^4(\mathrm{\Omega })}^2))||u||_{L^2(\mathrm{\Omega })}^2.$$
(4.3)
Taking $`ฯต=\frac{1}{\kappa }`$, we get, using also (2.9), the existence of $`\kappa _0`$ and $`C`$ such that, for $`\lambda [0,\lambda ^{(1)}+\eta ]`$ and for $`\kappa \kappa _0`$,
$$||(iA)u||^2(((1\frac{C}{\kappa })\lambda ^{(1)}\frac{C}{\kappa })||u||^2,$$
(4.4)
for any $`(u,A)`$ such that $`G_{\lambda ,\kappa }(u,A)0`$.
Coming back to (1.1), and, using again the negativity of the energy $`G_{\lambda ,\kappa }(u,A)`$ of the pair $`(u,A)`$, we get
$$\lambda _\mathrm{\Omega }|u|^4๐x(\eta +\frac{C}{\kappa })u_{L^2(\mathrm{\Omega })}^2.$$
(4.5)
But by Cauchy-Schwarz, we have
$$_\mathrm{\Omega }|u|^2๐x|\mathrm{\Omega }|^{\frac{1}{2}}(_\mathrm{\Omega }|u|^4๐x)^{\frac{1}{2}}.$$
(4.6)
So we get
$$u_{L^2(\mathrm{\Omega })}(\frac{|\mathrm{\Omega }|}{\lambda })^{\frac{1}{2}}(\eta +\frac{C}{\kappa })^{\frac{1}{2}}$$
(4.7)
We see that this becomes small with $`\eta `$ and $`\frac{1}{\kappa }`$. It is then also easy to control the norm of $`u`$ in $`H^1(\mathrm{\Omega })`$. We can indeed use successively (2.12), (2.13), (4.2) and the trivial inequality:
$$(iA)u_{L^2(\mathrm{\Omega })}^2\lambda u_{L^2(\mathrm{\Omega })}^2+G_\lambda (u,A).$$
(4.8)
The control of $`(AA_e)`$ in the suitable choice of gauge is also easy through (2.4) and (2.7).
Note also that if $`\lambda <\lambda ^{(1)}`$, we obtain the better
$$u_{L^2(\mathrm{\Omega })}\frac{C}{\kappa \lambda ^{\frac{1}{2}}}.$$
(4.9)
So we have shown in this section the following theorem:
###### Theorem 4.1
.
There exists $`\eta _0>0`$ such that, for $`0<\eta <\eta _0`$ and for $`\lambda \lambda ^{(1)}+\eta `$, then there exists $`\kappa _0`$ such that for $`\kappa \kappa _0`$, all the pairs $`(u,A)`$ with negative energy are in a suitable neighborhood $`๐ช(\eta ,\frac{1}{\kappa })`$ of the normal solution in $`H^1(\mathrm{\Omega },C\text{ }\text{ })\times H^1(\mathrm{\Omega },IR^2)`$ whose size tends to $`0`$ with $`\eta `$ and $`\frac{1}{\kappa }`$.
###### Remark 4.2
.
Using the same techniques as in \[Du\], one can also show that there are no solutions of the Ginzburg-Landau equations outside this neighborhood. This is discussed in Section 5.
## 5 A priori localization for solutions of Ginzburg-Landau equations
In this section, we give the proof of Remarks 2.7 and 4.2. The proof is adapted from Subsection 4.4 in \[Du\] which analyzes the Abrikosov situation. Similar estimates can also be found in \[GiPh\] (or in \[AfDa\]) but in a different asymptotical regime.
We assume that $`(u,A)`$ is a pair of solutions of the Ginzburg-Landau equations (1.3) and rewrite the second Ginzburg-Landau equation, with $`A=A_e+a`$ in the form :
$$La=\frac{\lambda }{\kappa ^2}\mathrm{Im}\left(\overline{u}(i(A_e+a))u\right).$$
(5.1)
Here $`L`$ is the operator defined on the space $`E^2(\mathrm{\Omega })`$, where, for $`kIN^{}`$,
$$E^k(\mathrm{\Omega }):=\{aH^k(\mathrm{\Omega };IR^2)|\mathrm{div}a=0,a\nu _{/\mathrm{\Omega }}=0\},$$
(5.2)
by
$$L=\mathrm{rot}^{}\mathrm{rot}=\mathrm{\Delta }.$$
(5.3)
One can easily verify that $`L`$ is an isomorphism from $`E^2(\mathrm{\Omega })`$ onto $`L^2(\mathrm{\Omega })`$. One first gets the following
###### Lemma 5.1
.
If $`(u,A_e+a)`$ is a solution of the GL-system (1.3) for some $`\lambda >0`$, then we have :
$$La\frac{|\mathrm{\Omega }|^{\frac{1}{2}}\lambda ^{\frac{3}{2}}}{\kappa ^2}.$$
(5.4)
Proof of Lemma.
We start from (5.1) and using Proposition 1.4, we obtain :
$$La^2\frac{\lambda ^2}{\kappa ^4}(iA)u^2.$$
(5.5)
Using the first GL-equation, we obtain :
$$La^2\frac{\lambda ^3}{\kappa ^4}_\mathrm{\Omega }|u|^2(1|u|^2)๐x$$
(5.6)
Using again Proposition 1.4, we obtain the lemma.
So Lemma 5.1 shows, together with the properties of $`L`$, that there exists a constant $`C_\mathrm{\Omega }`$ such that
$$a_{H^2(\mathrm{\Omega })}C_\mathrm{\Omega }\frac{\lambda ^{\frac{3}{2}}}{\kappa ^2}.$$
(5.7)
This permits to control the size of $`a`$ when $`\lambda `$ is small or $`\kappa `$ is large. In particular, using Sobolevโs injection Theorem, we get the existence of a constant $`C_\mathrm{\Omega }^{}`$ such that :
$$a_{L^{\mathrm{}}(\mathrm{\Omega })}C_\mathrm{\Omega }^{}\frac{\lambda ^{\frac{3}{2}}}{\kappa ^2}.$$
(5.8)
The second step consists in coming back to our solution $`(u,A)`$ of the Ginzburg-Landau equations. Let us rewrite the first one in the form :
$$\mathrm{\Delta }_{A_e}u=\lambda u(1|u|^2)2ia(iA_e)u|a|^2u.$$
(5.9)
Taking the scalar product with $`u`$ in $`L^2(\mathrm{\Omega })`$, we obtain:
$$\begin{array}{cc}\lambda |u|^2^2+\mathrm{\Delta }_{A_e}u,u\hfill & \lambda u^2+2a_L^{\mathrm{}}u\sqrt{\mathrm{\Delta }_{A_e}u,u}+a_L^{\mathrm{}}^2u^2\hfill \\ & (\lambda +(1+\frac{1}{ฯต})a_L^{\mathrm{}}^2)u^2+ฯต\mathrm{\Delta }_{A_e}u,u.\hfill \end{array}$$
We have finally obtained, for any $`ฯต]0,1[`$, and any pair $`(u,A)`$ solution of the GL-equations the folllowing inequality :
$$\lambda _\mathrm{\Omega }|u(x)|^4๐x+\mathrm{\Delta }_{A_e}u,u\frac{1}{1ฯต}(\lambda +(1+\frac{1}{ฯต})a_L^{\mathrm{}}^2)u^2.$$
(5.10)
Forgetting first the first term of the left hand side in (5.10), we get the following alternative :
* Either $`u=0`$,
* or
$$\lambda ^{(1)}\frac{1}{1ฯต}(\lambda +(1+\frac{1}{ฯต})a_L^{\mathrm{}}^2).$$
If we are in the first case, we obtain immediately (see (5.1), the equation $`La=0`$ and consequently $`a=0`$. So we have obtained that $`(u,A)`$ is the normal solution.
The analysis of the occurence or not of the second case depends on the assumptions done in the two remarks, through (5.8) and for a suitable choice of $`ฯต`$ ($`ฯต=\frac{1}{k}`$). So we get immediately the existence of $`\lambda _1(\kappa )`$ and its estimate when $`\kappa +\mathrm{}`$. If we now assume (see (4.1)) that $`\lambda ]\lambda ^{(1)}\eta ,\lambda ^{(1)}+\eta [`$, we come back to (5.10) and write :
$$\lambda _\mathrm{\Omega }|u(x)|^4๐x\left(\frac{1}{1ฯต}(\lambda +(1+\frac{1}{ฯต})a_L^{\mathrm{}}^2)\lambda ^1\right)u^2.$$
Using (4.6), this leads to
$$\lambda u^2\left(\frac{1}{1ฯต}(\lambda +(1+\frac{1}{ฯต})a_L^{\mathrm{}}^2)\lambda ^{(1)}\right)_+|\mathrm{\Omega }|.$$
(5.11)
This shows, as in (4.7), that $`u`$ is small in $`L^2`$ with $`\eta `$ and $`\frac{1}{\kappa }`$.
We can then conclude as in the proof of Theorem 4.1. The control of $`u`$ in $`H^1`$ is obtained through (5.10).
###### Theorem 5.2
.
There exists $`\eta _0>0`$ such that, for $`0<\eta <\eta _0`$ and for $`\lambda \lambda ^{(1)}+\eta `$, then there exists $`\kappa _0`$ such that for $`\kappa \kappa _0`$, all the pairs $`(u,A)`$ solutions of the (G-L)-equations are in a suitable neighborhood $`๐ช(\eta ,\frac{1}{\kappa })`$ of the normal solution in $`H^1(\mathrm{\Omega },C\text{ }\text{ })\times H^1(\mathrm{\Omega },IR^2)`$ whose size tends to $`0`$ with $`\eta `$ and $`\frac{1}{\kappa }`$.
## 6 About bifurcations and stability
### 6.1 Preliminaries
Starting from one normal solution, a natural way is to see, if, when increasing $`\lambda `$ from $`0`$, one can bifurcate for a specific value of $`\lambda `$. Proposition 2.4 shows that it is impossible before $`\lambda ^{(1)}`$. A necessary condition is actually that $`\lambda `$ becomes an eigenvalue of the Neumann realization of $`\mathrm{\Delta }_{A_e}`$ in $`\mathrm{\Omega }`$. We shall consider what is going on at $`\lambda ^{(1)}`$.
Note here that there is an intrinsic degeneracy to the problem related to the existence of an $`S^1`$ action. We have indeed the trivial lemma
###### Lemma 6.1
.
If $`(u,A)`$ is a solution, then $`(\mathrm{exp}i\theta u,A)`$ is a solution.
This degeneracy is independent of the gauge degeneracy.
In order to go further, we add the assumption
$$\lambda ^{(1)}\text{ is a simple eigenvalue}.$$
(6.1)
In this case, we denote by $`u_1`$ a corresponding normalized eigenvector.
Now, one can try to apply the general bifurcation theory due to Crandall-Rabinovitz. Note that, although, the eigenvalue is assumed to be simple, it is not exactly a simple eigenvalue in the sense of Crandall-Rabinowitz which are working with real spaces. Actually, this is only simple modulo this $`S^1`$-action. We are not aware of a general theory dealing with this situation in full generality (see however \[GoSc\]) but special cases involving Schrรถdinger operators with magnetic field are treated in \[Od\], \[BaPhTa\] and \[Du\]. The article \[BaPhTa\] is devoted to the case of the disk and \[Od\] (more recently \[Du\]) to the case of Abrikosovโs states.
All the considered operators are (relatively to the wave function or order parameter) suitable realizations of operators of the type
$$u\mathrm{\Delta }_Au\lambda f(|u|^2)u,$$
with $`f(0)=1`$.
The main theorem is the following :
###### Theorem 6.2
.
Under the assumptions (2.2) and (6.1), there exist $`ฯต_0`$ and a bifurcating family of solutions $`(u(;\alpha ),A(;\alpha ),\lambda (\alpha ))`$ in $`H^1(\mathrm{\Omega },C\text{ }\text{ })\times E^1(\mathrm{\Omega })\times IR^+`$, with $`\alpha D(0,ฯต_0)C\text{ }\text{ }`$ for the Ginzburg-Landau equations such that
$$\begin{array}{ccc}\hfill u(;\alpha )& =\alpha u_1+\alpha |\alpha |^2u^{(3)}(;\alpha ),\hfill & \text{ with }u_1,u^{(3)}=0,\hfill \\ \hfill A(,\alpha )& =A_e+|\alpha |^2a_2+|\alpha |^4a^{(4)}(;\alpha ),\hfill & \\ \hfill \lambda (\alpha )& =\lambda ^{(1)}+c(\kappa )|\alpha |^2+๐ช(\alpha ^4).\hfill & \end{array}$$
(6.2)
Here $`u^{(3)}(;\alpha )`$ and $`a^{(4)}(;\alpha )`$ are bounded in $`H^1`$.
This solution satisfies, $`sC\text{ }\text{ },|s|=1`$ :
$$u(;s\alpha )=su(;\alpha ),A(;s\alpha )=A(;\alpha ).$$
(6.3)
Moreover, if $`c(\kappa )0`$, all the solutions $`(u,A,\lambda )`$ of the Ginzburg-Landau equations lying in a sufficiently small neighborhood in $`H^1\times E^1\times IR^+`$ of $`(0,A_e,\lambda ^{(1)})`$ are described by the normal solutions $`(0,A_e,\lambda )`$ and the bifurcating solutions.
The constant $`c(\kappa )`$ will be explicited in the next subsection.
### 6.2 About the proof, construction of formal solutions.
The starting point is the GL-system written in the form
$$\begin{array}{cc}\hfill (\mathrm{\Delta }_{A_e}\lambda ^{(1)})u& =(\lambda \lambda ^{(1)})u\lambda u|u|^22ia(iA_e)ua^2u\hfill \\ \hfill La& =\frac{\lambda }{\kappa ^2}\mathrm{Im}\left(\overline{u}(iA)u\right)\hfill \end{array}$$
(6.4)
We then use the standard method. We look for a solution in the form
$$u=\alpha u_1+\alpha |\alpha |^2u_3+๐ช(\alpha ^5),$$
$$a=|\alpha |^2a_2+๐ช(\alpha ^4)$$
and
$$\lambda (\alpha )=\lambda ^{(1)}+c(\kappa )|\alpha |^2+๐ช(\alpha ^4).$$
We can eliminate the $`S^1`$-degeneracy by imposing $`\alpha `$ real (keeping only the parity). We refer to \[Du\] for details and just detail the beginning of the formal proof which gives the main conditions. We first obtain, using the second equation,
$$a_2=\frac{\lambda ^{(1)}}{\kappa ^2}b_2,$$
(6.5)
with
$$b_2:=L^1\mathrm{Im}\left(\overline{u}_1(iA_e)u_1\right).$$
(6.6)
Taking then the scalar product in $`L^2`$ with $`u_1`$, in the first equation, we get that
$$c(\kappa )=\lambda ^{(1)}\left(I_0\frac{2}{\kappa ^2}K_0\right),$$
(6.7)
with
$$I_0:=_\mathrm{\Omega }|u_1(x)|^4๐x,$$
(6.8)
and
$$K_0=ib_2(iA_e)u_1,u_1.$$
(6.9)
###### Remark 6.3
.
ยฟFrom this expression for $`c(\kappa )`$, we immediately see that there exists $`\kappa _1`$ such that, for $`\kappa \kappa _1`$, $`c(\kappa )>0`$. Moreover, the uniqueness statement in Theorem 6.2 is true in a neighborhood which can be chosen independently of $`\kappa [\kappa _1,+\mathrm{}[`$.
Let us now observe, that, $`b_2`$ being divergence free, it is immediate by integration by part that $`K_0`$ is real. Computing $`\mathrm{Re}K_0`$, we immediately obtain :
$$K_0=\mathrm{Re}K_0=L^1J_1,J_1,$$
(6.10)
where $`J_1`$ is the current :
$$J_1:=\mathrm{Im}\left(\overline{u}_1(iA_e)u_1\right).$$
(6.11)
We observe that $`K_0>0`$ if and only if $`J_1`$ is not identically $`0`$. In the non simply connected case, we shall find a case when $`J_1=0`$. (See Lemma 7.2).
Following the argument of \[Du\] (Lemme 3.4.9), let us analyze the consequences of $`J_1=0`$. By assumption $`u_1`$ does not vanish identically. If $`u_1(x_0)0`$, then we can perform in a sufficiently small ball $`B(x_0,r_0)`$ centered at $`x_0`$, the following computation in polar coordinates. We write $`u_1=r(x)\mathrm{exp}i\theta (x)`$ and get $`J_1=r(x)^2(A_e\theta )=0`$. So $`A_e=\theta `$ in this ball and this implies $`H_e=0`$ in the same ball. Using the properties of the zero set of $`u_1`$ in $`\mathrm{\Omega }`$ \[ElMaQi\] and the continuity of $`H_e`$, we then obtain $`H_e=0`$ in $`\mathrm{\Omega }`$. But we know that, if $`\mathrm{\Omega }`$ is simply connected, then this implies $`\lambda ^{(1)}=0`$. So we have the following lemma
###### Lemma 6.4
.
If $`\mathrm{\Omega }`$ is simply connected and $`\lambda ^{(1)}>0`$, then $`K_0>0`$.
.
Coming back to the first equation and projecting on $`u_1^{}`$, we get :
$$u_3=R_0v_3,$$
(6.12)
where $`v_3`$ is orthogonal to $`u_1`$ and given by :
$$v_3:=2a_2\left((iA_e)u_1\right),$$
(6.13)
and $`R_0`$ is the inverse of $`(\mathrm{\Delta }_{A_e}\lambda ^{(1)})`$ on the space $`u_1^{}`$ and satisfies
$$R_0u_1=0.$$
We emphasize that all this construction is uniform with the parameter $`\beta =\frac{1}{\kappa }`$ in $`]0,\beta _0]`$. One can actually extend analytically the equation in order to have a well defined problem in $`[\beta _0,\beta _0]`$.
### 6.3 About the energy along the bifurcating solution.
The proof is an adaptation of \[Du\]. Let us just present here the computation of the value of the GL-functional along the bifurcating curve. Although it is not the proof, this gives the right condition for the stability. For this, we observe that if $`(u,A_e+a)`$ is a solution of the GL-system, then we have:
$$G_{\lambda ,\kappa }(u,A)=\frac{\lambda }{2}_\mathrm{\Omega }|u|^4+\frac{\kappa ^2}{\lambda }_\mathrm{\Omega }|\mathrm{rot}a|^2๐x.$$
(6.14)
It is then easy to get the main term of the energy of the function for $`(u,A_e+a)`$ with $`a(;\alpha )=|\alpha |^2a_2()+๐ช(\alpha ^4)`$ and $`u(;\alpha )=\alpha u_1()+๐ช(\alpha ^3)`$.
$$G_{\lambda ,\kappa }(u(;\alpha ),A(;\alpha ))=|\alpha |^4\left(\frac{\lambda ^{(1)}}{2}_\mathrm{\Omega }|u_1|^4+\frac{\kappa ^2}{\lambda ^{(1)}}_\mathrm{\Omega }|\mathrm{rot}a_2|^2๐x\right)+๐ช(\alpha ^6).$$
(6.15)
Let us first analyze the structure of the term :
$$K_1:=\frac{\kappa ^2}{\lambda ^{(1)}}_\mathrm{\Omega }|\mathrm{rot}a_2|^2๐x=\frac{\kappa ^2}{\lambda ^{(1)}}La_2,a_2.$$
(6.16)
But we have :
$$K_1:=\frac{\lambda ^{(1)}}{\kappa ^2}Lb_2,b_2=\frac{\lambda ^{(1)}}{\kappa ^2}L^1J_1,J_1=\frac{\lambda ^{(1)}}{\kappa ^2}K_0.$$
(6.17)
With these expressions, we get
$$G_{\lambda (\alpha ),\kappa }(u(;\alpha ),A(;\alpha ))=|\alpha |^4\frac{\lambda ^{(1)}}{2}\left(I_0\frac{2}{\kappa ^2}K_0\right)+๐ช(\alpha ^6).$$
(6.18)
So we get that the energy becomes negative along the bifurcating solution for $`0<|\alpha |\rho _0`$, if the following condition is satisfied :
$$\kappa ^2>2\frac{K_0}{I_0}.$$
(6.19)
Another way of writing the result is :
###### Proposition 6.5
.
Under conditions (2.2) and (6.1), then, if
$$\kappa ^22\frac{K_0}{I_0},$$
(6.20)
there exists $`\alpha _0>0`$ such that, for all $`\alpha `$ satisfying $`0<|\alpha |\alpha _0`$,
$$(\lambda (\alpha )\lambda ^{(1)})G_{\lambda (\alpha ),\kappa }(u(;\alpha ),A(;\alpha ))<0.$$
(6.21)
In particular, we have shown, in conjonction with Theorem 4.1, the following theorem :
###### Theorem 6.6
.
There exists $`\eta >0`$ and $`\kappa _0`$, such that, for $`\kappa >\kappa _0`$ and $`\lambda \lambda ^{(1)}+\eta `$, the global minimum of $`G_{\lambda ,\kappa }`$ is realized by the normal solution for $`\lambda ]0,\lambda ^{(1)}]`$ and by the bifurcating solution for $`\lambda ]\lambda ^{(1)},\lambda ^{(1)}+\eta ]`$.
In particular, and taking account of Remark 2.10, we have
###### Corollary 6.7
.
There exists $`\kappa _c`$ such that the map $`\kappa \lambda _0^{opt}(\kappa )`$ is an increasing function from $`0`$ to $`\lambda ^{(1)}`$ for $`\kappa [0,\kappa _c]`$ and is constant and equal to $`\lambda ^{(1)}`$ for $`\kappa \kappa _c`$.
###### Remark 6.8
.
Note that Theorem 5.2 gives an additional information. For $`\eta `$ small enough and $`\kappa `$ large enough, there are actually no other solutions of the (GL)-equation.
### 6.4 Stability
The last point is to discuss the stability of the bifurcating solution. We expect that the bifurcating solution gives a local minimum of the GL-functional for $`\kappa `$ large enough, and more precisely under condition (6.19). The relevant notion is here the notion of strict stability. Following \[BaPhTa\], we say that $`(u,A)`$ (with $`u`$ not identically $`0`$) is strictly stable for $`G_{\lambda ,\kappa }`$ if it is a critical point, if its Hessian is positive and if its kernel in $`H^1\times E^1`$ is the one dimensional space $`IR(iu,\mathrm{\hspace{0.33em}0})`$.
We then have the following theorem :
###### Theorem 6.9
.
Under conditions (2.2),(6.1), and if (6.19) is satisfied, then there exists $`ฯต_0>0`$, such that, for $`0<|\alpha |ฯต_0`$, the solution $`(u(;\alpha ),A(;\alpha ))`$ is strictly stable.
We refer to \[Du\] for the detailed proof.
## 7 Bifurcation from normal solutions: special case of non simply connected models.
### 7.1 Introduction
In this section, we revisit the bifurcation problem in the case when $`\mathrm{\Omega }`$ is not simply connected and when the external field vanishes inside $`\mathrm{\Omega }`$. In this very particular situation which was considered by J. Berger and J. Rubinstein in \[BeRu\] (and later in \[HHOO1\], \[HHOO2\]), it is interesting to make a deeper analysis leading for example to the description of the nodal sets of the bifurcating solution. The situation is indeed quite different of the results obtained by \[ElMaQi\] in a near context (but with a simply connected $`\mathrm{\Omega }`$). We mainly follow here the presentation in \[HHOO2\] (for which we refer for other results or points of view) but emphasize on the link with the previous section.
### 7.2 The operator $`K`$
We shall now consider the specific problem introduced by \[BeRu\] and consider the case
$$suppH_e\overline{\mathrm{\Omega }}=\mathrm{},$$
(7.1)
and, in any hole $`๐ช_i`$, the flux of $`H_e`$ satisfies
$$\frac{1}{2\pi }_{๐ช_i}H_eZZ+\frac{1}{2}.$$
(7.2)
Here we recall that a hole associated to $`\mathrm{\Omega }`$ is a bounded connected component of the complementary of $`\mathrm{\Omega }`$.
We recall in this context, what was introduced in \[HHOO1\]. We observe that under conditions (7.1) and (7.2), there exists a multivalued function $`\varphi `$ such that $`\mathrm{exp}i\varphi C^{\mathrm{}}(\overline{\mathrm{\Omega }})`$ and
$$d\varphi =2\omega _A,$$
(7.3)
where $`\omega _A`$ is the $`1`$-form naturally attached to the vector $`A`$.
We also observe that, for the complex conjugation operator $`\mathrm{\Gamma }`$
$$\mathrm{\Gamma }u=\overline{u},$$
(7.4)
we have the general property
$$\mathrm{\Gamma }\mathrm{\Delta }_A=\mathrm{\Delta }_A\mathrm{\Gamma }.$$
(7.5)
Combining (7.3) and (7.5), we obtain, for the operator
$$K:=(\mathrm{exp}i\varphi )\mathrm{\Gamma },$$
(7.6)
which satisfies
$$K^2=Id,$$
(7.7)
the following commutation relation
$$K\mathrm{\Delta }_A=\mathrm{\Delta }_AK.$$
(7.8)
Let us also observe that the Neumann condition is respected by $`K`$.
As a corollary, we get
###### Lemma 7.1
.
If $`v`$ is an eigenvector of $`\mathrm{\Delta }_A^N`$, then $`Kv`$ has the same property.
This shows that one can always choose an orthonormal basis of eigenvectors $`u_j`$ such that $`Ku_j=u_j`$.
### 7.3 Bifurcation inside special classes.
Following \[BeRu\] (but inside our point of view), we look for solution of the GL equation in the form $`(u,A_e)`$ with $`Ku=u`$. Let us observe that
$$L_K^2(\mathrm{\Omega };C\text{ }\text{ }):=\{uL^2(\mathrm{\Omega };C\text{ }\text{ })|Ku=u\},$$
(7.9)
is a real Hilbert subspace of $`L^2(\mathrm{\Omega };C\text{ }\text{ })`$.
We denote by $`H_K^m`$ the corresponding Sobolev spaces :
$$H_K^m(\mathrm{\Omega };C\text{ }\text{ })=H^m(\mathrm{\Omega };C\text{ }\text{ })L_K^2.$$
(7.10)
We now observe the
###### Lemma 7.2
.
If $`uH_K^1`$, then $`\mathrm{Im}(\overline{u}(iA_e)u)=0`$ almost everywhere.
Proof of Lemma 7.2:
Let us consider a point where $`u0`$. Then we have $`u=\rho \mathrm{exp}i\theta `$ with $`2\theta =\varphi `$ modulo $`2\pi ZZ`$. Remembering that $`A_e=\frac{1}{2}\varphi `$, it is easy to get the property.
Once this lemma is proved, one immediately sees that $`(u,A_e)`$ (with $`Ku=u`$) is a solution of the GL system if and only if $`uH_K^1`$ and
$$\begin{array}{c}\mathrm{\Delta }_{A_e}u\lambda u(1|u|^2)=0,\hfill \\ (iA_e)u\nu =0,\text{ on }\mathrm{\Omega }.\hfill \end{array}$$
(7.11)
We shall call this new system the reduced GL-equation. But now we can apply the theorem by Crandall-Rabinowitz \[CrRa\]. By assumption (6.1), the kernel of $`(\mathrm{\Delta }_{A_e}\lambda ^{(1)})`$ is now a one-dimensional real subspace in $`L_K^2`$. Let us denote by $`u_1`$ a normalized โrealโ eigenvector. Note that $`u_1`$ is unique up to $`\pm 1`$. Therefore, we have the
###### Theorem 7.3
.
Under assumptions (6.1), (7.1) and (7.2), there exists a bifurcating family of solutions $`(u(;\alpha ),\lambda (\alpha ))`$ in $`H_K^1\times IR^+`$ with $`\alpha ]ฯต_0,+ฯต_0[`$, for the reduced GL-equation such that
$$\begin{array}{c}u(\alpha )=\alpha u_1+\alpha ^3v(\alpha ),\hfill \\ u_1,v(\alpha )_{L^2}=0,\hfill \\ v(\alpha )_{H^2(\mathrm{\Omega })}=๐ช(1),\hfill \end{array}$$
(7.12)
$$\lambda (\alpha )=\lambda ^{(1)}+c\alpha ^2+๐ช(\alpha ^4),$$
(7.13)
with
$$c=\lambda ^{(1)}_\mathrm{\Omega }|u_1|^4๐x.$$
(7.14)
Moreover
$$u(\alpha )=u(\alpha ),\lambda (\alpha )=\lambda (\alpha ).$$
(7.15)
###### Remark 7.4
.
Note that the property (7.15) is what remains of the $`S^1`$-invariance when one considers only โrealโ solutions.
Let us give here the formal computations of the main terms. If we denote by $`L_0`$ the operator $`L_0:=\mathrm{\Delta }_{A_e}\lambda ^{(1)}`$. Writing $`v(\alpha )=u_3+๐ช(\alpha )`$, we get :
$$(L_0c\alpha ^2)(\alpha u_1+\alpha ^3u_3)+(\lambda ^{(1)})\alpha ^3u_1|u_1|^2=๐ช(\alpha ^4).$$
Projecting on $`u_1`$, we get (7.14). Projecting on $`u_1^{}`$ and denoting by $`R_0`$ the operator equal to the inverse of $`L_0`$ on this subspace and to $`0`$ on $`\text{Ker }L_0`$, we get
$$u_3=\lambda ^{(1)}R_0(u_1|u_1|^2)=\lambda ^{(1)}R_0(u_1|u_1|^2cu_1).$$
(7.16)
###### Remark 7.5
.
By the uniqueness part in Theorem 6.2, we see that the solution $`(u(;\alpha ),A_e)`$ is actually the solution given in this theorem.
Another remark is that
$$G_{\lambda (\alpha ),\kappa }(u(\alpha ),A_e)=\frac{\lambda ^{(1)}}{2}\alpha ^4(_\mathrm{\Omega }|u_1(x)|^4๐x)+๐ช(|\alpha |^6),$$
(7.17)
so that when $`\alpha 0`$ the energy is decreasing. This is of course to compare with (6.16) (note that we have $`K_0=0`$). Once we have observed this last property, the local stability of the bifurcated solution near the bifurcation is clear.
The second result we would like to mention concerns the nodal sets. In the case when $`\mathrm{\Omega }`$ is simply connected, the analysis of the nodal set of $`u`$ when $`(u,A)`$ is a minimizer of the GL-functional is done in \[ElMaQi\], using the analyticity of the solutions of the GL-equation and techniques of Courant.
In the non simply connected case, very few results are known. The following theorem is true \[BeRu\], \[HHOO2\] :
###### Theorem 7.6
.
Under assumptions (2.9), (7.1) and (7.2), there exists $`ฯต_1>0`$ such that, for any $`\alpha ]0,ฯต_1]`$, the nodal set of $`u(\alpha )`$ in $`H_K^1`$ slits $`\overline{\mathrm{\Omega }}`$ in the sense of \[HHOO1\]. In particular, if there is only one hole, then the nodal set of $`u(\alpha )`$ consists exactly in one line joining the interior boundary and the exterior boundary.
An elegant way to recover these results (see \[HHOO1\], \[HHOO2\]) is to lift the situation to a suitable two-fold covering $`\mathrm{\Omega }^{}`$.
Acknowledgements: We would like to thank E. Akkermans, F. Bethuel, C. Bolley, G. Raugel, T. Riviรจre, S. Serfaty, M. and T. Hoffmann-Ostenhof for useful discussions. This work is partially supported by the TMR grant FMRX-CT 96-0001 of the European Union.
## Appendix A Analysis of the various scalings.
When considering asymptotical regimes, it is perhaps useful to have an interpretation in terms of the initial variables. According to the statistical interpretation of the Ginzburg-Landau functional (See for example \[BeRuSc\]), the starting point is the functional $`(\stackrel{~}{v},\stackrel{~}{A})(\stackrel{~}{v},\stackrel{~}{A})`$ with :
$$\begin{array}{cc}(\stackrel{~}{v},\stackrel{~}{A})\hfill & :=\frac{1}{8\pi }_{IR^2}|\mathrm{rot}\stackrel{~}{A}\stackrel{~}{H}_e|^2๐\stackrel{~}{x}\hfill \\ & +_\mathrm{\Omega }\frac{\mathrm{}^2}{4m}|(i\frac{2e}{c}\stackrel{~}{A})\stackrel{~}{u}|^2๐\stackrel{~}{x}\hfill \\ & +_\mathrm{\Omega }\left(a|\stackrel{~}{u}|^2+\frac{b}{2}|\stackrel{~}{u}|^4\right)๐\stackrel{~}{x}.\hfill \end{array}$$
Here $`a`$ is a parameter which is proportional to $`(TT_c)`$ (we are only interested in the case $`a<0`$) and $`b`$ is essentially independent of the temperature. The other parameters are standard : $`\mathrm{}=\frac{h}{2\pi }`$, $`h`$ is the Planck constant, $`e`$ is the charge of the electron and $`m`$ is the mass of the electron. With $`u=\frac{b}{|a|}\stackrel{~}{u}`$ and $`A=\frac{2e}{c}\stackrel{~}{A}`$, we obtain :
$$(\stackrel{~}{v},\stackrel{~}{A})=\frac{|a|\mathrm{}^2}{4mb}G_{\lambda ,\kappa }(u,A),$$
with $`H_e=\frac{2e}{\mathrm{}c}\stackrel{~}{H}_e`$, $`\lambda =\frac{4m|a|}{\mathrm{}^2}`$ and $`\kappa =\frac{mc}{e\mathrm{}}(\frac{b}{8\pi })^{\frac{1}{2}}`$. Here we emphasize that between the two functionals, no change of space variables is involved.
Let now compare with another standard representation of the Ginzburg-Landau functional. We make this time the change of variables $`x=\frac{\kappa }{\sqrt{\lambda }}\widehat{x}`$ and if we change $`u`$ and the $`1`$form corresponding to $`A`$ accordingly, we obtain the standard functional :
$$(\widehat{u},\widehat{A})=G_{\lambda ,\kappa }(u,A),$$
with
$$\begin{array}{cc}(\widehat{u},\widehat{A})\hfill & =\kappa ^2_{\widehat{\mathrm{\Omega }}}\left(|\widehat{u}|^2+\frac{1}{2}|\widehat{u}|^4\right)๐\widehat{x}\hfill \\ & +_{\widehat{\mathrm{\Omega }}}|(i\widehat{A})\widehat{u}|^2๐\widehat{x}\hfill \\ & +_{\widehat{\mathrm{\Omega }}}|\mathrm{rot}\widehat{A}\widehat{H}_e|^2๐\widehat{x},\hfill \end{array}$$
with
$$\begin{array}{cc}\widehat{H}_e\hfill & =\frac{\kappa ^2}{\lambda }H_e,\hfill \\ \widehat{\mathrm{\Omega }}\hfill & =\frac{\sqrt{\lambda }}{\kappa }\mathrm{\Omega }.\hfill \end{array}$$
Here we observe that the open set $`\mathrm{\Omega }`$ is not conserved in the transformation. We have to keep this in mind when comparing in the limit $`\kappa +\mathrm{}`$ the contributions of Sandier and Serfaty \[SaSe\] or \[LuPa1\] with the results presented in this paper. |
warning/0002/hep-th0002006.html | ar5iv | text | # References
Introduction.โ There are several reasons to study systems with interacting gravitational and non-Abelian gauge fields with the Lagrangian of the type
$$=\frac{1}{16\pi G}R\frac{1}{4g^2}F_{\mu \nu }^aF^{a\mu \nu }+\mathrm{}$$
(1)
Here dots stand for other possible fields, which can be scalars, Abelian vectors, etc. Also the gauge coupling constant $`g`$ can be position-dependent, which is the case when a dilaton is present. First, all modern models of particle physics include Yang-Mills fields as a basic element and, since the energies under consideration are climbing higher and higher, the effects of gravity should also be considered. Besides, as suggested by the recent brane-world scenarios, even at low energies gravity might be important . From the classical General Relativity point of view, the model (1) can be viewed as a natural generalization of the Einstein-Maxwell theory. This suggests studying its solutions and, in particular, checking whether the standard electrovacuum theorems still apply. Surprisingly, the answer to the latter question is negative โ the theory admits a large variety of new bizarre solutions, such as hairy black holes, say, whose existence manifestly contradicts the conventional wisdom like the no-hair conjecture . Unfortunately, most of these solutions are available only numerically, which is due to the high complexity of the equations.
Finally, string theory is perhaps the most natural place where models of the type (1) apply. The Lagrangian (1) arises in the low energy limit of heterotic string theory, and in type I and II string theories and M-theory compactified on manifolds with non-Abelian isometries. Supersymmetric string vacua play an important role in the analysis of string theory . However, apart from stringy monopoles and related solutions obtained via the heterotic five-brane construction (see also ), most of the literature is devoted to solutions with Abelian gauge fields. This is easily understood, since non-Abelian solutions are much more difficult to study than those from the Abelian sector. On the other hand, it is to be expected that also configurations with non-Abelian gauge fields will eventually play an important role. In particular, this is suggested by the AdS/CFT correspondence, which focuses on backgrounds of gauged supergravity models .
The tale of the EYMS sphaleron.โ As was mentioned above, only very few supergravity solutions with non-Abelian gauge fields are known explicitly. Of these we would like to describe one which is particularly interesting in view of the surprising developments it has caused. This solution was originally discovered numerically in the Einstein-Yang-Mills-Dilaton (EYMD) model for the gauge group SU(2) (in which case one has in Eq.(1) $`g=\mathrm{e}^\varphi `$, and there is also the standard kinetic term for the dilaton $`\varphi `$). The solution describes a static, spherically symmetric particle-like object with finite ADM mass and a globally regular geometry. The gauge field is essentially non-Abelian and purely magnetic with zero total charge โ due to the fast fall-off at infinity. It was realized that the solution is unstable , which however is not a drawback. On the contrary, it turns out that the solution admits an elegant interpretation as a sphaleron . Specifically, one can show that the solution relates to the top of the potential barrier between the topological vacua in the EYMD theory. This suggests that it might be responsible for some kind of fermion number non-conservation in the theory. Let us call this solution EYMD sphaleron.
A number of surprising regularities in the parameters of the solution have been observed . Specifically, the parameters fulfill simple algebraic relations with rational coefficients, which is usually not expected for numerically obtained solutions. It was therefore conjectured that the solution had a hidden symmetry, which hopefully might even help to obtain it analytically. A part of this symmetry was soon identified as the dilatational invariance , and this accounted for some of the relations. However, other relations remained unexplained, and all attempts to find the solution in a closed analytical form failed.
It was then natural to check whether supersymmetry could account for the remaining hidden symmetry. The immediate obstacle, however, was the instability of the solution, which made the existence of supersymmetry extremely unlikely. Nevertheless, it was still possible to look for gauged supergravities containing the EYMD theory or some deformation of it. Such theories might admit supersymmetric vacua similar to the EYMD sphaleron solution. The appropriate model was soon found and studied in โ it was the N=4 gauged SU(2)$`\times `$SU(2) supergravity, also known as the FreedmanโSchwarz (FS) model . This model is very similar to the EYMD theory, but contains in addition the dilaton potential
$$\mathrm{U}(\varphi )=\frac{1}{8}(g_{(1)}^2+g_{(2)}^2)\mathrm{e}^{2\varphi },$$
(2)
where $`g_{(1)}`$ and $`g_{(2)}`$ are the gauge coupling constants for the two non-Abelian gauge fields. It was shown in that the FS model admits essentially non-Abelian vacua of the BPS monopole type. These solutions are static, spherically symmetric, globally regular and have purely magnetic gauge field with unit magnetic charge โ even though there is no Higgs field in the theory. Next, they preserve 1/4 of the supersymmetries, and are probably stable, but not asymptotically flat โ due to the unbounded dilaton potential. The discovery of these solutions has enlarged somewhat the very narrow family of analytically known configurations for gravitating Yang-Mills fields. Moreover, it was shown in that the FS model itself can be obtained via dimensional reduction of the N=1, D=10 supergravity on $`S^3\times S^3`$ manifold. This established a direct link to string theory and showed that any solution of the FS model can be lifted to ten dimensions.
However, the situation with the EYMD sphaleron solution and its remaining hidden symmetry remained obscure. In order to obtain this solution within a gauged supergravity model it was necessary to get rid of the dilaton potential. However, the latter is present practically in all models. The following trick was therefore employed in : to truncate the FS model to the purely magnetic sector and then to pass to imaginary values of the gauge coupling constant $`g_{(2)}`$:
$$g_{(2)}ig_{(2)}.$$
(3)
For $`|g_{(2)}|=g_{(1)}`$ the potential vanishes. Surprisingly, such a formal trick does not destroy supersymmetry โ the replacement (3) in the FS fermion supersymmetry transformations leads to non-trivial Bogomolโnyi equations. These admit asymptotically flat solutions, one of which exactly coincides with the EYMD sphaleron ! To recapitulate, the EYMD sphaleron fulfills the first order Bogomolโnyi equations, and this explains all the remaining regularities in the parameters of the solution.
As the existence of Bogomolโnyi equations is usually related to supersymmetry, it was tempting to conclude that the solution is supersymmetric. However, it was unclear how to reconcile supersymmetry with the instability of the solution. Moreover, the replacement (3) is a rather formal trick, and it is unclear to what supergravity model it leads. The resolution of these puzzles is interesting. It was conjectured in and verified in that there is another, hitherto unknown consistent Euclidean gauged supergravity whose special truncation can be related to the purely magnetic sector of the FS model via the replacement (3).
The new supergravity can be obtained via a consistent dimensional reduction of the N=1, D=10 SUGRA on $`S^3\times AdS_3`$. Since the metric on the internal 6-space is not positive-definite, the timelike coordinate of the ten-dimensional space is viewed as one of the internal coordinates, while the remaining four-space becomes Euclidean. As a result, one obtains the Euclidean N=4, D=4 gauged supergravity called Euclidean Freedman-Schwarz (EFS) model. The fields of the new theory are similar to those of the FS model โ the bosonic sector contains gravitational field $`g_{\mu \nu }`$, two non-Abelian gauge fields $`A_\mu ^{(1)a}`$ and $`A_\mu ^{(2)a}`$ with coupling constants $`g_{(1)}`$ and $`g_{(2)}`$, the axion $`๐`$ and the dilaton $`\varphi `$. The principal new features of the theory is that it lives in Euclidean space, the gauge group is now SU(2)$`\times `$SU(1,1), and the dilaton potential is obtained from (2) via replacing the sum of the two gauge coupling constants by their difference, $`g_{(1)}^2g_{(2)}^2`$. The latter property is due to the opposite signs of the curvature of $`S^3`$ and $`AdS_3`$.
Now, one can consider โstatic and purely magneticโ (these notions do not have an invariant meaning in the Euclidean regime) sector of the EFS model. In this sector the SU(1,1) gauge field $`A_\mu ^{(2)a}`$ is zero, $`A_\mu ^{(1)a}`$ is โpurely magneticโ, and there is a hypersurface orthogonal โtimelikeโ Killing vector. Then one discovers that all bosonic equations and fermionic SUSY transformations are related to those of the static and purely magnetic sector of the FS model via the replacement (3). Next, one finds that the Euclidean theory admits a globally regular non-Abelian background, and this can be analytically continued to the Lorentzian sector, which gives precisely the EYMD sphaleron solution. This solves all puzzles: the EYMD sphaleron is not a supersymmetric solution. However, it becomes truly supersymmetric upon continuation to the Euclidean sector, simply via replacement time by imaginary time, $`tit`$. Both Euclidean and Lorentzian configurations fulfill the same system of Bogomolโnyi equations. The instability of the solution occurs only in the Lorentzian sector, where there is no supersymmetry.
Summarizing, the analysis of the EYMD sphaleron has led to the disclosure of the Kaluza-Klein origin of the FS model, to the supersymmetric monopoles, and to the new Euclidean gauged supergravity obtained from type I string theory via the dimensional reduction. The sphaleron can be analytically continued to the Euclidean sector to become the supersymmetric vacuum of the new supergravity. These features are briefly summarized in the following diagram.
Dimensional reduction of type I string theory.โ Let us sketch the essential steps leading to the above diagram. The details can be found in , and also in . The starting point is the bosonic part of the action of D=10, N=1 supergravity
$$S_{10}=\left(\frac{1}{4}\widehat{R}\frac{1}{2}_M\widehat{\varphi }^M\widehat{\varphi }\frac{1}{12}e^{2\widehat{\varphi }}\widehat{H}_{MNP}\widehat{H}^{MNP}\right)\sqrt{\widehat{g}}d^{10}\widehat{x}.$$
(4)
The idea is to find a parameterization of $`\widehat{g}_{MN}`$, $`\widehat{H}_{MNP}`$ and $`\widehat{\varphi }`$ in terms of four-dimensional variables which reduces the equations of motion for the action (4) to a consistent system of four-dimensional equations. As a first step, the tangent space metric is chosen in the form
$$\widehat{\eta }_{AB}=\mathrm{diag}(\stackrel{\eta _{\alpha \beta }}{\stackrel{}{\underset{4\mathrm{space}}{\underset{}{s^2,+1,+1,+1}}}},\underset{6\mathrm{space}}{\underset{}{\stackrel{\eta _{ab}^{(1)}}{\stackrel{}{+1,+1,+1}},\stackrel{\eta _{ab}^{(2)}}{\stackrel{}{+1,+1,s^2}}}}).$$
(5)
Here the parameter $`s`$ assumes two values: $`s=1`$ or $`s=i=\sqrt{1}`$. The two options correspond to the same theory in D=10 โ up to a renumbering of coordinates โ but to two different choices of the four-space. Keeping this parameter in all formulas, we shall be able to consider the two different cases together โ the reduction on $`S^3\times S^3`$ for $`s=i`$ and that on $`S^3\times AdS_3`$ for $`s=1`$.
The D=10 metric is expressed in terms of the vielbein
$$\widehat{g}^{MN}\frac{}{\widehat{x}^M}\frac{}{\widehat{x}^N}=\widehat{\eta }^{AB}\widehat{E}_A\widehat{E}_B=\eta ^{\alpha \beta }\widehat{E}_\alpha \widehat{E}_\beta +\underset{(\sigma )=1,2}{}\eta ^{(\sigma )ab}\widehat{E}_a^{(\sigma )}\widehat{E}_b^{(\sigma )},$$
(6)
where the vielbein vectors $`\widehat{E}_B=(\widehat{E}_\alpha ,\widehat{E}_a^{(\sigma )})`$ are
$$\widehat{E}_\alpha =\mathrm{e}^{3\varphi /4}\left(e_\alpha +\underset{(\sigma )=1,2}{}e_\alpha ^\mu A_\mu ^{(\sigma )a}\stackrel{~}{e}_a^{(\sigma )}\right),\widehat{E}_a^{(\sigma )}=\frac{g_{(\sigma )}}{\sqrt{2}}\mathrm{e}^{\varphi /4}\stackrel{~}{e}_a^{(\sigma )}.$$
(7)
Here the tetrad four-vectors $`e_\alpha e_\alpha ^\mu /x^\mu `$, the four-dilaton $`\varphi =\frac{1}{2}\widehat{\varphi }`$, and the fields $`A_\mu ^{(\sigma )a}`$ depend only on the four-coordinates $`x^\mu `$, the four-metric being $`g^{\mu \nu }=\eta ^{\alpha \beta }e_\alpha ^\mu e_\beta ^\nu `$. Vectors $`\stackrel{~}{e}_a^{(\sigma )}\stackrel{~}{e}_a^{(\sigma )i}/z^{(\sigma )i}`$ depend only on the internal six-coordinates $`z`$ and are assumed to be invariant vectors of the internal group space $`๐ข`$. Their commutation relations are $`[\stackrel{~}{e}_a^{(\sigma )},\stackrel{~}{e}_b^{(\sigma ^{})}]=\delta _{\sigma \sigma ^{}}f_{ab}^{(\sigma )c}\stackrel{~}{e}_c^{(\sigma )}`$ with constant $`f_{ab}^{(\sigma )c}`$. It is assumed that the internal space is a direct product of two group manifolds labeled by $`(\sigma )=1,2`$: $`๐ข`$=SU(2)$`\times `$SU(2) for $`s=i`$ and $`๐ข`$=SU(2)$`\times `$SU(1,1) for $`s=1`$. The structure constants for the two group factors are $`f_{ab}^{(\sigma )c}=\eta ^{(\sigma )cd}\epsilon _{dab}`$. The quantities $`A_\mu ^{(\sigma )a}`$ are regarded as the non-Abelian gauge fields with the field strength
$$F_{\mu \nu }^{(\sigma )a}=_\mu A_\nu ^{(\sigma )a}_\nu A_\mu ^{(\sigma )a}+f_{bc}^{(\sigma )a}A_\mu ^{(\sigma )b}A_\nu ^{(\sigma )c},$$
(8)
with $`g_{(\sigma )}`$ in (7) being the corresponding real gauge coupling constants.
The ten-dimensional antisymmetric tensor $`\widehat{H}_{MNP}`$ has the following non-vanishing components in the basis (7) (the value of $`(\sigma )`$ is the same for all internal indices):
$`\widehat{H}_{abc}^{(\sigma )}`$ $`=`$ $`{\displaystyle \frac{g_{(\sigma )}}{2\sqrt{2}}}\mathrm{e}^{3\varphi /4}\epsilon _{abc},`$
$`\widehat{H}_{a\alpha \beta }^{(\sigma )}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}g_{(\sigma )}}}\mathrm{e}^{5\varphi /4}\eta _{ab}^{(\sigma )}F_{\alpha \beta }^{(\sigma )b},`$
$`\widehat{H}_{\alpha \beta \gamma }`$ $`=`$ $`\mathrm{e}^{7\varphi /4}\epsilon _{\alpha \beta \gamma }^\delta e_\delta ^\mu _\mu ๐.`$ (9)
Here the axion $`๐`$ depends only on $`x^\mu `$, one has $`\epsilon ^{0123}=1`$, and $`F_{\alpha \beta }^{(\sigma )b}`$ are components of the gauge field strength with respect to the tetrad $`e_\alpha `$.
As a result, all D=10 fields are expressed in terms of the four-dimensional metric $`g_{\mu \nu }`$, two non-Abelian gauge fields $`A_\mu ^{(\sigma )a}`$, the axion $`๐`$ and the dilaton $`\varphi `$. These expressions are then inserted into the D=10 field equations for the action (4). It turns out that all these equations are fulfilled provided that the four-dimensional fields satisfy the equations derived from the following Lagrangian:
$`_4`$ $`=`$ $`{\displaystyle \frac{R}{4}}{\displaystyle \frac{1}{2}}_\mu \varphi ^\mu \varphi +{\displaystyle \frac{s^2}{2}}\mathrm{e}^{4\varphi }_\mu ๐^\mu ๐{\displaystyle \frac{1}{4}}\mathrm{e}^{2\varphi }{\displaystyle \underset{(\sigma )=1,2}{}}{\displaystyle \frac{1}{g_{(\sigma )}^2}}\eta _{ab}^{(\sigma )}F_{\mu \nu }^{(\sigma )a}F^{(\sigma )b\mu \nu }`$ (10)
$``$ $`{\displaystyle \frac{1}{2}}๐{\displaystyle \underset{(\sigma )=1,2}{}}{\displaystyle \frac{1}{g_{(\sigma )}^2}}\eta _{ab}^{(\sigma )}F_{\mu \nu }^{(\sigma )a}F^{(\sigma )b\mu \nu }+{\displaystyle \frac{1}{8}}(g_{(1)}^2s^2g_{(2)}^2)\mathrm{e}^{2\varphi }.`$
We thus obtain the four-dimensional theory upon a consistent dimensional reduction from D=10. Let us recall that the signature of the spacetime metric is $`(s^2,+1,+1,+1)`$, the internal metrics are $`\eta _{ab}^{(1)}=\delta _{ab}`$ and $`\eta _{ab}^{(2)}=\mathrm{diag}(+1,+1,s^2)`$. For $`s=i`$ the theory (10) exactly coincides with the bosonic sector of the Freedman-Schwarz model , which is the N=4 gauged SU(2)$`\times `$SU(2) supergravity. For $`s=1`$ we obtain the new theory โ Euclidean N=4 gauged SU(2)$`\times `$SU(1,1) supergravity. Notice that in the latter case the dilaton potential is proportional to the difference and not the sum of the two coupling constants.
The bosonic Lagrangian (10) is to be supplemented by the rules for computing the fermionic SUSY variations. These can be derived via the dimensional reduction of the D=10 SUSY variations
$`\delta \widehat{\psi }_M`$ $`=`$ $`\widehat{๐}_M\widehat{ฯต}{\displaystyle \frac{1}{48}}e^{\widehat{\varphi }}\left(\widehat{\mathrm{\Gamma }}_M^{SPQ}+9\delta _M^S\widehat{\mathrm{\Gamma }}^{PQ}\right)\widehat{H}_{SPQ}\widehat{ฯต},`$
$`\delta \widehat{\chi }`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}(\widehat{\mathrm{\Gamma }}^M_M\widehat{\varphi })\widehat{ฯต}{\displaystyle \frac{1}{12\sqrt{2}}}e^{\widehat{\varphi }}\widehat{\mathrm{\Gamma }}^{SPQ}\widehat{H}_{SPQ}\widehat{ฯต}.`$ (11)
Here $`\widehat{ฯต}`$ is the Majorana-Weyl spinor parameter of supersymmetry transformations, $`\widehat{๐}_M`$ being its covariant derivative, and $`\widehat{\mathrm{\Gamma }}^A`$ are D=10 gamma matrices. The idea now is to choose an explicit parameterization of $`\widehat{\mathrm{\Gamma }}^A`$ in terms of D=4 gamma matrices $`\gamma ^\mu `$ and the gauge group generators $`\text{T}_a^{(\sigma )}`$, where $`[\text{T}_a^{(\sigma )},\text{T}_b^{(\sigma )}]=\epsilon _{abc}\eta ^{(\sigma )cd}\text{T}_d^{(\sigma )}`$. Also 32-component Majorana-Weyl spinors $`\delta \widehat{\psi }_M`$, $`\delta \widehat{\chi }`$, and $`\widehat{ฯต}`$ are expressed in terms of the four-dimensional spinors $`\delta \psi _\mu `$, $`\delta \chi `$, and $`ฯต`$; details can be found in . Here each four-dimensional spinor, $`ฯต`$ say, can be thought of as a multiplet of D=4 Majorana spinors in the fundamental representation of the gauge group $`๐ข`$. One can write $`ฯตฯต_\kappa ^\mathrm{I}`$, where I$`=1,\mathrm{}4`$ is the group index and $`\kappa =1,\mathrm{}4`$ is the spinor index. The important moment is the D=4 Majorana condition, which is derived from the D=10 Majorana-Weyl condition:
$$ฯต^{}=ABฯต.$$
(12)
Here the asterisk denotes complex conjugation, and $`4\times 4`$ matrices $`A`$ and $`B`$ are such that
$$A\gamma ^\alpha A^1=\left(\gamma ^\alpha \right)^{},B\text{T}_a^{(\sigma )}B^1=\left(\text{T}_a^{(\sigma )}\right)^{},$$
(13)
and they fulfill also the condition of consistency of (12):
$$AA^{}BB^{}=1\mathrm{l}.$$
(14)
It is worth emphasizing that the Majorana condition (12) is consistent also in the Euclidean case. This is only possible due to the group degrees of freedom: because the group matrix $`B`$ is present in (13), (14). Specifically, in the Euclidean case each of the two factors in (14) has wrong sign, $`AA^{}=1`$ and $`BB^{}=1`$, but their product has correct sign.
Omitting further details, the reduction of (11) gives the following SUSY variation for the D=4 gaugino:
$`\delta \chi `$ $`=`$ $`\left({\displaystyle \frac{1}{\sqrt{2}}}\gamma ^\mu _\mu \varphi {\displaystyle \frac{1}{\sqrt{2}s}}\mathrm{e}^{2\varphi }\gamma _5\gamma ^\mu _\mu ๐\right)ฯต`$ (15)
$`+`$ $`{\displaystyle \frac{1}{2s}}\mathrm{e}^\varphi \left(s^{(1)}\gamma _5^{(2)}\right)ฯต+{\displaystyle \frac{1}{4s}}\mathrm{e}^\varphi \left(sg_{(1)}g_{(2)}\gamma _5\right)ฯต,`$
and for the gravitino:
$`\delta \psi _\mu =\left(_\mu +{\displaystyle \frac{1}{4}}\omega _{\alpha \beta ,\mu }\gamma ^\alpha \gamma ^\beta +{\displaystyle \underset{(\sigma )=1,2}{}}K_{ab}^{(\sigma )}\text{T}^{(\sigma )a}A_\mu ^{(\sigma )b}+{\displaystyle \frac{1}{2s}}\mathrm{e}^{2\varphi }\gamma _5_\mu ๐\right)ฯต`$
$`+{\displaystyle \frac{1}{2\sqrt{2}s}}\mathrm{e}^\varphi \left(s^{(1)}+\gamma _5^{(2)}\right)\gamma _\mu ฯต+{\displaystyle \frac{1}{4\sqrt{2}s}}\mathrm{e}^\varphi \left(sg_{(1)}+g_{(2)}\gamma _5\right)\gamma _\mu ฯต.`$ (16)
Here $`^{(\sigma )}=\frac{1}{g_{(\sigma )}}\eta _{ab}^{(\sigma )}\gamma ^\alpha \gamma ^\beta F_{\alpha \beta }^{(\sigma )a}\text{T}^{(\sigma )b}`$ and $`K_{ab}^{(1)}=\eta _{ab}^{(1)}`$, $`K_{ab}^{(2)}=s^2\eta _{ab}^{(2)}`$ with $`\text{T}_a^{(\sigma )}=\text{T}^{(\sigma )a}`$. These expressions complete the dimensional reduction procedure. For $`s=i`$ Eqs.(15) and (S0.Ex6) exactly coincide with the linearized fermionic SUSY variations in the Freedman-Schwarz model , while giving for $`s=1`$ the rules for the Euclidean Freedman-Schwarz model. It is important that the reduction is consistent in the sense that if the four-dimensional configuration is on-shell, then its uplifted version is a solution of the D=10 equations. Also, if the D=4 SUSY variations vanish, then the D=10 variations vanish too.
The FS and EFS models are related to each other via $`sis`$ and not via $`g_{(2)}ig_{(2)}`$. However, both models can be consistently truncated to the static, purely magnetic sector by setting $`A_\mu ^{(2)a}=A_0^{(1)a}=๐=0`$ and requiring $`\frac{}{x^0}`$ to be a hypersurface orthogonal Killing vector. The direct inspection of the field equations for the bosonic Lagrangian (10) and the SUSY variations (15) and (S0.Ex6) reveals then that $`sis`$ becomes equivalent to $`g_{(2)}ig_{(2)}`$. This explains the empirical relation in (3).
Bogomolโnyi equations.โ The D=4 models obtained above admit supersymmetric vacua. These are solutions of the bosonic field equations for which there are such non-trivial spinors $`ฯต`$ that $`\delta \chi =\delta \psi _\mu =0`$. It is rather difficult to directly solve the second order field equations for the bosonic Lagrangian (10), especially with non-trivial gauge fields. The procedure is therefore to focus on the conditions $`\delta \chi =\delta \psi _\mu =0`$, and this gives the system of equations for the $`ฯต`$. These equations are usually inconsistent, but one can look for the consistency conditions. The latter can be formulated as a system of first order equations for the bosonic background, usually called Bogomolโnyi equations. The Bogomolโnyi equations are compatible with the second order field equations, and are sometimes integrable.
To find the Bogomolโnyi equation we first truncate the system to the static, purely magnetic case as described above, and then impose in addition the spherical symmetry:
$$ds^2=s^2\mathrm{e}^{2(\varphi \varphi _{\mathrm{}})}dt^2+\mathrm{e}^{2\lambda }(dr^2+r^2d\mathrm{\Omega }^2),$$
(17)
$$AA^{(1)}=w(\text{T}_2d\theta +\text{T}_1\mathrm{sin}\theta d\phi )+\text{T}_3\mathrm{cos}\theta d\phi .$$
(18)
Here $`\varphi `$, $`\lambda `$ and $`w`$ are functions of the radial coordinate $`r`$, and one can show that the โmetric-dilaton relationโ $`g_{00}=s^2\mathrm{e}^{2(\varphi \varphi _{\mathrm{}})}`$ is consistent with the field equations. The procedure is then to insert (17), (18) into (15) and (S0.Ex6) and set the left-hand sides to zero, which gives a system of equations for $`ฯต`$. The next step is to restrict $`ฯต`$ to the sector with zero total (orbital+spin+isospin) angular momentum. The equations then essentially become algebraic, in which case the consistency conditions can be easily obtained . This gives the Bogomolโnyi equations
$`1+r{\displaystyle \frac{d\lambda }{dr}}`$ $`=`$ $`\sqrt{w^2+{\displaystyle \frac{1}{8}}\mathrm{e}^{2(\lambda +\mathrm{ln}r\varphi )}\left((B1)^2\xi ^2\right)}\nu ,`$
$`Ar{\displaystyle \frac{dw}{dr}}`$ $`=`$ $`2\xi w\nu +\xi ^2(w^21)2w^2(B+1),`$
$`Ar{\displaystyle \frac{d\varphi }{dr}}`$ $`=`$ $`(B+1)(\xi \nu +w(B1)).`$ (19)
Here $`A8w\nu \mathrm{e}^{2(\varphi \lambda \mathrm{ln}r)}+\xi (B1)`$ and $`B2\mathrm{e}^{2(\varphi \lambda \mathrm{ln}r)}(w^21)`$, and also $`\xi =g_{(2)}/s`$. These equations are compatible with the second order field equations for the Lagrangian (10), and for any solution there are two supersymmetry Killing spinors. If $`\xi =0`$ then the number of supersymmetries doubles.
Supersymmetric monopole.โ Let us set in (19) $`\xi =0`$. Notice that the equations then become the same both in the FS and EFS cases. After the substitution
$$w^2=\rho ^2e^{y(\rho )},\frac{1}{2}e^{2(\lambda +\mathrm{ln}r\varphi )}=\rho \frac{dy(\rho )}{d\rho }\rho ^2e^{y(\rho )}1,$$
(20)
the Bogomolโnyi equation (19) become equivalent to one equation
$$\frac{d^2y}{d\rho ^2}=2e^y,$$
(21)
which is integrable. This gives the solution in the closed form:
$$ds^2=2\mathrm{e}^{2\varphi }\left(s^2dt^2+d\rho ^2+Rd\mathrm{\Omega }^2\right),$$
(22)
$$R=2\rho \mathrm{coth}\rho \frac{\rho ^2}{\mathrm{sinh}^2\rho }1,w=\pm \frac{\rho }{\mathrm{sinh}\rho },\mathrm{e}^{2\varphi }=\frac{\mathrm{sinh}\rho }{\sqrt{R}}.$$
(23)
For $`s=1`$ (the EFS case) this solution describes a globally Euclidean background with an essentially non-Abelian gauge field. Since the configuration does not depend on $`t`$, the action is infinite. For $`s=i`$ (the FS case) the solution is of the BPS monopole type with unit magnetic charge. This is geodesically complete and globally hyperbolic . Unfortunately, the ADM mass is infinite and the solution is not asymptotically flat โ due to the dilaton potential. In view of its supersymmetry, it is very plausible that the Lorentzian solution is stable, while for its Euclidean counterpart the notion of dynamical stability makes no sense.
Supersymmetric sphaleron.โ Let us set in (19) $`\xi =1`$, in which case the theory is Euclidean and the dilaton potential vanishes. After some transformations described in the Bogomolโnyi equations can be reduced to one:
$$\frac{1}{2r}\frac{dw}{dr}=\frac{1w^2}{4r^2}\frac{(w+1)^3}{8}+\frac{(w1)^3}{8r^4},$$
(24)
which is invariant under $`r1/r`$, $`ww`$. When the solution $`w(r)`$ is found, the metric function $`\lambda `$ is obtained from
$$\lambda =\mathrm{ln}(2)+_0^r\left(\frac{U+w}{2}1\right)\frac{dr}{r},\mathrm{with}U=\frac{r^2(1+w)^2+(1w)^2}{r^2(1+w)^2(1w)^2},$$
(25)
while the dilaton is given by
$$\varphi =\lambda +\mathrm{ln}(r)+\frac{1}{2}\mathrm{ln}\left(\frac{(U+w)^22w^22}{2(w^21)^2}\right),$$
(26)
which is normalized such that $`\varphi (0)=0`$. Unfortunately, analytical solutions to Eq.(24) are unknown. The numerical integration reveals the existence of a globally regular solution in the interval $`r[0,\mathrm{})`$ which monotonically interpolates between the values specified by the local asymptotic solutions: $`w=1\frac{2}{3}r^2+O(r^4)`$ for $`r0`$ and $`w=1+2\sqrt{2}r^1+O(r^2)`$ as $`r\mathrm{}`$. This gives a globally regular supersymmetric Euclidean solution with non-trivial Yang-Mills field and infinite action.
Now, one can pass to the Lorentzian sector by simply changing the sign of $`dt^2`$ in the metric (17). Of course, this will not lead to the Freedman-Schwarz theory (we do not replace $`sis`$ but just $`tit`$), but rather to a solution of the Einstein-Yang-Mills-Dilaton model. This describes an unstable regular particle-like object with finite ADM mass and purely magnetic Yang-Mills amplitude โ the EYMD sphaleron. Despite its instability, this solution still bears an imprint of supersymmetry as it fulfills the same system of first order Bogomolโnyi equations as its Euclidean counterpart. Passing back to the Euclidean theory, the full supersymmetry is restored, and we call the Euclidean solution โsupersymmetric sphaleronโ.
The above considerations hopefully clarify somewhat the tale of the EYMD sphaleron. The story, however, is not yet finished, as the Bogomolโnyi equations (24) is still not solved analytically. Also it is unclear what happens when the Euclidean solution is uplifted to D=10. In the limits $`r0,\mathrm{}`$ the geometry of the solution is flat and the gauge field vanishes, in which case the uplifted solution is $`E^4\times S^3\times AdS_3`$. This is known to describe the near-horizon geometry of a system of parallel D5โD1 branes. If we consider the full range of $`r`$, then the uplifted solution presumably describes the near-horizon geometry of some deformation of the D5โD1 system. It is interesting to know what this configuration is. Finally, the EYMD sphaleron itself can also be lifted to ten dimensions as a non-supersymmetric solution of heterotic string theory. It is then interesting to know whether there is any relation to the recently discussed โD-sphaleronsโ . |
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